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# 1 Introduction
## 1 Introduction
The topological susceptibility $`\chi `$ plays a relevant role in understanding several low energy properties of QCD. It is defined as the zero–momentum two–point correlation function of the topological charge density $`Q(x)`$ ,
$`\chi `$ $``$ $`{\displaystyle \text{d}^4x_\mu 0|\text{T}\{K^\mu (x)Q(0)\}|0},`$
$`Q(x)`$ $``$ $`{\displaystyle \frac{g^2}{64\pi ^2}}ϵ^{\mu \nu \rho \sigma }F_{\mu \nu }^a(x)F_{\rho \sigma }^a(x),`$ (1)
where $`K^\mu (x)`$ is the Chern current
$$K^\mu (x)\frac{g^2}{16\pi ^2}ϵ^{\mu \nu \rho \sigma }A_\nu ^a(x)\left(_\rho A_\sigma ^a(x)\frac{1}{3}gf^{abc}A_\rho ^b(x)A_\sigma ^c(x)\right),$$
(2)
which satisfies $`Q(x)=_\mu K^\mu (x)`$. In these expressions $`g`$ is the QCD coupling constant and $`f^{abc}`$ are the structure functions of SU(3). In Eq. (1) we keep the derivative out of the expectation value in the definition of $`\chi `$ according to the continuum prescription .
The topological susceptibility in the quenched theory at zero temperature has been widely studied. In particular, a non–zero value of $`\chi `$ provides an explanation for the large mass of the $`\eta ^{}`$ particle, the would–be Goldstone boson of the axial symmetry . In Ref. we have found that the value of $`\chi `$ in pure Yang–Mills theory is $`\chi =(175(5)\text{ MeV})^4`$ in agreement with the phenomenological expectation . For a review of other calculations see and references therein.
In full QCD, Ward identities and current algebra relations imply in the chiral limit
$$\chi =\frac{m}{N_f}\overline{\psi }\psi ,$$
(3)
where $`m`$ is the fermion mass and $`N_f`$ the number of fermions. It is interesting to verify on the lattice the linear dependence on the mass displayed by Eq. (3).
The topological susceptibility at non–zero temperature for pure Yang–Mills theory has also been studied. In Ref. it has been shown that it is rather constant versus $`T`$ below the deconfinement temperature $`T_c`$ and it undergoes a sharp drop across the phase transition. It is also interesting to study the behaviour of $`\chi `$ for full QCD as a function of $`T`$. In particular, Eq. (3) suggests the existence of a drop in the signal of $`\chi `$ across $`T_c`$ also for the full QCD case , (see also ).
In the present paper we will present a lattice determination of $`\chi `$ in full QCD for $`N_f=2`$ degenerate flavours of staggered fermions at finite and zero temperature and for $`N_f=4`$ at finite temperature. The Monte Carlo simulation was done on the APE–Quadrics in Milano and Pisa.
The method used to extract $`\chi `$ from the lattice is described in Section 2. The results and technical details of the simulation for the $`N_f=4`$ and $`N_f=2`$ cases are shown in Section 3 and Section 4 respectively. The conclusions are drawn in the last section.
## 2 Method to determine $`\chi `$
We have simulated the theory on a space–time lattice. The topological charge density has been measured by making use of the lattice operator
$$Q_L(x)\frac{1}{2^9\pi ^2}\underset{\mu \nu \rho \sigma =\pm 1}{\overset{\pm 4}{}}\stackrel{~}{ϵ}_{\mu \nu \rho \sigma }\text{Tr}\left(\mathrm{\Pi }_{\mu \nu }(x)\mathrm{\Pi }_{\rho \sigma }(x)\right),$$
(4)
where $`\mathrm{\Pi }_{\mu \nu }(x)`$ is a plaquette in the $`\mu `$$`\nu `$ plane and $`\stackrel{~}{ϵ}_{\mu \nu \rho \sigma }`$ is a generalized Levi–Civita tensor which acquires an extra minus sign for each of its indices going negative. Green functions containing $`Q_L(x)`$ differ from the continuum counterparts by a finite renormalization $`Z`$ , in short
$$Q_L(x)=Za^4Q(x)+O(a^6),$$
(5)
where $`a`$ is the lattice spacing. Actually $`Q_L(x)`$ mixes with fermionic operators during renormalization but it has been shown that the off–diagonal mixings are negligible .
The lattice topological susceptibility $`\chi _L`$ is defined as
$$\chi _L\frac{Q_L^2}{V},$$
(6)
where $`V`$ is the lattice space–time volume and $`Q_L`$ is the total topological charge, $`Q_L_xQ_L(x)`$. The definition of $`\chi _L`$ differs from the continuum expression for $`\chi `$ in Eq. (1) by contact terms which must be subtracted. The relationship between $`\chi `$ and $`\chi _L`$ is
$$\chi _L=a^4Z^2\chi +M,$$
(7)
where $`M`$ contains mixings with operators with compatible quantum numbers.
We have measured $`Q_L(x)`$ after having applied two smearing steps with smearing parameter $`c=0.9`$ . The smearing process allows the operator to be less sensitive to the quantum fluctuations. Therefore the renormalization constant $`Z`$ approaches 1 and $`M`$ approaches 0. As a consequence the statistical errors are strongly diminished .
We evaluate $`Z`$ and $`M`$ by using the “heating method” which provides a non–perturba-tive determination of the renormalization constants . To calculate $`Z`$, several steps of an updating algorithm (the same used during the Monte Carlo simulation) are applied on a configuration containing a charge +1 classical instanton. These updatings thermalize first the quantum fluctuations, responsible for the renormalization effects, and due to the slowing down leave the topological content unchanged. The measurement of $`Q_L`$ on such updated configurations yields $`ZQ`$. As $`Q`$ is known, one can extract $`Z`$. Notice that this procedure is equivalent to imposing the continuum value for the 1–instanton charge (in the $`\overline{\mathrm{MS}}`$ scheme it is +1) and extracting the finite multiplicative renormalization $`Z`$ by evaluating a matrix element of $`Q_L`$.
The additive renormalization $`M`$ is obtained in a similar way. We apply a few heating steps with the same updating algorithm as above on a zero–field configuration. The topological susceptibility is then measured. This provides the value of $`M`$, if no instantons have been created during the few updating steps. As explained in , cooling tests must be done to check that the background topological charge has not been changed during the heating. Configurations where the topological content has been changed are discarded from the sample.
Once $`Z`$ and $`M`$ are known, we can extract $`a^4\chi `$ from Eq. (7). We have fixed $`a`$ in physical units by measuring the string tension $`\sigma `$ on a $`16^4`$ lattice and assuming that $`\sqrt{\sigma }=420`$ MeV. To this purpose, Wilson loops ranging from $`1\times 1`$ to $`8\times 8`$ have been evaluated by using smeared spatial links and a best fit has been performed for the interquark potential to the form $`V(r)=V_0+a_0/r+\sigma r`$ .
## 3 Full QCD with $`N_f=4`$
We have numerically simulated four flavours of quarks by using staggered fermions of bare mass $`am=0.05`$. For the pure gauge sector of the theory the Wilson action has been chosen. The $`\mathrm{\Phi }`$–type HMC algorithm has been used for the updating. We have simulated at the following values of $`\beta 6/g^2`$=5.00, 5.02, 5.04, 5.05 5.06 and 5.10 ($`g`$ is the lattice bare coupling).
The lattice volume was $`16^3\times 4`$. According to Ref. , at this temporal size the deconfinement transition occurs at $`\beta _c=5.04`$. We have checked this value by computing the Polyakov loop and the chiral condensate as shown in Fig. 1.
In the asymmetric lattice the temperature has been determined by the definition $`T=1/(aL_t)`$, where $`L_t`$ is the temporal size.
We have checked that the distribution of topological charge is thermalized enough. In Fig. 2 we show such a distribution for $`\beta =5.04`$ after 30 cooling steps. At smaller values of $`am`$ ($`am0.01`$) and higher $`\beta `$’s, simulations are affected by a slowing down in the sampling .
In Table 1 we list the numerical results. The lattice spacing $`a`$ has been determined at $`\beta =`$5.00, 5.04 and 5.10 on a $`16^4`$ lattice. For the other values of $`\beta `$, it was extrapolated by splines. From the value at $`\beta _c`$, $`a=0.30(2)`$ fm, we infer the critical temperature to be $`T_c=164(11)`$ MeV.
The behaviour of the ratio $`\chi (T)/\chi (T=0)`$ as a function of $`T/T_c`$ is shown in Fig. 3. The case $`N_f=2`$ will be discussed in the next section. For $`N_f=4`$ the value at zero temperature $`\chi (T=0)`$ is computed as the average of all results at $`TT_c`$. For comparison, we have included in the figure the data for the quenched theory taken from Ref. . There is clear evidence that the topological susceptibility in full QCD with $`N_f=4`$ presents a drop at $`T_c`$. Apparently the drop is sharper than in the case of pure Yang–Mills. However, the two theories have been simulated at rather different lattice volumes ($`32^3\times 8`$ in the pure Yang-Mills case). Moreover, at the small values of $`\beta `$ used in our present simulation, we expect violations of scaling in the ratio $`T/T_c`$. As a consequence we prefer not to draw definite conclusions about the relative slopes.
## 4 Full QCD with $`N_f=2`$
Full QCD with two flavours of quarks has been simulated by using staggered fermions of bare mass $`am=0.0125`$ and the usual Wilson action for the pure gauge sector. The updating has been performed with the R–type HMC algorithm , each trajectory consisting in 120 steps of $`\mathrm{\Delta }\tau =0.004`$ units of molecular dynamics time. We have simulated at $`\beta `$=5.40, 5.50, 5.55, 5.60 and 5.70 both at zero and finite temperature.
The lattice volume was $`32^3\times 8`$ for finite $`T`$ and $`16^4`$ for $`T=0`$. According to Ref. at this temporal size the deconfinement transition occurs at $`\beta _c=5.54(2)`$. We have checked this number by studying the Polyakov loop in a similar way as in section 3.
We have checked that the sets of configurations produced during the simulation are well decorrelated and the distribution of topological charge is well sampled. For the case of $`\beta =5.70`$ we needed a longer simulation to obtain a well sampled distribution of topological charge.
In Table LABEL:tab:nf2 we report the data obtained from the simulation at finite $`T`$. From the value extrapolated at $`\beta _c=5.54`$, we infer the critical temperature to be $`T_c=194(10)(15)`$ MeV where the first error is our statistical and the second one the error induced by the indetermination in the value of $`\beta _c`$.
In Fig. 3 the behaviour of the ratio $`\chi (T)/\chi (T=0)`$ as a function of $`T/T_c`$ is shown. For $`N_f=2`$ the value at zero temperature $`\chi (T=0)`$ has been obtained from an independent simulation on symmetric lattices $`16^4`$ at the same values of $`\beta `$ (see Table LABEL:tab:nf2mass). There is a clear drop in the signal when the deconfinement temperature is crossed.
In Fig. 4 we plot the topological susceptibility obtained at zero temperature for the five values of $`\beta `$. The corresponding set of data is listed in Table LABEL:tab:nf2mass. Having fixed $`am`$, the values of the bare mass vary as shown in the upper horizontal scale. A fit with a constant value of the topological susceptibility gives $`\chi =(163\pm 6)^4\mathrm{MeV}^4`$, with the statistical test $`\text{(chi–square)}/\mathrm{d}.\mathrm{o}.\mathrm{f}=0.37`$; a fit with a linear homogeneous dependence on $`m`$, like Eq. (3), gives $`\overline{\psi }\psi _{\mathrm{BARE}}=(6.2\pm 0.8)\times 10^7\mathrm{MeV}^3`$ and $`\text{(chi–square)}/\mathrm{d}.\mathrm{o}.\mathrm{f}=0.94`$. Therefore our present data are compatible with both behaviours, and we are not able, within our errors, to check Eq. (3).
## 5 Conclusions
We have studied the topological susceptibility $`\chi `$ in full QCD with 2 and 4 flavours of dynamical fermions. At zero temperature, we have not enough precision to check the chiral limit, Eq. (3).
At finite temperature, $`\chi `$ stays constant for values of $`T`$ up to the deconfinement temperature $`T_c`$. At $`T_c`$ it presents a sudden drop and $`\chi `$ becomes compatible with zero at $`T1.5T_c`$ for $`N_f=2`$ and $`T1.2T_c`$ for $`N_f=4`$. This behaviour is qualitatively similar to that found in Ref. for the quenched case.
## Acknowledgements
B.A. thanks the Theory Group in Pisa for the kind hospitality. Financial support from EC, Contract FMRX–CT97–0122, and from MURST is acknowledged.
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# 1 Introduction
## 1 Introduction
It is a long-standing problem to prove the quark confinement. There are some scenarios to explain this phenomenon qualitatively, one of which is widely known as a dual QCD vacuum scenario. However, these are not sufficient proofs and we need further efforts to investigate this phenomenon.
Recently several authors have proposed a novel reformulation of QCD as a perturbative deformation of a TQFT. By the use of this reformulation, a confining-deconfining phase transition has been investigated. This work is based on the Kugo-Ojima (KO) confinement criterion and the gauge fixing is performed in an $`OSp(4|2)`$ type Feynman gauge. Then, the TQFT sector becomes a 2D chiral model through the Parisi-Sourlas (PS) mechanism, equivalently a 2D $`O(4)`$ NLSM. When it is extended to the finite temperature in the real-time formalism, it has a deconfining phase under the twisted boundary conditions through the spontaneous symmetry breaking (SSB) of an $`O(4)`$ symmetry but not under the periodic ones. This result has been showed by the calculation of the effective potential. In general, the SSB in 2D systems is forbidden by the Coleman-Mermin-Wagner’s theorem. This theorem is based on the infrared divergence peculiar in 2D systems. But, if the twisted boundary conditions are imposed, this divergence is softened and the phase transition can take place by the SSB. Therefore, a deconfining (massless) phase can appear. This phase is retained as if pertubative parts are added although the boundary conditions are slightly modified.
The same analysis can be adopted to the case of the $`OSp(4|2)`$ type MAG. There, the role of the boundary conditions has not been clarified. The purpose of our study is to make it clear. In this case, the TQFT sector becomes a 2D coset model, equivalently a 2D $`O(3)`$ NLSM, and we can reach the similar conclusion about the phase structure in the TQFT sector. But the phase structure of the TQFT sector is not retained and only a deconfining phase can survive when the perturbative parts are added. The effect of these parts replaces the twisting factors in the twisted boundary conditions by the unit element and so the twisted boundary conditions become equivalent to the periodic ones. Moreover, we can show that the linear potential remains in the full $`\mathrm{QCD}_4`$ if we assume the $`U(1)`$ Abelian dominance. This linear potential means the quark confinement in the Wilson criterion. Then, a Polyakov loops’ correlator decays exponentially at large distance. This result also implies a confining phase. It may seem inconsistent to the case of the Feynman type gauge because the only difference is the gauge fixing. However, we notice that we have not considered the $`U(1)`$ background. It turns out that the information about the phase structure is encoded in it.
This paper is organized as follows. In section 2, the reformulation of the $`\mathrm{QCD}_4`$ at finite temperature is introduced. In section 3, we discuss the effect of the boundary conditions in the TQFT sector and investigate how it is modified when the perturbative parts are added. In section 4, we comment on the phase structure of the $`U(1)`$ background. This argument is based on the work and the topological object plays an important role. Finally, in section 5, we explain our results and discuss future problems.
## 2 Reformulation of $`\mathrm{QCD}_4`$ at Finite Temperature
We start with a finite temperature partition function
$$Z_{\mathrm{FT}}=_{\mathrm{periodic}}[dA_\mu ][dC][d\overline{C}][dB]\mathrm{exp}\left\{\frac{1}{2g_{\mathrm{YM}}^2}\mathrm{Tr}_GF_{\mu \nu }^2(A)i\delta _\mathrm{B}G_{\mathrm{GF}+\mathrm{FP}}[A_\mu ,C,\overline{C},B]\right\}$$
(2.1)
where $`F_{\mu \nu }\stackrel{def}{=}_\mu \mathrm{A}_\nu _\nu \mathrm{A}_\mu \mathrm{i}[\mathrm{A}_\mu ,\mathrm{A}_\nu ]`$ and $`\delta _\mathrm{B}`$ induces the BRST transformation
$`\delta _\mathrm{B}A_\mu `$ $`=`$ $`D_\mu [A]C,\delta _\mathrm{B}C=iC^2,`$ (2.2)
$`\delta _\mathrm{B}\overline{C}`$ $`=`$ $`iB,\delta _\mathrm{B}B=0.`$
The $`C(\overline{C})`$ is a (anti-)ghost of the system and the $`B`$ is a Nakanishi-Lautrup field. We consider the case of the gauge group $`G=SU(N)`$ without quark fields. The gauge field $`A_\mu `$ is expressed in terms of new fields $`V_\mu `$ and $`U`$
$$A_\mu =UV_\mu U^{}+\frac{i}{g_{\mathrm{YM}}}U_\mu U^{}.$$
(2.3)
We shall use the Faddeev-Popov (FP) trick and insert a unit in the path integral
$$1=\mathrm{\Delta }[A]\frac{1}{N}\underset{k=0}{\overset{N1}{}}_{B_k}[dU]\underset{x}{}\delta \left(^\mu A_\mu ^{U^1}(x)\right)$$
(2.4)
where the $`\mathrm{\Delta }[A]`$ is the FP determinant and the $`B_k`$ denotes the boundary condition of the $`U(x)`$
$$B_k:U(i\beta ,𝐱)=U(0,𝐱)e^{ik\pi /N}(\mathrm{for}k=0,\mathrm{},N1).$$
(2.5)
This condition is fully discussed in the next section. Eq.(2.4) can be rewritten as
$$1=\frac{1}{N}\underset{k=0}{\overset{N1}{}}_{B_k}[dU]_{\mathrm{periodic}}[d\gamma ][d\overline{\gamma }][d\beta ]\mathrm{exp}\left\{id^4x[i\stackrel{~}{\delta }_\mathrm{B}\stackrel{~}{G}_{\mathrm{GF}+\mathrm{FP}}(V_\mu ,\gamma ,\overline{\gamma },\beta )]\right\}$$
(2.6)
where the new BRST transformation $`\stackrel{~}{\delta }_\mathrm{B}`$ acts on the fields $`V_\mu ,\gamma ,\overline{\gamma }`$ and $`\beta `$ as
$`\stackrel{~}{\delta }_\mathrm{B}V_\mu `$ $`=`$ $`D_\mu [V]\gamma ,\stackrel{~}{\delta }_\mathrm{B}\gamma =i\gamma ^2,`$
$`\stackrel{~}{\delta }_\mathrm{B}\overline{\gamma }`$ $`=`$ $`i\beta ,\stackrel{~}{\delta }_\mathrm{B}\beta =0.`$ (2.7)
Here we used a formula of the $`\stackrel{~}{G}_{\mathrm{GF}+\mathrm{FP}}`$
$$\stackrel{~}{G}_{\mathrm{GF}+\mathrm{FP}}\stackrel{def}{=}\mathrm{Tr}_\mathrm{G}(\overline{\gamma }^\mu \mathrm{V}_\mu ).$$
(2.8)
Thus, we can obtain the following partition function
$`Z_{\mathrm{FT}}`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle _{B_k}}[dU]{\displaystyle _{\mathrm{periodic}}}[dV_\mu ][dC][d\overline{C}][dB][d\gamma ][d\overline{\gamma }][d\beta ]`$ (2.9)
$`\mathrm{exp}\{i{\displaystyle }d^4x[{\displaystyle \frac{1}{2g_{\mathrm{YM}}}}\mathrm{Tr}_GF_{\mu \nu }^2(V)i\stackrel{~}{\delta }_\mathrm{B}\stackrel{~}{G}_{\mathrm{GF}+\mathrm{FP}}[V_\mu ,\gamma ,\overline{\gamma },\beta ]`$
$`i\delta _\mathrm{B}G_{\mathrm{GF}+\mathrm{FP}}[{\displaystyle \frac{i}{g_{\mathrm{YM}}}}U_\mu U^{}+UV_\mu U^{},C,\overline{C},B]]\}.`$
Next, let us specify the gauge fixing term $`G_{\mathrm{GF}+\mathrm{FP}}`$. In the work, an $`OSp(4|2)`$ type MAG is used
$$G_{\mathrm{GF}+\mathrm{FP}}=\overline{\delta }_\mathrm{B}\mathrm{Tr}_{G/H}[A_\mu ^2+2iC\overline{C}].$$
(2.10)
The $`H`$ is the maximal Abelian subgroup of the G. On the other hand, an $`OSp(4|2)`$ type Feynman gauge is utilized in the work
$$G_{\mathrm{GF}+\mathrm{FP}}=\overline{\delta }_\mathrm{B}\mathrm{Tr}_G[A_\mu ^2+2iC\overline{C}].$$
(2.11)
These gauge fixing conditions lead to the following 2D TQFT sectors through the PS mechanism, respectively
$`S_{\mathrm{TQFT}}`$ $`=`$ $`{\displaystyle \frac{\pi }{g_{\mathrm{YM}}^2}}{\displaystyle d^2x\mathrm{Tr}_{G/H}[_\mu U_\mu U^{}]}(\mathrm{Coset}\mathrm{model}),`$ (2.12)
$`S_{\mathrm{TQFT}}`$ $`=`$ $`{\displaystyle \frac{\pi }{g_{\mathrm{YM}}^2}}{\displaystyle d^2x\mathrm{Tr}_G[_\mu U_\mu U^{}]}(\mathrm{Chiral}\mathrm{model}).`$ (2.13)
It should be noted that the difference of the two models is associated with the degrees of freedom in the maximal torus part $`H`$. In particular, the weak coupling limit ($`g_{\mathrm{YM}}0`$) of the finite temperature $`\mathrm{QCD}_4`$ is described by the TQFT sector with summing over all the boundary conditions
$$Z_{\mathrm{FTTQFT}}=\frac{1}{N}\underset{k=0}{\overset{N1}{}}_{B_k}[dU]_{\mathrm{periodic}}[dC][d\overline{C}][dB]\mathrm{exp}\{iS_{\mathrm{TQFT}}\}.$$
(2.14)
## 3 Study of Boundary Conditions
We restrict ourselves to the case of $`G=SU(2)`$ for simplicity. The real-time and imaginary-time formalisms are standard methods to deal with the finite temperature system. Both formalisms extend the time coordinate to a complex time and the fields obey the (anti-)periodic boundary conditions. In particular, the gauge field obeys a periodic condition
$$A_\mu (i\beta ,𝐱)=A_\mu (0,𝐱).$$
(3.1)
The gauge field in the TQFT sector has a pure gauge type configuration, and the twisted boundary conditions are allowed:
$$B_g:U(i\beta ,𝐱)=U(0,𝐱)g(\mathrm{for}g\mathrm{global}SU(2)).$$
(3.2)
Let us consider the TQFT sectors in eqs.(2.12), (2.13). It is proved that the TQFT sector (2.13), equivalently an $`O(4)`$ NLSM, has a deconfining phase under the twisted boundary conditions ($`g\mathrm{𝟏}`$) in the real-time formalism. But when one imposes usual periodic boundary conditions ($`g=\mathrm{𝟏}`$), this phase does not appear. Also, this phase transition cannot take place in the imaginary-time formalism. This point is discussed later in this section.
Next, we investigate the phase structure of the TQFT sector (2.12). Our main purpose is to study it. The coset model action (2.12) can be rewritten asWe omit the ghost term.
$$S_{\mathrm{TQFT}}=\frac{\pi }{g_{\mathrm{YM}}^2}d^2x[_\mu 𝐧_\mu 𝐧],𝐧𝐧=1.$$
(3.3)
This is the action of an $`O(3)`$ NLSM. Here we used the Euler angle representation of the $`SU(2)`$ matrix
$`U(x)`$ $`=`$ $`\mathrm{exp}\{i\chi (x)\sigma _3/2\}\mathrm{exp}\{i\theta (x)\sigma _2/2\}\mathrm{exp}\{i\phi (x)\sigma _3/2\}`$ (3.6)
$`=`$ $`\left(\begin{array}{cc}\mathrm{exp}\{\frac{i}{2}(\phi +\chi )\}\mathrm{cos}\frac{\theta }{2}& \mathrm{exp}\{\frac{i}{2}(\phi \chi )\}\mathrm{sin}\frac{\theta }{2}\\ \mathrm{exp}\{\frac{i}{2}(\phi \chi )\}\mathrm{sin}\frac{\theta }{2}& \mathrm{exp}\{\frac{i}{2}(\phi +\chi )\}\mathrm{cos}\frac{\theta }{2}\end{array}\right)`$
and we parameterize the unit vector field $`𝐧(x)`$ ($`𝐧:𝐑^2S^2`$) as
$$𝐧\stackrel{def}{=}\left(\begin{array}{c}\mathrm{n}^1(\mathrm{x})\\ \mathrm{n}^2(\mathrm{x})\\ \mathrm{n}^3(\mathrm{x})\end{array}\right)\stackrel{def}{=}\left(\begin{array}{c}\mathrm{sin}\theta (\mathrm{x})\mathrm{cos}\phi (\mathrm{x})\\ \mathrm{sin}\theta (\mathrm{x})\mathrm{sin}\phi (\mathrm{x})\\ \mathrm{cos}\theta (\mathrm{x})\end{array}\right).$$
(3.7)
Then, we have to determine the boundary condition on the $`𝐧(x)`$. Useful relations can be used in the calculationWe normalize the generators $`T^A`$’s of the SU(2) as $`\mathrm{Tr}_G(T^AT^B)=\frac{1}{2}\delta ^{AB}`$. :
$$n^A(x)T^A=U^{}(x)T^3U(x),n^A(x)=2\mathrm{T}\mathrm{r}_G[U(x)T^AU^{}(x)T^3](A=1,2,3).$$
(3.8)
We find that the $`𝐧(x)`$ is invariant under the $`U(1)`$ transformation generated by the $`T^3`$ and it can be rotated by generators associated with the coset $`SU(2)/U(1)`$. Also, eq.(2.12) has the following global $`SU(2)_\mathrm{L}`$ symmetry
$$UUh,hSU(2)_\mathrm{L}.$$
(3.9)
Then the $`𝐧(x)`$ transforms as
$$n^AT^An^A(h^{}T^Ah)\stackrel{def}{=}\mathrm{n}_{}^{}{}_{}{}^{\mathrm{A}}\mathrm{T}^\mathrm{A}$$
(3.10)
and we easily find an action of the element $`h`$ on $`n^A`$
$$n^An_{}^{}{}_{}{}^{A}=\underset{B=1}{\overset{3}{}}\mathrm{ad}(h^{})_B^An^B.$$
(3.11)
This means that the n$`(x)`$ is transformed under the $`SO(3)`$ rotation but it is invariant by an action of the center $`𝐙_2`$ of the $`SU(2)`$. By the use of eq.(3.11), the boundary condition on the $`U(x)`$ can be translated into that on the field n$`(x)`$
$$n^A(i\beta ,x_1)=\underset{B=1}{\overset{3}{}}\mathrm{ad}(g^{})_B^An^B(0,x_1),$$
(3.12)
where $`x_1`$ is a spatial coordinate. We can calculate the effective potential under this condition and show that the $`O(3)`$ NLSM has a deconfining (massless) phase in the real-time formalism under the twisted boundary conditions (ad($`g^{}`$)$`\mathrm{𝟏}`$). On the other hand, under the periodic boundary conditions (ad$`(g^{})=\mathrm{𝟏}`$), it has not this phase as in the case of the Feynman type gauge. Also, no massless phase appears in imaginary-time formalism.
The above discussion is limited to the TQFT sector. That is, the gauge field configuration is the pure gauge one. What happens when one incorporates the gauge adjoint part (i.e.the perturbative part)? This part replaces the twisted boundary conditions (3.2) by the following ones
$$B_k:U(i\beta ,𝐱)=U(0,𝐱)e^{ik\pi /2}(k=0,1).$$
(3.13)
Then, the $`𝐧(x)`$ obeys the periodic boundary conditions:
$$B_k:𝐧(i\beta ,x_1)=𝐧(0,x_1)(k=0,1).$$
(3.14)
Therefore, all the TQFT sectors have the periodic boundary conditions
$`{\displaystyle \frac{1}{N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle _{B_k}}[dU]`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{k=0}{\overset{N1}{}}}{\displaystyle _{B_k}}[dn]{\displaystyle _{\mathrm{periodic}}}[dn](N=2),`$
$`{\displaystyle \frac{1}{N}}(\mathrm{periodic}+\mathrm{twisted}+\mathrm{})\mathrm{periodic}.`$
That is, the sum over the boundary conditions becomes an integral of the periodic field $`𝐧(x)`$. Because a massless phase cannot exist under the periodic conditions, only a confining phase can exist. Once the gauge adjoint part is added, this phase of the TQFT sector is screened. We conclude that the phase structure of the TQFT sector is hidden in the full $`\mathrm{QCD}_4`$ with the MAG.
Moreover, if we assume the Abelian dominance, we can interpret this result in terms of the Polyakov loop. Then, we can see that the correlator of Polyakov loops $`P`$ and $`P^{}`$ is equivalent to the expectation value of an Abelian Wilson loop $`W`$ in the TQFT sector as shown in Figure 1
$$P^{}(𝐱)P(0)_{\mathrm{TQFT}}\stackrel{def}{=}\mathrm{W}_{\mathrm{TQFT}}=\mathrm{exp}(\beta \mathrm{V}).$$
(3.16)
Here the $`V`$ is a static potential or a free energy. We can reconstruct the expectation value in the full $`\mathrm{QCD}_4`$ $`P^{}(𝐱)P(0)`$ from eq.(3.16). Then, the Coulomb potential is added to the V.
In the $`|𝐱|\mathrm{}`$ limit, we obtain a decomposition
$$P^{}(𝐱)P(0)|P|^2.$$
(3.17)
From this formula, we conclude that if $`P=0`$, the free energy blows up for the large separation of the quarks($`|𝐱|`$$`\mathrm{}`$). We can interpret this result as a signal of the confinement:
$$P=0,(\mathrm{confinement}).$$
(3.18)
On the other hand, if $`P0`$, then the free energy of a static quark-antiquark pair approaches a constant for the limit $`|𝐱|`$$`\mathrm{}`$. We can interpret this as an evidence of the deconfinement:
$$P0,(\mathrm{deconfinement}).$$
(3.19)
The TQFT sector in the MAG, $`O(3)`$ NLSM has instanton solutions and it can be shown that the instanton gas induces the linear potential in the $`V`$. So we obtain the $`P=0`$ and this result implies the confinement in the full $`\mathrm{QCD}_4`$.<sup>§</sup><sup>§</sup>§ In this case, we would have to use the imaginary-time formalism.
Thus, we find that only a confining phase can appear in the analysis of the TQFT sector and so this point is different from the case of the Feynman type gauge. We will explain this point in the next section.
Here, we discuss the two methods to deal with the finite temperature system: the real-time formalism and the imaginary-time one. It is amazing that these methods lead to different results as commented in the work when they are applied to the same system. That is, the phase transition occurs under the twisted boundary conditions in the real-time formalism but not in the imaginary-time one. The KO confinement criterion is only applicable to the system with a continuous 4-momentum. Otherwise the imaginary-time formalism can not be applied and we have no contradiction. We may suspect to confront an issue in the MAG case. However, it would not so. It is the reason that the physical parts (gauge adjoint parts) really exist in the real world and only a confining phase can appear. Therefore, we cannot observe this difference of the two methods. We conclude that the difference of the two methods is not relevant to the phase structure of the $`\mathrm{QCD}_4`$ in the MAG case.
## 4 Comments on Gauge Fixing and $`U(1)`$ Background
In this section, we comment on the relation between the phase transition mechanism and the $`U(1)`$ background. First, let us recall that we have chosen the $`OSp(4|2)`$ type MAG. This gauge spoils the chiral $`SU(2)_\mathrm{L}SU(2)_\mathrm{R}`$ symmetry which exists in the Feynman type gauge. Eventually, only the $`SU(2)_\mathrm{L}`$ symmetry survives and the $`SU(2)_\mathrm{R}`$ symmetry is lost in the MAG case. Then, we obtain a topological object “monopole” peculiar to the MAG and it is absent in the Feynman gauge. This is an $`U(1)`$ Abelian mopole. Because the chiral symmetry is broken, we would not be able to apply the KO criterion. However, we have the topological object in our hand. Monopoles are interpreted as instantons in the 2D theory reduced by the PS mechanism in this case. These play an important role in the confinement of the Wilson criterion. It is also conjectured that these play a central role in the phase transition as well. In the above study, we saw that the phase structure of the TQFT sector was screened and only a confining phase can appear. However, in the case of the MAG, the information about the phase transition mechanism is encoded in the $`U(1)`$ background and the Berezinskii-Kosterlitz-Thouless (BKT) transition occurs by the condensation of the topological object. This would be clarified after we integrate out the non-diagonal gauge components. That manipulation is based on the $`U(1)`$ Abelian dominance and really was done in the work. Then, we can obtain an Abelian projected effective theory. This is just the $`U(1)`$ theory but is asymptotically free. When the scheme of the reformulation is applied to this $`U(1)`$ theory, the TQFT sector of it becomes a 2D $`O(2)`$ NLSM with a vortex solution. It is widely known that a BKT phase transition is induced by the condensation of the vortexes in this $`O(2)`$ NLSM. In this theory, by using parameterizations of the fields $`U`$ and $`𝐒`$
$$U(x)=\mathrm{e}^{i\phi (x)},𝐒(x)=(\mathrm{cos}\phi (x),\mathrm{sin}\phi (x)),$$
(4.1)
we obtain the TQFT sector of the effective $`U(1)`$ theory:
$$S_{\mathrm{TQFT}}=\frac{\pi }{g_{\mathrm{YM}}^2}d^2x_\mu 𝐒(x)_\mu 𝐒(x).$$
(4.2)
Also, the $`𝐒`$ obeys the boundary condition at finite temperature. But we can easily show that the boundary conditions on it becomes equivalent to the periodic ones when the perturbative parts are added. And so we cannot state the existence of the phase transition by the SSB in the TQFT sector as well. But, at least, we can find a BKT type phase transition by the $`U(1)`$ vortex condensation from the confining phase to the deconfining one. Moreover, it is commented in the work that the $`U(1)`$ vortex is closely relevant to the instanton in the $`O(3)`$ NLSM. This possibly implies that the vortex condensation induces an instanton condensation and the confining potential would vanish. This would just correspond to the phase transition to the deconfining phase.
## 5 Discussion and Conclusion
We have investigated the effect of the boundary conditions and found that it depends on the gauge fixings. Because the analysis based on the KO criterion does not depend on the complicated dynamical information about the $`\mathrm{QCD}_4`$, we could not derive the mechanism of the confinement directly. On the other hand, when we choose the MAG, we can understand the dynamical mechanism in terms of monopoles. But because the analysis depends on the special gauge, we cannot know how reliable it is. Also, we could not discuss the Higgs phase. It is an interesting problem to study this phase by the use of the reformulation of $`\mathrm{QCD}_4`$
The phase structure of the $`\mathrm{QCD}_4`$ has investigated in the $`OSp(4|2)`$ type MAG and Feynman gauge. These investigations are based on the toy models where the special $`OSp(4|2)`$ type gauge fixing is utilized. But we could explain the results based on the KO criterion in the framework of the Wilson criterion, and believe that our considerations shed light on an interrelation between the works based on the KO criterion and those on the Wilson criterion. We hope the paper proceeds the understanding of the quark confinement.
Acknowledgements
The author would like to thank W. Souma for valuable discussions and K. Sugiyama for careful reading the manuscript and useful comments.
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# Quantum Dot between Two Superconductors
## Abstract
Novel effects emerge from an interplay between multiple Andreev reflections and Coulomb interaction in quantum dot coupled to superconducting leads and subject to a finite potential bias $`V`$. Combining an intuitive physical picture with rigorous path integral formalism we evaluate the current $`I`$ through the dot and find that the interaction shifts the subharmonic pattern of the $`IV`$ curve is shifted toward higher $`V`$. For sufficiently strong interaction the subgap current (at $`eV<2\mathrm{\Delta }`$) is virtually suppressed.
Recent progress in nanotechnology enables the fabrication and experimental investigation of superconducting contacts of atomic size with few conducting channels . Transport properties of such systems are essentially determined by the mechanism of multiple Andreev reflections (MAR) which is responsible for Josephson current as well as for dissipative currents at subgap voltages. Theoretical analysis of MAR and current-voltage characteristics in small superconducting junctions is reported in a number of papers . In these works, an essential ingredient is the assumption that electron-electron interaction inside the contact can be neglected. It might indeed be justified provided a metallic contact is sufficiently large and/or strongly coupled to massive superconducting leads.
However, in very small contacts (quantum dots), the Coulomb interaction is not effectively screened, hence it is expected to substantially affect transport properties of the system. For instance, it is well known both from theory and experiment that tunneling through a quantum dot between superconductors can virtually be suppressed due to Coulomb effects. Thus, to the fascinating physics of SNS and SIS junctions, one should add that of an SAS junction composed of superconducting leads coupled by an interacting quantum dot.
In the present work, the physics of interplay between MAR and interaction effects in SAS junctions subject to a finite bias is exposed. It encodes the salient features of superconductivity, strong correlations and non-linear response. A simple intuitive physical picture is combined with a rigorous path integral technique by which irrelevant degrees of freedom are eliminated and an effective action is constructed (in the spirit of Feynman-Vernon influence functional ). Similar ideas proved to be useful elsewhere, see e.g. . In the present context they have been applied for the relatively simple case of an $`SAS`$ junction at zero bias, focusing on the equilibrium Josephson current . Our main achievements are: a) Derivation of a tractable expression for the non-linear tunneling current in the presence of interactions. b) prediction of novel physical effects pertaining to the $`IV`$ curve of an SAS junction at sub-gap bias.
Let us then commence with a simple and physically transparent picture of an interplay between MAR and Coulomb effects. Consider a quasiparticle (hole) which suffers $`n`$ Andreev reflections inside the superconducting junction, thereby gaining an energy $`neV`$, where $`V`$ is the voltage bias. As soon as $`nev=2\mathrm{\Delta }`$ the quasiparticle leaves the junction and does not contribute anymore to the subgap current. Hence, the number of Andreev reflections $`n`$ for a given voltage is $`n2\mathrm{\Delta }/eV`$.
Assume now that the Coulomb interaction inside the junction is switched on. For our qualitative discussion it suffices to account for it in terms of an effective capacitance $`C`$ and its related charging energy $`E_C=e^2/2C`$. At $`T=0`$ and for $`eVE_C`$ a single electron tunneling (and, hence, also MAR) is blocked, so in what follows we will consider the case $`eV>E_C>0`$. In order to leave the junction the quasiparticle should gain an energy $`neV`$ equal to $`2\mathrm{\Delta }+(n+1)E_C`$. The last term originates from the fact that during the MAR cycle with a given $`n`$ the charge $`(n+1)e`$ is transferred between the electrodes. Hence, an additional energy $`(n+1)E_C`$ should be paid. The above condition immediately fixes the number $`n`$ at a given voltage:
$$n=\left[\frac{2\mathrm{\Delta }+E_C}{eVE_C}\right].$$
(1)
Thus, in the presence of Coulomb interaction quasiparticles spend more time inside the junction and suffer more Andreev reflections. Another obvious observation is that at low temperature $`T`$, the transfer of the charge $`(n+1)e`$ is blocked by interaction at voltages $`eV(n+1)E_C`$. Combining this observation with (1) one arrives at the condition
$$eVeV_{th}=E_C\left(1+\sqrt{1+\frac{2\mathrm{\Delta }}{E_C}}\right),$$
(2)
under which the MAR current is suppressed due to Coulomb repulsion. For $`E_C\mathrm{\Delta }`$ the voltage threshold is $`eV_{th}\sqrt{2\mathrm{\Delta }E_C}E_C`$, i.e. in this case MAR should be blocked even at voltages much higher than $`E_C/e`$. For $`eV_{th}2\mathrm{\Delta }`$ eq. (2) yields
$$E_C2\mathrm{\Delta }/3.$$
For such values of $`E_C`$ one expects the subgap current to be fully suppressed due to Coulomb interaction.
Let us now recall that the subharmonic peaks occur on the $`IV`$ curves each time the MAR cycle with a given $`n`$ becomes impossible. Without interaction these peaks are located at voltages $`V_n=2\mathrm{\Delta }/en`$. It follows immediately from the above discussion that in the presence of Coulomb interaction the peaks should be shifted to higher voltages. From eq. (1) one finds
$$V_n=\frac{E_C}{e}+\frac{2\mathrm{\Delta }+E_C}{en},$$
(3)
i.e. one expects the subharmonic peaks to be shifted by $`\delta V_n=E_C(1+1/n)`$ towards larger $`V`$ as compared to the noninteracting case.
Thus, already a naive analysis of the interplay between MAR and Coulomb effects allows one to predict several novel effects which can be experimentally tested. To put these qualitative arguments on a firm basis we formulate below a realistic model of an SAS junction, and proceed with a rigorous calculation of the $`IV`$ characteristics.
The model and basic formalism. Consider, in two dimensions, a quantum dot at $`𝐫=0`$ weakly coupled to (half planar) superconducting electrodes. The Hamiltonian of the system is decomposed as,
$$𝑯=𝑯_L+𝑯_R+𝑯_{\mathrm{dot}}+𝑯_\mathrm{t}.$$
(4)
The Hamiltonians of the left ($`x<0`$) and right ($`x>0`$) superconducting electrodes have the standard BCS form
$`𝑯_j={\displaystyle }d𝒓[\mathrm{\Psi }_{j\sigma }^{}(𝒓)\xi ()\mathrm{\Psi }_{j\sigma }(𝒓)`$ (5)
$`\lambda \mathrm{\Psi }_j^{}(𝒓)\mathrm{\Psi }_j^{}(𝒓)\mathrm{\Psi }_j(𝒓)\mathrm{\Psi }_j(𝒓)].`$ (6)
Here $`\mathrm{\Psi }_{j\sigma }^{}`$ ($`\mathrm{\Psi }_{j\sigma }`$) are the electron creation (annihilation) operators, $`\xi ()=^2/2m\mu `$, and $`j=L,R`$ for left and right electrodes. The dot itself is modeled as an Anderson impurity center with Hamiltonian
$$𝑯_{\mathrm{dot}}=ϵ_0\underset{\sigma }{}C_\sigma ^{}C_\sigma +UC_{}^{}C_{}C_{}^{}C_{},$$
(7)
where $`C_\sigma ^{}`$ and $`C_\sigma `$ are dot electron operators. The impurity site energy $`ϵ_0`$ (counted from the Fermi energy $`\mu `$) is assumed to be far below the Fermi level $`ϵ_0<0`$. The presence of a strong Coulomb repulsion $`U>ϵ_0`$ between electrons in the same orbital guarantees that the dot is at most singly occupied.
Electron tunneling through the dot is described by the term,
$$𝑯_\mathrm{t}=𝒯\underset{j=L,R}{}\underset{\sigma }{}\mathrm{\Psi }_{j\sigma }^{}(0)C_\sigma +\mathrm{h}.\mathrm{c}.,$$
(8)
where $`𝒯`$ is an effective transfer amplitude.
The dynamics of the system is completely contained within the evolution operator on the Keldysh contour $`K`$ which consists of forward and backward oriented time paths. Its kernel $`J`$ is given by a path integral,
$$J=𝒟\overline{\mathrm{\Psi }}𝒟\mathrm{\Psi }𝒟\overline{C}𝒟C\mathrm{exp}(iS),$$
(9)
over Grassman fields corresponding to the fermion operators, with $`\overline{\mathrm{\Psi }}=(\mathrm{\Psi }_L^{},\mathrm{\Psi }_L^{},\mathrm{\Psi }_R^{},\mathrm{\Psi }_R^{})`$ with obvious definitions for $`\mathrm{\Psi }`$, $`\overline{C}`$ and $`C`$. Moreover, $`S=_KL𝑑t`$ is the action and $`L`$ is the Lagrangian pertaining to the Hamiltonian (4).
In order to avoid dealing with fields defined on both branches of the Keldysh contour one performs a rotation $`Cc`$ and $`\mathrm{\Psi }\psi `$ in Keldysh space:
$`\overline{c}=\overline{C}\sigma _z\widehat{Q}^1,c=\widehat{Q}C;\widehat{Q}`$ $`={\displaystyle \frac{1}{\sqrt{2}}}`$ $`\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)`$ (12)
and similarly for $`\overline{\psi }`$ and $`\psi `$. Here $`\sigma _z`$ is the third Pauli matrix operating in Keldysh space. The new Grassman variables $`\overline{c}`$, $`c`$, $`\overline{\psi }`$, $`\psi `$ are now defined solely on the forward time branch. Averages of the corresponding products of these fields determine the standard $`2\times 2`$ Green-Keldysh matrix composed of retarded ($`\widehat{G}^R`$), advanced ($`\widehat{G}^A`$) and Keldysh ($`\widehat{G}^K`$) Green functions which, in turn, are $`2\times 2`$ matrices in spin (Nambu) space.
The path integral (9) is expressed in terms of the new Grassman variables in the same way, and the action $`S`$ is now defined as $`S=S_{\mathrm{dot}}+S_0[\overline{\psi },\psi ]`$, where
$$S_{\mathrm{dot}}=𝑑t\left[\overline{c}\left(i\frac{}{t}\stackrel{~}{ϵ}\tau _z\right)c+\frac{U}{2}(\overline{c}c)^2\right],$$
(13)
$`S_0={\displaystyle }dt{\displaystyle \underset{j=L,R}{}}[{\displaystyle _j}d𝐫\overline{\psi }_j(𝐫,t)\widehat{G}_j^1\psi _j(𝐫,t)`$ (14)
$`+(𝒯\overline{\psi }_j(0,t)\tau _zc(t)+\mathrm{c}.\mathrm{c}.)],`$ (15)
where $`\stackrel{~}{ϵ}=ϵ_0+U/2`$ and the Pauli matrices $`\tau _{x,y,z}`$ act in Nambu space. The operator $`\widehat{G}_{L,R}^1`$ has the standard form
$$\widehat{G}_{L,R}^1(\xi )=i\frac{}{t}\tau _z\xi ()+\tau _+\mathrm{\Delta }_{L,R}+\tau _{}\mathrm{\Delta }_{L,R}^{},$$
(16)
where $`\tau _\pm =(\tau _x\pm i\tau _y)/2`$ and $`\mathrm{\Delta }_{L,R}`$ are the (spatially constant) BCS order parameters of the electrodes.
Effective action and transport current. The basic algorithm of our approach is to integrate out the electron variables in the superconducting electrodes which play the role of an effective environment for the dot. This procedure yields the influence functional $`F[\overline{c},c]`$ for the $`c`$-fields in the dot:
$$F\mathrm{exp}(iS_{\mathrm{env}}[\overline{c},c])=𝒟\overline{\psi }𝒟\psi \mathrm{exp}(iS_0[\overline{\psi },\psi ]),$$
(17)
which is evaluated exactly. Gaussian integration in (17) is carried out separately for $`L`$\- and $`R`$-electrodes.
Let us consider, say, the left superconductor and omit the subscript $`j=L`$ for the moment. The first step is to integrate out the fermion fields inside the superconductor thereby arriving at an intermediate effective action in terms of the fermion fields defined on the surface $`x=0`$. It is useful at this point to Fourier transform the fields $`\psi (x,y)`$ along the (translationally invariant) $`y`$ direction. The problem then reduces to a one dimensional one with fermion fields $`\psi _k(x)`$ where $`k`$ is the quasiparticle momentum in the direction normal to $`x`$. In order to evaluate the Gaussian integral we will look for a saddle point field $`\stackrel{~}{\psi }_k(x)`$ defined by $`\widehat{G}^1(\xi _x)\stackrel{~}{\psi }_k(x)=0,`$ where $`\xi _x=(1/2m)(^2/x^2)\mu _k`$ and $`\mu _k=\mu k^2/2m`$.
Decomposing $`\stackrel{~}{\psi }`$ into bulk and surface fields $`\stackrel{~}{\psi }_k(x)=\psi _k^b(x)+\psi _k(0)`$ and integrating out $`\psi _k^b(x)`$ we arrive at the intermediate effective action $`\stackrel{~}{S}`$ of a superconductor lead expressed only via the $`\psi `$-fields at the surface,
$$\stackrel{~}{S}=i𝑑t𝑑t^{}\underset{k}{}\frac{v_x}{2}\overline{\psi }_k(0,t)\tau _z\widehat{g}(t,t^{})\psi _k(0,t^{}).$$
(18)
Here $`v_x=\sqrt{2\mu _k/m}`$ and
$$\widehat{g}(t,t^{})=e^{\frac{i\phi (t)\tau _z}{2}}\widehat{g}(ϵ)e^{iϵ(tt^{})}\frac{dϵ}{2\pi }e^{\frac{i\phi (t^{})\tau _z}{2}},$$
(19)
is the Green-Keldysh-Eilenberger matrix of the (left) superconducting electrode
$`\widehat{g}`$ $`=`$ $`\left(\begin{array}{cc}\widehat{g}^R& \widehat{g}^K\\ \widehat{0}& \widehat{g}^A\end{array}\right),`$ (22)
$`\phi (t)=\phi _0+2e^tV(t_1)𝑑t_1`$ is the time-dependent phase of the superconducting order parameter and $`V(t)`$ is the electric potential of the electrode. The Fourier transformed retarded and advanced Eilenberger functions have the standard form
$$\widehat{g}^{R/A}(ϵ)=\frac{(ϵ\pm i0)\tau _z+i|\mathrm{\Delta }|\tau _y}{\sqrt{(ϵ\pm i0)^2|\mathrm{\Delta }|^2}},$$
(23)
and $`\widehat{g}^K=(\widehat{g}^R\widehat{g}^A)\mathrm{tanh}(ϵ/2T)`$ is the Keldysh function.
The second step in our derivation amounts to integrating out the $`\psi `$-fields on the surface. The integral
$$𝒟\overline{\psi }(0)𝒟\psi (0)\mathrm{exp}(i\stackrel{~}{S}+idt(𝒯\overline{\psi _k}(0)\tau _zc+\mathrm{c}.\mathrm{c}.))$$
(24)
can easily be evaluated. Carrying out exactly the same procedure for the right electrode, making use of the identity $`\widehat{g}_{L,R}^2=1`$ and adding up the results we obtain
$$S_{\mathrm{env}}=i\mathrm{\Gamma }𝑑t𝑑t^{}\overline{c}(t)\tau _z\widehat{g}_+(t,t^{})c(t^{}).$$
(25)
Here and below we define $`\mathrm{\Gamma }=4_k𝒯^2/v_x`$ and $`\widehat{g}_\pm =(\widehat{g}_L\pm \widehat{g}_R)/2`$.
Eq. (25) is one of our central results. It enables the expression of the kernel $`J`$ (9) solely in terms of the fields $`\overline{c}`$ and $`c`$:
$$J=𝒟\overline{c}𝒟c\mathrm{exp}(iS_{\mathrm{dot}}+iS_{\mathrm{env}}),$$
(26)
where $`S_{\mathrm{dot}}+S_{\mathrm{env}}S_{\mathrm{eff}}[\overline{c},c]`$ represents the effective action for a quantum dot between two superconductors.
In order to complete our derivation let us express the current through the dot in terms of the correlation function for the variables $`\overline{c}`$ and $`c`$. Starting from the general expression for the current and representing the correlator for the $`\psi `$-fields in terms of that for the $`c`$-fields we find,
$$I=\frac{e\mathrm{\Gamma }}{2}\mathrm{Tr}[\widehat{g}_{}\overline{c}c|_K+\mathrm{h}.\mathrm{c}.].$$
(27)
Thus, the problem of calculating the current through an interacting quantum dot is reduced to that of finding the correlator $`\overline{c}c`$ in the model defined by the effective action $`S_{\mathrm{eff}}=S_{\mathrm{dot}}+S_{\mathrm{env}}`$ (13), (25). It should be emphasized that our approach is appropriate for studying both equilibrium and nonequilibrium electron transport. In the noninteracting limit $`U0`$ the results of previous studies can be easily recovered within our formalism.
Mean field approximation. Consider now the case $`U0`$ and decouple the interacting term in (13) by means of Hubbard-Stratonovich transformation introducing additional scalar fields $`\gamma _\pm `$. The kernel $`J`$ now reads,
$$J=𝒟\overline{c}𝒟c𝒟\gamma _+𝒟\gamma _{}\mathrm{exp}(iS[\gamma ]+iS_{\mathrm{eff}}|_{U=0}),$$
(28)
$$S[\gamma ]=𝑑t\left(\overline{c}\gamma _+\sigma _xc+\overline{c}\gamma _{}c\frac{2}{U}\gamma _+\gamma _{}\right).$$
(29)
Here we will assume that the effective Kondo temperature $`T_K`$ =$`\sqrt{U\mathrm{\Gamma }}\mathrm{exp}[\pi |ϵ_0|/2\mathrm{\Gamma }]`$ is smaller than the superconducting gap $`\mathrm{\Delta }`$. In this case interactions can be accounted for within the mean field approximation. The fields $`\gamma _\pm `$ in (29) are considered as time-independent parameters determined self-consistently from the saddle point conditions $`\delta J/\delta \gamma _\pm =0`$:
$$\gamma _+=\frac{U}{2}𝑑t<\overline{c}c>,\gamma _{}=\frac{U}{2}𝑑t<\overline{c}\sigma _xc>.$$
(30)
As it turned out from our numerical analysis the effect of the parameter $`\gamma _+`$ is merely the renormalization of the tunneling rate $`\mathrm{\Gamma }`$. Absorbing $`\gamma _+`$-terms in $`\mathrm{\Gamma }`$ we arrive at the final effective action of our model
$`S_{\mathrm{eff}}[\gamma ]`$ $`=`$ $`{\displaystyle \frac{dϵ}{2\pi }𝑑ϵ^{}\overline{c}\widehat{M}(ϵ,ϵ^{})c},`$ (31)
$`\widehat{M}(ϵ,ϵ^{})`$ $`=`$ $`\delta (ϵϵ^{})(ϵ+\gamma _{}\tau _z\stackrel{~}{ϵ})+i\tau _z\mathrm{\Gamma }\widehat{g}_+(ϵ,ϵ^{}).`$ (32)
Subgap current in SAS junctions. In order to find the correlator $`\overline{c}c=i\widehat{M}^1`$ and the current (27) we numerically inverted the matrix (32) and simultaneously solved the self-consistency equation for $`\gamma _{}`$ (30). The resulting $`IV`$ characteristics for an $`SAS`$ junction in the presence of Coulomb interaction are displayed in Fig.1. One observes all the main features predicted within our simple picture of an interplay between MAR and Coulomb interaction: (i) at relatively low voltages $`VV_{th}`$ MAR current is essentially suppressed due to interaction, (ii) for higher voltages (but still $`eV<2\mathrm{\Delta }`$) MAR is possible and results in a nonzero subgap current which increases with $`V`$ and (iii) the subharmonic peaks in the differential conductance occur and are shifted to higher voltages as compared to the noninteracting case. An increase of $`U`$ results in a stronger current suppression and a more pronounced shift of the subharmonic peaks. Close to the gap edge $`eV=2\mathrm{\Delta }`$ the current shoots up sharply.
The parameters used in our numerical analysis are chosen in a way to observe all the key features (i), (ii) and (iii). It is interesting to quantitatively compare the results presented in Fig. 1 with the predictions (1)-(3) of an oversimplified “$`E_C`$-based” model. Let us estimate the effective value of $`E_C`$ (which is, strictly speaking, a function of $`U`$, $`ϵ_0`$ $`\mathrm{\Gamma }`$ and $`V`$ in our calculation) with the aid of eq. (2) and the $`IV`$ curves of Fig. 1. We find $`E_C0.2\mathrm{\Delta }`$ for $`U=2.4\mathrm{\Delta }`$ and $`E_C0.25\mathrm{\Delta }`$ for $`U=2.7\mathrm{\Delta }`$. Obviously, these values of $`E_C`$ are smaller than $`2\mathrm{\Delta }/3`$ and, hence, a finite subgap conductance is expected at $`eV\sqrt{2\mathrm{\Delta }E_C}E_C`$. This is precisely what we observe in Fig. 1. For such values of $`V`$ and $`E_C`$ eq. (1) yields $`n3`$, i.e. only two subharmonic peaks (with $`n=2,3`$) can occur. This is exactly the case in Fig. 1. Finally, substituting the above values of $`E_C`$ into eq. (3) we can estimate the magnitudes of the peak shifts $`\delta V_n`$. For $`U=2.4\mathrm{\Delta }`$ we find $`\delta V_n0.3\mathrm{\Delta }`$ for $`n=2`$ and $`0.26\mathrm{\Delta }`$ for $`n=3`$. Analogous values for $`U=2.7\mathrm{\Delta }`$ are respectively $`0.38\mathrm{\Delta }`$ and $`0.33\mathrm{\Delta }`$. These values are in a reasonably good agreement with our numerical results.
In summary, we presented a detailed analysis of an SAS junction at finite bias and derived its effective action using Keldysh path-integral techniques. Our approach applies for both equilibrium and nonequilibrium current transport in the presence of interactions. The repulsive Coulomb interaction leads to novel effects in the pattern of the subgap current. In particular, it shifts the peaks of the differential conductance toward larger bias. When the interaction is sufficiently strong the subgap current is highly suppressed. Our theoretical predictions can be directly tested in experiments with superconducting quantum dots.
We acknowledge useful discussions with J.C. Cuevas and J. von Delft. This research is supported by DIP German Israel Cooperation project, by the Israeli Science Foundation grant Center of Excellence and by the US-Israel BSF.
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# ON THE LOW-ENERGY LIMIT OF STRING AND M-THEORY**footnote *Lectures given at TASI97, Boulder, Colorado, USA June 1997. Published in: Supersymmetry, Supergravity and Supercolliders, Editor J. A. Bagger, World Scientific 1999, page 709
## 1 Introduction
The supersymmetric extension of the standard model of weak, eloctromagnetic and strong interactions still contains many parameters whose origin has to be understood. It is the general hope that an explanation of these parameters can be found in a more complete and fundamental theory. In such a theory one would hope to understand how the various coupling constants are unified. One also expects hints to understand the nature of supersymmetry breakdown and its consequences for the soft breaking parameters of the theory. Such a theory could be string theory. Especially in the framework of the heterotic string one is confident to have identified a candidate theory that could manifest itself as the supersymmetric standard model in its low energy limit. In these lectures, we shall try to discuss such questions.
TASI97 has been devoted to the various aspects of supersymmetric models. Not much has been said about string theory so far. So I have to introduce the basics of string theory in these lectures. I shall not repeat this discussion here in the written up version, since it has already been published elsewhere $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. A pretty complete picture of the status of string theory can been found in last years TASI lectures, which I encourage you to consult. To follow the lectures here, it might be useful to consider $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ for notations and conventions. Especially $`^\mathrm{?}`$ might be useful to learn the basics of constructing $`d=4`$ low energy effective supergravity theories from higher dimensional string theories. In our discussion we shall try to concentrate on the main points avoiding technical complications. We shall see that some of the aspects of the field are not so difficult to understand as one might have previously thought.
The part that is presented in this written version of the lectures concerns mainly two important aspects of the field that are at the center of current research: the question of the unification of gauge and gravitational coupling constants and the question of supersymmetry breakdown. In the first part we shall dicuss gauge coupling constants in string theory. In general, they are not constants but functions of so-called moduli fields, whose vacuum expectation values (vevs) are undetermined at the classical level. An understanding of the actual value of the couplings will thus require a determination of these vevs and require the consideration of supersymmetry breakdown. Still we shall first discuss these questions at the classical level. We start with the (weakly coupled) heterotic $`E_8\times E_8`$ string. The $`d=4`$ dimensional effective low energy theory is discussed in the approximation obtained through the method of reduction and truncation. The method is very simple, but allows us to understand the qualitative properties of these models. At the classical level the gauge coupling constants are determined through the so-called dilaton field and one obtains a definite relation between the gauge coupling and the gravitational coupling constant. We comment on this relation and discuss explicitely how quantum effects can modify this classical relation.
The absolute value of the coupling constants requires knowledge about the vev of the dilaton field and thus we have to consider the breakdown of supersymmetry. Here the mechanism of gaugino condensation is discussed in detail. While some aspects of this mechanism look very encouraging, there remain some problems, most notably the run-away behaviour of the dilaton. We discuss some attempts to shed light on that problem. After that we explicitely discuss various aspects of gauge coupling unification in the heterotic string.
With the recent progress in string dualities, a new picture of unification might emerge. It is closely connected to a conjectured $`d=11`$ dimensional theory known as M-theory. The $`E_8\times E_8`$ version of that theory on an 11-dimensional interval might be especially interesting in that respect. The second part of these lectures is devoted to a discussion of the possible low energy manifestations of that theory. Again, we first review its implications for the question of unification, that look even more promising than the one found in the previously considered version of the heterotic string. The question of the breakdown of supersymmetry, though, looks very similar in both cases, with one exception concerning the soft breaking parameters. It solves a long standing problem of the smallness of the observable sector gaugino masses. We comment on the various phenomenological properties of the models obtained in this framework. This will include a discussion of a possible nonuniverality of the soft scalar masses and its relevance for flavour changing neutral currents, the size of the neutralino masses and the question of the lightest supersymmetric particle in connection with the critical density of the universe.
## 2 Dynamical coupling constants
Most physics models contain coupling constants as free parameters that can be adjusted to fit the experimental values. In more complete theories one could imagine that the values of these parameters are determined dynamically by the theory itself. The question arises, how and why the coupling parameters take the values that are observed in nature. String theory $`^\mathrm{?}`$ provides an example for such a class of models. At first sight, one would then expect that it is rather easy to see whether a given string model has a chance to be realistic or not. Just compute the physical coupling constants and then check whether they conincide with the measured values. Looking more closely we find, however, that the situation is more complicated. First of all, it would be very difficult to do the computations needed. Secondly, the coupling constants might not be determined yet at the classical level. In string theory, the coupling constants are functions of so-called moduli fields and the actual values of the coupling parameters are determined through the vacuum expectation values (vev) of these fields. This is not only true for the gauge couplings, but also for the Yukawa couplings as well as the radii and other properties of the compactification parameters.
We shall concentrate here on the gauge coupling constants first. The heterotic string gives a definite prediction for the gauge coupling function as well as the functional relation between gauge couplings and gravitational couplings. At tree level the universal gauge coupling constant $`g_{\mathrm{string}}`$ is determined by the vev of the dilaton field $`^\mathrm{?}`$
$$S=\frac{4\pi }{g_{\mathrm{string}}^2}+i\frac{\theta }{2\pi }.$$
(1)
Nonuniversalities can appear at the one-loop level and depend on further moduli fields $`T`$, $`U`$ or $`B`$ that parametrize the properties of the compactification parameters as well as the Wilson lines.
$$\frac{1}{g_a^2(\mu )}=\frac{k_a}{g_{\mathrm{string}}^2}+\frac{b_a}{16\pi ^2}\mathrm{ln}\frac{M_{\mathrm{string}}^2}{\mu ^2}\frac{1}{16\pi ^2}\mathrm{}(T,U,B\mathrm{}).$$
(2)
Given this situation we have then to see how a theory with a realistic set of gauge coupling constants can emerge.
We would therefore like to connect such a theory with a low energy effective field theory describing the known particle physics phenomena. A prime candidate is a low energy supergravitational generalization of the standard model of strong and electroweak interactions, with supersymmetry broken in a hidden sector $`^{\mathrm{?},\mathrm{?}}`$. Unification of observable sector gauge couplings might appear as required in supersymmetric grand unified theories (SUSYGUTs) $`^\mathrm{?}`$. The vevs of the moduli fields $`S`$, $`T`$, $`\mathrm{}`$ should then determine gauge couplings in hidden and observable sector, including the correct values of the QCD coupling $`\alpha _s`$ and the weak mixing angle $`\mathrm{sin}^2\theta _\mathrm{W}`$.
In string theories with unbroken supersymmetry the vevs of the moduli fields are undetermined. A first step in the determination of gauge coupling constants requires therefore the discussion of supersymmetry breakdown. We then have to see how the vevs of the moduli are fixed. Of course, not any value of the moduli will lead to a satisfactory model. In fact we shall face some generic problems concerning the actual values of the coupling constants. We have to understand why the value of
$$\alpha _{\mathrm{string}}=\frac{g_{\mathrm{string}}^2}{4\pi }\frac{1}{20},$$
(3)
whereas a natural expectation for the vev of $`S`$ would be a number of order 1 or maybe 0 or even infinity, as happens in many simple models. A related question concerns the possibility of so-called string unification leading to the correct prediction of the weak mixing angle $`\mathrm{sin}^2\theta _\mathrm{W}`$. Naively one would expect string unification to appear at $`M_{\mathrm{string}}4\times 10^{17}`$GeV, while the correct prediction of $`\mathrm{sin}^2\theta _\mathrm{W}`$ seems to lead to a scale that is a factor of 20 smaller. We then have to face the question how such a situation can be achieved with natural values of the vevs of the moduli fields.
These are the questions we want to address in the next sections. We shall need to start with the discussion of supersymmetry breakdown. Here we shall concentrate in the framework of gaugino condensation. This will then lead us and include a discussion of the problem of a ”runaway” dilaton that any attempt of a dynamical determination of coupling constants has to face. The next question then concerns the value of $`<S>1`$ and its compatibility with weak coupling. Finally we shall discuss new results concerning string threshold corrections and the question of string unification $`^\mathrm{?}`$ in the framework of the weakly coupled heterotic string. After that we shall discuss alternative possibilities.
## 3 Gaugino Condensation
One of the prime motivations to consider the supersymmetric extension of the standard model is the stability of the weak scale ($`M_W`$) of order of a TeV in the presence of larger mass scales like a GUT-scale of $`M_X10^{16}GeV`$ or the Planck scale $`M_{Pl}10^{18}GeV`$. The size of the weak scale is directly related to the breakdown scale of supersymmetry, and a satisfactory mechanism of supersymmetry breakdown should explain the smallness of $`M_W/M_{Pl}`$ in a natural way. One such mechanism is based on the dynamical formation of gaugino condensates that has attracted much attention since its original proposal for a spontaneous breakdown of supergravity $`^{\mathrm{?},\mathrm{?}}`$. In the following we shall address some open questions concerning this mechanism in the framework of low energy effective superstring theories $`^{\mathrm{?},\mathrm{?}}`$.
Before discussing these detailed questions let us remind you of the basic facts of this mechanism. For simplicity we shall consider here a pure supersymmetric ($`N=1`$) Yang-Mills theory, with the vector multiplet $`(A_\mu ,\lambda )`$ containing gauge bosons and gauge fermions in the adjoint representation of the nonabelian gauge group. Such a theory is asymptotically free and we would therefore (in analogy to QCD) expect confinement and gaugino condensation at low energies $`^\mathrm{?}`$. We are then faced with the question whether such a nontrivial gaugino condensate $`\mathrm{<}\lambda \lambda \mathrm{>}0`$ leads to a breakdown of supersymmetry. A first look at the SUSY-transformation on the composite fermion $`\lambda \sigma ^\mu A_\mu `$ $`^\mathrm{?}`$
$$\{Q,\lambda \sigma ^\mu A_\mu \}=\lambda \lambda +\mathrm{}$$
(4)
might suggest a positive answer, but a careful inspection of the multiplet structure and gauge invariance leads to the opposite conclusion. The bilinear $`\lambda \lambda `$ has to be interpreted as the lowest component of the chiral superfield $`W^\alpha W_\alpha =(\lambda \lambda ,\mathrm{})`$ and therefore a non-vanishing vev of $`\lambda \lambda `$ does not break SUSY $`^\mathrm{?}`$. This suggestion is supported by index-arguments $`^\mathrm{?}`$ and an effective Lagrangian approach $`^\mathrm{?}`$. We are thus lead to the conclusion that in such theories gaugino condensates form, but do not break global (rigid) supersymmetry.
Not all is lost, however, since we are primarily interested in models with local supersymmetry including gravitational interactions. The weak gravitational force should not interfere with the formation of the condensate; we therefore still assume $`\mathrm{<}\lambda \lambda \mathrm{>}=\mathrm{\Lambda }^30`$. This expectation is confirmed by the explicit consideration of the effective Lagrangian of ref. $`^\mathrm{?}`$ in the now locally supersymmetric framework. We here consider a composite chiral superfield $`U=(u,\psi ,F_u)`$ with $`u=\mathrm{<}\lambda \lambda \mathrm{>}`$. In this toy model $`^{\mathrm{?},\mathrm{?}}`$ we obtain the surprising result that not only $`\mathrm{<}u\mathrm{>}=\mathrm{\Lambda }^30`$ but also $`\mathrm{<}F_u\mathrm{>}0`$, a signal for supersymmetry breakdown. In fact
$$\mathrm{<}F_u\mathrm{>}=M_S^2=\frac{\mathrm{\Lambda }^3}{M_{Pl}},$$
(5)
consistent with our previous result that in the global limit $`M_{Pl}\mathrm{}`$ (rigid) supersymmetry is restored. For a hidden sector supergravity model we would choose $`M_S10^{11}GeV`$ $`^\mathrm{?}`$.
Still more information can be obtained by consulting the general supergravity Lagrangian of elementary fields determined by the Kähler potential $`K(\mathrm{\Phi }_i,\mathrm{\Phi }_{}^{j}{}_{}{}^{})`$, the superpotential $`W(\mathrm{\Phi }_i)`$ and the gauge kinetic function $`f(\mathrm{\Phi }_i)`$ for a set of chiral superfields $`\mathrm{\Phi }_i=(\varphi _i,\psi _i,F_i)`$. Non-vanishing vevs of the auxiliary fields $`F_i`$ would signal a breakdown of supersymmetry. In standard supergravity notation these fields are given by
$$F_i=\mathrm{exp}(G/2)(G^1)_i^jG_j+\frac{1}{4}\frac{f}{\mathrm{\Phi }_k}(G^1)_i^k\lambda \lambda +\mathrm{},$$
(6)
where the gaugino bilinear appears in the second term $`^\mathrm{?}`$. This confirms our previous argument that $`\mathrm{<}\lambda \lambda \mathrm{>}0`$ leads to a breakdown of supersymmetry, however, we obtain a new condition: $`f/\mathrm{\Phi }_i`$ has to be nonzero, i.e. the gauge kinetic function $`f(\mathrm{\Phi }_i)`$ has to be nontrivial. In the fundamental action $`f(\mathrm{\Phi }_i)`$ multiplies $`W_\alpha W^\alpha `$ which in components leads to a form $`\text{Re}f(\varphi _i)F_{\mu \nu }F^{\mu \nu }`$ and tells us that the gauge coupling is field dependent. For simplicity we consider here one modulus field $`M`$ with
$$\mathrm{<}\text{Re}f(M)\mathrm{>}1/g^2.$$
(7)
This dependence of $`f`$ on the modulus $`M`$ is very crucial for SUSY breakdown via gaugino condensation. $`f/M0`$ leads to $`F_M\mathrm{\Lambda }^3/M_{Pl}`$ consistent with previous considerations. The goldstino is the fermion in the $`f(M)`$ supermultiplet. In the full description of the theory it might mix with a composite field, but the inclusion of the composite fields should not alter the qualitative behaviour discussed here. As we shall see in a moment, an understanding of the mechanism of SUSY breakdown via gaugino condensation is intimately related to the question of a dynamical determination of the gauge coupling constant as the vev of a modulus field. We would hope that in a more complete theory such questions could be clarified in detail.
The candidate at our disposal for such a theory is the $`E_8\times E_8`$ heterotic string. The second $`E_8`$ (or a subgroup thereof) could serve as the hidden sector gauge group and it was soon found $`^\mathrm{?}`$ that there we have nontrivial $`f=S`$ where $`S`$ represents the dilaton superfield. The heterotic string thus contains all the necessary ingredients for a successful implementation of the mechanism of SUSY breakdown via gaugino condensation $`^{\mathrm{?},\mathrm{?}}`$. Also the question of the dynamical determination of the gauge coupling constant can be addressed. A simple reduction and truncation $`^\mathrm{?}`$ from the $`d=10`$ theory leads to the following scalar potential $`^\mathrm{?}`$
$$V=\frac{1}{16S_RT_R^3}\left[|W(\mathrm{\Phi })+2(S_RT_R)^{3/2}(\lambda \lambda )|^2+\frac{T_R}{3}\left|\frac{W}{\mathrm{\Phi }}\right|^2\right],$$
(8)
where $`S_R=\text{Re}S`$, $`T_R=\text{Re}T`$ is the modulus corresponding to the overall radius of compactification and $`W(\mathrm{\Phi })`$ is the superpotential depending on the matter fields $`\mathrm{\Phi }`$. The gaugino bilinear appears via the second term in the auxiliary fields (6). To make contact with the dilaton field, observe that $`\mathrm{<}\lambda \lambda \mathrm{>}=\mathrm{\Lambda }^3`$ where $`\mathrm{\Lambda }`$ is the renormalization group invariant scale of the nonabelian gauge theory under consideration. In the one-loop approximation
$$\mathrm{\Lambda }=\mu \mathrm{exp}\left(\frac{1}{bg^2(\mu )}\right),$$
(9)
with an arbitrary scale $`\mu `$ and the $`\beta `$-function coefficient $`b`$. This then suggests
$$\lambda \lambda e^f=e^S$$
(10)
as the leading contribution (for weak coupling) for the functional $`f`$-dependence of the gaugino bilinear Relation (10) is of course not exact. For different implementations see $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. The qualitative behaviour of the potential remains unchanged..
In the potential (8) we can then insert (10) and determine the minimum. In our simple model (with $`W/T=0`$) we have a positive definite potential with vacuum energy $`E_{vac}=0`$. Suppose now for the moment that $`\mathrm{<}W(\mathrm{\Phi })\mathrm{>}0`$ In many places in the literature it is quoted incorrectly that $`\mathrm{<}W(\mathrm{\Phi })\mathrm{>}`$ is quantized in units of the Planck length since $`W`$ comes from $`H`$, the field strength of the antisymmetric tensor field $`B`$ and $`H=\text{d}B\omega _{3Y}+\omega _{3L}`$ ($`\omega `$ being the Chern-Simons form). Quantization is expected for $`\mathrm{<}\text{d}B\mathrm{>}`$ but not necessarily for $`H`$..
$`S`$ will now adjust its vev in such a way that $`|W(\mathrm{\Phi })+2(S_RT_R)^{3/2}(\lambda \lambda )|=0`$, thus
$$|W(\mathrm{\Phi })+2(S_RT_R)^{3/2}\mathrm{exp}(S)|=0.$$
(11)
This then leads to broken SUSY with $`E_{vac}=0`$ and a fixed value of the gauge coupling constant $`g^2\mathrm{<}\text{Re}S\mathrm{>}^1`$. For the vevs of the auxiliary fields we obtain $`F_S=0`$ and $`F_T0`$ with important consequences for the pattern of the soft SUSY breaking terms in phenomenologically oriented models $`^\mathrm{?}`$, which we shall discuss later.
Thus a satisfactory picture seems to emerge. However, we have just discussed a simplified example. In general we would expect also that the superpotential depends on the moduli, $`W/T0`$ and, including this dependence, the modified potential would no longer be positive definite and one would have $`E_{vac}<0`$.
But even in the simple case we have a further vacuum degeneracy. For any value of $`W(\mathrm{\Phi })`$ we obtain a minimum with $`E_{vac}=0`$, including $`W(\mathrm{\Phi })=0`$. In the latter case this would correspond to $`\mathrm{<}\lambda \lambda \mathrm{>}=0`$ and $`S\mathrm{}`$. This is the potential problem of the runaway dilaton. The simple model above does not exclude such a possibility. In fact this problem of the runaway dilaton does not seem just to be a problem of the toy model, but more general. It always appears in models that are continuously connected to a regime that is asymptotically free in the ultraviolet. The problem could then be avoided only when there are other minima for finite $`S`$. Alternatively one could consider a situation which is not connected to infinitely weak coupling (e.g. a theory that not asymptotically free). But such a situation we have excluded in this lectures from the beginning.
One attempt to avoid this problem with additional minima at finite $`S`$ was the consideration of several gaugino condensates $`^\mathrm{?}`$, but it still seems very difficult to produce satisfactory potentials that lead to a dynamical determination of the dilaton for reasonable values of $`\mathrm{<}S\mathrm{>}`$. In some cases it even seems impossible to fine tune the cosmological constant to zero. In absence of a completely satisfactory model it is then also difficult to investigate the detailed phenomenological properties of the approach. Here it would be of interest to know the actual size of the vevs of the auxiliary fields $`\mathrm{<}F_S\mathrm{>}`$, $`\mathrm{<}F_T\mathrm{>}`$ and $`\mathrm{<}F_U\mathrm{>}`$. In the models discussed so far one usually finds $`\mathrm{<}F_T\mathrm{>}`$ to be the dominant term, but it still remains a question whether this is true in general.
## 4 Fixing the dilaton
This is a very general problem that is unsolved at the moment. In the following we give a speculation of how the problem could be solved within the assumptions made above. It is not the only attempt to solve the problem, but it might point out some aspects of the problem that open a new way to look at it in the framework of the recent developments in string theory.
It seems that we need some new ingredient before we can understand the mechanism completely. The resolution of all these problems might comes with a better understanding of the form of the gauge kinetic function $`f`$ $`^{\mathrm{?},\mathrm{?}}`$. In all the previous considerations one assumed $`f=S`$. How general is this relation? Certainly we know that in one loop perturbation theory $`S`$ mixes with $`T`$ $`^\mathrm{?}`$, but this is not relevant for our discussion and, for simplicity, we shall ignore that for the moment. The formal relation between $`f`$ and the condensate is given through $`\mathrm{\Lambda }^3e^f`$ and we have $`f=S`$ in the weak coupling limit of string theory. In fact this argument only implies that
$$\underset{S\mathrm{}}{lim}f(S)=S.$$
(12)
Nonperturbative effects could lead to the situation that $`f`$ is a very complicated function of $`S`$. In fact a satisfactory incorporation of gaugino condensates in the framework of string theory might very well lead to such a complication. In $`^\mathrm{?}`$ we suggested that a nontrivial $`f`$-function is the key ingredient to better understand the mechanism of gaugino condensation. We still assume (12) to make contact with perturbation theory. How do we then control $`e^f`$ as a function of $`S`$? In absence of a determination of $`f(S)`$ by a direct calculation one might use symmetry arguments to make some progress. Let us here consider the presence of a symmetry called $`S`$-duality which in its simplest form is given by a $`SL(2,Z)`$ generated by the transformations
$$SS+i,S1/S.$$
(13)
Such a symmetry might be realized in two basically distinct ways: the gauge sector could close under the transformation (type I) or being mapped to an additional ‘magnetic sector’ with inverted coupling constant (type II). In the second case one would speak of strong-weak coupling duality, just as in the case of electric-magnetic duality $`^\mathrm{?}`$. Within the class of theories of type I, however, we could have the situation that the $`f`$-function is itself invariant <sup>§</sup><sup>§</sup>§More complicated choices of transformation properties for $`f`$ are possible and lead to similar results as obtained in our simple toy model. under $`S`$-duality; i.e. $`S1/S`$ does not invert the coupling constant since the gauge coupling constant is not given by $`\text{Re}S`$ but $`1/g^2\text{Re}f`$. In view of (12) we would call such a symmetry weak-weak coupling duality. The behaviour of the gauge coupling constant as a function of $`S`$ is shown in Fig. 1. Our assumption (12) implies that $`g^20`$ as $`\text{Re}S\mathrm{}`$ and by $`S`$-duality $`g^2`$ also vanishes for $`S0`$, with a maximum somewhere in the vicinity of the self-dual point $`S=1`$. Observe that $`S1`$ in this situation does not necessarily imply strong coupling, because $`g^21/\text{Re}f`$ and even for $`S1`$, $`\text{Re}f`$ could be large and $`g^2<<1`$, with perturbation theory valid in the whole range of $`S`$. Of course, nonperturbative effects are responsible for the actual form of $`f(S)`$.
Fig. 1 - Coupling constant $`g^2`$ as the function of $`S`$ in type-I models (dashed) vs $`g^2`$ given by $`f=S`$
To examine the behaviour of the scalar potential in this approach, let us consider a simple toy model, with chiral superfields $`U=Y^3=(\lambda \lambda ,\mathrm{})`$ as well as $`S`$ and $`T`$. We have to choose a specific example of a gauge kinetic function which is invariant under the $`S`$-duality transformations. Different choices are possible, the simplest is given by
$$f=\frac{1}{2\pi }\mathrm{ln}(j(S)744),$$
(14)
$`j(S)`$ being the usual generator of modular invariant functions. This function behaves like $`S`$ in the large $`S`$-limit. If we assume a type I-model where the gauge sector is closed under $`S`$-duality, then we also have to assume that the gaugino condensate does not transform under $`S`$-duality (because of the $`fW^\alpha W_\alpha `$-term in the Lagrangian) For type I-models it was shown in $`^\mathrm{?}`$ that one can always redefine the gauge kinetic function and condensate in such a way that this holds.. Under these conditions an obvious candidate for the superpotential is just the standard Veneziano-Yankielowicz superpotential (extended to take into account the usual $`T`$-duality, which we assume to be completely independent from $`S`$-duality) $`^{\mathrm{?},\mathrm{?}}`$
$$W=Y^3(f+3b\mathrm{ln}\frac{Y\eta ^2(T)}{\mu }+c).$$
(15)
This is clearly invariant under $`S`$-duality. Therefore we then cannot take the conventional form for the Kähler potential which would be given by
$$K=\mathrm{ln}(S+\overline{S})3\mathrm{ln}(T+\overline{T}Y\overline{Y}),$$
(16)
since it is not $`S`$-dual. To make it $`S`$-dual one could introduce an additional $`\mathrm{ln}|\eta (S)|^4`$ term, giving e.g.
$$K=\mathrm{ln}(S+\overline{S})3\mathrm{ln}(T+\overline{T}Y\overline{Y})\mathrm{ln}|\eta (S)|^4.$$
(17)
Because the only relevant quantity is
$$G=K+\mathrm{ln}|W|^2,$$
(18)
we can as well put this new term (which is forced upon us because of our demand for symmetry) into the superpotential and take the canonical Kähler function instead, which gives
$$K=\mathrm{ln}(S+\overline{S})3\mathrm{ln}(T+\overline{T}Y\overline{Y}),$$
(19)
$$W=\frac{Y^3}{\eta ^2(S)}(f+3b\mathrm{ln}\frac{Y\eta ^2(T)}{\mu }+c),$$
(20)
where the remarkable similarity to the effective potential for $`T`$-dual gaugino condensation $`W=W_{inv}/\eta ^6(T)`$ can be seen more clearly.
This model exhibits a well defined minimum at $`\mathrm{<}S\mathrm{>}=1`$, $`\mathrm{<}T\mathrm{>}=1.23`$ and $`\mathrm{<}Y\mathrm{>}\mu `$. Supersymmetry is broken with the dominant contribution being $`\mathrm{<}F_T\mathrm{>}\mu ^3`$. The cosmological constant is negative.
In contrast to earlier attempts $`^\mathrm{?}`$ this model fixes the problem of the runaway dilaton and breaks supersymmetry with only a single gaugino condensate. Previous models needed multiple gaugino condensates and (to get realistic vevs for the dilaton) matter fields in complicated representations. We feel that the concept of a nontrivial gauge kinetic function derived (or constrained) by a symmetry is a much more natural way to fix the dilaton and break supersymmetry, especially so because corrections to $`f=S`$ are expected in any case. Earlier models which included $`S`$-duality in different ways (both with and without gaugino condensates) $`^{\mathrm{?},\mathrm{?}}`$ were able to fix the vev of the dilaton but did not succeed in breaking supersymmetry. An alternative mechanism to fix the vev of the dilaton has been discussed in $`^\mathrm{?}`$.
Of course there are still some open questions not solved by this approach. The first is the problem of having a vanishing cosmological constant. Whereas early models of gaugino condensation often introduced ad hoc terms to guarantee a vanishing vacuum energy, it has been seen to be notoriously difficult to get this out of models based on string inspired supergravity. The only way out of this problem so far has been to introduce a constant term into the superpotential, parameterizing unknown effects. This approach does not even work in any arbitrary model, but at least in our model the cosmological constant can be made to vanish by adjusting such a constant.
Another question not addressed in this toy model is the mixing of $`S`$ and $`T`$ fields which happens at the one-loop level. It is still unknown whether one can keep two independent dualities in this case. In a consistent interpretation our toy model should describe an all-loop effective action. If it is considered to be a theory at the tree-level then the theory is not anomaly free. Introducing terms to cancel the anomaly which arises because of demanding $`S`$-duality will then destroy $`S`$-duality. At tree-level the theory therefore cannot be made anomaly free.
An additional interesting question concerns the vevs of the auxiliary fields, i.e. which field is responsible for supersymmetry breakdown. In all models considered so far (multiple gaugino condensates, additional matter, S-duality) it has always been $`F_T`$ which dominates all the other auxiliary fields. It has not been shown yet that this is indeed a generic feature. The question is an important one, since the hierarchy of the vevs of the auxiliary fields is mirrored in the structure of the soft SUSY breaking terms of the MSSM $`^\mathrm{?}`$. We want to argue that there is at least no evidence for $`F_T`$ being generically large in comparison to $`F_S`$, because all of the models constructed so far (including our toy model) are designed in such a way that $`\mathrm{<}F_S\mathrm{>}=0`$ by construction at the minimum (at least at tree-level for the other models). In fact, if one extends our model with a constant in the superpotential (see above), then $`\mathrm{<}F_S\mathrm{>}`$ increases with the constant (but does not become as large as $`\mathrm{<}F_T\mathrm{>}`$).
Of course there are still some assumptions we made by considering this toy model. We assumed that there is weak coupling in the large $`S`$ limit. This is an assumption because the nonperturbative effects are unknown (at tree-level it can be calculated that $`f=S`$). In addition it is clear that the standard form we take for the Kähler potential does not include nonperturbative effects and thus could be valid only in the weak coupling approximation (this is of course related to our choice of the superpotential). Of course, an equally valid assumption would be that nonperturbative effects destroy the calculable tree-level behaviour even in the weak coupling region. The model of ref. $`^\mathrm{?}`$ could be re-interpreted in that sense (they do not consider gaugino condensates and the gauge kinetic function, but their $`S`$-dual scalar potential goes to infinity for $`S\mathrm{}`$). We choose not to make this assumption, because it is equivalent to the statement that the whole perturbative framework developed so far in string theory is wrong.
Again it should be emphasized here that the $`S`$-duality considered is not a strong-weak coupling duality but a weak-weak coupling duality. Even if the present model does not seem fully satisfactory, we are convinced that the idea of weak-weak coupling duality might be of more general relevance. In type II-models one has a duality between strong and weak coupling $`^\mathrm{?}`$. At the moment it is not completely clear how to incorporate that in a realistic model.
## 5 $`S=1`$ and weak coupling
A problem could be the actual size of the gauge coupling constant. If $`f=S`$ and $`\mathrm{<}S\mathrm{>}=1`$ then the large value of the gauge coupling constant does not fit the low scale of gaugino condensation necessary for phenomenologically realistic supersymmetry breaking ($`10^{13}GeV`$). However if $`f=S`$ only in the weak coupling limit then one can have $`\mathrm{<}f\mathrm{>}>>1`$ and thus $`g^2<<1`$ even in the region $`S=O(1)`$. Therefore in our model $`\mathrm{<}S\mathrm{>}=1`$ is consistent with the demand for a small gauge coupling constant, whereas in models with $`f=S`$ a much larger (and therefore more unnatural) $`\mathrm{<}S\mathrm{>}`$ is needed. Of course, it still has to be understood how such a large value of $`f(S=1)`$ can appear.
To summarize we find that the choice of a nontrivial $`f`$-function (motivated by a symmetry requirement) gives rise to a theory where supersymmetry breaking is achieved by employing only a single gaugino condensate. The cosmological constant turns out to be negative, but can be adjusted by a simple additional constant in the superpotential. The vevs of all fields are at natural orders of magnitude and due to the nontrivial gauge kinetic function the gauge coupling constant can be made small enough to give a realistic picture.
Turning our attention to the observable sector we see that a small (grand unified) coupling constant is a necessity and the above mechanism is required for a satisfactory description of the size af the observed coupling constants like e.g. $`\alpha _{\mathrm{QCD}}`$. But this alone might not be sufficient for a realistic model. String theory should predict all low energy coupling constants correctly and should also give the correct ratio of electroweak and strong coupling constants.
LEP and SLC high precision electroweak data give for the minimal supersymmetric Standard Model (MSSM) with the lightest Higgs mass in the range $`60\mathrm{G}\mathrm{e}\mathrm{V}<M_\mathrm{H}<150\mathrm{G}\mathrm{e}\mathrm{V}`$
$$\begin{array}{ccc}\hfill \mathrm{sin}^2\widehat{\theta }_\mathrm{W}(M_\mathrm{Z})& =& 0.2316\pm 0.0003\hfill \\ \hfill \alpha _{em}(M_\mathrm{Z})^1& =& 127.9\pm 0.1\hfill \\ \hfill \alpha _\mathrm{S}(M_\mathrm{Z})& =& 0.12\pm 0.01\hfill \\ \hfill m_t& =& 160_{125}^{+11+6}\mathrm{GeV},\hfill \end{array}$$
(21)
for the central value $`M_\mathrm{H}=M_\mathrm{Z}`$ in the $`\overline{MS}`$ scheme $`^\mathrm{?}`$. This is in perfect agreement with the recent CDF/D0 measurements of $`m_t`$. Taking the first three values as input parameters leads to gauge coupling unification at $`M_{\mathrm{GUT}}210^{16}\mathrm{GeV}`$ with $`\alpha _{\mathrm{GUT}}\frac{1}{26}`$ and $`M_{\mathrm{SUSY}}1\mathrm{T}\mathrm{e}\mathrm{V}`$ $`^{\mathrm{?},\mathrm{?}}`$. Slight modifications arise from light SUSY thresholds, i.e. the splitting of the sparticle mass spectrum, the variation of the mass of the second Higgs doublet and two–loop effects. Whereas these effects are rather mild, huge corrections may arise from heavy thresholds due to mass splittings at the high scale $`M_{heavy}M_{\mathrm{GUT}}`$ arising from the infinite many massive string states $`^\mathrm{?}`$. In the following sections we shall discuss this question of string unification in detail.
## 6 Gauge coupling unification
In heterotic superstring theories all couplings are related to the universal string coupling constant $`g_{\mathrm{string}}`$ at the string scale $`M_{\mathrm{string}}1/\sqrt{\alpha ^{}}`$, with $`\alpha ^{}`$ being the inverse string tension. It is a free parameter which is fixed by the dilaton vacuum expectation value $`g_{\mathrm{string}}^2=\frac{S+\overline{S}}{2}`$. In general this amounts to string unification, i.e. at the string scale $`M_{\mathrm{string}}`$ all gauge and Yukawa couplings are proportional to the string coupling and are therefore related to each other. For the gauge couplings (denoted by $`g_a`$) we have $`^\mathrm{?}`$:
$$g_a^2k_a=g_{\mathrm{string}}^2=\frac{\kappa ^2}{2\alpha ^{}}.$$
(22)
Here, $`k_a`$ is the Kac–Moody level of the group factor labeled by $`a`$. The string coupling $`g_{\mathrm{string}}`$ is related to the gravitational coupling constant $`\kappa ^2`$. In particular this means that string theory itself provides gauge coupling and Yukawa coupling unification even in absence of a grand unified gauge group.
To make contact with the observable world one should construct the field theoretical low–energy limit of a string vacuum. This is achieved by integrating out all the massive string modes corresponding to excited string states as well as states with momentum or winding quantum numbers in the internal dimensions. The resulting theory then describes the physics of the massless string excitations at low energies $`\mu <M_{\mathrm{string}}`$ in field–theoretical terms. If one wants to state anything about higher energy scales one has to take into account threshold corrections $`\mathrm{}_a(M_{\mathrm{string}})`$ to the bare couplings $`g_a(M_{\mathrm{string}})`$ due to the infinite tower of massive string modes. They change the relations (22) to:
$$g_a^2=k_ag_{\mathrm{string}}^2+\frac{1}{16\pi ^2}\mathrm{}_a,$$
(23)
The corrections in (23) may spoil the string tree–level result (22) and split the one–loop gauge couplings at $`M_{\mathrm{string}}`$. This splitting could allow for an effective unification at a scale $`M_{\mathrm{GUT}}<M_{\mathrm{string}}`$ or destroy the unification.
The general expression of $`\mathrm{}_a`$ for heterotic tachyon–free string vacua is given in $`^\mathrm{?}`$. Various contributions to $`\mathrm{}_a`$ have been determined for several classes of models: First in $`^\mathrm{?}`$ for two $`\mathrm{𝖹𝖹}_3`$ orbifold models with a (2,2) world–sheet supersymmetry $`^\mathrm{?}`$. This has been extended to fermionic constructions in $`^\mathrm{?}`$. Threshold corrections for (0,2) orbifold models with quantized Wilson lines $`^\mathrm{?}`$ have been calculated in $`^\mathrm{?}`$. Threshold corrections for the quintic threefold and other Calabi–Yau manifolds $`^\mathrm{?}`$ with gauge group $`E_6\times E_8`$ can be found in $`^{\mathrm{?},\mathrm{?}}`$. In toroidal orbifold compactifications $`^\mathrm{?}`$ moduli dependent threshold corrections arise only from N=2 supersymmetric sectors. They have been determined for some orbifold compactifications in $`^{\mathrm{?},\mathrm{?}}`$ and for more general orbifolds in $`^\mathrm{?}`$. The full moduli dependence A lowest expansion result in the Wilson line modulus has been obtained in $`^{\mathrm{?},\mathrm{?}}`$. of threshold corrections for (0,2) orbifold compactifications with continuous Wilson lines has been derived in $`^{\mathrm{?},\mathrm{?}}`$. These models contain continuous background gauge fields in addition to the usual moduli fields $`^\mathrm{?}`$. In most of the cases these models are (0,2) compactifications. In all the above orbifold examples the threshold corrections $`\mathrm{}_a`$ can be decomposed into three parts:
$$\mathrm{}_a=\stackrel{~}{\mathrm{}}_ab_a^{N=2}\mathrm{}+k_aY.$$
(24)
Here the gauge group dependent part is divided into two pieces: The moduli independent part $`\stackrel{~}{\mathrm{}}_a`$ containing the contribution of the N=1 supersymmetric sectors as well as scheme dependent parts which are proportional to $`b_a`$. This prefactor $`b_a`$ is related to the one–loop $`\beta `$–function: $`\beta _a=b_ag_a^3/16\pi ^2`$. Furthermore the moduli dependent part $`b_a^{N=2}\mathrm{}`$ with $`b_a^{N=2}`$ being related to the anomaly coefficient $`b_a^{}`$ by $`b_a^{N=2}=b_a^{}k_a\delta _{\mathrm{GS}}`$. The gauge group independent part $`Y`$ contains the gravitational back–reaction to the background gauge fields as well as other universal parts $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. They are absorbed into the definition of $`g_{\mathrm{string}}`$: $`g_{\mathrm{string}}^2=\frac{S+\overline{S}}{2}+\frac{1}{16\pi ^2}Y`$. The scheme dependent parts are the IR–regulators for both field– and string theory as well as the UV–regulator for field theory. The latter is put into the definition of $`M_{\mathrm{string}}`$ in the $`\overline{\mathrm{DR}}`$ scheme $`^\mathrm{?}`$:
$$M_{\mathrm{string}}=2\frac{e^{(1\gamma _\mathrm{E})/2}3^{3/4}}{\sqrt{2\pi \alpha ^{}}}=0.527g_{\mathrm{string}}\times 10^{18}\mathrm{GeV}.$$
(25)
The constant of the string IR–regulator as well as the universal part due to gravity were recently determined in $`^\mathrm{?}`$.
The identities (23) are the key to extract any string–implication for low–energy physics. They serve as boundary conditions for our running field–theoretical couplings valid below $`M_{\mathrm{string}}`$ $`^\mathrm{?}`$. Therefore they are the foundation of any discussion about both low–energy predictions and gauge coupling unification. The evolution equations <sup>\**</sup><sup>\**</sup>\**We neglect the N=1 part of $`\stackrel{~}{\mathrm{}}_a`$ which is small compared to $`b_a^{N=2}\mathrm{}`$ $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. valid below $`M_{\mathrm{string}}`$
$$\frac{1}{g_a^2(\mu )}=\frac{k_a}{g_{\mathrm{string}}^2}+\frac{b_a}{16\pi ^2}\mathrm{ln}\frac{M_{\mathrm{string}}^2}{\mu ^2}\frac{1}{16\pi ^2}b_a^{N=2}\mathrm{},$$
(26)
allow us to determine $`\mathrm{sin}^2\theta _\mathrm{W}`$ and $`\alpha _\mathrm{S}`$ at $`M_Z`$. After eliminating $`g_{\mathrm{string}}`$ in the second and third equations one obtains
$`\mathrm{sin}^2\theta _\mathrm{W}(M_Z)`$ $`=`$ $`{\displaystyle \frac{k_2}{k_1+k_2}}{\displaystyle \frac{k_1}{k_1+k_2}}{\displaystyle \frac{\alpha _{em}(M_Z)}{4\pi }}\left[𝒜\mathrm{ln}\left({\displaystyle \frac{M_{\mathrm{string}}^2}{M_Z^2}}\right)𝒜^{}\mathrm{}\right],`$
$`\alpha _S^1(M_Z)`$ $`=`$ $`{\displaystyle \frac{k_3}{k_1+k_2}}\left[\alpha _{em}^1(M_Z){\displaystyle \frac{1}{4\pi }}\mathrm{ln}\left({\displaystyle \frac{M_{\mathrm{string}}^2}{M_Z^2}}\right)+{\displaystyle \frac{1}{4\pi }}^{}\mathrm{}\right],`$ (27)
with $`𝒜=\frac{k_2}{k_1}b_1b_2,=b_1+b_2\frac{k_1+k_2}{k_3}b_3`$ and $`𝒜^{},^{}`$ are obtained by exchanging $`b_ib_i^{}`$. For the MSSM one has $`𝒜=\frac{28}{5},=20`$. However to arrive at the predictions of the MSSM (21) one needs huge string threshold corrections $`\mathrm{}`$ due to the large value of $`M_{\mathrm{string}}`$ $`(3/5k_1=k_2=k_3=1)`$:
$$\mathrm{}=\frac{𝒜}{𝒜^{}}\left[\mathrm{ln}\left(\frac{M_{\mathrm{string}}^2}{M_{\mathrm{GUT}}^2}\right)+\frac{32\pi \delta _{\mathrm{sin}^2\theta _\mathrm{W}}}{5𝒜\alpha _{em}(M_Z)}\right].$$
(28)
At the same time, the N=2 spectrum of the underlying theory encoded in $`𝒜^{},^{}`$ which enters the threshold corrections has to fulfill the condition
$$\frac{^{}}{𝒜^{}}=\frac{}{𝒜}\frac{\mathrm{ln}\left(\frac{M_{\mathrm{string}}^2}{M_{\mathrm{GUT}}^2}\right)+\frac{32\pi }{3}\delta _{\alpha _\mathrm{S}^1}}{\mathrm{ln}\left(\frac{M_{\mathrm{string}}^2}{M_{\mathrm{GUT}}^2}\right)+\frac{32\pi }{5𝒜}\frac{\delta _{\mathrm{sin}\theta _\mathrm{W}^2}}{\alpha _{em}(M_Z)}},$$
(29)
where $`\delta `$ represents the experimental uncertainties appearing in (21). In addition $`\delta `$ may also contain SUSY thresholds.
For concreteness and as an illustration let us take the $`\mathrm{𝖹𝖹}_8`$ orbifold example of $`^\mathrm{?}`$ with $`𝒜^{}=2,^{}=6`$ and $`b_1^{}+b_2^{}=10`$. It is one of the few orbifolds left over after imposing the conditions on target–space duality anomaly cancellation $`^\mathrm{?}`$. To estimate the size of $`\mathrm{}`$ one may take in eq. (25) $`g_{\mathrm{string}}0.7`$ corresponding to $`\alpha _{\mathrm{string}}\frac{1}{26}`$, i.e. $`M_{\mathrm{string}}/M_{\mathrm{GUT}}20`$. Of course this is a rough estimate since $`M_{\mathrm{string}}`$ is determined by the first eq. of (26) together with (25). Nevertheless, the qualitative picture does not change. Therefore to predict the correct low–energy parameter (6) eq. (28) tells us that one needs threshold correction of considerable size:
$$17.1\mathrm{}16.3.$$
(30)
## 7 String thresholds
The construction of a realistic unified string model boils down to the question of how to achieve thresholds of that size. To settle the question we need explicit calculations within the given candidate string model. There we can encounter various types of threshold effects. Some depend continuously, others discretely on the values of the moduli fields. For historic reasons we also have to distinguish between thresholds that do or do not depend on Wilson lines. The reason is the fact that the calculations in the latter models are considerably simpler and for some time were the only available results. They were then used to estimate the thresholds in models with gauge group $`SU(3)\times SU(2)\times U(1)`$ and three families, although as a string model no such orbifold can be constructed without Wilson lines. Therefore, the really relevant thresholds are, of course, the ones found in the (0,2) orbifold models with Wilson lines $`^\mathrm{?}`$ which may both break the gauge group and reduce its rank. We will discuss the various contributions within the framework of our illustrative model. However the discussion can easily applied for all other orbifolds. The threshold corrections depend on the $`T`$ and $`U`$ modulus describing the size and shape of the internal torus lattice. In addition they may depend on non–trivial gauge background fields encoded in the Wilson line modulus $`B`$.
Moduli dependent threshold corrections $`\mathrm{}`$ can be of significant size for an appropriate choice of the vevs of the background fields $`T,U,B,\mathrm{}`$ which enter these functions. Of course in the decompactification limit $`Ti\mathrm{}`$ these corrections become always arbitrarily huge. This is in contrast to fermionic string compactifications or N=1 sectors of heterotic superstring compactifications. There one can argue that moduli–independent threshold corrections cannot become huge at all $`^\mathrm{?}`$. This is in precise agreement with the results found earlier in $`^{\mathrm{?},\mathrm{?}}`$. In field theory threshold corrections can be estimated with the formula $`^\mathrm{?}`$
$$\mathrm{}=\underset{n,m,k}{}\mathrm{ln}\left(\frac{M_{n,m,k}^2}{M_{\mathrm{string}}^2}\right),$$
(31)
with $`n,m`$ being the winding and momentum, respectively and $`k`$ the gauge quantum number of all particles running in the loop. The string mass in the $`N=2`$ sector of the $`\mathrm{𝖹𝖹}_8`$ model we consider later with a non–trivial gauge background in the internal directions is determined by $`^\mathrm{?}`$ :
$`\alpha ^{}M_{n,m,k}^2`$ $`=`$ $`4|p_R|^2`$
$`p_R`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{Y}}}\left[({\displaystyle \frac{T}{2\alpha ^{}}}UB^2)n_2+{\displaystyle \frac{T}{2\alpha ^{}}}n_1Um_1+m_2+Bk_2\right]`$
$`Y`$ $`=`$ $`{\displaystyle \frac{1}{2\alpha ^{}}}(T\overline{T})(U\overline{U})+(B\overline{B})^2.`$ (32)
In addition a physical state $`|n,m,k,l`$ has to obey the modular invariance condition $`m_1n_1+m_2n_2+k_1^2k_1k_2+k_2^2k_2k_3k_2k_4+k_3^2+k_4^2=1N_L\frac{1}{2}l_{E_8^{}}^2`$. Therefore the sum in (31) should be restricted to these states. This also guarantees its convergence after a proper regularization. In (31) cancellations between the contributions of various string states may arise. E.g. at the critical point $`T=i=U`$ where all masses appear in integers of $`M_{\mathrm{string}}`$ such cancellations occur. They are the reason for the smallness of the corrections calculated in $`^{\mathrm{?},\mathrm{?}}`$ and in all the fermionic models $`^\mathrm{?}`$. Let us investigate this in more detail. The simplest case ($`B=0`$) for moduli dependent threshold corrections to the gauge couplings was derived in $`^\mathrm{?}`$ :
$$\mathrm{}(T,U)=\mathrm{ln}\left[\frac{iT+i\overline{T}}{2\alpha ^{}}\left|\eta \left(\frac{T}{2\alpha ^{}}\right)\right|^4\right]+\mathrm{ln}\left[(iU+i\overline{U})\left|\eta (U)\right|^4\right].$$
(33)
Formula (33) can be used for any toroidal orbifold compactifications, where the two–dimensional subplane of the internal lattice which is responsible for the N=2 structure factorize from the remaining part of the lattice. If the latter condition does not hold, (33) is generalized $`^\mathrm{?}`$.
$$\begin{array}{cccccc}& & & & & \\ & & & & & \\ & T/2\alpha ^{}& U& M^2\alpha ^{}& ln(M^2\alpha ^{})& \mathrm{\Delta }^{II}\\ & & & & & \\ & & & & & \\ & & & & & \\ Ia& i& i& 1& 0& 0.72\\ & & & & & \\ & & & & & \\ & & & & & \\ Ib& 1.25i& i& \frac{4}{5}& 0.22& 0.76\\ & & & & & \\ & & & & & \\ & & & & & \\ Ic& 4.5i& 4.5i& \frac{4}{81}& 3.01& 5.03\\ & & & & & \\ & & & & & \\ & & & & & \\ Id& 18.7i& i& \frac{10}{187}& 2.93& 16.3\\ & & & & & \end{array}$$
Table 1: Lowest mass $`M^2`$ of particles charged
under $`G_A`$ and threshold corrections $`\mathrm{}(T,U)`$.
In Table 1 we determine the mass of the lowest massive string state being charged under the considered unbroken gauge group $`G_A`$ and the threshold corrections $`\mathrm{}(T,U)`$ for some values of $`T`$ and $`U`$.
The influence of moduli dependent threshold corrections to low–energy physics \[entailed in eqs. (6)\] has until now only been discussed for orbifold compactifications without Wilson lines by using (33). In these cases the corrections only depend on the two moduli $`T,U`$. However to obtain corrections of the size $`\mathrm{}16.3`$ one would need the vevs $`\frac{T}{2\alpha ^{}}=18.7,U=i`$ which are far away from the self–dual points $`^{\mathrm{?},\mathrm{?}}`$. It remains an open question whether and how such big vevs of $`T`$ can be obtained in a natural way in string theory.
A generalization of eq. (33) appears when turning on non–vanishing gauge background fields $`B0`$. According to (7) the mass of the heavy string states now becomes $`B`$–dependent and therefore also the threshold corrections change. This kind of corrections were recently determined in $`^\mathrm{?}`$. The general expression there is
$$\mathrm{}^{II}(T,U,B)=\frac{1}{12}\mathrm{ln}\left[\frac{Y^{12}}{1728^4}\left|𝒞_{12}(\mathrm{\Omega })\right|^2\right],$$
(34)
where $`B`$ is the Wilson line modulus, $`\mathrm{\Omega }=\left(\begin{array}{cc}\frac{T}{2\alpha ^{}}& B\\ B& U\end{array}\right)`$ and $`𝒞_{12}`$ is a combination of $`g=2`$ elliptic theta functions explained in detail in $`^\mathrm{?}`$. It applies to gauge groups $`G_A`$ which are not affected by the Wilson line mechanism. The case where the gauge group is broken by the Wilson line will be discussed later (those threshold corrections will be singular in the limit of vanishing $`B`$). Whereas the effect of quantized Wilson lines $`B`$ on threshold corrections has already been discussed in $`^\mathrm{?}`$ the function $`\mathrm{\Delta }^{II}(T,U,B)`$ now allows us to study the effect of a continuous variation in $`B`$.
Fig.2 – Dependence of the threshold corrections $`\mathrm{}^{II}`$
on the Wilson line modulus $`B=B_1+iB_2`$ for $`\frac{T}{2\alpha ^{}}=i=U`$.
We see in Fig.2 that the threshold corrections change very little with the Wilson line modulus $`B`$. They are comparable with $`\mathrm{}=0.72`$ corresponding to the case of $`B=0`$. In this case eq. (34) becomes eq. (33) for $`\frac{T}{2\alpha ^{}}=i=U`$.
So far all these calculations have been done within models where the considered gauge group $`G_A`$ is not broken by the Wilson line and its matter representations are not projected out. To arrive at SM like gauge groups with the matter content of the MSSM one has to break the considered gauge group with a Wilson line.
From the phenomenological point of view $`^\mathrm{?}`$, the most promising class of string vacua is provided by (0,2) compactifications equipped with a non–trivial gauge background in the internal space which breaks the $`E_6`$ gauge group down to a SM–like gauge group $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. Since the internal space is not simply connected, these gauge fields cannot be gauged away and may break the gauge group. Some of the problems present in (2,2) compactifications with $`E_6`$ as a grand unified group like e.g. the doublet–triplet splitting problem, the fine–tuning problem and Yukawa coupling unification may be absent in (0,2) compactifications. It is important that these properties can be studied in the full string theory, not just in the field theoretic limit $`^\mathrm{?}`$. The background gauge fields give rise to a new class of massless moduli fields again denoted by $`B`$ which have quite different low–energy implications than the usual moduli arising from the geometry of the internal manifold itself. In this framework the question of string unification can now be discussed for realistic string models. The threshold corrections for our illustrative model take the form $`^\mathrm{?}`$
$$\mathrm{}^I(T,U,B)=\frac{1}{10}\mathrm{ln}\left[Y^{10}\left|\frac{1}{128}\underset{k=1}{\overset{10}{}}\vartheta _k(\mathrm{\Omega })\right|^4\right],$$
(35)
where $`\vartheta _k`$ are the ten even $`g=2`$ theta–functions $`^\mathrm{?}`$. Equipped with this result we can now investigate the influence of the B–modulus on the thresholds and see how the conclusions of ref. $`^{\mathrm{?},\mathrm{?}}`$ might be modified. The results for a representative set of background vevs is displayed in Fig.3.
Fig.3 – Dependence of the threshold corrections $`\mathrm{}^I`$
on the Wilson line modulus $`B=B_1+iB_2`$ for $`\frac{T}{2\alpha ^{}}=4.5i=U`$.
From this picture we see that threshold corrections of $`\mathrm{}16.3`$ can be obtained for the choice of $`\frac{T}{2\alpha ^{}}4.5iU`$ and $`B=\frac{1}{2}`$. This has to be compared to the model in ref. $`^\mathrm{?}`$ where such a value was achieved with $`T=18.7i`$ and $`B=0`$. This turns out to be a general property of the models under consideration. With more moduli, sizeable threshold effects are achieved even with moderate values of the vevs of the background fields.
## 8 Heterotic string unification
Equipped with these explicit calculations of string threshold corrections we can now ask the question how string theory might lead to the correct prediction of gauge coupling constants. We also hope to deduce information on the spectrum of theories that lead to successful gauge coupling unification.
The modulus plays the rôle of an adjoint Higgs field which breaks e.g. the $`G_A=E_6`$ down to a SM like gauge group $`G_a`$. According to eq. (7) the vev of this field gives some particles masses between zero and $`M_{\mathrm{string}}`$. This is known as the stringy Higgs effect. Such additional intermediate fields may be very important to generate high scale thresholds. Sizeable threshold corrections $`\mathrm{}`$ can only appear if some particles have masses different from the string scale $`M_{\mathrm{string}}`$ and where cancellations between different states as mentioned above do not take place. In particular some gauge bosons of $`G_A`$ become massive receiving the mass:
$$\alpha ^{}M_I^2=\frac{4}{Y}|B|^2.$$
(36)
As before let us investigate the masses of the lightest massive particles charged under the gauge group $`G_a`$. For our concrete model we have $`M_{\mathrm{string}}=3.610^{17}\mathrm{GeV}`$.
$$\begin{array}{ccccccc}& & & & & & \\ & & & & & & \\ & T/2\alpha ^{}& U& B& M_I[\mathrm{GeV}]& ln(M_I^2\alpha ^{})& \mathrm{\Delta }^I\\ & & & & & & \\ & & & & & & \\ & & & & & & \\ IIa& i& i& \frac{1}{10^5}& 8.410^{12}& 23.0& 10.03\\ & & & & & & \\ & & & & & & \\ & & & & & & \\ IIb& i& i& \frac{1}{2}& 4.210^{17}& 1.39& 1.72\\ & & & & & & \\ & & & & & & \\ & & & & & & \\ IIc& 1.25i& i& \frac{1}{2}& 3.710^{17}& 1.61& 2.12\\ & & & & & & \\ & & & & & & \\ & & & & & & \\ IId& 4i& i& \frac{1}{2}& 2.110^{17}& 2.78& 7.86\\ & & & & & & \\ & & & & & & \\ & & & & & & \\ IIe& 4.5i& 4.5i& \frac{1}{2}& 9.310^{16}& 4.39& 16.3\\ & & & & & & \\ & & & & & & \\ & & & & & & \\ IIf& 18.7i& i& \frac{1}{2}& 1.110^{16}& 4.31& 43.3\\ & & & & & & \end{array}$$
Table 2: Lowest mass $`M_I`$ of particles charged
under $`G_a`$ and threshold corrections $`\mathrm{}^I`$ for $`B0`$.
Whereas $`\mathrm{\Delta }^{II}`$ describes threshold corrections w.r.t. to a gauge group which is not broken when turning on a vev of $`B`$, now the gauge group is broken for $`B0`$ and in particular this means that the threshold $`\mathrm{}^I`$ shows a logarithmic singularity for $`B0`$ when the full gauge symmetry is restored. This behaviour is known from field theory and the effects of the heavy string states can be decoupled from the former: Then the part of $`\mathrm{}_a^I`$ in (23) which is only due to the massive particles becomes $`^{\mathrm{?},\mathrm{?}}`$
$$\frac{b_Ab_a}{16\pi ^2}\mathrm{ln}\frac{M_{\mathrm{string}}^2}{|B|^2}\frac{b_A^{}}{16\pi ^2}\mathrm{ln}\left|\eta \left(\frac{T}{2\alpha ^{}}\right)\eta (U)\right|^4,$$
(37)
where the first part accounts for the new particles appearing at the intermediate scale of $`M_I`$ and the other part takes into account the contributions of the heavy string states. One of the questions of string unification concerns the size of this intermediate scale $`M_I`$. In a standard grand unified model one would be tempted to identify $`M_I`$ with $`M_{\mathrm{GUT}}`$. While this would also be a possibility for string unification, we have in string theory in addition the possibility to consider $`M_I>M_{\mathrm{GUT}}`$. The question remains whether the thresholds in that case can be big enough, as we shall discuss in a moment. Let us first discuss the general consequences of our results for the idea of string unification without a grand unified gauge group. Due to the specific form of the threshold corrections in eq. (26) unification always takes place if the condition $`𝒜^{}=𝒜^{}`$ is met within the errors arising from the uncertainties in (21). It guarantees that all three gauge couplings meet at a single point $`M_\mathrm{X}`$ $`^\mathrm{?}`$:
$$M_\mathrm{X}=M_{\mathrm{string}}e^{\frac{1}{2}\frac{𝒜^{}}{𝒜}\mathrm{}}.$$
(38)
For our concrete model this leads to $`M_\mathrm{X}210^{16}\mathrm{GeV}`$. Given these results we can now study the relation between $`M_I`$ and $`M_\mathrm{X}`$, which plays the rôle of the GUT–scale in string unified models. As a concrete example, consider the model $`IIe`$ in Table 2. It leads to an intermediate scale $`M_I`$ which is a factor 3.9 smaller than the string scale, thus $`10^{17}\mathrm{GeV}`$, although the apparent unification scale is as low as $`2\times 10^{16}\mathrm{GeV}`$. We thus have an explicit example of a string model where all the non–MSSM particles are above $`9.310^{16}\mathrm{GeV}`$, but still a correct prediction of the low energy parameters emerges. Thus string unification can be achieved without the introduction of a small intermediate scale.
Of course, there are also other possibilities which lead to the correct low–energy predictions. Instead of large threshold corrections one could consider a non–standard hypercharge normalization, i.e. a $`k_15/3`$ $`^\mathrm{?}`$. This would maintain gauge coupling unification at the string scale with the correct values of $`\mathrm{sin}^2\theta _\mathrm{W}(M_Z)`$ and $`\alpha _\mathrm{S}(M_Z)`$. However, it is very hard to construct such models. A further possibility would be to give up the idea of gauge coupling unification within the MSSM by introducing extra massless particles such as $`(\mathrm{𝟑},\mathrm{𝟐})`$ w.r.t. $`SU(3)\times SU(2)`$ in addition to those of the SM $`^{\mathrm{?},\mathrm{?}}`$. A careful choice of these matter fields may lead to sizable additional intermediate threshold corrections in (6) thus allowing for the correct low–energy data (21). Unfortunately the price for that is exactly an introduction of a new intermediate scale of $`M_I10^{1214}\mathrm{GeV}`$. It seems to be hard to explain such a small scale naturally in the framework of string theory. In some sense such a model can be compared to the model $`IIa`$ in table 2. Other possible corrections to (6) may arise from an extended gauge structure between $`M_\mathrm{X}`$ and $`M_{\mathrm{string}}`$. However this might even enhance the disagreement with the experiment $`^\mathrm{?}`$. Finally a modification of (6) appears from the scheme conversion from the string– or SUSY–based $`\overline{DR}`$ scheme to the $`\overline{MS}`$ scheme relevant for the low–energy physics data (21). However these effects are shown to be small $`^\mathrm{?}`$.
Let us stress here the important message that string unification can, in principle, be achieved with moduli dependent threshold corrections within (0,2) superstring compactification. The Wilson line dependence of these functions is comparable to that on the $`T`$ and $`U`$ fields thus offering the interesting possibility of large thresholds with background configurations of moderate size. All non–MSSM like states can e.g. be heavier than $`1/4`$ of the string scale, still leading to an apparent unification scale of $`M_\mathrm{X}=\frac{1}{20}M_{\mathrm{string}}`$. We do not need vevs of the moduli fields that are of the order 20 away from the natural scale, neither do we need to introduce particles at a new intermediate scale that is small compared to $`M_{\mathrm{string}}`$.
The situation could be even more improved with a higher number of moduli fields entering the threshold corrections: They may come from other orbifold planes giving rise to N=2 sectors or from additional Wilson lines. We think that the actual moderate vevs of the underlying moduli fields can be fixed by non–perturbative effects as e.g. gaugino condensation. Of course, unification can be achieved in different ways, as the introduction of an intermediate scale. This does not seem to be very natural, because it postulates a new scale that is a factor $`10^4`$ smaller than the GUT scale in order to explain a factor 20 discrepancy in the difference of the unification scale and the string scale. One might also argue that in the framework of string theory one should consider models with a grand unified group unbroken at the string scale but broken at the GUT scale. This might lead to interesting models and consequences, but it does not contribute to an explanation of the difference of the string scale and the GUT scale.
An alternative view of unification might arise according to recent developments in nonperturbative string theory. We shall discuss that in the following sections
## 9 Recent developments: M-Theory
From all the new and interesting results in string dualities, it is the heterotic M–theory of Hořava and Witten $`^\mathrm{?}`$ that seems to have immediate impact on the discussion of the phenomenological aspects of these theories. One of the results concerns the question of the unification of all fundamental coupling constants $`^\mathrm{?}`$ and the second one the properties of the soft terms (especially the gaugino masses) once supersymmetry is broken $`^{\mathrm{?},\mathrm{?}}`$. In both cases results that appeared problematic in the weakly coupled case get modified in a satisfactory way, while the overall qualitative picture remains essentially unchanged. In these lectures we shall therefore concentrate on these aspects of the new picture.
The heterotic M–theory is an 11–dimensional theory with the $`E_8\times E_8`$ gauge fields living on two 10–dimensional boundaries (walls), respectively, while the gravitational fields can propagate in the bulk as well. A $`d=4`$ dimensional theory with $`N=1`$ supersymmetry emerges at low energies when 6 dimensions are compactified on a Calabi–Yau manifold. The scales of that theory are $`M_{11}`$, the $`d=11`$ Planck scale, $`R_{11}`$ the size of the $`x^{11}`$ interval, and $`VR^6`$ the volume of the Calabi–Yau manifold. The quantities of interest in $`d=4`$, the Planck mass, the GUT–scale and the unified gauge coupling constant $`\alpha _{GUT}`$ should be determined through these higher dimensional quantities. The fit of ref. $`^\mathrm{?}`$ identifies $`M_{GUT}310^{16}`$ GeV with the inverse Calabi–Yau radius $`R^1`$. Adjusting $`\alpha _{GUT}=1/25`$ gives $`M_{11}`$ to be a few times larger than $`M_{GUT}`$. On the other hand, the fit of the actual value of the Planck scale can be achieved by the choice of $`R_{11}`$ and, interestingly enough, $`R_{11}`$ turns out to be an order of magnitude larger than the fundamental length scale $`M_{11}^1`$. A satisfactory fit of the $`d=4`$ scales is thus possible, in contrast to the case of the weakly coupled heterotic string, where naively the string scale seemed to be a factor 20 larger than $`M_{GUT}`$.
As we have seen before, otherwise the heterotic $`E_8\times E_8`$ string looks rather attractive from the point of view of phenomenological applications. One seems to be able to accommodate the correct gauge group and particle spectrum. The mechanism of hidden sector gaugino condensation leads to a breakdown of supersymmetry with vanishing cosmological constant to leading order. With a condensate scale $`\mathrm{\Lambda }10^{13}`$ GeV, one obtains a gravitino mass in the TeV range and soft scalar masses in that range as well. In the simplest models $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ this type of supersymmetry breakdown is characterized through the vacuum expectation value of moduli fields other than the dilaton, giving a small problem with the soft gaugino masses in the observable sector: they turn out to be too small, generically some two orders of magnitude smaller than the soft scalar masses. It is again in the framework of heterotic M–theory that this problem is solved $`^\mathrm{?}`$; gaugino masses are of the same size as (or even larger than) the soft scalar masses.
The mechanism of hidden sector gaugino condensation itself can be realized in a way very similar to the weakly coupled case. This includes the mechanism of cancellation of the vacuum energy, which in the weakly coupled case arises because of a cancellation of the gaugino condensate with a vacuum expectation value of the three index tensor field $`H`$ of $`d=10`$ supergravity. This cancellation is at the origin of the fact that supersymmetry breakdown is dominated by a $`T`$ modulus field rather than the dilaton ($`S`$). Hořava $`^\mathrm{?}`$ observed that this compensation of the vacuum expectation values of the condensate and $`H`$ carries over to the M–theory case. In $`^\mathrm{?}`$ this has been explicitly worked out for the mechanism of gaugino condensation in the heterotic M–theory and the similarity to the weakly coupled case was shown. Now the gaugino condensate forms at the hidden 4–dimensional wall and is cancelled locally at that wall by the vacuum expectation value (vev) of a Chern–Simons term. This also clarifies some questions concerning the nature of the vev of $`H`$ that arose in the weakly coupled case.
In the remainder of these lectures we want to discuss the phenomenological properties of the heterotic M–theory. This includes a presentation of the full effective four–dimensional $`N=1`$ supergravity action in leading and next–to–leading order, the mechanism of hidden sector gaugino condensation and its explicit consequences for supersymmetry breaking and the scalar potential and finally the resulting soft breaking terms in the 4–dimensional theory. Although some of the issues have already been discussed earlier, we shall at each step first explain the situation again for the weakly coupled theory and then compare it to the results obtained in the M–theory case.
These results are obtained using the method of reduction and truncation that has been successfully applied to the weakly coupled case $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. It is a simplified prescription that shows the main qualitative features of the effective $`d=4`$ effective theory. In orbifold compactification it would represent the fields and interactions in the untwisted sector. We compute Kähler potential ($`K`$), superpotential ($`W`$) and gauge kinetic function ($`f`$) both in the weakly and strongly coupled regime and explain similarities and differences.
The results in leading order had been obtained previously $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. These papers mainly focused on the breakdown of supersymmetry via a Scherk–Schwarz mechanism which we shall nor discuss here in detail. It remains to be seen, if and how such a mechanism can be related to the mechanism of gaugino condensation.
The remainder of the lectures will proceed as follows. First we discuss the scales and the question of unification as suggested in $`^\mathrm{?}`$ and compare the two cases. then we derive the effective $`d=4`$ action of M–theory using the method of reduction and truncation. In this case we have to deal with a nontrivial obstruction first encountered in $`^\mathrm{?}`$. It leads to an explicit $`x^{11}`$ dependence of certain fields, which is induced by vevs of antisymmetric tensor fields at the walls. To obtain the effective action in $`d=4`$ we have to integrate out this dependence. This then leads to corrections to $`K`$ and $`f`$ in next to leading order, which are very similar compared to those in the weakly coupled case. We also discuss the appearance and the size of a critical radius for $`R_{11}`$. The phenomenological fit presented in our discussion of unification implies that we are not too far from that critical radius. We then turn again to the question of supersymmetry breakdown. We start with the weakly coupled case and investigate the nature of the vev of the $`H`$–field (concerning some quantization conditions) and the cancellation of the vacuum energy. In the strongly coupled case we shall see that such a cancellation appears locally at one wall. This supports the interpretation that the gaugino condensate is matched by a nontrivial vev of a Chern–Simons term. We then explicitly identify the mechanism of supersymmetry breakdown and the nature of the gravitino. The goldstino turns out to be the fermionic component of the $`T`$ superfield that represents essentially the radius of the 11th dimension. It is a bulk field, with a vev of its auxiliary component on one wall. Integrating out the 11th dimension we then obtain explicitly the mass of the gravitino.
The remainder deals with the induced soft breaking terms in the observable sector: scalar and gaugino masses. We shall see a strong model dependence of the scalar masses and argue that they are not too different from the gravitino mass. This all is very similar to the situation in the weakly coupled case. We then compute the soft gaugino masses and see that in the strongly coupled case they are of the order of the gravitino mass. This comes from the fact that we are quite close to the critical radius and represents a decisive difference to the weakly coupled regime.
## 10 Scales and unification
As we have seen, models of particle physics that are derived as the low energy limit of the $`E_8\times E_8`$ heterotic string are particularly attractive. They seem to be able to accommodate the correct gauge group and particle spectrum to lead to the supersymmetric extension to the $`SU(3)\times SU(2)\times U(1)`$ standard model. It is exactly in this framework that a unification of the gauge coupling constants is expected to appear at a scale $`M_{GUT}=310^{16}`$ GeV. As we know, this heterotic string theory (weakly coupled at the string scale) gives a prediction for the relation between gauge and gravitational coupling constants. To see this explicitly let us have a look at the low energy effective action of the $`d=10`$–dimensional field theory:
$$L=\frac{4}{(\alpha ^{})^3}d^{10}x\sqrt{g}\mathrm{exp}(2\varphi )\left(\frac{1}{(\alpha ^{})}R+\frac{1}{4}\mathrm{tr}F^2+\mathrm{}\right),$$
(39)
where $`\alpha ^{}`$ is the string tension and $`\varphi `$ the dilaton field in $`d=10`$. A definite relation between gauge and gravitational coupling appears because of the universal behaviour of the dilaton term in eq. (39). The effective $`d=4`$–dimensional theory is obtained after compactification on a Calabi–Yau manifold with volume $`V`$:
$$L=\frac{4}{(\alpha ^{})^3}d^4x\sqrt{g}\mathrm{exp}(2\varphi )V\left(\frac{1}{(\alpha ^{})}R+\frac{1}{4}\mathrm{tr}F^2+\mathrm{}\right).$$
(40)
Thus a universal factor $`V\mathrm{exp}(2\varphi )`$ multiplies both the $`R`$ and $`F^2`$ terms. Newton’s and Einstein’s gravitational coupling constants are related as
$$G_N=\frac{1}{8\pi }\kappa _4^2=\frac{1}{M_{\mathrm{Planck}}^2},$$
(41)
with $`M_{\mathrm{Planck}}1.210^{19}`$ GeV. From eq. (40) we then deduce:
$$G_N=\frac{\mathrm{exp}(2\varphi )(\alpha ^{})^4}{64\pi V},$$
(42)
as well as
$$\alpha _{GUT}=\frac{\mathrm{exp}(2\varphi )(\alpha ^{})^3}{16\pi V},$$
(43)
leading to the relation
$$G_N=\frac{\alpha _{GUT}\alpha ^{}}{4}.$$
(44)
Putting in the value for $`M_{\mathrm{Planck}}`$ and $`\alpha _{GUT}1/25`$ one obtains a value for the string scale $`M_{\mathrm{string}}=(\alpha ^{})^{1/2}`$ that is in the region of $`10^{18}`$ GeV. This is apparently much larger than the GUT–scale of $`310^{16}`$ GeV, while naively one would like to identify $`M_{\mathrm{string}}`$ with $`M_{GUT}`$. The discrepancy of the scales is sometimes called the unification problem in the framework of the weakly coupled heterotic string. We have discussed it in the previous sections. There we have seen that the above argumentation is rather simple minded and that more sophisticated (threshold) calculations are needed to settle this issue. In any case, the natural appearance of $`M_{\mathrm{string}}M_{GUT}`$ would have been desirable. Let us now see how the situation looks in the case of heterotic string theory at stronger coupling.
We now consider the $`E_8\times E_8`$ M–theory. The effective action of the strongly coupled $`E_8\times E_8`$$`M`$–theory in the “downstairs” approach is given by $`^\mathrm{?}`$ (we take into account the numerical corrections found in $`^\mathrm{?}`$)
$`L`$ $`=`$ $`{\displaystyle \frac{1}{\kappa ^2}}{\displaystyle _{M^{11}}}d^{11}x\sqrt{g}[{\displaystyle \frac{1}{2}}R{\displaystyle \frac{1}{2}}\overline{\psi }_I\mathrm{\Gamma }^{IJK}D_J\left({\displaystyle \frac{\mathrm{\Omega }+\widehat{\mathrm{\Omega }}}{2}}\right)\psi _K{\displaystyle \frac{1}{48}}G_{IJKL}G^{IJKL}`$
$`{\displaystyle \frac{\sqrt{2}}{384}}\left(\overline{\psi }_I\mathrm{\Gamma }^{IJKLMN}\psi _N+12\overline{\psi }^J\mathrm{\Gamma }^{KL}\psi ^M\right)\left(G_{JKLM}+\widehat{G}_{JKLM}\right)`$
$`{\displaystyle \frac{\sqrt{2}}{3456}}ϵ^{I_1I_2\mathrm{}I_{11}}C_{I_1I_2I_3}G_{I_4\mathrm{}I_7}G_{I_8\mathrm{}I_{11}}]`$
$`+`$ $`{\displaystyle \frac{1}{4\pi (4\pi \kappa ^2)^{2/3}}}{\displaystyle _{M_i^{10}}}d^{10}x\sqrt{g}[{\displaystyle \frac{1}{4}}F_{iAB}^aF_i^{aAB}{\displaystyle \frac{1}{2}}\overline{\chi }_i^a\mathrm{\Gamma }^AD_A(\widehat{\mathrm{\Omega }})\chi _i^a`$
$`{\displaystyle \frac{1}{8}}\overline{\psi }_A\mathrm{\Gamma }^{BC}\mathrm{\Gamma }^A(F_{iBC}^a+\widehat{F}_{iBC}^a)\chi _i^a+{\displaystyle \frac{\sqrt{2}}{48}}\left(\overline{\chi }_i^a\mathrm{\Gamma }^{ABC}\chi _i^a\right)\widehat{G}_{ABC11}]`$
where $`M^{11}`$ is the “downstairs” manifold while $`M_i^{10}`$ are its 10–dimensional boundaries. In the lowest approximation $`M^{11}`$ is just a product $`M^4\times X^6\times S^1/Z_2`$. Compactifying to $`d=4`$ in such an approximation we obtain $`^{\mathrm{?},\mathrm{?}}`$
$$G_N=\frac{\kappa _4^2}{8\pi }=\frac{\kappa ^2}{8\pi R_{11}V},$$
(46)
$$\alpha _{GUT}=\frac{(4\pi \kappa ^2)^{2/3}}{V}$$
(47)
with $`V`$ the volume of the Calabi–Yau manifold $`X^6`$ and $`R_{11}=\pi \rho `$ the $`S^1/Z_2`$ length.
The fundamental mass scale of the 11–dimensional theory is given by $`M_{11}=\kappa ^{2/9}`$. Let us see which value of $`M_{11}`$ is favoured in a phenomenological application. For that purpose we identify the Calabi–Yau volume $`V`$ with the GUT–scale: $`V(M_{GUT})^6`$. From (47) and the value of $`\alpha _{GUT}=1/25`$ at the grand unified scale, we can then deduce the value of $`M_{11}`$
$$V^{1/6}M_{11}=(4\pi )^{1/9}\alpha _{GUT}^{1/6}2.3,$$
(48)
to be a few times larger than the GUT–scale. In a next step we can now adjust the gravitational coupling constant by choosing the appropriate value of $`R_{11}`$ using (46). This leads to
$$R_{11}M_{11}=\left(\frac{M_{Planck}}{M_{11}}\right)^2\frac{\alpha _{GUT}}{8\pi (4\pi )^{2/3}}2.910^4\left(\frac{M_{Planck}}{M_{11}}\right)^2.$$
(49)
This simple analysis tells us the following:
* In contrast to the weakly coupled case ( where we had a prediction (44)), the correct value of $`M_{\mathrm{Planck}}`$ can be fitted by adjusting the value of $`R_{11}`$.
* The numerical value of $`R_{11}^1`$ turns out to be approximately an order of magnitude smaller than $`M_{11}`$.
* Thus the 11th dimension appears to be larger than the dimensions compactified on the Calabi–Yau manifold, and at an intermediate stage the world appears 5–dimensional with two 4–dimensional boundaries (walls).
We thus have the following picture of the evolution and unification of coupling constants. At low energies the world is 4–dimensional and the couplings evolve accordingly with energy: a logarithmic variation of gauge coupling constants and the usual power law behaviour for the gravitational coupling. Around $`R_{11}^1`$ we have an additional 5th dimension and the power law evolution of the gravitational interactions changes. Gauge couplings are not effected at that scale since the gauge fields live on the walls and do not feel the existence of the 5th dimension. Finally at $`M_{GUT}`$ the theory becomes 11–dimensional and both gravitational and gauge couplings show a power law behaviour and meet at the scale $`M_{11}`$, the fundamental scale of the theory. It is obvious that the correct choice of $`R_{11}`$ is needed to achieve unification. We also see that, although the theory is weakly coupled at $`M_{GUT}`$, this is no longer true at $`M_{11}`$. The naive estimate for the evolution of the gauge coupling constants between $`M_{GUT}`$ and $`M_{11}`$ goes with the sixth power of the scale. At $`M_{11}`$ we thus expect unification of the couplings at $`\alpha O(1)`$. In that sense, the M–theoretic description of the heterotic string gives an interpolation between weak coupling and moderate coupling. In $`d=4`$ this is not strong–weak coupling duality in the usual sense. We shall later come back to these questions when we discuss the appearance of a critical limit on the size of $`R_{11}`$.
These are, of course, rather qualitative results. In order to get a more quantitative feeling for the range of $`M_{11}`$ and $`R_{11}`$, let us be a bit more specific and write the relation of the unification scale $`M_{GUT}`$ to the characteristic size of the Calabi–Yau space as:
$$V^{1/6}=aM_{GUT}^1.$$
(50)
The above formula corresponds to the situation in which we identify the unification scale with the radius, $`R`$, of $`X^6`$ which volume is given by $`V=(aR)^6`$. We expect the parameter $`a`$ to be somewhere in the range from 1 to $`2\pi `$. Using the above identification and the value of $`M_{GUT}=310^{16}`$ GeV we obtain:
$$M_{11}\frac{2.3}{a}M_{GUT}.$$
(51)
As said before, the scale $`M_{11}`$ occurs to be of the order of the unification scale $`M_{GUT}`$. However, we do not expect $`M_{11}`$ to be smaller than $`M_{GUT}`$ because we need the ordinary logarithmic evolution of the gauge coupling constants up to $`M_{GUT}`$. In fact, $`M_{11}`$ should be somewhat bigger in order to allow for the evolution of $`\alpha `$ from its unification value 1/25 to the strong regime. Thus, we expect the parameter $`a`$ to be quite close to 1. Putting the above value of $`M_{11}`$ into eq. (49) we get the length of $`S^1/Z_2`$:
$$R_{11}9.2a^2M_{11}^14a^3M_{GUT}^1.$$
(52)
It is about one order of magnitude bigger than the scale characteristic for the 11–dimensional theory. This is the reason for the relatively large value of the $`d=4`$ Planck Mass. Of course $`R_{11}`$ can not be too large. For $`a`$ between 1 and 2.3 (values corresponding to $`M_{11}>M_{GUT}`$) we obtain $`R_{11}^1`$ in the range $`(6.210^{14}7.410^{15})`$ GeV (as we discussed, the parameter $`a`$ should not be too different from 1 so the upper part of the above range is favoured). Smaller values of $`R_{11}^1`$ seem to be very unnatural. Trying to push $`R_{11}^1`$ to smaller values would need a redefinition of $`M_{11}`$. For that purpose in $`^\mathrm{?}`$ a definition $`m_{11}=2\pi (4\pi \kappa ^2)^{1/9}`$ was used. This allows then to push $`a`$ to the extreme limit of $`2\pi `$. With these extreme choices of both $`a`$ and $`m_{11}`$ one would then be able to obtain $`R_{11}^1`$ as small as $`310^{13}`$ GeV. Values smaller than that (like values of $`10^{12}`$ GeV as sometimes quoted in the literature) cannot be obtained. In any case, even values in the lower $`10^{13}`$ GeV range seem to be in conflict with the critical value of $`R_{11}`$, as we shall see later.
## 11 The effective action in $`d=4`$
We now want to work out more explicitely the effective action in $`d=4`$ as obtained using the method of reduction and truncation.
We shall first consider again the $`d=10`$ effective field theory for the heterotic string (in more detail as given previously):
$$L=\frac{4}{(\alpha ^{})^3}d^{10}x\sqrt{g}\mathrm{exp}(2\varphi )\left(\frac{1}{(\alpha ^{})}R+\frac{1}{4}\mathrm{tr}F^2+\frac{1}{12}\alpha ^{}H^2+\mathrm{}\right),$$
(53)
where we have included the three index tensor field strength
$$H=dB+\omega ^{YM}\omega ^L.$$
(54)
$`B`$ is the two–index antisymmetric tensor while
$$\omega ^{YM}=\mathrm{Tr}(AF\frac{2}{3}A^3)$$
(55)
and
$$\omega ^L=\mathrm{Tr}(\omega R\frac{2}{3}\omega ^3)$$
(56)
are the Yang–Mills and Lorentz–Chern–Simons terms, respectively. The addition of these terms in the definition of $`H`$ is needed for supersymmetry and anomaly freedom of the theory.
To obtain the effective theory in $`d=4`$ dimensions we use as an approximation the method of reduction and truncation explained in ref. $`^\mathrm{?}`$. It essentially corresponds to a torus compactification, while truncating states to arrive at a $`d=4`$ theory with $`N=1`$ supersymmetry. In string theory compactified on an orbifold this would describe the dynamics of the untwisted sector. We retain the usual moduli fields $`S`$ and $`T`$ as well as matter fields $`C_i`$ that transform nontrivially under the observable sector gauge group. In this approximation, the Kähler potential is given by $`^{\mathrm{?},\mathrm{?}}`$
$$G=\mathrm{log}(S+S^{})3\mathrm{log}(T+T^{}2C_i^{}C_i)+\mathrm{log}\left|W\right|^2$$
(57)
with superpotential originating from the Chern–Simons terms $`\omega ^{YM}`$ $`^\mathrm{?}`$
$$W(C)=d_{ijk}C_iC_jC_k$$
(58)
and the gauge kinetic function is given by the dilaton field
$$f=S.$$
(59)
For a detailed discussion of this method and the explicit definition of the fields see the review $`^\mathrm{?}`$. These expressions for the $`d=4`$ effective action look quite simple and it remains to be seen whether this simplicity is true in general or whether it is an artifact of the approximation. Our experience with supergravity models tells us that the holomorphic functions $`W`$ and $`f`$ might be protected by nonrenormalization theorems, while the Kähler potential is strongly modified in perturbation theory. In addition we have to be aware of the fact that the expressions given above are at best representing a subsector of the theory. In orbifold compactification this would be the untwisted sector, and we know that the Kähler potential for twisted sectors fields will look quite different. Nonetheless the used approximation turned out to be useful for the discussion of those aspects of the theory that determine the dynamics of the $`T`$– and $`S`$–moduli. When trying to extract, however, detailed masses and other properties of the fields one should be aware of the fact, that some results might not be true in general and only appear as a result of the simplicity of the approximation.
So far the classical action. What about loop corrections? Not much can be said about the details of the corrections to the Kähler potential. This has to be discussed on a model by model basis. The situation with the superpotential is quite easy. There we expect a nonrenormalization theorem to be at work. The inclusion of other sectors of the theory will lead to new terms in the superpotential that in general have $`T`$–dependent coefficients. Such terms can be computed in simple cases by using e.g. methods of conformal field theory $`^\mathrm{?}`$.
The situation for $`f`$, the gauge–kinetic function is more interesting. Symmetries and holomorphicity lead us to believe, that although there are nontrivial corrections at one–loop, no more perturbative corrections are allowed at higher orders $`^{\mathrm{?},\mathrm{?}}`$. The existence of such corrections at one loop seems to be intimately connected to the mechanism of anomaly cancellation in the $`d=10`$ theory $`^{\mathrm{?},\mathrm{?}}`$. To see this consider one of the anomaly cancellation counter–terms introduced by Green and Schwarz $`^\mathrm{?}`$:
$$ϵ^{NEGAWSKLOV}B_{VO}\mathrm{Tr}F_{LKSW}^2F_{AGEN}^2.$$
(60)
We are interested in a $`d=4`$ theory with $`N=1`$ supersymmetry, and thus expect nontrivial vacuum expectation values for the curvature terms $`\mathrm{Tr}R^2`$ and field strengths $`\mathrm{Tr}F^2`$ in the extra six dimensions. Consistency of the theory requires a condition for the 3–index tensor field strength. For $`H`$ to be well defined, the quantity
$$dH=\mathrm{Tr}F^2\mathrm{Tr}R^2$$
(61)
has to vanish cohomologically $`^\mathrm{?}`$. In the simplest case (the so–called standard embedding leading to gauge group $`E_6\times E_8`$) one chooses equality pointwise $`\mathrm{Tr}R^2=\mathrm{Tr}F^2`$. Let us now assume that $`\mathrm{Tr}F_{agen}^2`$ is nonzero. The Green–Schwarz term given above by eq. (60) then leads to
$$ϵ^{mn}B_{mn}ϵ^{\mu \nu \rho \sigma }\mathrm{Tr}F_{\mu \nu }F_{\rho \sigma }$$
(62)
in the four–dimensional theory. An explicit inspection of the fields tells us that $`ϵ^{mn}B_{mn}`$ is the pseudoscalar axion that belongs to the $`T`$–superfield. Upon supersymmetrization the term in eq. (62) will then correspond to a one–loop correction to the holomorphic $`f`$–function (59) that is proportional to $`T`$ with the coefficient fixed entirely by the anomaly considerations. This is, of course, nothing else than a threshold correction. In the simple case of the standard embedding with gauge group $`E_6\times E_8`$ one obtains e.g.
$$f_6=S+ϵT;f_8=SϵT.$$
(63)
respectively, where $`ϵ`$ is the constant fixed by the anomaly. These results can be backed up by explicit calculations in string theory. In cases where such an explicit calculation is feasible, many more details about these corrections can be deduced. The above result (63) obtained in $`d=10`$ field theory represents an approximation of the exact result in the large $`T`$–limit. For a detailed discussion of these calculations and the limiting procedure see $`^\mathrm{?}`$. We have here mainly concentrated on that limit, because it represents a rather model independent statement.
Thus we have seen that there are corrections to the gauge–kinetic function at one loop. Their existence is found to be intimately related to the mechanism of anomaly cancellation. The corrections found are exactly those that are expected by general symmetry considerations $`^\mathrm{?}`$. In (63) we have given the result for the standard embedding. Coefficients might vary for more general cases, but the fact that they have opposite sign for the two separate groups is true in all known cases.
Superpotential and $`f`$–function should not receive further perturbative corrections beyond one loop. This implies that the knowledge of $`f`$ at one loop represents the full perturbative result. Combined with the fact that the coefficients are fixed by anomaly considerations one would then expect that this result for the $`f`$–function might be valid even beyond the weakly coupled limit. Not much can be said about the Kähler potential beyond one loop.
We now turn to the calculation in the M-theory case $`^\mathrm{?}`$. In the strongly coupled case we have to perform a compactification from $`d=11`$ to $`d=4`$. again we use the method of reduction and truncation. For the metric we write
$$g_{MN}^{(11)}=\left(\begin{array}{ccc}c_4e^\gamma e^{2\sigma }g_{\mu \nu }& & \\ & e^\sigma g_{mn}& \\ & & e^{2\gamma }e^{2\sigma }\end{array}\right)$$
(64)
with $`M,N=1\mathrm{}11`$; $`\mu ,\nu =1\mathrm{}4`$; $`m,n=5\mathrm{}10`$ and det($`g_{mn}`$)=1. This is the frame in which the 11–dimensional Einstein action gives the ordinary Einstein action after the reduction do $`d=4`$:
$$\frac{1}{2\kappa ^2}d^{11}x\sqrt{g^{(11)}}R^{(11)}=\frac{c_4\widehat{V_7}}{2\kappa ^2}d^4x\sqrt{g}R+\mathrm{}$$
(65)
where $`\widehat{V_7}=d^7x`$ is the coordinate volume of the compact 7–manifold and the scaling factor $`c_4`$ describes our freedom to choose the units in $`d=4`$. The most popular choice in the literature is $`c_4=1`$. This however corresponds to the unphysical situation in which the 4–dimensional Planck mass is determined by the choice of $`\widehat{V_7}`$ which is just a convention. With $`c_4=1`$ one needs further rescaling of the 4–dimensional metric. We instead prefer the choice
$$c_4=V_7/\widehat{V_7}$$
(66)
where $`V_7=d^7x\sqrt{g^{(7)}}`$ is the physical volume of the compact 7–manifold. This way we recover eq. (46) in which the 4–dimensional Planck mass depends on the physical, and not coordinate, volume of the manifold on which we compactify. As a result, if we start from the product of the 4–dimensional Minkowski space and some 7–dimensional compact space (in the leading order of the expansion in $`\kappa ^{2/3}`$) as a ground state in $`d=11`$ we obtain the Minkowski space with the standard normalization as the vacuum in $`d=4`$.
To find a more explicit formula for $`c_4`$ we have to discuss the fields $`\sigma `$ and $`\gamma `$ in some detail. In the leading approximation $`\sigma `$ is the overall modulus of the Calabi–Yau 6–manifold. We can divide it into a sum of the vacuum expectation value, $`\sigma `$, and the fluctuation $`\stackrel{~}{\sigma }`$. In general both parts could depend on all 11 coordinates but in practice we have to impose some restrictions. The vacuum expectation value can not depend on $`x^\mu `$ if the 4–dimensional theory is to be Lorentz–invariant. In the fluctuations we drop the dependence on the compact coordinates corresponding to the higher Kaluza–Klein modes. Furthermore, we know that in the leading approximation $`\sigma `$ is just a constant, $`\sigma _0`$ , while corrections depending on the internal coordinates, $`\sigma _1`$, are of the next order in $`\kappa ^{2/3}`$. Thus, we obtain
$$\sigma (x^\mu ,x^m,x^{11})=\sigma (x^m,x^{11})+\stackrel{~}{\sigma }(x^\mu )=\sigma _0+\sigma _1(x^m,x^{11})+\stackrel{~}{\sigma }(x^\mu ).$$
(67)
To make the above decomposition unique we define $`\sigma _0`$ by requiring that the integral of $`\sigma _1`$ over the internal space vanishes. The analogous decomposition can be also done for $`\gamma `$. With the above definitions the physical volume of the compact space is
$$V_7=d^7xe^{2\sigma }e^\gamma =e^{2\sigma _0}e^{\gamma _0}\widehat{V_7}$$
(68)
up to corrections of order $`\kappa ^{4/3}`$. Thus, the parameter $`c_4`$ can be written as
$$c_4=e^{2\sigma _0}e^{\gamma _0}.$$
(69)
The choice of the coordinate volumes is just a convention. For example in the case of the Calabi–Yau 6–manifold only the product $`e^{3\sigma }\widehat{V_6}`$ has physical meaning. For definiteness we will use the convention that the coordinate volumes are equal 1 in $`M_{11}`$ units. Thus, $`e^{3\sigma }`$ describes the Calabi–Yau volume in these units. Using eqs. (48,49) we obtain $`e^{3\sigma _0}=VM_{11}^6(2.3)^6`$, $`e^{\gamma _0}e^{\sigma _0}=R_{11}M_{11}9.2a^2`$. The parameter $`c_4`$ is equal to the square of the 4–dimensional Planck mass in these units and numerically $`c_4(35a)^2`$.
At the classical level we compactify on $`M^4\times X^6\times S^1/Z_2`$. This means that the vacuum expectation values $`\sigma `$ and $`\gamma `$ are just constants and eq. (67) reduces to
$$\sigma =\sigma _0+\stackrel{~}{\sigma }(x^\mu ),\gamma =\gamma _0+\stackrel{~}{\gamma }(x^\mu ).$$
(70)
In such a situation $`\sigma `$ and $`\gamma `$ are 4–dimensional fields. We introduce two other 4–dimensional fields by the relations
$`{\displaystyle \frac{1}{4!c_4}}e^{6\sigma }G_{11\lambda \mu \nu }`$ $`=`$ $`ϵ_{\lambda \mu \nu \rho }\left(^\rho D\right),`$ (71)
$`C_{11a\overline{b}}`$ $`=`$ $`C_{11}\delta _{a\overline{b}}`$ (72)
where $`x^a`$ ($`x^{\overline{b}}`$) is the holomorphic (antiholomorphic) coordinate of the Calabi–Yau manifold. Now we can define the dilaton and the modulus fields by
$`𝒮`$ $`=`$ $`{\displaystyle \frac{1}{\left(4\pi \right)^{2/3}}}\left(e^{3\sigma }+i24\sqrt{2}D\right),`$ (73)
$`𝒯`$ $`=`$ $`{\displaystyle \frac{1}{\left(4\pi \right)^{2/3}}}\left(e^\gamma +i6\sqrt{2}C_{11}+C_i^{}C_i\right)`$ (74)
where the observable sector matter fields $`C_i`$ originate from the gauge fields $`A_M`$ on the 10–dimensional observable wall (and $`M`$ is an index in the compactified six dimensions). The Kähler potential takes its standard form as in eq. (57)
$$K=\mathrm{log}(𝒮+𝒮^{})3\mathrm{log}(𝒯+𝒯^{}2C_i^{}C_i).$$
(75)
The imaginary part of $`𝒮`$ (Im$`𝒮`$) corresponds to the model independent axion, and with the above normalization the gauge kinetic function is $`f=𝒮`$. We have also
$$W(C)=d_{ijk}C_iC_jC_k$$
(76)
Thus the action to leading order is very similar to the weakly coupled case.
Before drawing any conclusion from the formulae obtained above we have to discuss a possible obstruction at the next to leading order. For the 3–index tensor field $`H`$ in $`d=10`$ supergravity to be well defined one has to satisfy $`dH=\mathrm{tr}F_1^2+\mathrm{tr}F_2^2\mathrm{tr}R^2=0`$ cohomologically. In the simplest case of the standard embedding one assumes $`\mathrm{tr}F_1^2=\mathrm{tr}R^2`$ locally and the gauge group is broken to $`E_6\times E_8`$. Since in the M–theory case the two different gauge groups live on the two different boundaries (walls) of space–time such a cancellation point by point is no longer possible $`^\mathrm{?}`$. We expect nontrivial vacuum expectation values (vevs) of
$$(dG)\underset{i}{}\delta (x^{11}x_i^{11})\left(\mathrm{tr}F_i^2\frac{1}{2}\mathrm{tr}R^2\right)$$
(77)
at least on one boundary ($`x_i^{11}`$ is the position of $`i`$–th boundary). In the case of the standard embedding we would have $`\mathrm{tr}F_1^2\frac{1}{2}\mathrm{tr}R^2=\frac{1}{2}\mathrm{tr}R^2`$ on one and $`\mathrm{tr}F_2^2\frac{1}{2}\mathrm{tr}R^2=\frac{1}{2}\mathrm{tr}R^2`$ on the other boundary. This might pose a severe problem since a nontrivial vev of $`G`$ might be in conflict with supersymmetry ($`G_{11ABC}=H_{ABC}`$). The supersymmetry transformation law in $`d=11`$ reads
$$\delta \psi _M=D_M\eta +\frac{\sqrt{2}}{288}G_{IJKL}\left(\mathrm{\Gamma }_M^{IJKL}8\delta _M^I\mathrm{\Gamma }^{JKL}\right)\eta +\mathrm{}$$
(78)
Supersymmetry will be broken unless e.g. the derivative term $`D_M\eta `$ compensates the nontrivial vev of $`G`$. Witten has shown $`^\mathrm{?}`$ that such a cancellation can occur and constructed the solution in the linearized approximation (linear in the expansion parameter $`\kappa ^{2/3}`$). This solution requires some modification of the metric on $`M^{11}`$:
$$g_{MN}^{(11)}=\left(\begin{array}{ccc}(1+b)\eta _{\mu \nu }& & \\ & (g_{ij}+h_{ij})& \\ & & (1+\gamma ^{})\end{array}\right).$$
(79)
$`M^{11}`$ is no longer a direct product $`M^4\times X^6\times S^1/Z_2`$ because $`b`$, $`h_{ij}`$ and $`\gamma ^{}`$ depend now on the compactified coordinates. The volume of $`X^6`$ depends on $`x^{11}`$ $`^\mathrm{?}`$:
$$\frac{}{x^{11}}V=\frac{\sqrt{2}}{8}d^6x\sqrt{g}\omega ^{AB}\omega ^{CD}G_{ABCD}$$
(80)
where the integral is over the Calabi–Yau manifold $`X^6`$ and $`\omega `$ is the corresponding Kähler form. The parameter $`(1+b)`$ is the scale factor of the Minkowski 4–manifold and depends on $`x^{11}`$ in the following way
$$\frac{}{x^{11}}b=\frac{1}{2}\frac{}{x^{11}}\mathrm{log}v_4=\frac{\sqrt{2}}{24}\omega ^{AB}\omega ^{CD}G_{ABCD}$$
(81)
where $`v_4`$ is the physical volume for some fixed coordinate volume in $`M^4`$. In our simple reduction and truncation method with the metric $`g_{MN}^{(11)}`$ given by eq. (64) we can reproduce the $`x^{11}`$ dependence of $`V`$ and $`v_4`$. The volume of $`X^6`$ is determined by $`\sigma `$:
$$\frac{}{x^{11}}\mathrm{log}V=\frac{}{x^{11}}\left(3\sigma \right)=3\frac{}{x^{11}}\sigma $$
(82)
while the scale factor of $`M^4`$ can be similarly expressed in terms of $`\sigma `$ and $`\gamma `$ fields:
$$\frac{}{x^{11}}\mathrm{log}v_4=\frac{}{x^{11}}\left(2\gamma +4\sigma \right)=\frac{}{x^{11}}(2\gamma +4\sigma ).$$
(83)
Substituting $`\sigma `$ with $`\sigma `$ in the above two equations is allowed because, due to our decomposition (67), only the vev of $`\sigma `$ depends on the internal coordinates (the same is true for $`\gamma `$). The scale factor $`b`$ calculated in ref. $`^\mathrm{?}`$ depends also on the Calabi–Yau coordinates. Such a dependence can not be reproduced in our simple reduction and truncation compactification so we have to average eq. (81) over $`X^6`$. Using equations (8083) after such an averaging we obtain (to leading order in the expansion parameter $`\kappa ^{2/3}`$) $`^\mathrm{?}`$
$$\frac{\gamma }{x^{11}}=\frac{\sigma }{x^{11}}=\frac{\sqrt{2}}{24}\frac{d^6x\sqrt{g}\omega ^{AB}\omega ^{CD}G_{ABCD}}{d^6x\sqrt{g}}.$$
(84)
Substituting the vacuum expectation value of $`G`$ found in $`^\mathrm{?}`$ we can rewrite it in the form
$$\frac{\gamma }{x^{11}}=\frac{\sigma }{x^{11}}=\frac{2}{3}\alpha \kappa ^{2/3}V^{2/3}$$
(85)
where
$$\alpha =\frac{\pi c}{2(4\pi )^{2/3}}$$
(86)
and c is a constant of order unity given for the standard embedding of the spin connection by
$$c=V^{1/3}\left|\frac{\omega \mathrm{tr}(RR)}{8\pi ^2}\right|.$$
(87)
Our calculations, as those of Witten, are valid only in the leading nontrivial order in the $`\kappa ^{2/3}`$ expansion. The expression (85) for the derivatives of $`\sigma `$ and $`\gamma `$ have explicit factor $`\kappa ^{2/3}`$. This means that we should take the lowest order value for the Calabi–Yau volume in that expression. An analogous procedure has been used in obtaining all formulae presented in this paper. We always expand in $`\kappa ^{2/3}`$ and drop all terms which are of higher order than our approximation. Taking the above into account and using our units in which $`M_{11}=1`$ we can rewrite eq. (85) in the simple form:
$$\frac{\gamma }{x^{11}}=\frac{\sigma }{x^{11}}=\frac{2}{3}\alpha e^{2\sigma _0}.$$
(88)
Eqs. (8488) have been derived in ref. $`^\mathrm{?}`$. As we will see in the following, these results contain all the information to deduce the effective action, i.e. Kähler potential, superpotential and gauge kinetic function of the 4–dimensional effective supergravity theory.
It is the above dependence of $`\sigma `$ and $`\gamma `$ on $`x^{11}`$ that leads to these consequences. One has to be careful in defining the fields in $`d=4`$. It is obvious, that the 4–dimensional fields $`𝒮`$ and $`𝒯`$ can not be any longer defined by eqs. (73, 74) because now $`\sigma `$ and $`\gamma `$ are 5–dimensional fields. We have to integrate out the dependence on the 11th coordinate. In the present approximation, this procedure is quite simple: we have to replace $`\sigma `$ and $`\gamma `$ in the definitions of $`𝒮`$ and $`𝒯`$ with their averages over the $`S^1/Z_2`$ interval $`^\mathrm{?}`$. With the linear dependence of $`\sigma `$ and $`\gamma `$ on $`x^{11}`$ their average values coincide with the values taken at the middle of the $`S^1/Z_2`$ interval
$$\overline{\sigma }=\sigma \left(\frac{\pi \rho }{2}\right)=\sigma _0+\stackrel{~}{\sigma }(x^\mu ),$$
(89)
$$\overline{\gamma }=\gamma \left(\frac{\pi \rho }{2}\right)=\gamma _0+\stackrel{~}{\gamma }(x^\mu ).$$
(90)
When we reduce the boundary part of the Lagrangian of M–theory to 4 dimensions we find exponents of $`\sigma `$ and $`\gamma `$ fields evaluated at the boundaries. Using eqs. (67) and (88) we get
$`e^\gamma |_{M_i^{10}}`$ $`=`$ $`e^{\gamma _0}\pm {\displaystyle \frac{1}{3}}\alpha e^{3\sigma _0},`$ (91)
$`e^{3\sigma }|_{M_i^{10}}`$ $`=`$ $`e^{3\sigma _0}\pm \alpha e^{\gamma _0}.`$ (92)
The above formulae have very important consequences for the definitions of the Kähler potential and the gauge kinetic functions. For example, the coefficient in front of the $`D_\mu C_i^{}D^\mu C_i`$ kinetic term is proportional to $`e^\gamma `$ evaluated at the $`E_6`$ wall where the matter fields propagate. At the lowest order this was just $`e^{\gamma _0}`$ or $`𝒯^1`$ up to some numerical factor. From eq. (91) we see that at the next to leading order also $`𝒮^1`$ is involved with relative coefficient $`\alpha /3`$. Taking such corrections into account we find that at this order the Kähler potential is given by
$$K=\mathrm{log}(𝒮+𝒮^{})+\frac{2\alpha C_i^{}C_i}{𝒮+𝒮^{}}3\mathrm{log}(𝒯+𝒯^{}2C_i^{}C_i)$$
(93)
with $`𝒮`$ and $`𝒯`$ now defined by
$`𝒮`$ $`=`$ $`{\displaystyle \frac{1}{\left(4\pi \right)^{2/3}}}\left(e^{3\overline{\sigma }}+i24\sqrt{2}\overline{D}+\alpha C_i^{}C_i\right),`$ (94)
$`𝒯`$ $`=`$ $`{\displaystyle \frac{1}{\left(4\pi \right)^{2/3}}}\left(e^{\overline{\gamma }}+i6\sqrt{2}\overline{C}_{11}+C_i^{}C_i\right)`$ (95)
where bars denote averaging over the 11th dimension. It might be of some interest to note that the combination $`𝒮𝒯^3`$ is independent of $`x^{11}`$ even before this averaging procedure took place. The solution above is valid only for terms at most linear in $`\alpha `$. Keeping this in mind we could write the Kähler potential also in the form
$$K=\mathrm{log}(𝒮+𝒮^{}2\alpha C_i^{}C_i)3\mathrm{log}(𝒯+𝒯^{}2C_i^{}C_i).$$
(96)
Equipped with this definition the calculation of the gauge kinetic function(s) from eqs. (88, 92) becomes a trivial exercise $`^\mathrm{?}`$. In the five–dimensional theory $`f`$ depends on the 11–dimensional coordinate as well, thus the gauge kinetic function takes different values at the two walls. The averaging procedure allows us to deduce these functions directly. For the simple case at hand (the so–called standard embedding) eq. (92) gives $`^\mathrm{?}`$
$$f_6=𝒮+\alpha 𝒯;f_8=𝒮\alpha 𝒯.$$
(97)
It is a special property of the standard embedding that the coefficients are equal and opposite. The coefficients vary for more general cases. This completes the discussion of the $`d=4`$ effective action in next to leading order, noting that the superpotential does not receive corrections at this level.
The nontrivial dependence of $`\sigma `$ and $`\gamma `$ on $`x^{11}`$ can also enter definitions and/or interactions of other 4–dimensional fields. Let us next consider the gravitino. After all we have to show that this field is massless to give the final proof that the given solution respects supersymmetry. Its 11–dimensional kinetic term
$$\frac{1}{2}\sqrt{g}\overline{\psi }_I\mathrm{\Gamma }^{IJK}D_J\psi _K$$
(98)
remains diagonal after compactification to $`d=4`$ if we define the 4–dimensional gravitino, $`\psi _\mu ^{(4)}`$, and dilatino ,$`\psi _{11}^{(4)}`$, fields by the relations
$`\psi _\mu `$ $`=`$ $`e^{(\sigma \sigma _0)/2}e^{(\gamma \gamma _0)/4}\left(\psi _\mu ^{(4)}+{\displaystyle \frac{1}{\sqrt{6}}}\mathrm{\Gamma }_\mu \psi _{11}^{(4)}\right),`$ (99)
$`\psi _{11}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{6}}}e^{(\sigma \sigma _0)/2}e^{(\gamma \gamma _0)/4}\mathrm{\Gamma }^{11}\psi _{11}^{(4)}.`$ (100)
The $`d=11`$ kinetic term (98) gives after the compactification also a mass term for the $`d=4`$ gravitino of the form
$$\frac{3}{8}e^{\sigma _0}e^{\gamma _0}\frac{\gamma }{x^{11}}=\frac{\sqrt{2}}{64}e^{\sigma _0}e^{\gamma _0}\frac{d^6x\sqrt{g}\omega ^{AB}\omega ^{CD}G_{ABCD}}{d^6x\sqrt{g}}=\frac{1}{4}\alpha e^{\sigma _0}e^{\gamma _0}.$$
(101)
The sources of such a term are nonzero values of the spin connection components $`\omega _\mu ^{\alpha 11}`$ and $`\omega _m^{a11}`$ resulting from the $`x^{11}`$ dependence of the metric. It is a constant mass term from the 4–dimensional point of view. This, however, does not mean that the gravitino mass is nonzero. There is another contribution from the 11–dimensional term
$$\frac{\sqrt{2}}{384}\sqrt{g}\overline{\psi }_I\mathrm{\Gamma }^{IJKLMN}\psi _N\left(G_{JKLM}+\widehat{G}_{JKLM}\right).$$
(102)
After redefining fields according to (99,100) and averaging the nontrivial vacuum expectation value of $`G`$ over $`X^6`$ we get from eq. (102) a mass term which exactly cancels the previous contribution (101). The gravitino is massless – the result which we expect in a model with unbroken supersymmetry and vanishing cosmological constant. Thus, we find that our simple reduction and truncation method (including the correct $`x^{11}`$ dependence in next to leading order) reproduces the main features of the model.
The factor $`\mathrm{exp}(3\sigma )`$ represents the volume of the six–dimensional compact space in units of $`M_{11}^6`$. The $`x^{11}`$ dependence of $`\sigma `$ then leads to the geometrical picture that the volume of this space varies with $`x^{11}`$ and differs at the two boundaries:
$$V_{E_8}=V_{E_6}2\pi ^2\rho \left(\frac{\kappa }{4\pi }\right)^{2/3}\left|\omega \frac{\mathrm{tr}(FF)\frac{1}{2}\mathrm{tr}(RR)}{8\pi ^2}\right|$$
(103)
where the integral is over $`X^6`$ at the $`E_6`$ boundary. In the given approximation, this variation is linear, and for growing $`\rho `$ the volume on the $`E_8`$ side becomes smaller and smaller. At a critical value of $`\rho `$ the volume will thus vanish and this will provide us with an upper limit on $`\rho `$:
$$\rho <\rho _{crit}=\frac{(4\pi )^{2/3}}{c\pi ^2}M_{11}^3V_{E_6}^{2/3}$$
(104)
where $`c`$ was defined in eq. (87). To estimate the numerical value of $`\rho _{crit}`$ we first recall that from eq. (46) we obtained <sup>††</sup><sup>††</sup>††With $`V`$ depending on $`x^{11}`$ we have to specify which values should be used in eqs. (46,47,50). The appropriate choice in the expression for $`G_N`$ is the average value of $`V`$ while in the expressions for $`\alpha _{GUT}`$ and for the $`V`$$`M_{GUT}`$ relation we have to use $`V`$ evaluated at the $`E_6`$ wall.
$$M_{11}V_{E_6}^{1/6}=\left(\alpha _{GUT}(4\pi )^{2/3}\right)^{1/6}2.3.$$
(105)
Thus, we get
$$\rho ^1>\rho _{crit}^10.16cV_{E_6}^{1/6}.$$
(106)
The numerical value of $`V`$ at the $`E_6`$ boundary depends on what we identify with the unification scale $`M_{GUT}`$ via eq. (50):
$$V_{E_6}^{1/6}=aM_{GUT}^1$$
(107)
with $`a`$ somewhere between 1 and about 2. Thus, the bound (106) can be written in the form
$$R_{11}^1>0.05\frac{c}{a}M_{GUT}.$$
(108)
For the phenomenological applications we have to check whether our preferred choice of $`6.210^{14}\mathrm{GeV}<R_{11}^1<7.410^{15}`$ GeV that fits the correct value of the $`d=4`$ Planck mass satisfies the bound (108). In a rather extreme case of $`c=1`$ and $`a=2.3`$ we find that the upper bound on $`R_{11}^1`$ is of the order of $`6.510^{14}`$ GeV. Even for $`c=1`$ this bound goes up to about $`1.510^{15}`$ GeV if we identify $`V^{1/6}`$ with $`M_{GUT}`$. Although some coefficients are model dependent we find in general that the bound can be satisfied, but that $`R_{11}`$ is quite close to its critical value. Values of $`R_{11}^1`$ about $`10^{12}`$ GeV as necessary in $`^\mathrm{?}`$ seem to be beyond the critical value, even with the modifications discussed before. In any case, models where supersymmetry is broken by a Scherk–Schwarz mechanism seem to require the absence of the next to leading order corrections in (97), i.e. $`\alpha =0`$. It remains to be seen whether such a possibility can be realized.
Inspection of (63) and (97) reveals a close connection between the strongly and weakly coupled case $`^{\mathrm{?},\mathrm{?}}`$. The variation of the Calabi–Yau manifold volume as discussed above is the analogue of the one loop correction of the gauge kinetic function (63) in the weakly coupled case and has the same origin, namely a Green–Schwarz anomaly cancellation counterterm. In fact, also in the strongly coupled case this leads to a correction for the gauge coupling constants at the $`E_6`$ and $`E_8`$ side. As seen, gauge couplings are no longer given by the (averaged) $`𝒮`$–field, but by that combination of the (averaged) $`𝒮`$ and $`𝒯`$ fields which corresponds to the $`𝒮`$–field before averaging at the given boundary leading to
$$f_{6,8}=𝒮\pm \alpha 𝒯$$
(109)
at the $`E_6`$ ($`E_8`$) side respectively. The critical value of $`R_{11}`$ will correspond to infinitely strong coupling at the $`E_8`$ side $`𝒮\alpha 𝒯=0`$. Since we are here close to criticality a correct phenomenological fit of $`\alpha _{\mathrm{GUT}}=1/25`$ should include this correction $`\alpha _{\mathrm{GUT}}^1=𝒮+\alpha 𝒯`$ where $`𝒮`$ and $`\alpha 𝒯`$ give comparable contributions. This is a difference to the weakly coupled case, where in $`f=S+ϵT`$ the latter contribution was small compared to $`S`$. This stable result for the corrections to $`f`$ when going from weak coupling to strong coupling is only possible because of the rather special properties of $`f`$. $`f`$ does not receive further perturbative corrections beyond one loop $`^{\mathrm{?},\mathrm{?}}`$, and the one loop corrections are determined by the anomaly considerations. The formal expressions for the corrections are identical, the difference being only that in the strongly coupled case these corrections are as important as the classical value.
## 12 Supersymmetry breaking at the hidden wall
We shall now discuss the question of supersymmetry breakdown within this framework. We consider the breakdown of supersymmetry in a hidden sector, transmitted to the observable sector via gravitational interactions. Such a scenario was suggested in $`^\mathrm{?}`$ after having observed that gaugino condensation can break supersymmetry in $`d=4`$ supergravity models. As we have seen, a nontrivial gauge kinetic function $`f`$ seems to be necessary for such a mechanism to work $`^\mathrm{?}`$. In the heterotic string both ingredients, a hidden sector $`E_8`$ and a nontrivial $`f`$, were present in a natural way and a coherent picture of supersymmetry breakdown via gaugino condensation emerged $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. In the strongly coupled case, such a mechanism can be realized as well $`^{\mathrm{?},\mathrm{?}}`$. In fact the notion of the hidden sector acquires a geometrical interpretation: the gaugino condensate forms at one boundary (the hidden wall) of spacetime. We shall now discuss this mechanism in detail. First we remind you of some relevant formulae in the weakly coupled case. Our aim then is to compare the strong coupling regime with the weak coupling regime and clarify similarities as well as differences. For the weakly coupled case we start with the action of $`d=10`$ supergravity. Supersymmetry transformation laws for the $`d=10`$ gravitino fields $`\psi _M`$ and the dilatino field $`\lambda `$ are written <sup>‡‡</sup><sup>‡‡</sup>‡‡Here we use the conventions of $`^\mathrm{?}`$, where the Lagrangian is given in the Einstein frame. To recover the effective action (53) in the string frame, one has to make a proper Weyl transformation and identify $`\phi =\mathrm{exp}(\varphi /3)`$.
$`\delta \lambda `$ $`=`$ $`{\displaystyle \frac{1}{8}}\phi ^{3/4}\mathrm{\Gamma }^{MNP}H_{MNP}+{\displaystyle \frac{\sqrt{2}}{384}}\mathrm{\Gamma }^{MNP}\overline{\chi }^a\mathrm{\Gamma }_{MNP}\chi ^a+\mathrm{},`$
$`\delta \psi _M`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{32}}\phi ^{3/4}(\mathrm{\Gamma }_M^{NPQ}9\delta _M^N\mathrm{\Gamma }^{PQ})H_{NPQ}`$ (110)
$`+{\displaystyle \frac{1}{256}}(\mathrm{\Gamma }_{MNPQ}5g_{MN}\mathrm{\Gamma }_{PQ})\overline{\chi }^a\mathrm{\Gamma }^{NPQ}\chi ^a+\mathrm{},`$
implying that a condensate of gauginos $`\overline{\chi }\chi `$ and/or non–vanishing vevs of the $`H`$ fields may break supersymmetry. Here we assume the appearance of the gaugino condensate in the hidden sector
$$\overline{\chi }^a\mathrm{\Gamma }_{mnp}\chi ^a=\mathrm{\Lambda }^3ϵ_{mnp},$$
(111)
with $`\mathrm{\Lambda }`$ being the gaugino condensation scale and $`ϵ_{mnp}`$ the covariantly constant holomorphic three–form. The perfect square structure seen in the Lagrangian $`^\mathrm{?}`$
$$\frac{3}{4}\phi ^{3/2}(H_{MNP}\sqrt{2}\phi ^{3/4}\overline{\chi }^a\mathrm{\Gamma }_{MNP}\chi ^a)^2$$
(112)
will be a very important ingredient to discuss the quantitative properties of the mechanism. When reducing to the $`d=4`$ effective action we will find a cancellation of the vevs of the $`H`$ field and the gaugino condensate at the minimum of the potential such that the term in eq. (112) vanishes. Before we look at this in detail, let us first comment on such a possible vev of $`H`$ and a possible quantization condition of the antisymmetric tensor. In $`^\mathrm{?}`$ it was shown, that an antisymmetric tensor field $`H=dB`$ has a quantized vacuum expectation value. In many subsequent papers this has been incorrectly taken as an argument for the quantization of the vev of $`H=dB+\omega ^{YM}\omega ^L`$ as given in eq. (54). The correct way to interpret this situation is to have a cancellation of the gaugino condensate with the vev of a Chern–Simons term $`^\mathrm{?}`$, for which such a quantization condition does not hold. After all the Chern–Simons term $`\omega ^{YM}`$ contains the superpotential of the $`d=4`$ effective theory $`^\mathrm{?}`$. This cancellation leads to a certain combination of $`\psi _M`$ and $`\lambda `$ as the candidate goldstino that will provide the longitudinal component of the gravitino. While in $`d=10`$ this looks rather complicated, it simplifies tremendously once one reduces to $`d=4`$. Qualitatively the scalar potential takes the following form at the classical level (for the detailed factors see $`^\mathrm{?}`$):
$$V=\frac{1}{ST^3}\left[W2(ST)^{3/2}(\overline{\chi }\chi )^2+\frac{T}{3}\frac{W}{C}^2\right].$$
(113)
We observe the important fact that the potential is positive and vanishes at the minimum. Thus we have broken supersymmetry with a vanishing cosmological constant at the classical level. The first term in the brackets of eq. (113) corresponds to the contribution from eq. (112) once reduced to $`d=4`$ and vanishes at the minimum. In the $`d=4`$ theory it represents the auxiliary component $`F_S`$ of the dilaton superfield $`S`$. Thus we have $`F_S=0`$ and supersymmetry is broken by a nonvanishing vev of $`F_T`$ $`^\mathrm{?}`$. The goldstino is then the fermion in the $`T`$–multiplet and we are dealing with a situation that has later been named moduli–dominated supersymmetry breakdown. This fact has its origin in the special properties of the $`d=10`$ action (the term in eq. (112)) and seems to be of rather general validity. The statement $`F_S=0`$ is, of course, strictly valid only in the classical theory. The corrections discussed in section 3, eq. (63) will slightly change these results as we shall discuss later.
Having minimized the potential and identified the goldstino we can now compute the gravitino mass according to the standard procedure. The result has a direct physical meaning because we are dealing with a theory with vanishing vacuum energy. We obtain
$$m_{3/2}\frac{F_T}{M_{Planck}}\frac{\mathrm{\Lambda }^3}{M_{Planck}^2}.$$
(114)
A value of $`\mathrm{\Lambda }10^{13}`$ GeV will thus lead to a gravitino mass in the TeV region.
Next we turn to supersymmetry breaking in the strongly coupled case ($`d=11`$ M–theory picture) and start with the $`d=11`$ action. Supersymmetry transformation laws for the gravitino fields in this case are given by
$`\delta \psi _A`$ $`=`$ $`D_A\eta +{\displaystyle \frac{\sqrt{2}}{288}}G_{IJKL}\left(\mathrm{\Gamma }_A^{IJKL}8\delta _A^I\mathrm{\Gamma }^{JKL}\right)\eta `$ (115)
$`{\displaystyle \frac{1}{1152\pi }}\left({\displaystyle \frac{\kappa }{4\pi }}\right)^{2/3}\delta (x^{11})\left(\overline{\chi }^a\mathrm{\Gamma }_{BCD}\chi ^a\right)\left(\mathrm{\Gamma }_A^{BCD}6\delta _A^B\mathrm{\Gamma }^{CD}\right)\eta +\mathrm{}`$
$`\delta \psi _{11}`$ $`=`$ $`D_{11}\eta +{\displaystyle \frac{\sqrt{2}}{288}}G_{IJKL}\left(\mathrm{\Gamma }_{11}^{IJKL}8\delta _{11}^I\mathrm{\Gamma }^{JKL}\right)\eta `$ (116)
$`+{\displaystyle \frac{1}{1152\pi }}\left({\displaystyle \frac{\kappa }{4\pi }}\right)^{2/3}\delta (x^{11})\left(\overline{\chi }^a\mathrm{\Gamma }_{ABC}\chi ^a\right)\mathrm{\Gamma }^{ABC}\eta +\mathrm{}`$
where gaugino bilinears appear in the right hand side of both expressions. Again we consider gaugino condensation at the hidden $`E_8`$ boundary
$$\overline{\chi }^a\mathrm{\Gamma }_{ijk}\chi ^a=g_8^2\mathrm{\Lambda }^3ϵ_{ijk}.$$
(117)
The $`E_8`$ gauge coupling constant appears in this equation because the straightforward reduction and truncation leaves a non–canonical normalization for the gaugino kinetic term. An important property of the weakly coupled case (d=10 Lagrangian) was the fact that the gaugino condensate and the three–index tensor field $`H`$ contributed to the scalar potential in a full square. Hořava made the important observation that a similar structure appears in the M–theory Lagrangian as well $`^\mathrm{?}`$:
$$\frac{1}{12\kappa ^2}_{M^{11}}d^{11}x\sqrt{g}\left(G_{ABC11}\frac{\sqrt{2}}{32\pi }\left(\frac{\kappa }{4\pi }\right)^{2/3}\delta (x^{11})\overline{\chi }^a\mathrm{\Gamma }_{ABC}\chi ^a\right)^2$$
(118)
with the obvious relation between $`H`$ and $`G`$. Let us now have a closer look at the form of $`G`$. At the next to leading order we have
$`G_{11ABC}`$ $`=`$ $`(_{11}C_{ABC}+\text{permutations})`$ (119)
$`+{\displaystyle \frac{1}{4\pi \sqrt{2}}}\left({\displaystyle \frac{\kappa }{4\pi }}\right)^{2/3}{\displaystyle \underset{i}{}}\delta (x^{11}x_i^{11})(\omega _{ABC}^{YM}{\displaystyle \frac{1}{2}}\omega _{ABC}^L).`$
Observe, that in the bulk we have $`G=dC`$ with the Chern–Simons contributions confined to the boundaries. Formula (118) suggests a cancellation between the gaugino condensate and the $`G`$–field in a way very similar to the weakly coupled case, but the nature of the cancellation of the terms becomes much more transparent now. In the former case we had to argue via the quantization condition for $`dB`$ that the gaugino condensate is cancelled by one of the Chern–Simons terms. Here this becomes obvious. The condensate is located at the wall as are the Chern–Simons terms, so this cancellation has to happen locally at the wall and $`dC`$ should vanish for $`G`$ not to have a vev in the bulk. In any case there is a quantization condition for $`dC`$ as well $`^\mathrm{?}`$.
So this cancellation is very similar to the one in the weakly coupled case. At the minimum of the potential we obtain $`G_{ABCD}=0`$ everywhere and
$$G_{ABC11}=\frac{\sqrt{2}}{32\pi }\left(\frac{\kappa }{4\pi }\right)^{2/3}\delta (x^{11})\overline{\chi }^a\mathrm{\Gamma }_{ABC}\chi ^a$$
(120)
at the hidden wall. Eqs. (115) and (116) then become
$`\delta \psi _A`$ $`=`$ $`D_A\eta +\mathrm{}`$ (121)
$`\delta \psi _{11}`$ $`=`$ $`D_{11}\eta +{\displaystyle \frac{1}{384\pi }}\left({\displaystyle \frac{\kappa }{4\pi }}\right)^{2/3}\delta (x^{11})\left(\overline{\chi }^a\mathrm{\Gamma }_{ABC}\chi ^a\right)\mathrm{\Gamma }^{ABC}\eta +\mathrm{}.`$ (122)
An inspection of the potential shows that $`\delta \psi _{11}`$ is nonvanishing and supersymmetry is spontaneously broken. Because of the cancellation in eq. (118), the cosmological constant vanishes to leading order. Recalling supersymmetry transformation law for the elfbein
$$\delta e_I^m=\frac{1}{2}\overline{\eta }\mathrm{\Gamma }^m\psi _I,$$
(123)
one finds that the superpartner of the $`𝒯`$ field plays the role of the goldstino. Again we have a situation where $`F_𝒮=0`$ (due to the cancellation in (118)) with nonvanishing $`F_𝒯`$. But here we find the novel and interesting situation that $`F_𝒯`$ differs from zero only at the hidden wall, although the field itself is a bulk field. In general it would be interesting to consider also situations where the goldstino is not a bulk but a wall field.
At that wall our discussion is completely 4–dimensional although we are still dealing effectively with a $`d=5`$ theory. To reach the effective theory in $`d=4`$ we have to integrate out the dependence of the $`x^{11}`$ coordinate. As in the previous section this can be performed by the averaging procedure explained there. With the gaugino condensation scale $`\mathrm{\Lambda }`$ sufficiently small compared to the compactification scale $`M_{GUT}`$, the low–energy effective theory is well described by four dimensional $`N=1`$ supergravity in which supersymmetry is spontaneously broken. In this case, the modes which remain at low energies will be well approximated by constant modes along the $`x^{11}`$ direction. This observation justifies our averaging procedure to obtain four dimensional quantities. Averaging $`\delta \psi _{11}`$ over $`x^{11}`$, we thus obtain the vev of the auxiliary field $`F_𝒯`$
$$F_𝒯=\frac{1}{2}𝒯\frac{𝑑x^{11}\sqrt{g_{1111}}\delta \psi _{11}}{𝑑x^{11}\sqrt{g_{1111}}}.$$
(124)
Note that this procedure allows for a nonlocal cancellation of the vev of the auxiliary field in $`d=4`$. A condensate with equal size and opposite sign at the observable wall could cancel the effect and restore supersymmetry. Using $`𝑑x^{11}\sqrt{g_{1111}}\delta (x^{11})=1`$, the auxiliary field is found to be
$$F_𝒯=𝒯\frac{1}{32\pi (4\pi )^{2/3}}\frac{g_8^2\mathrm{\Lambda }^3}{R_{11}M_{11}^3}$$
(125)
Similarly one can easily show that $`F_𝒮`$ as well as the vacuum energy vanish. This allows us then to unambiguously determine the gravitino mass, which is related to the auxiliary field in the following way:
$$m_{3/2}=\frac{F_𝒯}{𝒯+𝒯^{}}=\frac{1}{64\pi (4\pi )^{2/3}}\frac{g_8^2\mathrm{\Lambda }^3}{R_{11}M_{11}^3}=\frac{\pi }{2}\frac{\mathrm{\Lambda }^3}{M_{Planck}^2}.$$
(126)
As a nontrivial check one may calculate the gravitino mass in a different way. A term in the Lagrangian
$$\frac{\sqrt{2}}{192\kappa ^2}𝑑x^{11}\sqrt{g}\overline{\psi }_I\mathrm{\Gamma }^{IJKLMN}\psi _NG_{JKLM},$$
(127)
becomes the gravitino mass term when compactified to four dimensions. Using the vevs of the $`G_{IJK11}`$ given by eq. (120), one can obtain the same result as eq. (126). This is a consistency check of our approach and the fact that the vacuum energy vanishes in the given approximation.
It follows from eq. (126), that the gravitino mass tends to zero when the radius of the eleventh dimension goes to infinity. When the four–dimensional Planck scale is fixed to be the measured value, however, the gravitino mass in the strongly coupled case is expressed in a standard manner, similar to the weakly coupled case as can be seen by inspecting (126) and (114). To obtain the gravitino mass of the order of 1 TeV, one has to adjust $`\mathrm{\Lambda }`$ to be of the order of $`10^{13}`$ GeV when one constructs a realistic model by appropriately breaking the $`E_8`$ gauge group at the hidden wall.
In the minimization of the potential we have implicitly used the leading order approximation. As was explained in a previous section, the next to leading order correction gives the non–trivial dependence of the background metric on $`x^{11}`$. Then the Einstein–Hilbert action in eleven dimensions gives additional contribution to the scalar potential in the four–dimensional effective theory, which shifts the vevs of the $`G_{IJKL}`$. As a consequence, $`F_S`$ will no longer vanish. Though this may be significant when we discuss soft masses, it does not drastically change our estimate of the gravitino mass (126) and our main conclusion drawn here is still valid after the higher order corrections are taken into account.
## 13 Soft supersymmetry breaking terms
In the previous section, we have shown that the gaugino condensation breaks supersymmetry both in the weakly coupled heterotic string and in the heterotic $`M`$–theory. We chose $`\mathrm{\Lambda }`$ in such a way that the gravitino mass appeared in the TeV–range. In this section we shall discuss the soft supersymmetry breaking terms that appear in the low–energy effective theory as a consequence of this nonzero gravitino mass.
We first give the relevant formulae for gaugino and scalar masses in the observable sector. Given the gauge kinetic function $`f_6`$ in the observable sector, the gaugino mass is calculated to be
$$m_{1/2}=\frac{f_6}{\varphi ^i}\frac{F^i}{2\text{Re}f_6},$$
(128)
where $`\varphi ^i`$ symbolically denote hidden sector fields responsible for supersymmetry breakdown. Writing the Kähler potential
$$K=\widehat{K}(\varphi ^i,\varphi _i^{})+Z(\varphi ^i,\varphi _i^{})C^{}C+(\text{higher orders in }C,C^{}),$$
(129)
one can also calculate the mass of a matter field $`C`$ $`^{\mathrm{?},\mathrm{?}}`$
$$m_0^2=m_{3/2}^2F^iF_j^{}\frac{Z_i^jZ_iZ^1Z^j}{Z}.$$
(130)
Here a vanishing cosmological constant is assumed.
Using the classical approximation naively, these formulae lead to a surprising result. All soft masses vanish. At the basis of this fact it had been suggested that the gravitino mass could be arbitrarily high, still leading to softly broken supersymmetry in the TeV range. It has been observed meanwhile that this surprising result is an artifact of the approximation and it is now commonly accepted that generically the soft masses tend to be of the order of the gravitino mass or at least not arbitrarily small compared to it. In general the result for the soft scalar masses is strongly model dependent. We shall see in the following that the situation concerning the gaugino mass is less model dependent but varies when we go from the weakly to the strongly coupled case $`^\mathrm{?}`$.
We start again with the weakly coupled case. At the leading order (tree level), the gauge kinetic function for the observable sector is simply $`f_6=S`$, whereas the gaugino condensation gives $`F_S=0`$, $`F_T=m_{3/2}(T+T^{})`$. Thus, at this level, the gaugino mass vanishes. As was discussed earlier in these lectures, the gauge kinetic function receives corrections at one–loop order. Using eq. (63), the gaugino mass is explicitly written as
$$m_{1/2}=\frac{F_S+ϵF_T}{2\text{Re}(S+ϵT)}.$$
(131)
Note that $`F_T/(T+T^{})m_{3/2}`$. Also we expect $`F_S`$ to be of the order of $`ϵTm_{3/2}`$ due to the one–loop corrections. Plugging them into the above expression, we obtain
$$m_{1/2}\frac{ϵT}{S}m_{3/2}.$$
(132)
Since in the weakly coupled case the ratio $`ϵT/S`$ is small, the gaugino becomes much lighter than the gravitino.
Let us now consider the scalar masses. At the tree level, the Kähler potential is
$$K=\mathrm{ln}(S+S^{})3\mathrm{ln}(T+T^{})+(T+T^{})^nC^{}C+(\text{higher orders in }C^{}C),$$
(133)
where $`n`$ denotes the modular weight of a field $`C`$. For a field with $`n=1`$ (untwisted sector in an orbifold construction), which naturally appears in the simple truncation procedure, we recover the previous formula (57). From eq. (130), it follows that
$$m_0^2=m_{3/2}^2+\frac{|F_T|^2}{(T+T^{})^2}=(1+n)m_{3/2}^2.$$
(134)
A scalar field with the modular weight $`1`$ has a vanishing supersymmetry breaking mass at the leading order. It is an artifact of the approximation of reduction and truncation (i.e. torus compactification) that the fields have modular weight $`1`$. A field whose modular weight is different from $`1`$ has a mass comparable to the gravitino mass. Though, as discussed in section 3, corrections at the one–loop level are model dependent, one expects they are of the order of $`ϵT/Sm_{3/2}^2`$. Summarizing these contributions, one obtains
$$m_0^2=(1+n)m_{3/2}^2+O(\frac{ϵT}{S}m_{3/2}^2),$$
(135)
where the actual value of the second term depends on the model one considers. A conclusion we can draw from eqs. (132) and (135) is that the gaugino masses tend to be much smaller than the scalar masses:
$$m_{1/2}<<m_0O(m_{3/2}).$$
(136)
Phenomenologically this relation might be problematic. Requiring that the gaugino masses are at the electro–weak scale, eq. (136) would then imply that the masses of the squarks and sleptons should be well above the 1 TeV region, which raises the fine–tuning problem to reproduce the Fermi scale. Another potential problem is the relic abundance of the lightest superparticles (LSPs) which are likely the lightest neutralinos in the present case. With the parameters characterized by (136), the standard computation of the relic abundances shows that too many LSPs would (if stable) still be around today, resulting in the overclosure of the Universe.
Thus in the weak coupling regime, one can conclude that, though the gaugino condensation realizes supersymmetry breaking, it tends to lead to a picture where gaugino masses are generically smaller than gravitino and scalar masses. A satisfactory situation might only be achieved, if one fine–tunes the scalar masses in a way that they become comparable to the gaugino masses.
Next we want to discuss how the situation changes when one considers the strongly coupled case (heterotic $`M`$–theory).
As in the weakly coupled heterotic string theory, the gaugino mass vanishes at the leading order of the $`\kappa ^{2/3}`$ expansions, because $`f_6=𝒮`$ and $`F_𝒮=0`$. Again the next to the leading order is important. The analogue of eq. (131) in the strongly coupled case is
$$m_{1/2}=\frac{F_𝒮+\alpha F_𝒯}{2\text{Re}(𝒮+\alpha 𝒯)}.$$
(137)
Thus we obtain, as before
$$m_{1/2}\frac{\alpha 𝒯}{𝒮}m_{3/2}.$$
(138)
A crucial difference in this case, however, is the fact that the ratio $`\alpha 𝒯/𝒮`$ is not a small number, but can be as large as unity. This is because the values of $`𝒮`$ and $`𝒯`$ inferred from our input variables (see section 2.2) suggests that we are rather close to criticality (in which case the ratio becomes unity). Thus we can conclude that, unlike the weakly coupled case, the gaugino mass in the strongly coupled regime is comparable to the gravitino mass. This observation confirms the expectation that the gravitino mass should be in the TeV–region and the gaugino condensation scale $`\mathrm{\Lambda }10^{13}`$ GeV. Because of the simplicity of the mass formula (128) and the fact that the gauge–kinetic function $`f`$ is stable in higher order perturbation theory, the statement concerning the soft gaugino masses is rather model independent.
The situation is more complicated in the case of the scalar masses which we consider now in the framework of heterotic $`M`$–theory. At the leading order we arrive at the same conclusions as in the weak coupling case, since the Kähler potential is identical in both cases. In section 4, we calculated the corrections to the Kähler potential at the next to leading order, which reads
$`\widehat{K}`$ $`=`$ $`\mathrm{ln}(𝒮+𝒮^{})3\mathrm{ln}(𝒯+𝒯^{})`$ (139)
$`Z`$ $`=`$ $`{\displaystyle \frac{6}{𝒯+𝒯^{}}}+{\displaystyle \frac{2\alpha }{𝒮+𝒮^{}}}`$ (140)
where the latter is valid for a field with the modular weight $`1`$. Now using the formula (130) one may be able to calculate the scalar masses, with the result
$`m_0^2=m_{3/2}^2{\displaystyle \frac{2\frac{1}{1+\delta }}{1+\delta }}{\displaystyle \frac{|F_𝒯|^2}{(𝒯+𝒯^{})^2}}{\displaystyle \frac{\delta (2\delta \frac{1}{1+\delta })}{1+\delta }}{\displaystyle \frac{|F_𝒮|^2}{(𝒮+𝒮^{})^2}}`$
$`{\displaystyle \frac{\delta }{(1+\delta )^2}}(F_𝒮F_𝒯^{}+F_𝒮^{}F_𝒯)`$ (141)
where
$$\delta \frac{\alpha }{3}\frac{𝒯+𝒯^{}}{𝒮+𝒮^{}}.$$
(142)
We can clearly see from this expression that the structure obtained in the leading order is badly violated. Given the fact that the expansion parameter $`\alpha (𝒯+𝒯^{})/(𝒮+𝒮^{})`$ is of order unity it is no longer possible to fine tune the scalar masses (by choosing modular weight $`1`$ for all of them) to a small value and then hope that the corrections respect this fine tuning. In addition the scalar masses depend strongly on the form of the Kähler potential which, in contrast to the gauge kinetic function, receives further corrections in higher order. Thus detailed statements about the scalar masses are very model dependent. It remains to be seen whether any sensible quantitative statement can be made about the scalar masses with the formulae given above. The results for the gaugino masses are more reliable since $`f`$ does not receive corrections in higher order.
In summary we can, however, conclude with the qualitative statement that in the strong coupling regime,
$$m_{1/2}m_0m_{3/2}.$$
(143)
This contrasts with the relation (136) for the weak coupling regime and represents an important improvement concerning phenomenological applications. In the strongly coupled case, the difference between dilaton– and moduli–dominated supersymmetry breakdown seems less pronounced than it is in the weakly coupled case.
## 14 Some phenomenological consequences
We have presented a consistent framework of supersymmetry breaking and soft breaking terms triggered by the gaugino condensate at the hidden wall. In the strongly coupled case, in complete analogy to the weakly coupled case, the gravitino mass $`m_{3/2}`$ is related to the gaugino condensation scale $`\mathrm{\Lambda }`$ as
$$m_{3/2}\frac{\mathrm{\Lambda }^3}{M_{Planck}^2}.$$
(144)
Furthermore, as explained in detail, the soft masses are of the order of the gravitino mass. This implies that these masses should be in the TeV range in order to solve the naturalness problem of the Higgs boson mass in the supersymmetric framework. This requires that $`\mathrm{\Lambda }`$ should be around $`10^{13}`$ GeV, three orders of magnitude smaller than the GUT scale (the compactification scale) and thus the 11D Planck scale as well. The gauge coupling constant at the $`E_8`$ wall, where the gaugino condensate is supposed to occur, is larger than the one at the $`E_6`$ wall. If the eleventh dimensional radius $`\rho `$ approaches the critical radius $`\rho _{crit}`$, the $`E_8`$ gauge coupling constant becomes strong at a scale as large as the GUT scale, and the running coupling constant will blow up at that scale already. Then the gaugino condensation scale $`\mathrm{\Lambda }`$, which is approximately identified with the blow–up energy scale, would become too large. For a value of $`\mathrm{\Lambda }10^{13}`$ GeV, $`\rho `$ should (although close) not be too close to the critical value so that the gauge coupling constant does not blow up immediately. This gives a constraint on the constant $`\alpha `$ (defined in (86)), which depends on the detailed properties of the Calabi–Yau manifold under consideration. In any case it is probably necessary to break the hidden $`E_8`$ to a smaller group to obtain a smaller coefficient of the $`\beta `$–function. These considerations should be kept in mind when one attempts to construct a realistic model.
The fact that the gravitino mass cannot be arbitrarily large, but should lie in the TeV range in the heterotic $`M`$–theory regime suggests that the theory might share a problem already encountered in the weakly coupled case $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. Late time decay of the gravitinos would upset the success of the standard big–bang nucleosynthesis scenario. This problem is rather universal in most of the supergravity models where breakdown of supersymmetry is mediated through gravity. Indeed this is not really a serious difficulty, but just implies that the universe underwent inflationary expansion followed by reheating at a relatively low temperature ($`T<10^9`$ GeV for $`m_{3/2}=1`$ TeV $`^\mathrm{?}`$), in which the gravitino number density is diluted by the inflation and the low reheat temperature suppresses gravitino production after that.
A main difference between the weakly and the strongly coupled case manifests itself when we consider phenomenological issues associated with the soft masses. In the weakly coupled string case, the gaugino condensation scenario gives a very small gaugino mass compared to the scalar masses. For a typical size of the compactification radius of the 6D manifold, the gaugino mass is shown to be more than one order of magnitude smaller than the scalar mass (see for example eqs. (7.20) and (7.24) (with $`\mathrm{sin}\theta 0`$ limit) of ref. $`^\mathrm{?}`$ for more detail). This hierarchy among the soft masses obviously raises a naturalness problem. With gaugino masses of the order of 100 GeV, the scalar masses would be far above 1 TeV, requiring fine tuning to obtain the electroweak symmetry breaking scale. This causes problems for explicit model building. Another phenomenological difficulty caused by the small gaugino mass arises in the context of relic abundances of the lightest superparticles (LSPs). Under the assumption of $`R`$–parity conservation, the LSP is stable and remains today as a dark matter candidate. Given the superparticle spectrum in the weak coupling regime, the bino, the superpartner of the $`U(1)_Y`$ gauge boson, is most likely to be the LSP. To evaluate the relic abundances of the bino, one has to know its annihilation cross section (see ref. $`^\mathrm{?}`$ and references therein). In our case, the bino pair annihilates into fermion (quarks and leptons) pairs via t–channel scalar (squarks and sleptons) exchange. The cross section is roughly proportional to
$$\sigma \frac{m_{\stackrel{~}{B}}^2}{m_{\stackrel{~}{f}}^4}$$
(145)
where $`m_{\stackrel{~}{B}}`$ is the bino mass and $`m_{\stackrel{~}{f}}`$ represents a scalar mass. As the scalar becomes heavier, the cross section is suppressed, yielding a larger relic abundance. Indeed when the scalar mass is more than an order of magnitude larger than the gaugino mass, a standard calculation shows that the relic abundance exceeds the critical value of the universe. This overclosure is a serious problem in the weakly coupled case.
In the strong coupling regime, the gaugino acquires a mass comparable to the gravitino mass and the scalar masses. Thus the above two problems do not appear. All the soft masses are in the same range. If this is not far from the electroweak scale, one can naturally realize the electroweak symmetry breaking at the correct scale without fine tuning. Moreover in this scenario, the annihilation cross section of the bino becomes larger, and thus we can obtain a relic abundance compatible with the observations. In some regions of parameter space we may even realize a situation where the LSP is the dominant component of the dark matter of the universe.
A characteristic of the mechanism of gaugino condensation is the fact that it is the $`T`$ field that plays the dominant role in the breakdown of supersymmetry. In this scenario scalar fields with different modular weight will have different masses, which may cause problems with flavor changing neutral currents (FCNC). In the strong coupling case, the situation may be improved through the presence of a large gaugino mass which contributes to the scalar masses at low energies through radiative corrections that can be computed via renormalization group methods. In a situation where scalar masses at the GUT scale are small enough, this universal radiative contribution might wash out nonuniversalities and avoid problems with FCNC. Details of the superparticle phenomenology in the strongly coupled case, including the issues outlined above, will be discussed elsewhere $`^\mathrm{?}`$.
Eqs. (63) (in the weak coupling case) and (97) (in the strong coupling case) show that the imaginary part of the complex scalar fields, $`S`$ and $`T`$, has an axion–like coupling to the gluon fields. In the weakly coupled case, world–sheet instanton effects $`^\mathrm{?}`$ and possibly other non–perturbative effects give non–negligible contributions to the potential. Then the axion candidates receive masses comparable to the gravitino mass, and they do not solve the strong $`CP`$ problem. However, in the strongly coupled case, it has been argued that these non–perturbative contributions originated at high energy physics might be suppressed to a negligible level $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. If this is the case, a linear combination of the Im$`𝒮`$ and Im$`𝒯`$ will play a role of the axion, whose potential is dominated by the QCD contribution.. Then this axion, referred to as the $`M`$–theory axion, will be able to solve the strong $`CP`$ problem. A word of caution should be added here, since a reliable calculation of these world sheet nonperturbative effects has only been performed in the weakly coupled case $`^\mathrm{?}`$. The above argumentation in the M–theory framework uses the implicit assumption that those Yukawa couplings remain as weak as in the case of the weakly coupled string, an assumption that might not be necessarily correct. Apart from that, the axion decay constant in this case becomes as as large as $`10^{16}`$ GeV, which leads to the potential problem that the energy density of the coherent oscillation of the axion field exceeds the critical energy density of the universe. This problem could be solved if the entropy production occurs after the QCD phase transition when the axion gets massive, or if this world is almost $`CP`$ conserving and the initial displacement of the axion field is very small. The direct detection of the relic axions with such a large decay constant would be extremely difficult. However the $`M`$–theory axion may give a significant contribution to the isocurvature density fluctuations during the inflationary epoch, which may be detectable in future satellite observations $`^\mathrm{?}`$. It remains to be seen whether this mechanism leads to a satisfactory solution of the strong CP–problem.
## 15 Summary and outlook
In any case we have seen that the M–theoretic version of the heterotic string shows some highly satisfactory phenomenological properties concerning the unification of fundamental coupling constants as well as the nature of the soft supersymmetry breaking parameters.
Still there remain some problems that still resist attempts for a satisfactory solution. Certainly one of them is the question of fixing the vev of the dilaton. One would like to see whether the M-theoretic approach to the problem might give us some new hints in that direction.
In the last years there has been revolutionary progress in the understanding of nonperturbative aspects of string theory. Here we have discussed the first consequences of phenomenological interest that could be derived from this new insights. Let us hope that other aspects of that field might also be of relevance for this questions and increase our understanding of the low-energy effective actions that could be derived from string theory.
## Acknowledgements
I would like to thank J. Conrad, Z. Lalak, A. Niemeyer, M. Olechowski, S. Stieberger, and M. Yamaguchi for useful discussions and collaboration. This work was partially supported by the European Commission programs ERB FMRX–CT96–0045 and CT96–0090.
## References
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# Strongly n-trivial Knots
## 1 Introduction
We start with a little background.
###### Definition 1.1.
A knot $`k`$ is called “$`(n`$-$`1)`$-trivial” if there exists a projection of $`k`$, such that one can choose $`n`$ disjoint sets of crossings of the projection with the property that making the crossing changes corresponding to any of the $`2^n1`$ nontrivial combination of the sets of crossings turns the original knot into the unknot. The collection of sets of crossing changes is said to be an “$`(n`$-$`1)`$-trivializer for $`k`$”.
###### Conjecture 1.2.
The unknot is the only knot that is $`n`$-trivial for all $`n`$.
Note: A knot that is $`n`$-trivial is automatically $`m`$-trivial for all $`mn`$.
Work of Gusarov \[Gu\] and Ng and Stanford \[NS\] shows that this question equates to showing that the only knot with vanishing Vassiliev invariants for all $`n`$ is the unknot. Thus, Conjecture 1.2 is at the heart of the study of Vassiliev invariants.
One reason why this question is interesting is that it takes a geometric approach to Vassiliev invariants, instead of the traditional algebraic approach and therefore is relatively unexplored. Vassiliev invariants measure geometric properties of knots, which in turn are geometric objects, so it is reasonable to hope that the geometry might play an integral role in their study.
The following definition derives its motivation from $`n`$-triviality.
###### Definition 1.3.
A knot $`k`$ is called “strongly $`(n`$-$`1)`$-trivial.” if there exists a projection of $`k`$, such that one can choose $`n`$ crossings of the projection with the property that making the crossing changes corresponding to any of the $`2^n1`$ nontrivial combination of the selected crossings turns the original knot into the unknot. The collection of crossing changes is said to be a “strong $`(n`$-$`1)`$-trivializer for $`k`$”.
Note: The expression “$`n`$ adjacent to the unknot” is used interchangeably with “strongly $`(n`$-$`1)`$-trivial.” We will stick with the latter throughout this paper.
In Section 6 we show that for any $`n`$ there is a non-trivial knot that is strongly $`n`$-trivial. On the other hand in Section 5 we prove the main result of this paper:
###### Theorem 1.4.
Any non-trivial knot $`k`$ of genus $`g`$ fails to be strongly $`n`$-trivial for all $`n`$, $`n3g1`$.
Note: A knot that is strongly $`n`$-trivial is automatically strongly $`m`$-trivial for all $`mn`$. Also any knot that is strongly $`n`$-trivial is clearly $`n`$-trivial, too.
In analogy with Conjecture 1.2 we have
###### Corollary 1.5.
The unknot is the only knot that is strongly $`n`$-trivial for all $`n`$.
Theorem 1.4 is proven by repeated use of the following theorem of Gabai
###### Theorem 1.6.
(Corollary 2.4 \[G\]) Let $`M`$ be a Haken manifold whose boundary is a nonempty union of tori. Let $`F`$ be a Thurston norm minimizing surface representing an element of $`H_2(M,M)`$ and let $`P`$ be a component of $`M`$ such that $`PF=\mathrm{}`$. Then with at most one exception (up to isotopy) $`F`$ remains norm minimizing in each manifold $`M(\alpha )`$ obtained by filling $`M`$ along an essential simple closed curve $`\alpha `$ in $`P`$ In particular $`F`$ remains incompressible in all but at most one manifold obtained by filling $`P`$.
## 2 Notation
Let $`k`$ be a knot that is strongly $`(n`$-$`1)`$-trivial. Let $`p:kR^2`$ be a projection with crossings $`\{a_1,\mathrm{}a_n\}`$ demonstrating the strong $`(n`$-$`1)`$-triviality. For each $`a_i`$ let $`c_i`$ be the small vertical circle that bounds a disk $`D_i`$ that intersects $`k`$ geometrically twice, but algebraically 0 times. We call the $`c_i`$ linking circles of $`k`$ and call $`D_i`$ a crossing disk after \[ST\]. Let $`M`$ be the link exterior of $`kc_1\mathrm{}c_n`$ and $`P_i`$ be the torus boundary component in $`M`$ corresponding to $`c_i`$. Either $`+1`$ or $`1`$ filling of $`P_i`$ will result in the desired crossing change depending on orientation. We adopt the convention that each $`P_i`$ will be oriented so that $`+1`$ filling of $`P_i`$ corresponds to the appropriate crossing change dictated by $`a_i`$.
## 3 Irreducibility
###### Lemma 3.1.
Let $`k`$ be a nontrivial knot. Let $`\{c_1,\mathrm{}c_n\}`$ be linking circles for $`k`$ that show $`k`$ is strongly $`(n`$-$`1)`$-trivial, then $`M`$, the exterior of $`kc_1\mathrm{}c_n`$, is irreducible and $`M`$ is therefore Haken.
###### Proof.
Assume $`M`$ is reducible. Let $`S`$ be a sphere that does not bound a ball on either side. $`S`$ cannot be disjoint from $`D_1\mathrm{}D_n`$ or else it would bound a ball on the side that does not contain $`k`$. Assume $`S`$ intersects $`D_1\mathrm{}D_n`$ minimally and transversally. The intersection will consist of a union of circles. Let $`r`$ be one of these circles that is innermost on $`S`$ (any circle that bounds a disk on $`S`$ disjoint from all the other circles of intersection). Without loss of generality assume $`rD_1`$. $`r`$ cannot be trivial on $`D_1(D_1k)`$ since $`SD_1`$ is minimal. $`r`$, however must be trivial in $`M`$ so must divide $`D_1`$ into two pieces, one containing both points of $`D_1k`$ and the other consisting of an annulus running from $`r`$ to $`c_1`$. This disk on $`S`$ bounded by $`r`$ shows that $`c_1`$ bounds a disk in the exterior of $`k`$. This, however, means that $`+1`$ surgery on $`c_1`$ leaves $`k`$ unchanged instead of turning it into an unknot, yielding the desired contradiction. ∎
## 4 Minimal genus Seifert surfaces
This section is dedicated to proving the following theorem.
###### Theorem 4.1.
If $`k`$ has a strong $`(n`$-$`1)`$-trivializer $`\{c_1,\mathrm{}c_n\}`$ and $`F`$ is a Seifert surface for $`k`$ disjoint from $`\{c_1,\mathrm{}c_n\}`$ which is minimal genus among all such surfaces, then $`F`$ is a minimal genus Seifert surface for $`k`$.
###### Proof.
Because the $`c_i`$ have linking number 0 with $`k`$ we can find a Seifert surface for $`k`$ disjoint from the $`c_i`$. Let $`F`$ be a minimal genus Seifert surface for $`k`$ in the link complement.
We supplement the notation introduced in Section 2. Recall $`M`$ is the link exterior of $`kc_1\mathrm{}c_n`$. Let $`L`$ be the corresponding link of $`n+1`$ components in $`S^3`$. $`P_i`$ is the torus boundary component in $`M`$ corresponding to $`c_i`$. Let $`M(\alpha )`$ refer to the manifold obtained by filling $`M`$ along an essential simple closed curve of slope $`\alpha `$ in $`P_n`$. When $`\alpha =1/m,mZ`$, $`M(\alpha )`$ is a link exterior. Let $`L_\alpha `$ be the corresponding link in $`S^3`$. Let $`k_\alpha `$ be the image of $`k`$ in $`L_\alpha `$ and $`F_\alpha `$ be the image of $`F`$ in $`L_\alpha `$.
We now prove Theorem 4.1 by induction on $`n`$. If $`F`$ is ever a disk then Theorem 4.1 is clearly true, so we will assume that $`F`$ is not a disk throughout the proof.
The base case: Let $`k`$ be a strongly 0-trivial knot. This means that $`k`$ is unknotting number 1 and there is one linking circle $`c_1`$ that dictates a crossing change that unknots $`k`$.
By Lemma 3.1 if $`M`$ is reducible, then $`k`$ is the unknot. As in the proof of Lemma 3.1 $`c_1`$ bounds a disk in the complement of $`k`$, so $`kc_1`$ is the unlink on two components. Therefore, $`F`$ being least genus must be a disk, which is a contradiction, verifying the claim for $`M`$ reducible and $`n=1`$. We may assume $`M`$ is irreducible to complete the base case. $`k_1`$ is the unknot. Since $`F_1`$ is not a disk, it is no longer norm minimizing after the filling. Thus by Theorem 1.6 $`F`$ is norm minimizing under any other filling of $`P_1`$. In particular $`F_{\mathrm{}}`$ is Thurston norm minimizing for $`L_{\mathrm{}}`$, which is just $`k`$. Thus, $`F`$ is a least genus Seifert surface for $`k`$.
The inductive step: Now we assume that if $`k`$ has a strong $`(n`$-$`2)`$-trivializer $`\{c_1,\mathrm{}c_{(n1)}\}`$ and $`F`$ is a Seifert surface for $`k`$ disjoint from $`\{c_1,\mathrm{}c_{n1}\}`$, which is minimal genus among all such surfaces, then $`F`$ is also a minimal genus Seifert surface for $`k`$ and show that the same must be true for any strong $`(n`$-$`1)`$-trivializer for $`k`$.
Again by Lemma 3.1 if $`M`$ is reducible, $`k`$ must be the unknot. As in previous arguments, the separating sphere $`S`$ must intersect at least one $`D_i`$ in a curve that is essential on $`D_ik`$. Without loss of generality, we may assume that $`D_n`$ is such a disk. Then $`c_n`$ bounds a disk in the complement of $`k\{c_1,\mathrm{}c_{n2}\}`$. Since $`\{c_1,\mathrm{}c_{n1}\}`$ forms a strong $`(n`$-$`2)`$-trivializer for $`k`$, the induction assumption implies $`k`$ bounds a disk $`\mathrm{\Delta }`$ disjoint from $`c_1\mathrm{}c_{n1}`$. Since $`c_n`$ bounds a disk disjoint from $`kc_1\mathrm{}c_{n1}`$, $`\mathrm{\Delta }`$ can clearly be chosen to be disjoint from $`c_n`$, too, but this contradicts the assumption that $`F`$ was minimal genus, but not a disk.
We now may finish the proof of Theorem 4.1 knowing that $`M`$ is irreducible. $`k_1`$ is an unknot in the link $`L_1`$. $`\{c_1,\mathrm{}c_{n1}\}`$ is a strong $`(n`$-$`2)`$-trivializer for $`k`$ in $`L_1`$. The inductive assumption means that $`k_1`$ bounds a disk in the exterior of $`L_1`$. This disk is in the same class as $`F_1`$ in $`H_2(M(1),M(1))`$, thereby showing that $`F_1`$ is not Thurston norm minimizing. Thus, by Theorem 1.6 $`F`$ remains norm minimizing under any other filling of $`P_n`$. In particular $`F_{\mathrm{}}`$ is Thurston norm minimizing in $`L_{\mathrm{}}`$. Thus, $`F`$ is a least genus Seifert surface for $`k`$ in the complement of $`\{c_1,\mathrm{}c_{n1}\}`$. $`\{c_1,\mathrm{}c_{n1}\}`$, however, forms a strong $`(n`$-$`2)`$-trivializer for $`k`$ in $`L_{\mathrm{}}`$. By the inductive assumption, $`F`$ must be Thurston norm minimizing for $`k`$ in the knot complement as well as the link complement. ∎
## 5 Arcs on a Seifert surface
We now prove Theorem 1.4: Any non-trivial knot $`k`$ of genus $`g`$ fails to be strongly $`n`$-trivial for all $`n`$, $`n3g1`$.
###### Proof.
Let $`k`$ be strongly $`n`$-trivial with $`n`$-trivializers $`\{c_1,\mathrm{}c_{n+1}\}`$. Let $`F`$ be a minimal genus Seifert surface for $`k`$ disjoint from $`\{c_1,\mathrm{}c_{n+1}\}`$ as in Theorem 4.1. $`F`$ has genus $`g`$.
Each linking circle $`c_i`$ bounds a disk $`D_i`$ that intersects $`F`$ in an arc running between the two points of $`kD_i`$ and perhaps some simple closed curves. Simple closed curves inessential in $`D_ik`$ can be eliminated since $`F`$ is incompressible. Any essential curves $`s_j`$ must be parallel to $`c_i`$ in $`D_ik`$. These curves can be removed one at a time using the annulus running from $`c_i`$ to the outermost $`s_j`$ to reroute $`F`$, decreasing the number of intersections. Thus, if $`F`$ is assumed to have minimal intersection with each of the $`D_i`$ then it intersects each one in an arc which we shall call $`a_i`$ as in Figure 1. Each $`a_i`$ is essential in $`F`$. Otherwise $`c_i`$ would bound a disk disjoint from $`F`$ and the crossing change along $`c_i`$ would fail to unknot $`k`$.
###### Lemma 5.1.
$`a_i`$ is never parallel on $`F`$ to $`a_j`$ for $`ij`$.
###### Proof.
If $`a_i`$ is parallel on $`F`$ to $`a_j`$ there must be an annulus running from $`P_i`$ to $`P_j`$ in the link exterior. Recall that we adopted the convention that $`P_i`$ and $`P_j`$ are each oriented so that $`+1`$ surgery results in the appropriate crossing changes. The two tori cannot have opposite orientations or else $`+1`$ fillings on both $`P_i`$ and $`P_j`$ is the same as $`\mathrm{}`$ fillings on both, thus, instead of unknotting $`k`$ changing both crossings leaves $`k`$ knotted. If the two tori have the same orientation we could replace $`P_iP_j`$ with a single torus $`T`$ obtained by cutting and pasting of the two tori along the annulus. Now $`+1`$ filling for $`P_i`$ and $`\mathrm{}`$ filling for $`P_j`$ is equivalent to $`+1`$ filling on $`T`$, while $`+1`$ filling on both $`P_i`$ and $`P_j`$ is equivalent to $`\frac{1}{2}`$ filling on $`T`$. This implies that $`F`$ fails to be norm-minimizing after both $`+1`$ and $`\frac{1}{2}`$ filling of $`T`$. This contradicts Theorem 1.6 completing the proof of the Lemma. ∎
Then $`\{a_1,\mathrm{}a_n\}`$ is a collection of embedded arcs on $`F`$, no two of which are parallel. An Euler characteristic argument shows that $`m3g1`$. Since the arcs are in one to one correspondence with the linking circles, a strong $`n`$-trivializer can produce at most $`3g1`$ linking circles for $`k`$ finally proving Theorem 1.4. ∎
We note that Theorem 1.4 predicts that a genus one knot can be at most strongly $`1`$-trivial. Given standard projections of the trefoil and the figure eight knot it is easy to find a pair of crossing changes that show the knots are strongly $`1`$-trivial. The theorem is therefore sharp at least for genus one knots. It is possible, but unlikely, that the theorem remains as precise for higher genus knots.
Finally as noted in the introduction, Theorem 1.4 implies Corollary 1.5: The unknot is the only knot that is strongly $`n`$-trivial for all $`n`$.
## 6 Constructing strongly $`n`$-trivial knots
One might think that there exists a bound $`n`$ such that no nontrivial knot is strongly $`n`$-trivial. Given any $`n`$, this section gives one way to construct strongly $`n`$-trivial knots.
In Figure 3 we will give projections of graphs that show how to turn an unknot into a strongly $`n`$-trivial knot. The circle running around the outside of the graph should be viewed as an unknot $`k^{}`$. Each arc $`a_i`$ suggests a linking circle $`c_i`$ and a crossing disk $`D_i`$. If we alter the link $`kc_1\mathrm{}c_n`$ in $`S^3`$ by twisting -1 times along each of the disks $`D_1,D_2,\mathrm{}D_n`$ $`k^{}`$ becomes a new knot $`k`$ see Figure 2. The linking circles remain fixed, so we get a new link in $`S^3`$, $`kc_1\mathrm{}c_n`$.
Figure 3 gives graphs that generate examples of strongly $`1`$-trivial and strongly $`2`$-trivial knots. Note that the figure on the right is obtained from the figure on the left by replacing one arc by two new arcs that follow along the original arc, clasp, return along the original arc, and then, close to the boundary, clasp once again. This process could be iterated indefinitely by choosing an arc of the new graph and repeating the construction. It is modeled on doubling one component of a link. Given a Brunnian link of $`n`$ components (a nontrivial link for which any $`n1`$ components is the unlink), doubling one of the components yields a Brunnian link of $`n+1`$ components. The graph on the left in Figure 3 has a Brunnian link of 2 components as a subgraph and the one on the right has the double of that link as a subgraph. Let $`\mathrm{\Gamma }_n`$ be the graph after $`n2`$ iterations ($`n2`$).
###### Theorem 6.1.
$`\mathrm{\Gamma }_n`$ contains a Brunnian link, $`l_{n+2}`$, of $`n+2`$ components and yields $`k`$ a non-trivial, strongly $`(n+1)`$-trivial knot.
###### Proof.
The link consists of the arcs $`\{a_1,\mathrm{}a_{n+2}\}`$, together with short segments from $`k^{}`$ connecting the end points of the segments (and disjoint from the end points of the other segments). The base case is trivial because, $`\mathrm{\Gamma }_0`$ contains a Brunnian link of $`2`$ components: the Hopf link. $`\mathrm{\Gamma }_n`$ is obtained from $`\mathrm{\Gamma }_{n1}`$ by doubling one of the components of a Brunnian link of $`n+1`$ components. This yields a Brunnian link of $`n+2`$ components.
As a result of the Brunnian structure in $`\mathrm{\Gamma }_n`$ any $`n+1`$ edges from $`\{a_1,\mathrm{}a_{n+2}\}`$ can be disjointly embedded on a disk bounded by $`k^{}`$. So $`k^{}`$ forms an unlink with any proper subset of $`\{c_1,\mathrm{}c_{n+2}\}`$.
We can use this fact to show that $`\{c_1,\mathrm{}c_{n+2}\}`$ are an $`n`$-trivializer for $`k`$. Let $`S`$ be any nontrivial subset of $`\{c_1,\mathrm{}c_{n+2}\}`$. Let $`S^c`$ be the complement of $`S`$. If we take $`k`$ together with $`S`$, and do $`+1`$ surgery on each component of $`S`$ the resulting knot is an unknot. This is because it is exactly the same as if we took $`k^{}`$ and did $`1`$ surgery on each of the components of $`S^c`$. Since $`S`$ is a nontrivial subset, $`S^c`$ is a proper subset. $`k^{}`$ together with a the linking circles in $`S^c`$, therefore form an unlink, so each of the components of $`S^c`$ bounds a disk disjoint from $`k^{}`$ and doing $`1`$ surgery on these linking circles leave $`k^{}`$ unchanged.
Now that we know that $`k`$ is strongly $`(n+1)`$-trivial, we need only show $`k`$ is a non-trivial knot. Assume $`k`$ is trivial. By Theorem 4.1, $`k`$ bounds a disk $`\mathrm{\Delta }`$ in the complement of $`c_1\mathrm{}c_{n+2}`$. Since $`kc_1\mathrm{}c_{n+2}`$ was obtained from $`k^{}c_1\mathrm{}c_{n+2}`$ by spinning along the $`D_i`$’s, the exteriors of the two links are homeomorphic, and therefore $`k^{}`$ must bound a disk $`\mathrm{\Delta }^{}`$ also disjoint from $`c_1\mathrm{}c_{n+2}`$ (note that one could even prove that both $`kc_1\mathrm{}c_{n+2}`$ and $`k^{}c_1\mathrm{}c_{n+2}`$ are unlinks). $`k^{}`$ intersects each $`D_i`$ in 2 points, so as before we may assume $`\mathrm{\Delta }^{}D_i`$ is an arc for each $`i`$, but these arcs must, of course, be isotopic to the $`a_is`$ which in turn shows that the $`a_is`$ can be disjointly embedded on $`\mathrm{\Delta }^{}`$, proving that $`l_{n+2}`$ is planar and not Brunnian, the desired contradiction. Thus, $`k`$ is a strongly $`(n+1)`$-trivial knot, but not the unknot.
## 7 References
\[G\] Gabai, David, Foliations and the topology of 3-manifolds. II, Journal of Differential Geometry 26 (1987), pp. 461–478.
\[Gu\] Gusarov, M., On $`n`$-equivalence of knots and invariants of finite degree. Topology of manifolds and varieties, pp. 173–192, Adv. Soviet Math., 18, Amer. Math. Soc., Providence, RI, 1994.
\[NS\] Ng, Ka Yi; Stanford, Ted, On Gusarov’s groups of knots, to appear in Math Proc Camb Phil.
\[S\] Scharlemann, Martin, Unknotting number one knots are prime. Invent. Math. 82 (1985), no. 1, 37–55.
\[ST\] Scharlemann, Martin; Thompson, Abigail, Link genus and the Conway moves. Comment. Math. Helv. 64 (1989), no. 4, 527–535.
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# Interplay between Zeeman Coupling and Swap Action in Spin-based Quantum Computer Models: Error Correction in Inhomogeneous Magnetic Fields
\[
## Abstract
We consider theoretically the interplay between Zeeman coupling and exchange-induced swap action in spin-based quantum dot quantum computer models in the presence of inhomogeneous magnetic fields, which are invariably present in real systems. We estimate quantitatively swap errors caused by the inhomogeneous field, establishing that error correction would, in principle, be possible in the presence of non-uniform magnetic fields in realistic structures.
\]
In recent years quantum computing has attracted widespread attention. The computers based on the principles of quantum mechanics, such as quantum parallelism and entanglement, promise to deliver results much faster than classical computers in certain tasks such as factoring and searching . The impetus for constructing real quantum computer architectures arose from the seminal results establishing that quantum error correction is theoretically possible and therefore decoherence is not an insurmountable barrier as was assumed earlier. There have been many proposals for quantum computer (QC) architectures based on various physical two-level systems, such as those using laser-cooled trapped ions , photons or atoms trapped in cavities , nuclear spins in bulk solutions and crystal lattices , spins of electrons trapped in quantum dots or donors , spins of donor nuclei in Si , superconducting devices , and others. The minimal requirements for a QC architecture are the existence of fundamental quantum bits (qubit) and the ability to carry out single- and two-qubit operations such as swap and controlled-NOT, as well as suitable hardwares for reading (input) and writing (output or measurement) operations. Among the many obstacles facing the successful demonstration of nontrivial quantum computation in specific QC hardwares, the most daunting are the problem of quantum decoherence and the difficulty in achieving precise control over the various unitary operations necessary for quantum computation. In this Letter we consider the theoretical issue of controlling the swap operation in the proposed solid state QC architecture involving quantum dot spin entanglement, with externally-controlled electrostatic gates and magnetic fields providing the unitary operation.
Initial work on QC hardwares concentrated on atomic/molecular systems partly because well-defined single-qubit operations are comparatively easier to control in atomic systems such as trapped ions . Successful QC architectures will eventually require many qubits working in parallel, which would be difficult to achieve in atomic systems. Proposed solid state QC architectures are, in principle, scalable to many qubits, although no one has yet successfully demonstrated a controlled single-qubit operation in solid state QC systems. This is a curious dichotomy in the current QC hardware research—the architectures with demonstrated single-qubit operations are difficult to scale up while the presumably scalable solid state architectures have not yet been able to demonstrate single-qubit control. It is thus crucial to explore the challenges facing coherent control of qubits in solid state structures, particularly the issue of possible error corrections in realistic systems. Among various microscopic degrees of freedom that have been considered for the role of qubits in solid state QC architectures, spins of electrons or nuclei are natural candidates because of their well-defined Hilbert spaces and their relatively long decoherence time compared to the orbital degrees of freedom. In the proposed spin-based QC architectures, exchange interaction plays a fundamental role of establishing two-qubit entanglement , while Zeeman coupling to an external magnetic field provides various single-qubit operations. Generally Zeeman coupling is treated separately from the exchange interaction because the latter originates from the Coulomb interaction (due to Pauli principle) while the former is purely a spin effect, and the two interaction terms in the Hamiltonian commute with each other if the Zeeman term is homogeneous. However, solid state heterostructures are intrinsically inhomogeneous and magnetic imperfections or impurities are likely to be present, leading to inhomogeneous stray magnetic fields. Furthermore, parallel pulse schemes, in which exchange interaction and inhomogeneous Zeeman coupling are present simultaneously, have been proposed to expedite operations of spin-based QCs . Therefore, it is necessary to explore the interplay between exchange interaction and Zeeman coupling in relation to spin-based QC architectures, as we do here.
In this Letter we focus on the effects of Zeeman coupling when the external magnetic field is spatially inhomogeneous. Zeeman coupling is generally neglected in these studies based partly on the assumption that the effects are small and do not lead to any qualitative changes. While it is true in uniform fields, this assumption may not hold in an inhomogeneous magnetic field. The key motivations for our study are the observations that spatially inhomogeneous fields are intrinsic to spin-based solid state QC architectures and that for parallel schemes the magnetic fields at the locations of the two spins are generally different , which means the total field must be inhomogeneous. If only the single-spin evolution is involved, effects of local magnetic field can be corrected in such single-qubit operations by techniques such as spin echoes . However, if two-spin entangled evolution is also involved, it is not clear whether errors caused by inhomogeneous fields can still be eliminated, i.e. whether such errors lie within the current QC error correction constraints.
As an example we study how Zeeman coupling may affect the proposed operations of spin-based quantum dot quantum computers (QDQC). Our results, however, are quite general and can be applied to other spin-based models with finite magnetic fields. In particular, we study the effect of a finite inhomogeneous magnetic field on swap actions. For most of the spin-based schemes, swap action $`U_{sw}`$, in which two spins exchange their states, is one of the most basic operations. It is used to construct conditional phase shifts (CPS), $`U_{CPS}=e^{i\frac{\pi }{4}\sigma _{1z}}e^{i\frac{\pi }{4}\sigma _{2z}}U_{sw}^{\frac{1}{2}}e^{i\frac{\pi }{2}\sigma _{1z}}U_{sw}^{\frac{1}{2}}`$ (with $`\sigma `$ being the Pauli matrices for spins), which can then be converted to controlled-NOT (CNOT) easily . In addition, swaps are used to move spin states around so that an arbitrary pair of spins can be brought into controlled entanglement, which is essential to quantum computation . It is therefore important to investigate how Zeeman splitting will affect the swap gate and what are the consequences if these effects are non-trivial.
In a single envelope function approach , the Hamiltonian for two electrons trapped in a lateral double-quantum-dot is
$`H`$ $`=`$ $`{\displaystyle \underset{i}{\overset{2}{}}}[{\displaystyle \frac{1}{2m^{}}}(𝐩+{\displaystyle \frac{e}{c}}𝐀(𝐫_i))^2+V(𝐫_i)`$ (2)
$`+g^{}\mu _B𝐁(𝐫_i)𝐒_i]+{\displaystyle \frac{e^2}{ϵr_{12}}}.`$
Notice that here the magnetic field has a spatial dependence. The Hamiltonian can be expanded in a basis of two-spin eigenstates (spin singlet and triplet states). If the magnetic field is uniform, the Zeeman coupling depends only on the total spin along the field direction and commutes with the total spin operator $`𝐒^2`$, so that singlet and triplet states remain eigenstates of the two-electron system. On the other hand, if the magnetic field is inhomogeneous, it destroys the symmetry between the two spins, and the singlet state and one of the triplet states (the one with $`S_𝐧=0`$, where $`𝐧`$ is the field direction) couple. For example the coupling between the lowest singlet and triplet states (in the Heitler-London approximation) is $`T|V_Z|S=g\mu _B/2\times \{L|B_z(𝐫)|LR|B_z(𝐫)|R+L|B_z(𝐫)|RR|L`$$`R|B_z(𝐫)|LL|R\}`$, where the singlet state is $`|S=[|L(1)|R(2)+|L(2)|R(2)](||)/2`$, the triplet state is $`|T=[|L(1)|R(2)|L(2)|R(2)](|`$$`+|)/2`$, and the Zeeman coupling is $`V_Z=g\mu _B\left[B_z(𝐫_1)S_{1z}+B_z(𝐫_2)S_{2z}\right]`$. Here we have assumed that the external field is along the z direction. $`|L`$ and $`|R`$ are the left and right quantum dot ground states. Therefore, the inhomogeneous magnetic field couples the singlet and the $`S_z=0`$ triplet states so that two of the eigenstates of the system are no longer the eigenstates of the total spin. When the overlap between the two ground state wavefunctions is small, the main contribution to this coupling comes from the average field difference between the two quantum dots. Below we explore the consequences of this loss of symmetry.
When the overlap of the electronic wavefunctions is small, we can assume that the electron orbital degrees of freedom are frozen, so that an effective spin Hamiltonian quite faithfully describes the two-spin system :
$`H_s`$ $`=`$ $`J𝐒_1𝐒_2+\gamma _1S_{1z}+\gamma _2S_{2z},`$ (3)
where $`\gamma _1=g\mu _BL|B_z(𝐫)|L`$ and $`\gamma _2=g\mu _BR|B_z(𝐫)|R`$ are local Zeeman couplings due to applied or stray magnetic fields. Here we have implicitly assumed, based on the small interdot wavefunction overlap, that the two spins are distinguishable, with spin 1 on the left dot and spin 2 on the right dot. We have also assumed that the field is entirely along the z direction. Whether such a choice is reasonable will be discussed later.
Hamiltonian (3) can be expressed in the basis of four two-spin states $`|`$, $`|`$, $`|`$, and $`|`$, and its eigenstates can be easily obtained. The two polarized states are decoupled from the other two, which are mixtures of singlet and $`S_z=0`$ triplet states: $`|\psi _1=|`$, $`|\psi _2=|`$, $`|\psi _3=c_1|+c_2|`$, and $`|\psi _4=c_2|c_1|`$. Here the coefficients $`c_1`$ and $`c_2`$ satisfy the relations $`c_2/c_1=\sqrt{1+(\delta /2J)^2}\delta /2J`$ and $`c_1^2+c_2^2=1`$, where $`\delta =\gamma _1\gamma _2`$ represents the field inhomogeneity. The energies of the latter two states are also shifted from those of singlet and triplet states (with $`\stackrel{~}{\delta }`$ being the shift): $`E_1=J+\mathrm{\Delta }`$, $`E_2=J\mathrm{\Delta }`$, $`E_3=J+\sqrt{4J^2+\delta ^2}=J+\stackrel{~}{\delta }`$, and $`E_4=J\sqrt{4J^2+\delta ^2}=3J\stackrel{~}{\delta }`$, where $`\mathrm{\Delta }=\gamma _1+\gamma _2`$ represents the average magnetic field. An important question here is whether these mixtures and shifts will cause any error in quantum computation in the schemes based on the exchange interaction. After all, the swap action in these models depends on the perfect phase matching in the evolution of singlet and triplet states, as we will show below. Since swap operation is an essential component of a spin-based QDQC and several other architectures, we need to precisely quantify the effects of mixtures in singlet and triplet states on the swap action.
To determine whether swap is affected, we explore whether a product state of spins 1 and 2 will evolve into a product state (pure states for both spins) again. Our strategy here is to calculate the Schmidt number of the two-spin state. We will evolve our two-spin state, calculate its density matrix, trace out the second spin so that the density matrix of the first spin is left. We can then look for the eigenvalues of this density matrix. If at sometime it has only one non-vanishing eigenvalue, the state of the first spin is pure. We can then find out whether this pure state corresponds to a swapped state.
Our initial state is a product state given by
$$|\varphi (0)=(\alpha _1|+\alpha _2|)(\beta _1|+\beta _2|).$$
(4)
If the two electrons are located in two well-separated quantum dots in the beginning, the above product state does not violate the antisymmetry requirement of a two-fermion state. This state can be expanded in the basis of the eigenstates of Hamiltonian (3). It then evolves under Hamiltonian (3). The two-spin state at time $`t`$ takes the form $`|\varphi (t)=\alpha _1\beta _1e^{iE_1t/\mathrm{}}|+\alpha _2\beta _2e^{iE_2t/\mathrm{}}|+(\alpha _1\beta _2c_1+\alpha _2\beta _1c_2)e^{iE_3t/\mathrm{}}|+(\alpha _1\beta _2c_2\alpha _2\beta _1c_1)e^{iE_4t/\mathrm{}}|`$. The corresponding density matrix for the first spin can then be calculated straightforwardly: $`\rho _1=\mathrm{Tr}_2\{\rho _{12}\}=\mathrm{Tr}_2\{|\varphi (t)\varphi (t)|\}`$. The matrix elements of $`\rho _1`$ are all functions of the constants $`\alpha _1`$, $`\alpha _2`$, $`\beta _1`$, $`\beta _2`$, $`c_1`$, $`c_2`$, field inhomogeneity $`\delta `$, exchange coupling $`J`$, and time $`t`$ ($`\rho _1`$ is also a function of the average field $`\mathrm{\Delta }`$). The eigenvalue equation for $`\rho _1`$ is
$$\lambda ^2(\rho _1+\rho _1)\lambda +(\rho _1\rho _1|\rho _1|^2)=0.$$
(5)
To have a pure state for the first spin, which means that only one eigenvalue of the density matrix $`\rho _1`$ is non-vanishing, the last term on the left-hand-side of the above equation has to vanish:
$$\rho _1\rho _1|\rho _1|^2=0.$$
(6)
The explicit expression for Eq. (6) is quite complicated, so let us first look at the simple uniform field situation, when $`\delta =0`$ and $`c_1=c_2=1/\sqrt{2}`$. Equation (6) then becomes $`|\alpha _1\beta _1\alpha _2\beta _2\frac{1}{4}[(\alpha _1\beta _2+\alpha _2\beta _1)+(\alpha _1\beta _2\alpha _2\beta _1)e^{i\theta }][(\alpha _1\beta _2+\alpha _2\beta _1)(\alpha _1\beta _2\alpha _2\beta _1)e^{i\theta }]|^2=0`$, where $`\theta =4Jt`$. The solutions here are $`e^{i\theta }=\pm 1`$. When $`e^{i\theta }=1`$, $`|\varphi _1(t)=\alpha _1|+\alpha _2e^{i\mathrm{\Delta }t}|`$, the state of the first spin returns to its initial state with a phase shift between the two coefficients. When $`e^{i\theta }=1`$, $`|\varphi _1(t)=\beta _1e^{i\mathrm{\Delta }t}|+\beta _2|`$, the swap is achieved with the exception of an additional phase that can be corrected easily with a single-spin operation. For example, a pulse sequence can be constructed based on the phase-shifted swap $`U_{psw}`$ to produce a CNOT gate
$`U_{CNOT}`$ $`=`$ $`e^{i\frac{\pi }{4}\sigma _{2y}}e^{i(\frac{\pi }{4}+\frac{\stackrel{~}{\delta }t}{2})\sigma _{1z}}e^{i(\frac{\pi }{4}\frac{\stackrel{~}{\delta }t}{2})\sigma _{2z}}`$ (8)
$`\times U_{psw}^{\frac{1}{2}}e^{i\frac{\pi }{2}\sigma _{1z}}U_{psw}^{\frac{1}{2}}e^{i\frac{\pi }{4}\sigma _{2y}}.`$
Physically, a uniform field means that the Zeeman coupling couples to the total spin (including both electron spins), so that the Zeeman term commutes with the exchange term in the Hamiltonian (3), and therefore does not change the eigenstates. The shifts in the energy levels of the polarized states cause additional phase shift, but can be corrected by applying an opposite magnetic field with the same pulse shape, magnitude, and length. In summary, an external uniform magnetic field does not qualitatively change the proposed QC algorithm. Logistically it makes the QC operation more difficult because of the necessary correction pulses.
If the magnetic field is inhomogeneous, $`\delta 0`$. Let us look at a simple situation when the initial state is $`|\varphi (0)=|`$, i.e. $`\alpha _1=\beta _2=1`$ and $`\alpha _2=\beta _1=0`$. Eq. (6) now takes the simple form of
$$|c_1c_2(1e^{i\theta })(c_1^2+c_2^2e^{i\theta })|^2=0,$$
(9)
where $`\theta =(4J+2\stackrel{~}{\delta })t`$. When $`c_1c_2`$, as is the case when $`\delta 0`$, the only solution of this equation is $`e^{i\theta }=1`$, which corresponds to the “return” operation with an additional phase: $`|\varphi _1(t)=\alpha _1|+\alpha _2e^{i\mathrm{\Delta }t}|`$. The condition for an exact swap does not exist anymore. To find the best approximation to a swap (so that the state of first spin is as close as possible to being “down” in the current special case), we calculate the minima of the expression in Eq. (9) and find that $`e^{i\theta }=1`$ does still produce a minimum in $`|(1e^{i\theta })(c_1^2+c_2^2e^{i\theta })|^2`$ (which is, however, no longer zero). When this condition is satisfied, the density matrix of the first spin is
$$\rho _1|_{e^{i\theta }=1}=\frac{1}{1+x^2}||+\frac{x^2}{1+x^2}||,$$
(10)
where $`x=\delta /(2J)`$. Therefore, the state of the first spin cannot be simultaneously pure and the same as the initial state of the second spin. The state of the first spin will remain mixed (when it is close to a swap) with the state of the second spin. According to Eq. (10), there would be an error of the magnitude $`x^2`$ if we perform such a “swap” operation, which needs to be corrected.
For GaAs-based QDQC, the Zeeman splitting is $`g\mu _B2.55\times 10^2`$ meV/Tesla, while a single Bohr magneton produces a field of about 5 Gauss at 1 nm distance and almost nothing at a distance of 100 nm in GaAs. This difference in the magnetic field leads to a $`x^2`$ in the order of $`10^6`$, which is right within the capability of currently available error correction schemes. Since the field involved in $`\delta `$ is the average value over an entire quantum dot, local magnetic field inhomogeneity caused by impurities should not cause intractable errors.
Mathematically, inhomogeneity in the magnetic field means that a part of the Hamiltonian (3) does not commute with the exchange term $`J𝐒_1𝐒_2`$, thus the eigenstates will change, leading to errors in swap, which will cause errors in CNOT in the proposed sequential scheme . In the parallel pulse scheme , conditional phase shifts can be produced directly from the Hamiltonian (3) (with precisely controlled field inhomogeneity), circumventing the swap actions. However, stray fields, which may be invariably present in realistic solid state QDQC architectures due to magnetic impurities/imperfections, can still produce errors which need to be corrected (even in solely exchange-based models which have no applied magnetic field ). Furthermore, the impossibility of exact swap does make the transfer of spins through swap an error-prone process in the presence of an inhomogeneous field.
Finally, inhomogeneity in the external magnetic field is not only a possible error source in swap and other two-qubit operations, but also an important factor in the single-qubit operations determined solely by the Zeeman coupling. Here the main concern is the precise definition of the direction of the quantization axis for the electron spins. For example, if spin-up and -down along z direction correspond to the two states of our qubit, then the z axis is our quantization axis which provides a reference frame for all the single-qubit operations. In our calculations the magnetic field is purely along the z direction. The only spatial variation is in its magnitude. However, according to Maxwell equations, in a steady state, $`\times 𝐁=0`$. Thus, if $`B_z/x0`$, $`B_x/z=B_z/x0`$. If we have a perfect two dimensional QD system lying entirely in the x-y plane, we can assume that $`B_x`$ happens to vanish in the plane of the quantum dots. However, using the example of GaAs quantum well QD’s, the thickness of the well is generally around $`510`$ nm, which characterizes the typical z-width of the QD’s. Therefore, to have a vanishing Zeeman coupling along x direction, the $`B_x`$ field has to vanish along the middle plane of the quantum well and be an odd function in the z (growth) direction, which is a quite stringent constraint. If we choose z-direction as our quantizing axis, a finite $`B_x`$ field, which may be unavoidable, would tend to flip the spins. For a 1 Gauss $`B_x`$ field, this flipping rate is about 0.65 MHz, corresponding to a flipping time of 1.5 $`\mu `$s. Considering that the gates should be operated as slowly as possible to satisfy the adiabatic condition, this is a constraint one has to take into consideration. In fact, “decoherence” induced by a fluctuating inhomogeneous transverse field may turn out to be an important extrinsic decoherence channel for QDQC operations.
This work is supported by ARDA, LPS, and US-ONR.
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# Brane versus shell cosmologies in Einstein and Einstein-Gauss-Bonnet theories
## I Introduction
Motivated by recent developments in high energy physics there is at present a considerable increase of activity in the domain of cosmology with extra dimensions. In these models gravity is assumed to act in a $`n`$-dimensional ’bulk’ while the standard model interactions are confined to a 4-dimensional slice (’brane’ worldsheet) of this multi-dimensional spacetime. Randall and Sundrum have recently proposed two models in which all the matter is confined to a 4-dimensional brane worldsheet embedded in a 5-dimensional anti-de Sitter ($`AdS_5`$) spacetime with imposed $`Z_2`$ metric symmetry, $`ww`$ ($`w`$ denotes the fifth dimension) . In the first model, the fifth dimension is compact and bounded by two branes with positive and negative tension, the visible universe being modelled by the negative tension brane. Their second model assumes a single brane embedded in a $`AdS_5`$ with a non-compact fifth dimension.
More recently, several authors have found exact cosmological solutions describing ’4+1’ brane universes. Binétruy et al. solved the Einstein equations for a single brane embedded in a bulk governed by a negative cosmological constant. They impose the $`Z_2`$ symmetry to obtain a global cosmological solution in the Gaussian normal coordinates (’BDL’ solutions). Later, Ida and Kraus considered in Schwarzschild coordinates a spherically symmetric hypersurface moving in a 5-dimensional Schwarzschild-$`AdS_5`$ spacetime. By performing an explicit coordinate transformation from the Schwarzschild coordinates used in to the Gaussian normal coordinates used in , Mukhoyama et al. showed that all these solutions represent the same spacetime described in different coordinate systems.
Now, what is usually done in General Relativity, when, e.g., studying the gravitational collapse of spherical bodies is to join two metrically different solutions of the Einstein field equations. Integration of the Einstein equations across the surface separating the two regions leads to the Israel junction conditions (, see also Appendix A for a review) relating the surface stress-energy tensor to the discontinuity of the extrinsic curvature across the surface. The sign of the extrinsic curvature, and thus the form of the junction conditions, depend on the definitions of normal vectors in the neighbourhood of the surface. However, once the directions of the normal vectors on each side of the surface are fixed, e.g. pointing in a defined positive sense, the formalism becomes unambiguous and yields in general a nonzero stress-energy tensor for the surface (massive ’shell’).
On the other hand matter of $`Z_2`$ symmetric $`AdS_5`$ branes, which connect two metrically identical solutions of the Einstein equations, arises formally by a flip of the normal vectors at one or the other side of the boundary surface while preserving the form of the Israel junction conditions. We illustrate in section II how this flip of the normals across the brane describes the connection, via the Israel junction conditions, of a bulk region and its mirror image. Were no formal flip of the normal vectors performed, i.e. dropping out of the $`Z_2`$ symmetry, one would join in a topologically trivial way two complementary parts of the $`AdS_5`$ spacetime across a non-massive boundary surface.
As pointed out by Binétruy et al., brane cosmology leads to Friedmann-like equations different from the standard ones. They showed, however, that the standard cosmological evolution can be obtained in a $`Z_2`$ symmetric model à la Randall-Sundrum if one finely tunes the tension of the brane. In section III we show that the standard behaviour at late times can also beobtained when one drops out the $`Z_2`$ symmetry: we present a cosmological model which is a shell separating two metrically different anti-de Sitter regions. Just as Binétruy et al. we fine-tune the tension of the shell (that is we impose a particular equation of motion for the shell) in order to recover standard cosmology at late times in the shell. On the other hand the evolution of the scale factor at early times is very different from what give standard as well as brane cosmologies, as the scale factor grows from a non-zero value at the big bang singularity.
If these brane or shell models are to be the low energy limit of string theory, it is likely that the field equations include in particular the Gauss-Bonnet term, which, in five dimensions, is the only non-linear term in the curvature which yields second order field equations (see e.g. and references therein). In section IV we show how the Israel matching conditions across the membrane (either brane or shell) have to be modified to take into account the Gauss-Bonnet correction. We conclude that in this case the thin membrane approximation fails. The complete evolution of the Einstein-Gauss-Bonnet universe must be studied more carefully taking into account the internal structure of the membrane. We shall see however that the microphysics of the thick membrane may, at late times, be hidden in a renormalisation of Newton’s constant so that the Einstein and Einstein-Gauss-Bonnet membranes become identical.
## II The geometry and topology of the BDL brane solutions
In this section we analyse the first BDL brane model , which shares the same topological properties as the models discussed in and, although simpler, is not included in the analysis of Ida and Mukhoyama et al.. Consider, in a coordinate system $`x^A=\{\tau ,x^i,w\}`$ ($`i=1,2,3`$), the 5-dimensional metric of the first BDL brane model
$$ds^2=n^2(\tau ,w)d\tau ^2+S^2(\tau ,w)\delta _{ij}dx^idx^j+dw^2$$
(1)
with
$`S(\tau ,w)`$ $`=`$ $`a(\tau )\dot{a}(\tau )|w|`$ (2)
$`n(\tau ,w)`$ $`=`$ $`1{\displaystyle \frac{\ddot{a}(\tau )}{\dot{a}(\tau )}}|w|`$ (3)
where $`a(\tau )`$ is an arbitrary function of time $`\tau `$, where a dot denotes a time derivative and $`|\mathrm{}|`$ an absolute value.
It is easy to check that the Riemann tensor of the metric (1) is zero everywhere, except on the brane $`\mathrm{\Sigma }`$ defined by $`w=0`$, where it exhibits a $`\delta `$-like discontinuity. Now, were the $`Z_2`$ symmetry of the metric dropped out and $`|w|`$ replaced by $`w`$ or $`w`$, the metric (1) would be flat everywhere, including on $`\mathrm{\Sigma }`$. The presence of matter on the brane is therefore due to the reflexion symmetry imposed on the metric.
In order to illustrate how matter on the brane arises from a topologically non trivial pasting of two flat regions, we transform the metric (1) into its minkowskian form
$$ds^2=(dX^0)^2+(dX^1)^2+(dX^2)^2+(dX^3)^2+(dX^4)^2.$$
(4)
A simple way to obtain the transformation $`X^B=X^B(x^A)`$ is to integrate the equations which give the Christoffel symbols $`\mathrm{\Gamma }_{BC}^A`$ for the metric (1) as
$$\frac{X^A}{x^E}\mathrm{\Gamma }_{BC}^E=\frac{^2X^A}{x^Bx^C}$$
(5)
and to impose
$$g_{AB}=\frac{X^C}{x^A}\frac{X^D}{x^B}\eta _{CD}.$$
(6)
After some algebra, one gets (up to a Lorentz transformation) :
$`X^0`$ $`=`$ $`S(\tau ,w)\left({\displaystyle \frac{r^2}{4}}+1{\displaystyle \frac{1}{4h^2a^2}}\right){\displaystyle \frac{1}{2}}{\displaystyle 𝑑\tau \frac{\dot{h}}{ah^3}}`$ (7)
$`X^i`$ $`=`$ $`S(\tau ,w)x^i`$ (8)
$`X^4`$ $`=`$ $`S(\tau ,w)\left({\displaystyle \frac{r^2}{4}}1{\displaystyle \frac{1}{4h^2a^2}}\right){\displaystyle \frac{1}{2}}{\displaystyle 𝑑\tau \frac{\dot{h}}{ah^3}}.`$ (9)
where $`h\dot{a}/a`$ and $`r^2\delta _{ij}x^ix^j`$.
Inverting (8) gives $`|w|`$ as some function of the $`X^A`$ coordinates : $`|w|=f(X^A)`$. In the simple case $`a(\tau )=\tau `$, we have for example
$$|w|=\frac{1}{4}(3X^0+5X^4)\frac{R^2}{X^0X^4}$$
(10)
(with $`R^2\delta _{ij}X^iX^j`$). As for the brane $`\mathrm{\Sigma }`$, it is represented by the hypersurface $`f(X^A)=0`$. If $`a(\tau )=\tau ^q`$ for example, its equation reads
$$\frac{\eta _{AB}X^AX^B}{(X^0X^4)^{2q}}=\frac{q^2}{2^{2q}(2q1)}.$$
(11)
The point to note is that the coordinates $`x^A=\{\tau ,x^i,w\}`$ cover only the portion $`f(X^A)0`$ of the full Minkowski spacetime spanned by the $`X^A`$ coordinates, and that each point of this ’half’-space represents two points of the bulk, one corresponding to, say, $`w=a`$ and the other to $`w=a`$. In order to have a one-to-one correspondence between the $`x^A`$ and the $`X^A`$ coordinates, one must unfold this half Minkowski space along the brane $`\mathrm{\Sigma }`$. In doing so the normal vector $`n^A`$, chosen to point to, say, increasing $`f(X^A)`$, will point to positive, increasing, $`w`$ on one side of the brane (the ’$`+`$’ side), and to negative, decreasing, $`w`$ on the ’$``$’ side. In other words, $`n^A`$ flips across the brane. As a consequence the extrinsic curvature $`K_{\mu \nu }^\pm `$ of the $`w0_\pm `$ surfaces ($`x^\mu `$ being four coordinates on the brane and $`ds_\mathrm{\Sigma }^2=\gamma _{\mu \nu }dx^\mu dx^\nu `$ its induced metric) becomes discontinuous across the brane in such a way that
$$K_{\mu \nu }^+=K_{\mu \nu }^{}.$$
(12)
The standard Israel junctions conditions then give the stress-energy tensor $`𝒯_{\mu \nu }`$ of the matter on the brane as
$$\kappa \left(𝒯_\nu ^\mu \frac{1}{3}\delta _\nu ^\mu 𝒯\right)=\widehat{K}_\nu ^\mu \text{with}\widehat{K}_{\mu \nu }K_{\mu \nu }^+K_{\mu \nu }^{},$$
(13)
where indices are raised with the inverse metric $`\gamma ^{\mu \nu }`$, where $`𝒯\gamma ^{\mu \nu }𝒯_{\mu \nu }`$ and where $`\kappa `$ is Einstein’s constant.
If one had simply considered, in a flat, simply oriented, 5-dimensional universe in Minkowskian coordinates $`X^A`$, the surface $`\mathrm{\Sigma }`$ defined by e.g. (10), there would have been (with our convention on the directions of normal vectors) no discontinuity in the extrinsic curvature across that surface so that no matter would have been required on $`\mathrm{\Sigma }`$ for the Israel matching conditions (13) to be satisfied. Equivalently, as we have already remarked, if one replaced in the BDL metric (1) $`|w|`$ by $`w`$ or $`w`$, the coordinates $`x^A`$ would cover the full Minkowski spacetime and one would describe a completely flat universe, so that there would be no discontinuity in the extrinsic curvature across the surface $`w=0`$ and no matter on the brane.
The procedure to obtain (13) is reminiscent of what is sometimes done in the Kerr spacetime, where a region between two planes $`z=\pm z_0`$ is excised in order to create a discontinuity in the extrinsic curvature and hence a source for the Kerr solution (cf. e.g. ).
The more general, second, BDL brane cosmological model is a solution, everywhere except on the brane $`\mathrm{\Sigma }`$, of $`G_{AB}+\mathrm{\Lambda }g_{AB}=0`$, where $`G_{AB}`$ is the five-dimensional Einstein tensor of the metric $`g_{AB}`$ and $`\mathrm{\Lambda }(<0)`$ a cosmological constant. Mukhoyama et al. showed that the bulk metric is that of a Schwarzschild-anti-de Sitter space time: they performed a coordinate transformation $`\{\tau ,x^i,w\}\{t,r,\chi ,\theta ,\varphi \}`$ in order to bring the BDL metric in the bulk into the more familiar form
$$ds^2=\mathrm{\Phi }(r)dt^2+\frac{dr^2}{\mathrm{\Phi }(r)}+r^2[d\chi ^2+f_k^2(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)]$$
(14)
with
$$\mathrm{\Phi }(r)=k\frac{M}{r^2}\frac{\mathrm{\Lambda }}{6}r^2$$
(15)
where $`f_1(\chi )=\mathrm{sin}\chi `$, $`f_0(\chi )=\chi `$, $`f_1(\chi )=\mathrm{sinh}\chi `$, and where $`M`$ is the Schwarzschild mass parameter. (Note that the case $`k=M=\mathrm{\Lambda }=0`$ considered above is not included in this transformation.) In the coordinates $`\{t,r,\chi ,\theta ,\varphi \}`$ the brane $`\mathrm{\Sigma }`$ is a four dimensional sphere, $`r=a(t)`$ (and hence is geometrically simpler than when the bulk is flat, cf. e.g. (11)).
Again, the presence of matter on $`\mathrm{\Sigma }`$ arises from the fact that the BDL spacetime is not a simply oriented Schwarzschild-anti-de Sitter space time but is obtained by disregarding the region, say, inside the hypersphere $`\mathrm{\Sigma }`$ and unfolding the region outside along it, which implies a flip of the normal vectors across $`\mathrm{\Sigma }`$. In a simply oriented spacetime (14) there would be no need for the matter on the hypershpere $`\mathrm{\Sigma }`$ for the junction conditions to be satisfied.
## III The universe as a shell separating two anti-de Sitter regions
Let us now consider, in a spirit more akin to what is usually done in General Relativity, a simply oriented five dimensional spacetime, consisting of two metrically different anti-de Sitter regions, one solution of $`G_{AB}+\mathrm{\Lambda }_+g_{AB}=0`$, the other of $`G_{AB}+\mathrm{\Lambda }_{}g_{AB}=0`$, separated by a shell $`\mathrm{\Sigma }`$. The bulk regions, with Riemann tensors $`R_{ABCD}=L_\pm (g_{AC}g_{BD}g_{AD}g_{BC})`$, where $`L_\pm =\mathrm{\Lambda }_\pm /6`$, are both described by coordinates $`(t,r,\chi ,\theta ,\varphi )`$. The metrics outside, resp. inside, $`\mathrm{\Sigma }`$ have the following forms
$`ds^2|_+`$ $`=`$ $`\mathrm{\Phi }_+(r)dt^2+\mathrm{\Phi }_+^1(r)dr^2+r^2[d\chi ^2+f_k^2(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)],`$ (16)
$`ds^2|_{}`$ $`=`$ $`\mathrm{\Phi }_{}(r)\alpha ^2(t)dt^2+\mathrm{\Phi }_{}^1(r)dr^2+r^2[d\chi ^2+f_k^2(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)],`$ (17)
with
$$\mathrm{\Phi }_\pm (r)=kL_\pm r^2,$$
(18)
and where the lapse function $`\alpha (t)`$ is introduced in order to match the two coordinate grids through the shell <sup>*</sup><sup>*</sup>*One can consider two Schwarzschild-Anti de Sitter regions corresponding to the BDL solution with $`\mathrm{\Lambda }<0`$, $`𝒞0`$ . Ida showed that the Schwarzschild mass parameter $`M=M_+=M_{}`$ generates an effective radiative term in the Friedmann-like equation, $`Ma^4`$, and thus corresponds to the constant $`𝒞`$ of the BDL solution. We shall not consider this generalisation here (cf. also )..
The equation of motion of the shell is given parametrically by
$$r=a(\tau );t=t(\tau ).$$
(19)
Choosing $`\tau `$ as the proper time on the shell and requiring the continuity of the metrics (16) and (17) at $`\mathrm{\Sigma }`$ then gives
$$\dot{t}=\frac{\sqrt{\mathrm{\Phi }_+(a)+\dot{a}^2}}{\mathrm{\Phi }_+(a)};\alpha (\tau )=\frac{\mathrm{\Phi }_+(a)}{\mathrm{\Phi }_{}(a)}\sqrt{\frac{\dot{a}^2+\mathrm{\Phi }_{}(a)}{\dot{a}^2+\mathrm{\Phi }_+(a)}},$$
(20)
where a dot denotes $`d/d\tau `$. The metric on $`\mathrm{\Sigma }`$ then takes the Friedmann-Lemaître form
$$ds_\mathrm{\Sigma }^2\gamma _{\mu \nu }dx^\mu dx^\nu =d\tau ^2+a^2(\tau )[d\chi ^2+f_k^2(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)].$$
(21)
In the interior and exterior regions, the independent tangent vectors at the shell are given as $`e_\tau ^{A\pm }=(\dot{t},\dot{a},0,0,0)`$ and $`e_i^{A\pm }=\delta _i^A`$ where the index $`i`$ stands for $`(\chi ,\theta ,\varphi )`$. The normal vectors to the shell both pointing in the positive direction, by definition from ’$``$’ to ’$`+`$’, normalized to unity, are
$$n_A^{}=\alpha (t)(\dot{a},\dot{t},0,0,0),n_A^+=(\dot{a},\dot{t},0,0,0).$$
The extrinsic curvature is defined by
$$K_{\mu \nu }^\pm e_\mu ^Ae_\nu ^B_An_B|^\pm ,$$
where the indices $`(\mu ,\nu )`$ stand for $`(\tau ,\chi ,\theta ,\varphi )`$ and are raised and lowered by the induced Friedmann-Lemaître metric (21). $`_A`$ is the covariant derivative associated with the metric (16) or (17).
The non-vanishing components of the extrinsic curvature read
$`K_\tau ^{\tau \pm }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{h^2+\frac{k}{a^2}L_\pm }}}\left[{\displaystyle \frac{\ddot{a}}{a}}L_\pm \right],`$ (22)
$`K_\chi ^{\chi \pm }`$ $`=`$ $`K_\theta ^{\theta \pm }=K_\varphi ^{\varphi \pm }=\sqrt{h^2+{\displaystyle \frac{k}{a^2}}L_\pm },`$ (23)
with $`h\dot{a}/a`$.
Israel’s junction conditions give the stress-energy tensor of the shell in terms of the jump in its extrinsic curvature. The stress-energy tensor of the shell has the perfect-fluid form, in the coordinates of metric (21) it reads $`𝒯_\nu ^\mu =diag(\rho ,p,p,p)`$ where $`\rho `$ and $`p`$ are the energy density and pressure of the shell respectively. The junction condition (13) gives
$$\rho =\frac{3}{\kappa }\widehat{K}_\chi ^\chi .$$
(24)
As for $`p`$, it is given by the conservation law
$$\dot{\rho }+3(\dot{a}/a)(\rho +p)=0.$$
(25)
Putting together (23) and (24) we get
$$\frac{\kappa \rho }{3}=\sqrt{h^2+\frac{k}{a^2}L_{}}\sqrt{h^2+\frac{k}{a^2}L_+}.$$
(26)
By contrast the brane cosmological model is obtained by imposing the $`Z_2`$ symmetry: $`LL_+=L_{}`$ and changing the direction of one of the normal vectors. The reflexion $`n_+^A=n_{}^A`$ then gives in the brane
$$\frac{\kappa \rho }{6}=\sqrt{h^2+\frac{k}{a^2}L}.$$
(27)
In ’shell cosmology’ on the other hand to impose $`L_+=L_{}`$ would yield a massless boundary surface.
In order to recover standard cosmology at late times in the shell we have to consider $`L_+0`$, $`L_{}0`$, otherwise $`h^2+k/a^2`$ can never go to zero. We also choose $`L_+<L_{}<0`$ (the particular case $`L_{}=0`$ is considered below). Furthermore, as in , we must decompose the stress-energy tensor in the shell into:
$$\rho \rho _m+\sigma ;pp_m\sigma ,$$
(28)
where $`\rho _m`$, $`p_m`$ can be interpreted as the energy density and pressure of ordinary matter and where the tension $`\sigma `$ must be fine-tuned as
$$\sigma \frac{3}{\kappa }\left(\sqrt{L_+}\sqrt{L_{}}\right).$$
Equation (26) then reduces at late times to the Friedman-like equation for the energy density of ordinary matter
$$h^2+\frac{k}{a^2}=\frac{8\pi G}{3}\rho _m+O(\rho _m^2),$$
(29)
where
$$\frac{8\pi G}{3}\frac{2\kappa }{3}\frac{\sqrt{L_+L_{}}}{\sqrt{L_+}\sqrt{L_{}}}>0,$$
(30)
and $`G`$ can be interpreted as Newton’s constant. Thus, at late times one recovers as in brane cosmology the standard FL behaviour.
The complete time evolution of the shell universe can be analysed easily for $`k=0`$. The conservation law (25) for the equation of state
$$p_m=v\rho _m,v=constant,$$
(31)
gives $`\rho _ma^q`$, $`q3(1+v)`$. Proceeding as in standard cosmology, we then obtain from (26) and (28):
$$\tau =\frac{2}{q\sigma }_{\frac{\rho _m}{\sigma }}^1\frac{(1x)dx}{x^{2/3}\sqrt{2x}}\frac{1}{\sqrt{x^2+2x+\gamma }}\text{with}\gamma \frac{4\sqrt{L_+L_{}}}{\sigma ^2}.$$
(32)
The behaviour of the shell in the neighbourhood of the big bang is different from that in both standard $`(a\tau ^{2/q})`$ and brane $`(a\tau ^{1/q})`$ cosmologies. Indeed, in the passage to the limit $`\tau 0`$ in (32) we have
$$\frac{a}{a_{bb}}1+\frac{1}{q}\sqrt{\frac{\tau }{d}}\text{with}\frac{1}{d}q\sigma \sqrt{1+\gamma }.$$
(33)
Hence, the scale factor grows from a non-zero value $`a_{bb}`$ at the big bang as $`\sqrt{\tau }`$ (for all equations of state $`p_m=v\rho _m`$, $`v0`$). Nevertheless, the singularity of curvature is present at the big bang: $`R_{\nu \kappa \lambda }^\mu \mathrm{}`$, because $`\dot{a},\ddot{a}\mathrm{}`$ when $`\tau 0`$.
In the particular case $`L_{}=0`$ the shell connects a flat and a $`AdS_5`$ regions. It can therefore represent the ’negative-tension brane’ in the GRS model of ’quasi-localized’ gravity (cf. also ). As one can see from (29-30), $`G=0`$ and standard cosmology cannot be recovered at late times in such a shell; we have $`a\tau ^{1/q}`$ instead of $`a\tau ^{2/q}`$ when $`\tau \mathrm{}`$.
## IV Einstein-Gauss-Bonnet membrane cosmologies
Up to now we have considered membrane cosmological models which are solutions of the Einstein equations: $`G_{AB}+\mathrm{\Lambda }g_{AB}=0`$ everywhere except on the membrane $`\mathrm{\Sigma }`$. In brane cosmology the cosmological constant $`\mathrm{\Lambda }`$ is the same on each side of $`\mathrm{\Sigma }`$, in shell cosmology, it jumps from $`\mathrm{\Lambda }_+`$ to $`\mathrm{\Lambda }_{}`$. Now, if these models are to be the low energy limit of string theory, it is likely that the field equations should be generalized and include the Gauss-Bonnet term (cf. e.g. and ref. therein). In this section we shall therefore consider the gravitational action
$$S_g=d^5x\sqrt{g}(2\mathrm{\Lambda }+R+\alpha L_2)$$
(34)
with
$$L_2=R_{ABCD}R^{ABCD}4R_{AB}R^{AB}+R^2$$
(35)
where $`\alpha `$ is a coupling constant, and where $`R_{ABCD}`$, $`R_{AB}`$ and $`R`$ are the Riemann tensor, the Ricci tensor and the scalar curvature of the five dimensional metric $`g_{AB}`$ with determinant $`g`$. The corresponding field equations, outside the membrane, are (see e.g. )
$$\mathrm{\Lambda }g_{AB}+G_{AB}+\alpha H_{AB}=0$$
(36)
with
$$H_{AB}2R_{ALMN}R_B^{LMN}4R_{AMBN}R^{MN}4R_{AM}R_B^M+2RR_{AB}\frac{1}{2}g_{AB}L_2.$$
(37)
Contrarily to Einstein’s equations, the equations (36) possess, for $`\alpha 0`$ and a given value of the cosmological constant $`\mathrm{\Lambda }`$, two (anti) de Sitter solutions
$`R_{ABCD}`$ $`=`$ $`L_\pm (g_{AC}g_{BD}g_{AD}g_{BC})`$ (38)
$`\text{with}L_\pm `$ $`=`$ $`{\displaystyle \frac{1}{4\alpha }}\left(1\pm \sqrt{1+{\displaystyle \frac{4\alpha \mathrm{\Lambda }}{3}}}\right).`$ (39)
In brane cosmology we shall choose one or the other solution everywhere in the bulk. In shell cosmology, we shall choose the solution $`L_+`$ on one side of the shell and the solution $`L_{}`$ on the other side. (Shell cosmologies are therefore more satisfactory in Einstein-Gauss-Bonnet theory as one does not have to impose different cosmological constants on each side of the shell.)
In order to get the stress-energy tensor on the membrane, one proceeds along Israel’s line. Like Binetruy et al., we first choose a Gaussian coordinate system $`(w,x^\mu )`$ such that the metric reads
$`ds^2`$ $`=`$ $`dw^2+\gamma _{\mu \nu }dx^\mu dx^\nu `$ (40)
$`=`$ $`dw^2n^2(\tau ,w)d\tau ^2+S^2(\tau ,w)[d\chi ^2+f_k^2(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)],`$ (41)
$`w=0`$ being the equation of the membrane $`\mathrm{\Sigma }`$. In this coordinate system the extrinsic curvature of the surfaces $`w=constant`$ is simply given by
$$K_{\mu \nu }=\frac{1}{2}\frac{\gamma _{\mu \nu }}{w}.$$
(42)
It jumps across the membrane from $`K_{\mu \nu }^+`$ to $`K_{\mu \nu }^{}`$ (with $`K_{\mu \nu }^+=K_{\mu \nu }^{}`$ in the case of branes) and this discontinuity can be described in terms of the Heaviside distribution.
Expressing now the Riemann tensor (39) in terms of $`K_{\mu \nu }`$ and the four dimensional Riemann tensor of the metric $`\gamma _{\mu \nu }`$ we then obtain from (36), everywhere outside the membrane (see Appendices A and B for the ’4+1’ decompositions of $`G_B^A`$ and $`H_B^A`$)
$$\mathrm{\Lambda }\delta _\mu ^\nu +G_\mu ^\nu +\alpha H_\mu ^\nu =(1+4\alpha L)\left(\frac{K_\mu ^\nu }{w}\delta _\mu ^\nu \frac{K}{w}\right)+\mathrm{}(=0)$$
(43)
where $`K\gamma ^{\alpha \beta }K_{\alpha \beta }`$, where $`L=L_+`$ or $`L_{}`$, and where the dots stand for terms containing at most first order $`w`$-derivatives of $`\gamma _{\mu \nu }`$.
In Einstein’s theory, $`\alpha =0`$, and (43), in the vicinity of $`\mathrm{\Sigma }`$, is well defined in a distributional sense: $`K_\mu ^\nu /w`$ can be expressed in terms of the Dirac distribution and the integration of (43) across the membrane gives Israel’s junction conditions, that is the stress-energy tensor on the membrane in terms of the jump in the extrinsic curvature, eq. (13).
When $`\alpha 0`$ on the other hand, (43) is not well defined in a distributional sense, as $`L`$ cannot be considered as an infinitely $`w`$-differentiable function. Indeed, in shell cosmology, $`L`$ jumps from $`L_+`$ to $`L_{}`$ across $`\mathrm{\Sigma }`$, and in brane cosmology, $`L=L_+=L_{}`$ is continuous across $`\mathrm{\Sigma }`$, but, because of the reflexion symmetry, has a discontinuous $`w`$-derivative. This mathematical obstruction simply means that, in Einstein-Gauss-Bonnet theory, membranes cannot be treated in the thin wall approximation: the jumps in the extrinsic curvature and in $`L`$ or its derivative have to be described in detail within specific microphysical models.
When the thickness of the membrane is taken into account, the distributions $`K_\mu ^\nu /w`$ and $`L`$ are replaced by rapidly varying but $`C^{\mathrm{}}`$ functions. Supposing that the metric keeps the form (41), we can define from the $`\tau `$-$`\tau `$ component of (43) the sharply peaked function (cf. (B19) in Appendix B)
$$\kappa \rho 3(1+4\alpha L)\frac{K_\chi ^\chi }{w}$$
(44)
(The $`\chi `$-$`\chi `$ component of (43) is redundant thanks to the conservation equation, that is the Bach-Lanczos identity (B6).) In the vicinity and inside the membrane $`K_\chi ^\chi `$ can be written as
$$K_\chi ^\chi =\frac{1}{2}\overline{K}_\chi ^\chi +\frac{1}{2}\widehat{K}_\chi ^\chi f(\tau ,w)$$
(45)
where $`\overline{K}_{\chi \chi }K_{\chi \chi }^++K_{\chi \chi }^{}`$ and where the function $`f(\tau ,w)`$, which varies rapidly from $`1`$ to $`+1`$ across $`\mathrm{\Sigma }`$, encapsulates its microphysics. Similarly, in the case of brane cosmology, we can write
$$L=\stackrel{~}{L}g_b(\tau ,w)$$
(46)
where $`\stackrel{~}{L}=L_+`$ or $`L_{}`$ and where $`g_b(\tau ,w)`$ is some even function of $`w`$ which varies rapidly from $`+1`$ to $`+1`$ across the brane. In shell cosmology on the other hand
$$L=\frac{1}{2}\overline{L}+\frac{1}{2}\widehat{L}g_s(\tau ,w)$$
(47)
where $`g_s(\tau ,w)`$ varies rapidly from $`1`$ to $`+1`$ and $`\overline{L}=\frac{1}{2\alpha }`$, $`\widehat{L}=\frac{1}{2\alpha }\sqrt{1+\frac{4\alpha \mathrm{\Lambda }}{3}}`$.
Integrating (44) across $`\mathrm{\Sigma }`$ we therefore get the energy density of the membrane as
$$\kappa \varrho _\eta ^{+\eta }𝑑w\kappa \rho =3\widehat{K}_\chi ^\chi \left[1+2\alpha \stackrel{~}{L}_\eta ^{+\eta }𝑑wg_b\frac{f}{w}\right]$$
(48)
in the case of branes, and
$$\kappa \varrho _\eta ^{+\eta }𝑑w\kappa \rho =\frac{3}{2}\sqrt{1+\frac{4\alpha \mathrm{\Lambda }}{3}}\widehat{K}_\chi ^\chi _\eta ^{+\eta }𝑑wg_s\frac{f}{w}$$
(49)
in the case of shells (the fact that one does not recover the results of Einstein’s theory when $`\alpha =0`$ is not surprising as $`L_\pm `$ is divergent in that case) . Now, if $`f`$ or $`g_{b/s}`$ depend on $`\tau `$, the integrals in (48)-(49) are some functions of $`\tau `$. But, if $`f`$ and $`g_{b/s}`$ do not depend on time, which is probably to be expected whence the brane has reached a stationary state, that is at late times, then the integrals in (48)-(49) are just numbers. In this case then, the microphysics of the membrane is simply hidden in a renormalization of the Einstein constant $`\kappa `$ and Einstein and Einstein-Gauss-Bonnet branes are indistinguishable.
We therefore conclude that Einstein and Einstein-Gauss-Bonnet membranes cosmologies are probably indistinguishable (and Friedmannian) at late times, but that the early time behaviour of the scale factor needs to be studied more carefully, taking into account the microphysics of the membrane, as is done in e.g. .
## Acknowledgements
We thank Thibaut Damour and David Langlois for fruitful discussions.
## A Junction conditions for non-null surfaces in General Relativity
This appendix summarizes the junction conditions in the theory of Einstein (Lanczos , Darmois , Misner and Sharp , Israel ). Suppose we are given a 4-dimensional hypersurface ($`\mathrm{\Sigma }`$) in a 5-dimensional spacetime (metric $`g_{AB}`$) which can be imagined as the element of a family of surfaces. The normal vectors $`n^A`$ to this family of surfaces are not null; $`n_An^Aϵ=\pm 1`$. They are all oriented in a positive direction defined in the bulk. Let the surface be either spacelike ($`ϵ=1`$) or timelike ($`ϵ=+1`$). As an aid in deriving junction conditions we introduce Gaussian normal coordinates in the neighbourhood of $`\mathrm{\Sigma }`$. The metric $`g_{AB}`$ has the form
$$ds^2=ϵdw^2+\gamma _{\mu \nu }dx^\mu dx^\nu ,$$
(A1)
and the extrinsic curvature of the surfaces $`w=constant`$ is
$$K_{\mu \nu }=\frac{1}{2}\frac{\gamma _{\mu \nu }}{w}.$$
(A2)
The curvature tensor of the metric $`g_{AB}`$ can be expressed in terms of the intrinsic curvature of 4-dimensional hypersurface (metric $`\gamma _{\mu \nu }`$) and of its extrinsic curvature; one gets the so-called Gauss-Codazzi equations. In the special case of Gaussian normal coordinates the equations simplify to
$`R_{w\mu w\nu }`$ $`=`$ $`{\displaystyle \frac{K_{\mu \nu }}{w}}+K_{\rho \nu }K_\mu ^\rho ,`$ (A3)
$`R_{w\mu \nu \rho }`$ $`=`$ $`_\nu K_{\mu \rho }_\rho K_{\mu \nu },`$ (A4)
$`R_{\lambda \mu \nu \rho }`$ $`=`$ $`{}_{}{}^{4}R_{\lambda \mu \nu \rho }^{}+ϵ\left[K_{\mu \nu }K_{\lambda \rho }K_{\mu \rho }K_{\lambda \nu }\right],`$ (A5)
where $`_\rho `$ is the covariant derivative with respect to the 4-dimensional metric $`\gamma _{\mu \nu }`$. From (A3)-(A5) we obtain the decomposition of the Ricci tensor ($`R_{AB}=g^{CD}R_{CADB}`$) and of the scalar curvature ($`R=g^{AB}R_{AB}`$) as:
$`R_{ww}`$ $`=`$ $`\gamma ^{\mu \nu }{\displaystyle \frac{K_{\mu \nu }}{w}}+Tr(K^2),`$ (A6)
$`R_{w\mu }`$ $`=`$ $`_\mu K_\nu K_\mu ^\nu ,`$ (A7)
$`R_{\mu \nu }`$ $`=`$ $`{}_{}{}^{4}R_{\mu \nu }^{}+ϵ\left[{\displaystyle \frac{K_{\mu \nu }}{w}}+2K_\mu ^\rho K_{\rho \nu }KK_{\mu \nu }\right],`$ (A8)
$`R`$ $`=`$ $`{}_{}{}^{4}R+ϵ\left[2\gamma ^{\mu \nu }{\displaystyle \frac{K_{\mu \nu }}{w}}+3Tr(K^2)K^2\right],`$ (A9)
where we defined $`KK_\mu ^\mu `$ and $`Tr(K^2)K^{\mu \nu }K_{\mu \nu }`$.
In terms of the intrinsic and extrinsic curvature of the 4-dimensional hypersurfaces $`w=constant`$, the Einstein tensor ($`G_A^B=R_A^B(1/2)\delta _A^BR`$) and the field equations have components
$`G_w^w`$ $`=`$ $`{\displaystyle \frac{1}{2}}{}_{}{}^{4}R+{\displaystyle \frac{1}{2}}ϵ\left[K^2Tr(K^2)\right]=\kappa T_w^w,`$ (A10)
$`G_\mu ^w`$ $`=`$ $`ϵ\left[_\mu K_\nu K_\mu ^\nu \right]=\kappa T_\mu ^w,`$ (A11)
$`G_\nu ^\mu `$ $`=`$ $`{}_{}{}^{4}G_{\nu }^{\mu }+ϵ\left[{\displaystyle \frac{K_\nu ^\mu }{w}}\delta _\nu ^\mu {\displaystyle \frac{K}{w}}\right]`$ (A13)
$`+ϵ\left[KK_\nu ^\mu +{\displaystyle \frac{1}{2}}\delta _\nu ^\mu Tr(K^2)+{\displaystyle \frac{1}{2}}\delta _\nu ^\mu K^2\right]=\kappa T_\nu ^\mu .`$
If the stress-energy tensor $`T_B^A`$ contains a ’delta-function contribution’ at $`\mathrm{\Sigma }`$, the integral of $`T_B^A`$ with respect to the proper distance $`w`$ measured perpendicularly through $`\mathrm{\Sigma }`$,
$$𝒯_B^A\underset{\eta 0}{lim}\left[_\eta ^\eta T_B^A𝑑w\right],$$
(A14)
is non-zero and represents the surface stress-energy tensor. In this case the extrinsic curvature must be a distribution of ’Heaviside type’ at $`\mathrm{\Sigma }`$. Integral (A14) applied on equations (A10)-(A13) yields the junction conditions relating the stress-energy tensor of $`\mathrm{\Sigma }`$ to the discontinuity of the extrinsic curvature at $`\mathrm{\Sigma }`$. In the passage to the limit $`\eta 0`$ only the terms $`(K_\nu ^\mu /w)`$ contribute to yield
$`\kappa 𝒯_w^w`$ $`=`$ $`0,`$ (A15)
$`\kappa 𝒯_\mu ^w`$ $`=`$ $`0,`$ (A16)
$`\kappa 𝒯_\nu ^\mu `$ $`=`$ $`\widehat{K}_\nu ^\mu \delta _\nu ^\mu \widehat{K},\text{where}\widehat{K}_\nu ^\mu K_\nu ^\mu (0+)K_\nu ^\mu (0).`$ (A17)
It is useful to denote $`K_{\mu \nu }^+K_{\mu \nu }(0+)`$, $`K_{\mu \nu }^{}K_{\mu \nu }(0)`$.
As for the intrinsic geometry of $`\mathrm{\Sigma }`$, it must be continuous across $`\mathrm{\Sigma }`$; this is the second junction contition completing equations (A17). If there are no ’delta singularities’ contained in $`T_B^A`$, the bulk is sliced by massless ’boundary surfaces’ (see for details).
## B Junction conditions for the theory of Einstein-Gauss-Bonnet
The theory of Einstein-Gauss-Bonnet is based on the following LagrangianWe analyse the particular case of a 5-dimensional spacetime sliced by 4-dimensional hypersurfaces. (cf. e.g. and ref. therein)
$$=\sqrt{g}\left[2\mathrm{\Lambda }+R+\alpha L_2\right],$$
(B1)
where $`g`$ is the determinant of the 5-dimensional metric $`g_{AB}`$ and $`\alpha `$ is a constant of the dimension of $`[length]^2`$. $`L_2`$ is the Gauss-Bonnet Lagrangian which reads
$$L_2=R_{ABCD}R^{ABCD}4R_{AB}R^{AB}+R^2.$$
(B2)
The Euler variation of $``$ gives the following field equations:
$$\mathrm{\Lambda }g_{AB}+G_{AB}+\alpha H_{AB}=0.$$
(B3)
$`G_{AB}`$ is the Einstein tensor and $`H_{AB}`$ is its analogue stemmed from the Gauss-Bonnet part of the Lagrangian, $`L_2`$,
$`G_{AB}`$ $``$ $`R_{AB}{\displaystyle \frac{1}{2}}g_{AB}R,`$ (B4)
$`H_{AB}`$ $``$ $`2\left[R_{ALMN}R_B^{LMN}2R_{AMBN}R^{MN}2R_{AM}R_B^M+RR_{AB}\right]{\displaystyle \frac{1}{2}}g_{AB}L_2.`$ (B5)
$`G_B^A`$ and $`H_B^A`$ satisfy the Bianchi and Bach-Lanczos identities respectively
$$_AG_B^A=0,_AH_B^A=0.$$
(B6)
In order to derive the junction condition in the theory of Einstein-Gauss-Bonnet we need to express $`H_{AB}`$ in terms of the intrinsic curvature of hypersurfaces $`w=constant`$ and their extrinsic curvatures. We adopt the notation used in Appendix A in which the ’4+1’ decomposition of the Einstein tensor $`G_{AB}`$ is shown (expressions (A10)-(A13)).
Inserting the decomposition of the curvature tensors (A3)-(A9) into (B5) one finds the following results: $`H_B^A`$ does not contain terms $`(K_\nu ^\mu /w)^2`$ as one would expect since $`H_B^A`$ contains terms $`(R_{ABCD})^2`$. Further, there are no terms linear in ($`K_\nu ^\mu /w`$) in $`H_w^w`$ and $`H_\mu ^w`$ so that these junction conditions correspond to those in the theory of Einstein: $`𝒯_w^w=𝒯_\mu ^w=0`$. The components $`H_\nu ^\mu `$ are
$`H_\nu ^\mu `$ $`=`$ $`\left\{{\displaystyle \frac{K_\nu ^\mu }{w}}\right\}\left(2Tr(K^2)2K^2\right)+\left\{{\displaystyle \frac{K_\lambda ^\mu }{w}}\right\}\left(4KK_\nu ^\lambda 4K_{\nu \beta }K^{\beta \lambda }\right)`$ (B7)
$`+`$ $`\left\{{\displaystyle \frac{K_{\nu \lambda }}{w}}\right\}\left(4KK^{\lambda \mu }4K_\beta ^\mu K^{\beta \lambda }\right)+\left\{{\displaystyle \frac{K}{w}}\right\}\left(4K_\beta ^\mu K_\nu ^\beta 4KK_\nu ^\mu \right)`$ (B8)
$`+`$ $`\left\{{\displaystyle \frac{K_{\alpha \beta }}{w}}\right\}\left(4K_\nu ^\mu K^{\alpha \beta }4K^{\alpha \mu }K_\nu ^\beta \right)+\left\{{\displaystyle \frac{K_{\alpha \beta }}{w}}\right\}\delta _\nu ^\mu \left(4K_\gamma ^\alpha K^{\gamma \beta }4KK^{\alpha \beta }\right)`$ (B9)
$`+`$ $`\left\{{\displaystyle \frac{K}{w}}\right\}\delta _\nu ^\mu \left(2K^22Tr(K^2)\right)`$ (B10)
$`+`$ $`ϵ\left(4{}_{}{}^{4}R_{\alpha \nu }^{\mu \beta }{\displaystyle \frac{K_\beta ^\alpha }{w}}4{}_{}{}^{4}R_{\nu }^{\alpha }{\displaystyle \frac{K_\alpha ^\mu }{w}}4{}_{}{}^{4}R_{}^{\alpha \mu }{\displaystyle \frac{K_{\nu \alpha }}{w}}\right)`$ (B11)
$`+`$ $`ϵ\left(4{}_{}{}^{4}R_{\nu }^{\mu }{\displaystyle \frac{K}{w}}+2{}_{}{}^{4}R{\displaystyle \frac{K_\nu ^\mu }{w}}+4\delta _\nu ^\mu {}_{}{}^{4}R_{}^{\alpha \beta }{\displaystyle \frac{K_{\alpha \beta }}{w}}2\delta _\nu ^\mu {}_{}{}^{4}R{\displaystyle \frac{K}{w}}\right)`$ (B12)
$`+`$ $`\mathrm{},`$ (B13)
where ’…’ includes terms of zeroth order in $`(K_\nu ^\mu /w)`$ which disappear in the passage to the limit $`\eta 0`$ of the integration (A14).
In what follows we confine the analysis to timelike surfaces ($`ϵ=+1`$). The previous expression can be simplified by using the Leibnitz rule that holds for both the function and the distributions
$`H_\nu ^\mu `$ $`=`$ $`4{\displaystyle \frac{}{w}}\left\{KK_\alpha ^\mu K_\nu ^\alpha K_\alpha ^\mu K^{\alpha \beta }K_{\beta \nu }+{\displaystyle \frac{1}{2}}K_\nu ^\mu Tr(K^2){\displaystyle \frac{1}{2}}K_\nu ^\mu K^2\right\}`$ (B14)
$`+`$ $`4{\displaystyle \frac{}{w}}\left\{\delta _\nu ^\mu {\displaystyle \frac{1}{2}}KTr(K^2)+\delta _\nu ^\mu {\displaystyle \frac{1}{3}}Tr(K^3)+\delta _\nu ^\mu {\displaystyle \frac{1}{6}}K^3\right\}`$ (B15)
$`+`$ $`4\left({}_{}{}^{4}R_{\alpha \nu }^{\mu \beta }{\displaystyle \frac{K_\beta ^\alpha }{w}}{}_{}{}^{4}R_{\nu }^{\alpha }{\displaystyle \frac{K_\alpha ^\mu }{w}}{}_{}{}^{4}R_{}^{\alpha \mu }{\displaystyle \frac{K_{\nu \alpha }}{w}}\right)`$ (B16)
$`+`$ $`4\left({}_{}{}^{4}R_{\nu }^{\mu }{\displaystyle \frac{K}{w}}+{\displaystyle \frac{1}{2}}{}_{}{}^{4}R{\displaystyle \frac{K_\nu ^\mu }{w}}+\delta _\nu ^\mu {}_{}{}^{4}R_{}^{\alpha \beta }{\displaystyle \frac{K_{\alpha \beta }}{w}}{\displaystyle \frac{1}{2}}\delta _\nu ^\mu {}_{}{}^{4}R{\displaystyle \frac{K}{w}}\right)`$ (B17)
$`+`$ $`\mathrm{},`$ (B18)
where $`Tr(K^3)K_\beta ^\alpha K_\gamma ^\beta K_\alpha ^\gamma `$. In the case the metric has the form
$$ds^2=dw^2n^2(\tau ,w)d\tau ^2+S^2(\tau ,w)[d\chi ^2+f_k^2(\chi )(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)]$$
we have
$$H_\tau ^\tau =12L(\tau ,w)\frac{K_\chi ^\chi }{w}+\mathrm{}\text{with}L(\tau ,w)(K_\chi ^\chi )^2+\frac{\dot{S}^2+kn^2}{n^2S^2}$$
(B19)
and where $`K_\chi ^\chi =S^{}/S`$, a prime denotes $`/w`$. Outside the membrane, spacetime is anti-de Sitter, $`n`$ and $`S`$ are given by BDL , and $`L(\tau ,w)L_\pm `$ (as an explicit calculation shows). In the vicinity and inside the membrane the function $`L(\tau ,w)`$ is either continous with discontinuous $`w`$-derivative (case of branes) or discontinous (case of shells) and can be modelled by the expressions (46)-(47) in the text.
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# Remarks on scaling properties inherent to the systems with hierarchically organized supplying network
## Abstract
We study the emergence of a power law distribution in the systems which can be characterized by a hierarchically organized supplying network. It is shown that conservation laws on the branches of the network can, at some approximation, impose power law properties on the systems. Some simple examples taken from economics, biophysics etc. are considered.
The existence of power law distribution in different systems of nature is a fascinating property still waiting for sufficient explanation. Despite a very long history of investigation of different scaling laws in systems of outworld this problem is very attractive up to the present time. A high interest to this problem is associated primarily with the fact that the systems, different in their nature, possess the similar power law distributions. This in its turn indicates that complex systems which consist of tremendous amount of objects have universal behavior, whose investigation gives the possibility to penetrate deeply into the essence of organization of nature.
V.Pareto (1897, ) seems to be first who discussed the power law distributions in wealth of the individuals in a stable economy. Since that time more than one century passed and the scaling properties are already discovered in many systems. Now we meet power law distributions in turbulence , in biological systems , economics and finance , social phenomena , linguistic phenomena , systems which can be characterized by network organization like river network , etc.
In this paper we single out a certain subclass of systems widely met in the nature which obey power law. Main characteristic feature of such systems is their hierarchical organization. The hierarchical organization allows such system to be divided into sets of subsystems (which are called levels) involving many elements that are similar in their properties. The elements of the different levels are substantially different in their characteristics.
We consider systems that can exist only due to the permanent flow through the system of some transport agents joining different subsystems and elements into a whole organization. Such interconnection of systems elements is formed by hierarchically organized supplying network. The specific function of the network is to drain the systems or supply them with some transport agents needed for systems activity. The characteristic feature of such system consists of conservation laws which should be fulfilled on each hierarchy level. This denotes that total flow of the transport agents through the system is constant.
We represent the system under consideration by a simple model schematically drawn on Fig. (1,a). Each hierarchy level $`i`$ consists of set of $`n_i`$ independent elements. All these elements take part in transferring transport agent flow from one level to another, namely taking it from the elements belonging to the lower level, and transferring it to the elements of higher level and
$$n_{i+1}>n_i>n_{i1};i.$$
Due to the transport agent flow being constant through each hierarchy level $`i`$ the total flow $`X_{\text{total}}`$ is constant and
$$X_{\text{total}}=\underset{j}{}X_{i,j}n_iX_i,i,$$
(1)
here $`n_i`$-is the number of system elements on the $`i`$-th hierarchy level, $`X_{i,j}`$-the flow through the $`j`$-th element of the $`i`$-th hierarchy level, $`X_i`$-the average over level flow through each element of level $`i`$. It was supposed that all the elements of the same hierarchy have similar characteristics and it is possible to sum up their flows directly. In the framework of this article we pay attention to the fact that the power law dependence is inherent to different natural systems and their occurrence is naturally connected to the function of systems and is the intrinsic property of the systems.
The flow $`X_i`$ though each element is determined by characteristics of the properties of the element and system organization. We denote all these characteristics by some parameters $`A_i`$. Namely these parameters are observable features of systems under consideration. It is conventional way to investigate power law distribution of these parameters. In this case from (1) we get an obvious relation
$$n_i1/X_i(A_i).$$
(2)
We will touch some examples of model to determine $`A_i`$.
Now consider a simple example taken from economics. As the economic systems exist only with the streams of commodity and money, they possess all the above-listed properties.
The number of parameters, which define the economic market state, at first glance seems to be actually endless because there is a tremendous amount of goods on the market. At the same time a relatively small number of raw materials in the industry shows that there should be a large hierarchical systems in the market structure. So it is possible to present an economic system as a hierarchical system which hierarchy levels consist of the firms of different power. On the first level there are firms processing row materials, on the last level – the shops of retail trade. In this context we point out that after reaching the consumers goods the flows transform into money flows are going in the opposite direction. The schematic representation of such system is shown in fig. (1,b). We shall consider the situations, when the market of the considered goods exist, i.e. all the streams should be more than zero.
For simplicity we confine our consideration to some commodity market, in which producers, dealers and consumers of some kinds of goods will form a common connected network. For example, the market of steel or meat items and the market of furniture can interact with each other mainly through the changes in financial states of the whole set of the consumers. In this case we can consider some branch of industry as a single one which practically does not interact with any other branches of industry . <sup>1</sup><sup>1</sup>1It should be noted that this restriction is not so strong. Main conclusion of this article will be true for the whole economic system (for example, the economic system of the state), including many sets of material flows twisted with each other. Indeed, in this case we can expend total material flow onto the basis of different material flows. For each of them situation will be practically the same as listed below.
The level $`i`$ of this system is a set of $`n_i`$ independent firms, bringing out the products of some sort, which buy the products from firms of the lower level of the hierarchy, and sell the results of their activity to the firms of the higher level. Firms, forming the common level $`i`$, are considered as identical from the viewpoint of their power. The bottom of the system (the firms of the first level) is formed by the firms, obtaining and processing raw material. The market participants of the last level are the “points” of retail trade, supplying the consumers with required goods.
The level of the firm production on each level $`i`$ is measured in the units of the initial raw material flow $`X_i`$, “flowing” through the given firm. It should be noted, that in this system of units the value $`X`$ having dimensionality (material)/(time). Particular interconnection between different firms can occur and disappear during the formation and evolution of the market under consideration. Their interaction is governed by trade with firms belonging to the neighbor levels.
For simplicity, consider some hierarchical system which can be represented as a homogeneous tree. Each firm is represented by the node on the tree. One branch enter this node and $`a`$ branches go out. Such representation denotes, that the given firm connected with trade relations with one firm of the lower level and $`a`$firms of the higher level. It is assumed that tree consist of $`N`$ levels. In this case number of firms on $`i`$-th level - $`n_i`$ \- equal
$$n_i=a^{i1}.$$
(3)
The fraction of the system, which make up the $`i`$-th level is
$$p(i)=\frac{n_i}{S_N},$$
(4)
where $`S_N`$ \- number of nodes in the tree (firms in our system)
$$S_N=\underset{i=1}{\overset{N}{}}n_i=\frac{a^N1}{a1}.$$
(5)
The fraction $`p(i)`$ may be regarded as the probability that chosen at random firm belong to the $`i`$-th level. Cumulative probability defined as the probability that chosen at random firm belong to one of levels from the 1-st to the $`i`$-th
$$P(i)=\frac{S_i}{S_N}=\frac{an_i1}{a^N1},$$
(6)
where $`S_i=_{k=1}^in_k`$.
As it was mentioned above $`n_i=X_{\text{total}}/X_i`$ due to this fact probability to choose node with flow $`x`$ less or equal to $`X_i`$ (cumulative distribution function) is
$$P(xX_i)=\frac{a(X_{\text{total}}/X_i)1}{a^N1}\frac{1}{X_i}.$$
(7)
Finally, it is important to express flow $`X_i`$using the dependence of the parameter $`A_i`$ from the flow $`X_i`$. There is a question what the parameter $`A_i`$ means. For example let us consider the parameter $`A_i`$as a firm’s profit $`\pi _i`$. It is natural to suppose that profit of each firm is proportional to the material or financial flow through the firm, namely $`\pi _iX_i^\alpha `$ $`(\alpha >0)`$. Then we can write cumulative probability distribution of firm’s profit for $`n_i1`$
$$P(\pi _i)\frac{1}{\pi _i^{1/\alpha }}.$$
(8)
In the most intuitive case $`\alpha =1`$ and we obtain Zipf’s law distribution as it was observed in . It seems to be reasonable to note that distribution of the firms from their income analyzed in the is a direct consequence of the hierarchical organization of the firms into corresponding industrial branches, which are connected to each other by material and financial flows.
As an another typical example of the systems under consideration we may regard living tissue where blood, flowing through the vascular network, which involves arterial and venous beds, supplies cellular tissue with oxygen, nutritious products, etc. At the same time blood withdraws carbon dioxide and products resulting from life activities of the cellular tissue. Both the arterial and venous beds being of the tree form contain a large number of hierarchy levels and are similar in structure. A similar situation takes place in respiratory systems where oxygen going through hierarchical system of bronchial tubes reaches small vessels (capillaries).
A microcirculatory bed of living tissue can be reasonably regarded as a space-filling fractal, being a natural structure for ensuring that all cells are serviced by capillaries . The vessel network must branch so that every small group of cells, referred below to as “elementary tissue domain”, is supplied by at least one capillary. In this case a vessel network generated by arteries should contain a sufficiently large number of hierarchy levels At each level $`i`$ of the vascular network the tissue domain supplied by a given microcirculatory bed as a whole can be approximated by the union of the tissue subdomains whose mean size is about the typical length $`l_i`$ of the $`i`$-th level vessels. Thus, the individual volume of these subdomains is estimated as $`V_il_i^3`$. From this simple consideration immediately follows the sufficiently obvious power law distribution for vascular network. Indeed, for the space-filling vascular network this artery supplies with blood the tissue region of volume about $`l_i^3`$ and, so, under normal conditions the blood flow rate $`X_i`$ in it should be equal to
$$X_ijl_i^3,$$
(9)
where $`j`$ is the blood perfusion rate (the volume of blood, flowing through the tissue domain of volume unit per time unit) assumed to be the same at all the points of the given microcirculatory bed. In living tissue the ratio $`l_i/a_i`$ takes a certain fixed value, $`l_i\mathrm{constant}a_i`$ and, thus, $`X_i\mathrm{constant}^{}a_i^3`$. Due to the blood conservation at branching nodes the scaling law of such system of vessels controlling finally the blood flow redistribution over the microcirculatory beds can be represented as
$$n_iconst/a_i^3const/l_i^3$$
(10)
In the same way we can obtain approximate power law dependance for river network organization. Recapitulating practically one-to-one the explanation listed above, for the hierarchically river network as two-dimensional space-filling fractal we can write
$$X_ijl_i^2.$$
In this case from equation (1) we get
$$n_iconst/l_i^2const/l_i^2.$$
(11)
The obtained power law dependence is approximately valid for qualitative understanding of scaling properties so as the results obtained in framework of this approach correspond to that one, obtained by more rigorous consideration .
From the described examples we can develop a general mathematical model of the system under consideration consisting from the distributed basic medium $`𝕄`$ and the transport hierarchical network $`S`$ (Fig. 1,b). The transport network can supply basic medium by some transport agent or can drain basic medium. The medium $`𝕄`$ is a $`d`$ \- dimensional homogeneous continuum. The transport network can be reasonably regarded as a space-filling fractal embedded in $`d`$-dimensional homogeneous continuum.
$$X_ijl_i^d$$
where $`j`$ is the flow of transport agents (the volume of transport agent through some domain of unit volume per unit time) assumed to be the same at all the points of the given bed. In this case $`n_iconst/l_i^d,`$ and the dimensionality of homogeneous continuum should be determined. In the given examples such basic medium has dimensions $`d=1`$, $`d=3`$ and $`d=2`$ respectively to (9),(10),(11). In general case such dimensions can be noninteger and their properties seem to determine the properties of basic medium $`𝕄`$ as a whole.
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# Twin wall of cubic-tetragonal ferroelastics
\[
## Abstract
We derive solutions for the twin wall linking two tetragonal variants of the cubic-tetragonal ferroelastic transformation, including for the first time the dilatational and shear energies and strains. Our solutions satisfy the compatibility relations exactly and are obtained at all temperatures. They require four non-vanishing strains except at the Barsch-Krumhansl temperature $`T_{BK}`$ (where only the two deviatoric strains are needed). Between the critical temperature and $`T_{BK}`$, material in the wall region is dilated, while below $`T_{BK}`$ it is compressed. In agreement with experiment and more general theory, the twin wall lies in a cubic 110-type plane. We obtain the wall energy numerically as a function of temperature and we derive a simple estimate which agrees well with these values.
\]
Ferroelastic transformations are diffusionless, first-order, shape-changing phase changes in the solid state. In cubic-tetragonal (C-T) systems like Nb<sub>3</sub>Sn, V<sub>3</sub>Si, In-Tl alloys, Fe-Pd alloys and Ni<sub>2</sub>MnGa, the cubic unit cell elongates (or contracts) along one of three axes to form a tetragonal unit cell; below the transition temperature $`T_c`$ there are three possible homogeneous products (variants) differing only in orientation.
Barsch and Krumhansl (BK) obtained an analytical solution for the twin wall linking two tetragonal variants of C-T ferroelastics. The dilatational and shear strains vanish identically, and the two remaining strains are functions of a single coordinate. The wall lies in a cubic 110-type plane, in agreement with experiment and as already known on more general grounds.
The BK solution is valid, however, at only a single temperature ($`T=T_{BK}`$). At any other $`T`$, the dilatational and shear strains are not zero, and the only known method to find the C-T twin wall structure requires solving the full three-dimensional (3D) partial differential equations (previous attempts at a 1D solution omitted the non-deviatoric strains). Some numerical solutions of these equations have been obtained, but no results have been given for the wall structure.
That is, 16 years after BK, the wall structure is still unknown except at $`T_{BK}`$.
The following solves this long-standing problem, which is of considerable physical interest given the large strains and large magnetic-field effects in Ni<sub>2</sub>MnGa. Specifically, we present a 1D solution for the C-T twin wall at all $`T`$. Our solutions, which include dilatational and shear energies and strains, satisfy the compatibility relations by virtue of analytical relations that we derive between the strains; these relations allow us to reduce the problem to the solution of ordinary, rather than partial, differential equations. We recover the BK solution at $`T_{BK}`$ and we present results at both higher and lower $`T`$.
The three paragraphs immediately following define the strains and the two parts of the free-energy density. The next two paragraphs obtain the key new results of our analysis, namely two relations between the strains, and the resulting Euler-Lagrange equations. The remaining paragraphs discuss the results from solution of these equations. We find that the twin-wall region is dilated near $`T_c`$, and compressed below $`T_{BK}`$.
Six strains are required to describe C-T ferroelastics, the dilatational strain $`e_1`$, the deviatoric strains $`e_2`$ and $`e_3`$, and the shear strains $`e_4`$, $`e_5`$ and $`e_6`$. With the coordinate axes along the four-fold axes, and in the small-strain approximation, these are
$$\begin{array}{cc}e_1=\left(u_{1,1}+u_{2,2}+u_{3,3}\right)/\sqrt{3},\hfill & e_2=\left(u_{1,1}u_{2,2}\right)/\sqrt{2},\hfill \\ e_3=\left(u_{1,1}+u_{2,2}2u_{3,3}\right)/\sqrt{6},\hfill & e_6=\left(u_{1,2}+u_{2,1}\right)/2\hfill \end{array}$$
(1)
plus obvious expressions for $`e_4`$ and $`e_5`$. Here $`𝐮=(u_1,u_2,u_3)`$ is the displacement of the material point originally at $`𝐱`$, and $`u_{i,j}=_ju_i=u_i/x_j`$. We need also the components $`\omega _3=(u_{1,2}u_{2,1})/2,\mathrm{𝑒𝑡𝑐}.`$ of the local rotation $`𝝎`$.
The free energy $`F`$ is the integral $`F=_Vd^3x`$ of the free-energy density $``$ over the undeformed volume $`V`$. For proper ferroelastics (where the strain is the primary order parameter), $``$ is the sum of strain and strain-gradient parts. The strain part is
$$_s=\frac{A_1}{2}e_1^2+\frac{A_2}{2}\left(e_2^2+e_3^2\right)\frac{B_2}{3}\left(e_3^33e_2^2e_3\right)$$
$$+\frac{C_2}{4}\left(e_2^2+e_3^2\right)^2+\frac{A_4}{2}\left(e_4^2+e_5^2+e_6^2\right).$$
(2)
In Voigt notation, the dilatational and shear constants are $`A_1=C_{11}+C_{12}`$ and $`A_4=4C_{44}`$ respectively, with both $`>0`$ for stability; the corresponding terms in the density were omitted in previous treatments. The coefficient $`A_2`$ depends on temperature as $`A_2=A_2^{}(TT_0)`$, where $`T_0`$ is the stability limit of the cubic phase and $`A_2^{}`$ is a material-dependent constant; above $`T_c`$, $`A_2=C_{11}C_{12}`$. For $`A_2>\frac{1}{4}B_2^2/C_2`$, the energy has only the cubic minimum (all strains zero). For $`A_2<\frac{1}{4}B_2^2/C_2`$, there are in addition three degenerate minima symmetrically located in the plane of the deviatoric strains:
$`e_2=0,e_3`$ $`=`$ $`e_{30}`$ (3a)
$`e_2=\sqrt{3}e_{30}/2,e_3`$ $`=`$ $`e_{30}/2`$ (3b)
$`e_2=\sqrt{3}e_{30}/2,e_3`$ $`=`$ $`e_{30}/2`$ (3c)
with $`e_1=e_4=e_5=e_6=0`$ and
$$e_{30}=\left[B_2+\left(B_2^24A_2C_2\right)^{1/2}\right]/\left(2C_2\right);$$
(4)
the tetragonal four-fold axes are in the 3, 1 and 2 directions respectively. The C-T transition, which is first-order, occurs at $`A_2=\frac{2}{9}B_2^2/C_2`$ where $`e_{30}=\frac{2}{3}B_2/C_2`$.
The free-energy density requires also strain-gradient terms, so that energy is required to introduce variant-variant walls (otherwise the system can subdivide into arbitrarily fine variants). We keep only the two invariants quadratic in the deviatoric strain derivatives:
$$_{sg}=\frac{d_2}{2}\left[\left(e_{2,1}^{}\right)^2+\left(e_{2,2}^{\prime \prime }\right)^2+\left(e_{2,3}\right)^2\right]$$
$$+\frac{d_3}{2}\left[\left(e_{3,1}^{}\right)^2+\left(e_{3,2}^{\prime \prime }\right)^2+\left(e_{3,3}\right)^2\right]$$
(5)
where $`e_2^{}`$, $`e_2^{\prime \prime }`$, $`e_3^{}`$ and $`e_3^{\prime \prime }`$ are obtained from $`e_2`$ and $`e_3`$ by $`2\pi /3`$ rotations about the cubic 111 axis:
$$\begin{array}{cc}\hfill e_2^{}=\left(u_{2,2}u_{3,3}\right)/\sqrt{2}=& \left(e_2+\sqrt{3}e_3\right)/2,\hfill \\ \hfill e_2^{\prime \prime }=\left(u_{3,3}u_{1,1}\right)/\sqrt{2}=& \left(e_2\sqrt{3}e_3\right)/2,\hfill \\ \hfill e_3^{}=\left(u_{2,2}+u_{3,3}2u_{1,1}\right)/\sqrt{6}=& \left(e_3\sqrt{3}e_2\right)/2,\hfill \\ \hfill e_3^{\prime \prime }=\left(u_{3,3}+u_{1,1}2u_{2,2}\right)/\sqrt{6}=& \left(e_3+\sqrt{3}e_2\right)/2.\hfill \end{array}$$
(6)
Both terms are transparently invariant and non-negative, and so we have the stability requirements $`d_20`$ and $`d_30`$ (which differ from those in Ref.); contact with previous treatments is made by writing $`d_2=g_2g_3`$ and $`d_3=g_2+g_3`$, resulting in
$$_{sg}=\frac{1}{2}g_2[(\stackrel{}{}e_2)^2+(\stackrel{}{}e_3)^2]+\frac{1}{2}g_3\{\frac{1}{2}[(e_{2,1})^2(e_{3,1})^2]$$
$$+\frac{1}{2}[(e_{2,2})^2(e_{3,2})^2][(e_{2,3})^2(e_{3,3})^2]$$
$$+\sqrt{3}(e_{2,1}e_{3,1}e_{2,2}e_{3,2})\}.$$
(7)
We seek the solution linking the variants with four-fold axes in the 1 and 2 directions, Eqs. (3b) and (3c); the results for this pair are simpler than for the others. The method is easily extended to treat a twin band. We assume a solution with $`e_4`$, $`e_5`$, $`\omega _1`$ and $`\omega _2=0`$ all identically zero, and with the other four strains (and $`\omega _3`$) independent of $`x_3`$. The strains are not independent for physical settings like a twin wall where the strains depend on position; rather, they are linked by the compatibility relations (necessary and sufficient conditions that the strains be derivable from the displacement). The nine first-order relations, of the form $`u_{i,jk}=u_{i,kj}`$, involve the first derivatives of the strains and the rotation components $`\omega _i`$. The more familiar second-order relations, which involve the second derivatives of the strains, are easily obtained by differentiation to eliminate the $`\omega _i`$. By virtue of the above assumptions, only the relations
$$\begin{array}{cc}\hfill _2\left(\sqrt{2}e_1+\sqrt{3}e_2+e_3\right)/\sqrt{6}& =_1\left(e_6+\omega _3\right),\hfill \\ \hfill _1\left(\sqrt{2}e_1\sqrt{3}e_2+e_3\right)/\sqrt{6}& =_2\left(e_6\omega _3\right),\hfill \\ \hfill _1\left(e_1\sqrt{2}e_3\right)& =0,\hfill \\ \hfill _2\left(e_1\sqrt{2}e_3\right)& =0,\hfill \end{array}$$
(8)
need be considered. We try for functions of $`X=x_1\mathrm{cos}\beta +x_2\mathrm{sin}\beta `$ alone; a similar 1D solution is not possible for a cubic-tetragonal soliton. The compatibility relations are then
$`\mathrm{sin}\beta \left(\sqrt{2}e_1+\sqrt{3}e_2+e_3\right)/\sqrt{6}`$ $``$ $`\mathrm{cos}\beta \left(e_6+\omega _3\right)=K_1,`$ (9)
$`\mathrm{cos}\beta \left(\sqrt{2}e_1\sqrt{3}e_2+e_3\right)/\sqrt{6}`$ $``$ $`\mathrm{sin}\beta \left(e_6\omega _3\right)=K_2,`$ (10)
$`e_1`$ $``$ $`\sqrt{2}e_3=K_3.`$ (11)
The constants $`K_1`$, $`K_2`$ and $`K_3`$ are evaluated from the boundary conditions at $`X=\pm \mathrm{}`$, namely $`e_2=\pm \frac{1}{2}\sqrt{3}e_{30}`$, $`e_3=\frac{1}{2}e_{30}`$, $`e_1=e_6=0`$, and $`\omega _3=\pm \mathrm{\Omega }`$. One finds easily that a solution is possible only if $`\mathrm{cos}2\beta =0`$, and so $`X=(x_1\pm x_2)/\sqrt{2}`$; that is, the walls lie in the $`110`$ or $`1\overline{1}0`$ planes. In this way, we find the key new results relating $`e_1`$ and $`e_6`$ to the deviatoric strains:
$`e_1(X)`$ $`=`$ $`\sqrt{2}\left[e_3(X)+e_{30}/2\right],`$ (12)
$`e_6(X)`$ $`=`$ $`\pm \sqrt{3/2}\left[e_3(X)+e_{30}/2\right],`$ (13)
plus $`\omega _3(X)=\pm e_2(X)/\sqrt{2}`$. These results are independent of the details of the free-energy density (they apply whether or not the dilatational and shear energies appear in $``$).
Since the compatibility relations are satisfied, we can use the density
$`_{twin}=`$ $`{\displaystyle \frac{A_{ds}}{2}}\left(e_3+{\displaystyle \frac{e_{30}}{2}}\right)^2+{\displaystyle \frac{A_2}{2}}\left(e_2^2+e_3^2\right){\displaystyle \frac{B_2}{3}}\left(e_3^33e_2^2e_3\right)`$ (15)
$`+{\displaystyle \frac{C_2}{4}}\left(e_2^2+e_3^2\right)^2+{\displaystyle \frac{D_2}{2}}\left({\displaystyle \frac{de_2}{dX}}\right)^2+{\displaystyle \frac{D_3}{2}}\left({\displaystyle \frac{de_3}{dX}}\right)^2`$
and minimize $`F`$ with respect to $`e_2`$ and $`e_3`$. Here $`A_{ds}=2A_1+\frac{3}{2}A_4`$, $`D_2=\frac{1}{4}d_2+\frac{3}{4}d_3`$, and $`D_3=\frac{3}{4}d_2+\frac{1}{4}d_3`$; adding the diagonal invariants $`\frac{1}{2}d_1(\stackrel{}{}e_1)^2`$ and $`\frac{1}{2}d_4[(\stackrel{}{}e_4)^2+(\stackrel{}{}e_5)^2+(\stackrel{}{}e_6)^2]`$ (with $`d_1`$ and $`d_4`$ $`>0`$) to the density of Eq. (5) leaves $`D_2`$ unchanged while adding $`2d_1+\frac{3}{2}d_4`$ to $`D_3`$. The corresponding Euler-Lagrange equations are
$$A_2e_2+2B_2e_2e_3+C_2e_2(e_2^2+e_3^2)=D_2d^2e_2/dX^2,$$
(12a)
$`A_{ds}(e_3+e_{30}/2)+A_2e_3+B_2(e_2^2e_3^2)+C_2e_3(e_2^2+e_3^2)`$ (12b)
$`=D_3d^2e_3/dX^2.`$ (12c)
The term $`A_{ds}(e_3+e_{30}/2)`$ in Eq. (12c), new with this article, results from satisfying the compatibility relations. The same equations are obtained, after integrations, on using the density of Eqs.(2) plus (3), demanding that $`F`$ be stationary with respect to the displacement, and only then using Eqs.(10). The boundary conditions are
$$e_2(\pm \mathrm{})=\pm \sqrt{3}e_{30}/2,e_3(\pm \mathrm{})=e_{30}/2.$$
(13)
At $`A_2=2B_2^2/C_2`$, which defines the temperature $`T_{BK}`$, the solutions are
$$e_2=(\sqrt{3}e_{30}/2)\mathrm{tanh}(\kappa X),e_3=e_{30}/2$$
(14)
with $`\kappa ^2=3B_2^2/(2C_2D_2)`$; the $`A_{ds}`$ term and the dilatational and shear strains all vanish identically. We note that $`T_{BK}`$ may possibly be identified experimentally as the temperature where $`e_{30}`$ ($`=2B_2/C_2`$) is three times the value at $`T_c`$.
At other temperatures, the strains $`e_1`$ and $`e_6`$ are not zero, and the equations must be solved numerically. We define the reduced temperature $`\tau =(TT_0)/(T_cT_0)`$; then $`\tau =1`$ at the transition and $`\tau =9`$ at $`T=T_{BK}`$. We also take $`D_3=D_2`$.
Figure 1 shows $`e_2`$ and $`e_3`$ as functions of $`X`$ at three different temperatures, $`T_c`$ (which is $`>T_{BK}`$), $`T_{BK}`$ and $`T<T_{BK}`$, as determined from numerical solution of Eqs. (12) and (13). The value $`A_{ds}=A_2^{}(T_cT_0)=A_2(T=T_c)`$, which is rather soft, was chosen for display purposes; for larger values, $`e_3`$ remains close to $`e_{30}/2`$. One sees that $`e_3+\frac{1}{2}e_{30}`$ is $`>0`$ or $`<0`$ for $`T>T_{BK}`$ or $`<T_{BK}`$ respectively. From Eqs. (10) then, the wall region is dilated near $`T_c`$ and compressed below $`T_{BK}`$.
To estimate the size of these effects, we use data for FePd alloys. The data quoted in Ref. give $`T_c=268.6`$K, $`T_0=265`$K and $`T_{BK}=233`$K. Combining these with data from Refs. , we find that $`A_{ds}/A_2(T=T_c)`$ has a lower limit of $`200`$ based on the smallest observed values for $`A_2`$; we use a more conservative estimate of 400. Between $`T_c`$ and $`T_{BK}`$, the volume change $`\mathrm{\Delta }V/V`$ at the centre of the wall reaches a maximum at $`T258`$ K, with a very small value $`10^4`$. At $`T=0`$ we find $`\mathrm{\Delta }V/V10^2`$.
Figure 2 shows twin-wall trajectories in the $`(e_3,e_2)`$ plane. The trajectories bow toward the third variant for $`T_{BK}<T<T_c`$, and away for $`T<T_{BK}`$. They shift toward the vertical for larger $`A_{ds}`$.
Reference, in effect, assumed that $`A_1=A_6=0`$ and so the term $`A_{ds}(e_3+e_{30}/2)`$ was absent from their differential equations. Solution of these equations, done only near $`T_c`$, gave trajectories which passed close to the origin. We point out however that the origin is not the cubic state but rather a highly dilated, highly sheared state with $`e_1=e_{30}/\sqrt{2}`$ and $`e_6=\pm \sqrt{3/8}e_{30}`$, as seen from Eqs.(10). Reference, which also assumed in effect that $`A_1=A_6=0`$, proposed a trajectory that is a $`2\pi /3`$ circular arc centered at the origin. As evident from Figure 2, this is possibly useful only for $`TT_{BK}`$.
The wall energy $`W`$ (per unit area) is the energy required to form an interface between two variants:
$$W=_{\mathrm{}}^{\mathrm{}}\left(_{twin}_h\right)𝑑X$$
(15)
where $`_{twin}`$ is the density (11) for the twin-wall solution of Eqs. (12) and (13), and
$$_h=\frac{1}{2}A_2e_{30}^2\frac{1}{3}B_2e_{30}^3+\frac{1}{4}C_2e_{30}^4$$
(16)
is the density for a single variant. Although not directly observable, $`W`$ has physical content and so we provide the following.
Figure 3 shows the wall energy $`W`$ as a function of temperature, as determined numerically for three different values of $`A_{ds}`$; it shows also two variational approximations which we now obtain.
If we consider Eq.(14) as a trial function, with $`\kappa `$ an adjustable parameter, we find
$$We_{30}^3\sqrt{3C_2D_2/8}$$
(17)
at the optimal $`\kappa `$, namely $`\kappa =e_{30}\sqrt{3C_2/(8D_2)}`$; the $`T`$ dependence is in $`e_{30}`$. From Figure 3, this does very well over the temperatures examined, not the least because it is an equality at $`T_{BK}`$ for all $`A_{ds}`$. The coefficients $`A_{ds}`$ and $`D_3`$ do not appear because of the form of our trial function; explicitly, in obtaining Eq.(17), we did not drop the corresponding terms in the density, and so Eq.(17) is valid independent of the magnitude of the dilatational and shear energies and strains. Equation (17) seems to be an equality also in the limit $`A_{ds}\mathrm{}`$ for all $`T`$.
The wall energy was estimated previously, using the circular trajectory described above. On setting $`A_1=A_6=A_{ds}=0`$, taking $`D_3=D_2`$, and using $`e_3=e_{30}\mathrm{cos}\varphi `$, $`e_2=e_{30}\mathrm{sin}\varphi `$, with $`\frac{2\pi }{3}\varphi \frac{4\pi }{3}`$, we find
$$W\sqrt{64B_2D_2e_{30}^5/27}.$$
(18)
Unlike Eq.(17), this assumes that $`A_1=A_6=0`$; including dilatation and shear only increases the value relative to Eq.(17). Figure 3 shows that Eq.(18) is considerably poorer than Eq.(17) for the temperatures examined; it is a factor of 3 larger at $`T_c`$. But the trajectories bend toward the Hong-Olson arc at lower $`T`$, and so Eq.(18) increases with decreasing $`T`$ less strongly than Eq.(17) (as $`e_{30}^{5/2}`$ rather than as $`e_{30}^3`$). Indeed, for $`\tau <151`$, Eq.(18) is actually better than Eq.(17); but this temperature may not be accessible since the strain there is $`10`$ times that at $`T_c`$. For Fe<sub>70</sub>Pd<sub>30</sub>, $`\tau =151`$ is inaccessible ($`T<0`$).
In summary, we have found for the first time the solution for the C-T twin wall, at all $`T`$, in the physical case with dilatational and shear energies and strains. The dilatation and shear strains, which are localized near the interface, change sign at $`T_{BK}`$. The magnitudes (at most 1% in Fe-Pd alloys) may however be too small to detect.
###### Acknowledgements.
This research was supported by the Natural Sciences and Engineering Research Council (Canada).
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# 1 Introduction
## 1 Introduction
To describe the motion of mechanical systems there is a variety of mathematical models which are based on different principles. Most of the physical models are obtained using an appropriate variational principle in a certain evolution space. But variational principles are not only important in physics, but also in many branches of engineering or economics , where one is interested in optimizing a given functional, possibly subject to some restrictions. In fact, constraints are ubiquitous in many mechanical systems and much more different situations.
In this paper we are going to study Lagrangian systems, i.e., dynamical systems in which the equations of motion are obtained by finding the critical paths of a functional
$$_{t_1}^{t_2}Ldt,$$
where $`L`$ is a function defined on the tangent bundle $`\mathrm{T}Q`$ of a given differentiable manifold $`Q`$, the configuration space. We will not consider arbitrary variations, but only variations satisfying certain conditions. These conditions arise from some given constraints on the dynamics of the system. We will analize the case when the constraints are defined by a certain submanifold $`C`$ of the tangent bundle $`\mathrm{T}Q`$. Such constraints are usually called nonholonomic.
There are two different approaches when dealing with constraints. The first one is based on the idea of understanding the constraints as constraint forces. This point of view, that seems very natural in a physical context, gives rise to the classical d’Alembert-Lagrange principle. Mechanics of Lagrangian systems with nonholonomic constraints based on this principle is called nonholonomic mechanics . But there is another different point of view, that seems more natural when one is interested in optimizing a given functional defined as above when there are constraints. For example, if we wish to change the state of a given system minimizing a cost functional (a typical problem in engineering or economics), it is not natural to understand the constraints as forces acting on the system. In this case one is interested in minimizing the functional considering only the variations allowed by the constraints. Mechanics of Lagrangian systems with nonholonomic constraints based on this idea is often called vakonomic mechanics (mechanics of variational axiomatic kind ). For example, optimal control theory is a typical example of vakonomic mechanics. It is interesting to notice that both mechanics do not coincide in general, but they agree when the constraints are integrable. Several references on these topics are , , and . See also .
The paper is organized as follows. In section 2 we give a generalized notion of a variational problem. Section 3 is devoted to study vakonomic mechanics from a geometric point of view, with special attention to the existence of admissible variations in order to obtain the equation of motion. In section 4 we do the same for nonholonomic mechanics. Both mechanics are understood as generalized variational problems. In section 5 it is proved that, when the constraints are integrable, vakonomic and nonholonomic mechanics coincide. In section 6, as an example of vakonomic mechanics, a geometric formulation of optimal control theory is studied. Finally, in section 7, we give a geometric formulation of vakonomic mechanics when the constraints are defined by a distribution (linear constraints in the velocities).
Basic knowledge of differential geometric structures is assumed. The presentation is almost self-contained but the interested reader may consult the bibliography for more specific topics as the vertical lift, the fibre derivative and the canonical involution , or the Euler-Lagrange operator .
## 2 Variational problems
### 2.1 Elements of a variational problem
First, we are going to define what we mean by a variational problem. A variational problem consists of the data $`(Q,L,C,𝒞,𝒲)`$ where:
* $`Q`$ is a $`n`$-dimensional differentiable manifold, the configuration space.
* $`L`$ is the lagrangian function defined on the tangent bundle, $`L:\mathrm{T}Q\text{}`$.
* $`C`$ is the constraint submanifold, and it is a submanifold of $`\mathrm{T}Q`$.
* $`𝒞`$ is the family of admissible paths. Given two points $`q_1,q_2Q`$ and a compact interval $`I=[t_1,t_2]`$, we will say that a path $`\gamma :IQ`$ of class $`C^2`$ is admissible if:
$`\gamma (t_1)=q_1`$, $`\gamma (t_2)=q_2`$ and
$`\dot{\gamma }(t)C`$, for all $`tI`$.
* $`𝒲`$ are the admissible variation fields (or infinitesimal variations). For a given admissible path $`\gamma `$, $`𝒲_\gamma `$ consists on a certain set of $`C^1`$ vector fields along $`\gamma `$.
Notice that we do not consider variations of a path $`\gamma `$, but variation fields along $`\gamma `$. Now we are ready to define the variational problem associated to $`(Q,L,C,𝒞,𝒲)`$. The action of $`L`$ along a path $`\gamma `$ is the functional $`S:𝒞\text{}`$ given by the integral
$$S[\gamma ]=_IL(\dot{\gamma }(t))dt.$$
(2.1)
A variational problem consists in finding the critical admissible paths of the functional $`S`$, in a sense that will be precised later.
### 2.2 Variations and variation fields
Let $`\gamma :IQ`$, $`\gamma (t_1)=q_1`$, $`\gamma (t_2)=q_2`$, be an admissible path. A variation of $`\gamma `$ is a $`\mathrm{C}^2`$ function $`\mathrm{\Gamma }:(\delta ,\delta )\times IQ`$ such that:
1. $`\mathrm{\Gamma }_\epsilon =\mathrm{\Gamma }(\epsilon ,)`$ is a one-parameter family of paths defined on $`I`$ with fixed end-points, $`\mathrm{\Gamma }(\epsilon ,t_i)=q_i`$, $`\epsilon (\delta ,\delta )`$, $`i=1,2`$, and
2. $`\mathrm{\Gamma }(0,t)=\gamma (t)`$, $`tI`$ (if there is no variation, $`\epsilon =0`$, we obtain the original path $`\gamma `$).
Given a function $`\mathrm{\Gamma }(\epsilon ,t)`$ of two real variables, we will denote its derivatives with respect to $`\epsilon `$ and $`t`$ as $`\mathrm{\Gamma }^{}`$ and $`\dot{\mathrm{\Gamma }}`$, respectively. It is clear that $`\mathrm{\Gamma }^{}`$ and $`\dot{\mathrm{\Gamma }}`$ are vector fields along $`\mathrm{\Gamma }`$.
Now, we are in conditions to define a variation field of $`\mathrm{\Gamma }`$.
###### Definition 1
Given a variation $`\mathrm{\Gamma }(\epsilon ,t)`$ of an admissible path $`\gamma `$, the variation field of $`\mathrm{\Gamma }`$ is the vector field $`𝐰`$ along $`\gamma `$ defined by
$$𝐰(t)=\mathrm{T}_{(0,t)}\mathrm{\Gamma }\frac{\mathrm{d}}{\mathrm{d}\epsilon }|_{(0,t)}=\mathrm{\Gamma }^{}(0,t).$$
Notice that
$$𝐰(t_1)=0,𝐰(t_2)=0,$$
since the $`\mathrm{\Gamma }_\epsilon `$ have fixed end-points.
Therefore, given a family of variations of $`\gamma `$, we can associate to them a family of variation vector fields along $`\gamma `$. We will say that a variation $`\mathrm{\Gamma }`$ of $`\gamma `$ is admissible if its associated variation vector field along $`\gamma `$, $`𝐰(t)=\mathrm{\Gamma }^{}(0,t)`$, is an admissible variation field of $`\gamma `$, i.e., $`𝐰𝒲_\gamma `$.
We finish this description about variations and variation fields with a useful lemma, whose proof is straighforward in local coordinates.
###### Lemma 1
For any $`𝐰`$ vector field along $`\gamma `$ and any function $`\lambda :I\text{}`$,
$$\text{s}(\lambda 𝐰)^.=(\mathrm{D}\lambda )\mathrm{vl}(\dot{\gamma },𝐰)+\lambda \text{s}\dot{𝐰}.$$
$`\mathrm{}`$
Notice that, if $`𝐰`$ is a vector field along $`\gamma `$, then $`\dot{𝐰}`$ is a vector field along $`𝐰`$ and $`\text{s}\dot{𝐰}`$ is a vector field along $`\dot{\gamma }`$. The function $`\lambda :I\text{}`$ denotes a function of time, and it is clear that $`\lambda 𝐰`$ is another vector field along $`\gamma `$. $`\mathrm{D}`$ is the usual derivative with respect to the time. The map $`\mathrm{vl}`$ denotes the vertical lift $`\mathrm{vl}:\mathrm{T}Q\times _Q\mathrm{T}Q\mathrm{T}(\mathrm{T}Q)`$. Its local expression is $`\mathrm{vl}(q,v,u)=(q,v;0,u)`$. Finally, $`\text{s}:\mathrm{T}(\mathrm{T}Q)\mathrm{T}(\mathrm{T}Q)`$ denotes the canonical involution, which is an isomorphism between the two vector bundle structures of $`\mathrm{T}(\mathrm{T}Q)`$. Its local expression is $`\text{s}(q,v;u,a)=(q,u;v,a)`$.
### 2.3 Critical admissible paths
###### Definition 2
An admissible path $`\gamma `$ is said to be critical if, for each admissible variation $`\mathrm{\Gamma }_\epsilon `$, the first variation of $`S[\mathrm{\Gamma }_\epsilon ]`$ is zero; i.e.,
$$\frac{\mathrm{d}}{\mathrm{d}\epsilon }S[\mathrm{\Gamma }_\epsilon ]|_{\epsilon =0}=0$$
for each admissible variation $`\mathrm{\Gamma }_\epsilon `$ of $`\gamma `$.
The main purpose of this paper consists in discussing the criticity conditions for different variational problems and describing the solutions. First of all, we are going to describe the criticity condition for a general problem.
It is clear that, if $`g:Q\text{}`$ is a function, then, for any function $`\mathrm{\Gamma }(\epsilon ,t)`$ ($`\mathrm{\Gamma }:U\text{}^2Q`$) of two real variables,
$$\frac{}{\epsilon }g(\mathrm{\Gamma }(\epsilon ,t))=\mathrm{d}g(\mathrm{\Gamma }(\epsilon ,t)),\mathrm{\Gamma }^{}(\epsilon ,t),$$
and similarly for $`/t`$.
Now, suppose that $`\mathrm{\Gamma }(\epsilon ,t)`$ is a variation of a path $`\gamma `$. Let us consider $`\dot{\mathrm{\Gamma }}:(\delta ,\delta )\times I\mathrm{T}Q`$. Derivation of $`\dot{\mathrm{\Gamma }}`$ with respect to $`\epsilon `$ and $`t`$ yields $`(\dot{\mathrm{\Gamma }})^{}`$ and $`(\dot{\mathrm{\Gamma }})^.`$, which are now vector fields along $`\dot{\mathrm{\Gamma }}`$. Taking $`\epsilon =0`$ yields two vector fields along $`\dot{\gamma }`$, which are $`\text{s}\dot{𝐰}`$ and $`\ddot{\gamma }`$.
Then, if $`f:\mathrm{T}Q\text{}`$ is a function, we have
$$\frac{}{\epsilon }|_{\epsilon =0}f(\dot{\mathrm{\Gamma }}(\epsilon ,t))=\mathrm{d}f(\dot{\gamma }(t)),\text{s}(\dot{𝐰}(t)).$$
(2.2)
(Remember that $`\dot{𝐰}`$ is a vector field along $`𝐰`$, and $`\text{s}\dot{𝐰}`$ is a vector field along $`\dot{\gamma }`$, so the contraction makes sense.)
Now, we can characterize the criticity condition in a more manageable way.
###### Proposition 1
Given a variational problem $`(Q,L,C,𝒞,𝒲)`$, an admissible path $`\gamma `$ is critical if and only if
$$_I\mathrm{d}L(\dot{\gamma }(t)),\text{s}(\dot{𝐰}(t))dt=0,$$
for each admissible vector field $`𝐰𝒲_\gamma `$.
Proof: Using (2.2) in (2.1), we obtain
$$\frac{\mathrm{d}}{\mathrm{d}\epsilon }S[\mathrm{\Gamma }_\epsilon ]|_{\epsilon =0}=_I\mathrm{d}L(\dot{\gamma }(t)),\text{s}(\dot{𝐰}(t))dt,$$
and the result follows. $`\mathrm{}`$
Observe that this condition does not depend on the full variation $`\mathrm{\Gamma }(\epsilon ,t)`$, but only on its variation field (see also ). Therefore, in our study of variational calculus, we will shift our attention to infinitesimal variations rather to finite variations.
### 2.4 The Euler-Lagrange operator
To obtain a more manageable condition of criticity, it is convenient to perform an integration by parts. First, let us define the Euler-Lagrange operator of $`L`$.
###### Definition 3
The Euler-Lagrange operator associated with a function $`L:\mathrm{T}Q\text{}`$ is a mapping $`_L:\mathrm{T}^2Q\mathrm{T}^{}Q`$ defined by the relation
$$_L\ddot{\gamma },𝐰=\mathrm{d}L\dot{\gamma },\text{s}\dot{𝐰}\mathrm{D}𝒟L\dot{\gamma },𝐰,$$
for any path $`\gamma :IQ`$ and vector field $`𝐰`$ along $`\gamma `$.
Here, the map $`𝒟L:\mathrm{T}Q\mathrm{T}^{}Q`$ is the fibre derivative of $`L`$. Recall that, given a vector bundle $`EB`$, if $`f:E\text{}`$ is a function, then the derivatives of the restrictions of $`f`$ to the fibres define the fibre derivative of $`f`$, which is a map $`𝒟f:EE^{}`$. Its local expression is $`𝒟f(b;a)=(b;f/a)`$.
It is easy to check (in coordinates) that the Euler-Lagrange operator is well-defined by this relation.
Therefore, the Euler-Lagrange operator is a one-form along the projection $`\mathrm{T}^2QQ`$, and also an affine bundle map along $`\mathrm{T}QQ`$. The expression in local coordinates of $`_L`$ is the usual one,
$$_L=\left(\frac{L}{q}\frac{\mathrm{d}}{\mathrm{d}t}\left(\frac{L}{v}\right)\right)\mathrm{d}q,$$
where $`\mathrm{d}/\mathrm{d}t`$ is the total time-derivative operator. The Euler-Lagrange operator can be extended in the same way to a time dependent Lagrangian.
Using the definition of the Euler-Lagrange operator, we can characterize the criticity condition in the usual form.
###### Theorem 1
Given a variational problem $`(Q,L,C,𝒞,𝒲)`$, an admissible path $`\gamma `$ is critical if and only if
$$_I_L(\ddot{\gamma }(t)),𝐰(t)dt=0,$$
for each admissible variation field $`𝐰𝒲_\gamma `$.
Proof: From proposition 1 and the definition 3 of the Euler-Lagrange operator we obtain that
$$\frac{\mathrm{d}}{\mathrm{d}\epsilon }S[\mathrm{\Gamma }_\epsilon ]|_{\epsilon =0}=_I_L(\ddot{\gamma }(t)),𝐰(t)dt+\left[𝒟L(\dot{\gamma }(t)),𝐰(t)\right]_{t_1}^{t_2}.$$
The result follows observing that, since $`𝐰`$ is a variation field, the last term vanishes ($`𝐰(t_1)=0`$, $`𝐰(t_2)=0`$). $`\mathrm{}`$
For the case when there are no constraints, $`C=\mathrm{T}Q`$ and $`𝒲_\gamma `$ is the set of all the vector fields along $`\gamma `$, we obtain the well-known Euler-Lagrange equation.
###### Corollary 1
Given a unconstrained variational problem, a path $`\gamma 𝒞`$ is critical if and only if
$$_L\ddot{\gamma }=0.$$
In this paper we will be interested in variational problems when $`C\mathrm{T}Q`$, $`C\mathrm{T}Q`$. Hence, given a set of admissible paths, it is necessary to select a set of admissible variation fields (or infinitesimal variations) along the admissible paths. We will consider two different approaches to this problem. The first one is vakononomic mechanics, which can be considered as a strictly variational approach. The second one is nonholonomic mechanics, which is variational in our generalized sense, but not in the classical one. Nonholonomic mechanics is the usual way to describe the dynamics of a mechanical system with constraints. In the next two sections we will describe the dynamical equations obtained in each case. It is interesting to remark that both approaches are equivalent when the constraints are integrable (holonomic constraints).
We finish this section with a useful property of the Euler-Lagrange operator that will be used in many calculations in the following. The proof is straighforward in local coordinates.
###### Lemma 2
For any $`f:\mathrm{T}Q\text{}`$ and $`\mu :I\text{}`$ (a function of time),
$$_{\mu f}=\mu _f(\mathrm{D}\mu )𝒟f\tau _{\mathrm{\hspace{0.17em}1}}^2,$$
where $`\tau _{\mathrm{\hspace{0.17em}1}}^2:\mathrm{T}^2Q\mathrm{T}Q`$ is the canonical projection. $`\mathrm{}`$
## 3 Vakonomic mechanics
Roughly speaking, vakonomic mechanics is the result of variational calculus when the variations are restricted by some constraints on the positions and also the velocities.
Our initial setting is therefore a submanifold $`C\mathrm{T}Q`$ of codimension $`m<n`$; let us denote by $`j`$ the inclusion of $`C`$ in $`\mathrm{T}Q`$. A constraint is any function $`\varphi `$ vanishing on $`C`$. Locally $`C`$ is defined by the vanishing of some constraints $`\varphi ^i:\mathrm{T}Q\text{}`$ ($`i=1,\mathrm{},m`$) whose differentials $`\mathrm{d}\varphi ^i`$ are linearly independent at each point of $`C`$.
We will assume that the projection of $`C`$ to $`Q`$, $`\tau _Qj:CQ`$, is a submersion. It can be easily proved that this statement is equivalent to say that the constraints $`\varphi ^i`$ can be chosen such that their fibre derivatives $`𝒟\varphi ^i`$ are linearly independent at every point of $`C`$. In local coordinates, this means that $`\varphi ^i/v^k`$ has maximal rank. That is, the constraints restrict the velocities, not the positions.
With the assumptions above, the image $`(\tau _Qj)(C)Q`$ is open, so we may assume that the projection $`CQ`$ is a surjective submersion. Then there exists the vertical subbundle $`\mathrm{V}(C)\mathrm{T}(C)`$, which has rank $`nm`$ (the dimension of the fibres of the submersion). Indeed, at each $`v_qC`$ we have $`\mathrm{V}_{v_q}(C)=\mathrm{T}_{v_q}(C)\mathrm{V}_{v_q}(\mathrm{T}Q)`$.
To obtain the equations of motion of vakonomic mechanics, we need first to describe which are the admissible variations.
### 3.1 The variations of vakonomic mechanics
We remember that an admissible path is a mapping $`\gamma :IQ`$ such that $`\dot{\gamma }`$ takes its values in the submanifold $`C\mathrm{T}Q`$. Due to our assumptions on $`C`$, there exist vector fields locally defined on $`Q`$ taking values in $`C`$. Their integral curves have their derivatives in
$`C`$, so there are many admissible paths.
###### Definition 4
Let $`\mathrm{\Gamma }`$ be a variation of an admissible path $`\gamma `$. The variation $`\mathrm{\Gamma }`$ is called a strongly admissible variation of $`\gamma `$ if every path $`\mathrm{\Gamma }_\epsilon `$ is admissible.
If $`\mathrm{\Gamma }`$ is a strongly admissible variation, then $`\varphi (\dot{\mathrm{\Gamma }}(\epsilon ,t))=0`$, for any constraint $`\varphi `$. Taking the derivative with respect to $`\epsilon `$ at $`\epsilon =0`$ and using (2.2) we have
$$\mathrm{d}\varphi \dot{\gamma },\text{s}\dot{𝐰}=0$$
for every constraint $`\varphi `$. This can also be expressed as
$$\text{s}(\dot{𝐰}(t))\mathrm{T}_{\dot{\gamma }(t)}(C)$$
for each $`tI`$.
###### Definition 5
A variation field $`𝐰`$ of an admissible path $`\gamma `$ is called an admissible variation field for a given vakonomic problem if
$$\text{s}(\dot{𝐰}(t))\mathrm{T}_{\dot{\gamma }(t)}(C),$$
that is, $`\mathrm{d}\varphi \dot{\gamma },\text{s}\dot{𝐰}=0`$.
From the definition of the Euler-Lagrage operator (definition 3), we obtain that $`𝐰`$ is a variation field if and only if
$$_\varphi \ddot{\gamma },𝐰=\mathrm{D}𝒟\varphi \dot{\gamma },𝐰$$
for every constraint $`\varphi `$.
It is important to remark that an admissible path may not have any nontrivial strongly admissible variation, and so an admissible variation field may not arise from a strongly admissible variation. One may say that the variations having admissible variation fields are the variations that preserve the constraints up to first order in $`\epsilon `$. These variations may be called weakly admissible variations.
Next, we are going to give a more detailed description of admissible variation fields. Among all the vector fields $`𝐰`$ along $`\gamma `$, we consider a particular submodule. Take the subbundle
$$\text{L}_\gamma ^C\gamma ^{}\mathrm{T}(Q)=I\times _\gamma \mathrm{T}(Q)$$
whose sections are the vector fields $`𝐰`$ along $`\gamma `$ of class $`C^1`$ whose vertical lifts $`\mathrm{vl}(\dot{\gamma },𝐰)`$ are tangent to $`C`$. Using this subbundle we can express the admissible variation fields in a more manageable way. First, notice that, since $`I`$ is an interval, both $`\text{L}_\gamma ^C`$ and $`\gamma ^{}\mathrm{T}(Q)`$ are trivializable. Therefore there exists a global frame for $`\gamma ^{}\mathrm{T}(Q)`$, $`(𝐰_k)`$ ($`k=1,\mathrm{},n`$). Since $`\text{L}_\gamma ^C`$ is a subbundle of rank $`nm`$, we can assume that the last $`nm`$ of the $`𝐰_k`$ span this subbundle.
Any vector field along $`\gamma `$ can be thus uniquely written $`𝐰=_{k=1}^n\lambda ^k𝐰_k`$, where $`\lambda ^k`$ are functions of time. Then $`𝐰`$ is a variation field if and only if the coefficients $`\lambda ^k`$ vanish at the end-points of $`I`$. Moreover, according to definition 5, $`𝐰`$ is an admissible variation field if it is a variation field and $`\mathrm{d}\varphi ^i\dot{\gamma },\text{s}\dot{𝐰}=0`$, for $`i=1,\mathrm{},m`$. Taking into account lemma 1, this condition can be written
$$\underset{k=1}{\overset{n}{}}𝒟\varphi ^i\dot{\gamma },𝐰_k\mathrm{D}\lambda ^k+\underset{k=1}{\overset{n}{}}\mathrm{d}\varphi ^i\dot{\gamma },\text{s}\dot{𝐰}_k\lambda ^k=0.$$
(3.1)
Notice that, since the fibre derivatives $`𝒟\varphi ^i`$ are linearly independent at each point, the matrix with entries $`𝒟\varphi ^i\dot{\gamma },𝐰_k`$ has maximal rank, $`m`$. By the special choice of the $`𝐰_k`$, the last $`nm`$ of them vanish under the $`𝒟\varphi ^i`$, and hence the square matrix $`A=(𝒟\varphi ^i\dot{\gamma },𝐰_j)_{i,j=1,\mathrm{},m}`$ is invertible. So, writing the equation as
$$\underset{j=1}{\overset{m}{}}A_j^i\mathrm{D}\lambda ^j+\underset{j=1}{\overset{m}{}}B_j^i\lambda ^j+\underset{l=m+1}{\overset{n}{}}C_l^i\lambda ^l=0,$$
(3.2)
we can isolate the $`\mathrm{D}\lambda ^j`$ ($`j=1,\mathrm{},m`$) linearly in terms of all the $`\lambda ^k`$ ($`k=1,\mathrm{},n`$). This determines uniquely $`\lambda ^j`$, $`j=1,\mathrm{},m`$ as functions of $`\lambda ^l`$, $`l=m+1,\mathrm{},n`$, due to the initial condition $`\lambda ^j(t_1)=0`$. However, not any $`\lambda ^l`$, $`l=m+1,\mathrm{},n`$ vanishing on $`t_1,t_2`$ are admissible. Notice that the solutions $`\lambda ^j`$, $`j=1,\mathrm{},m`$ must vanish also in $`t_2`$. In fact, the existence of solutions of (3.2) satisfying $`\lambda ^k(t_1)=\lambda ^k(t_2)=0`$, $`k=1,\mathrm{}n`$, is not guaranteed in principle. If we write (3.2) as
$$A\dot{\lambda }_{(1)}=B\lambda _{(1)}C\lambda _{(2)},$$
the solution satisfying $`\lambda _{(1)}(t_1)=0`$ is
$$\lambda _{(1)}(t)=\nu (t)_{t_1}^t[\nu (s)]^1A^1(s)C(s)\lambda _{(2)}(s)ds,$$
where $`\nu (t)`$ is the fundamental matrix of the homogeneous system $`A\dot{\lambda }_{(1)}=B\lambda _{(1)}`$ satisfying the initial condition $`\nu _i^j(t_2)=\delta _i^j`$. If $`\lambda _{(1)}(t_2)=0`$, then necessarily
$$_{t_1}^{t_2}[\nu (s)]^1A^1(s)C(s)\lambda _{(2)}(s)ds=0.$$
(3.3)
If the system is homogeneous ($`C=0`$) we obtain the trivial solution $`\lambda ^j(t)=0`$, $`j=1,\mathrm{},m`$, remaining $`\lambda ^l`$, $`l=m+1,\mathrm{},n`$ as arbitrary functions satisfying the boundary conditions $`\lambda ^l(t_1)=\lambda ^l(t_2)=0`$. As for the general case, in the following section we will show that admissible variations always exist in vakonomic mechanics.
Condition (3.1) is very useful to study variation fields in vakonomic mechanics, as we show in the following example.
Example Let $`Q=\text{}^2`$ be the configuration space, with coordinates $`(x,y)`$, and consider a Lagrangian function and a constraint both depending only on the velocities, i.e., $`L=L(\dot{x},\dot{y})`$ and $`\varphi =\varphi (\dot{x},\dot{y})`$.
From our assumptions on the constraints we can write locally $`\varphi =\dot{y}f(\dot{x})0`$. In this case, using theorem 2, it is a simple calculus to show that the equations of motion of the vakonomic problem are $`\ddot{x}=0`$, that is $`x(t)=at+b`$ and $`y(t)=f(a)t+c`$. The parameters $`a`$, $`b`$ and $`c`$ are obtained from the boundary conditions $`x(t_1)`$, $`y(t_1)`$, $`x(t_2)`$ and $`y(t_2)`$. If $`\gamma (t)=(x(t),y(t))`$ is a straight line satisfying the boundary conditions, using (3.1), the reader can check that there exist admissible variation fields, and they are vector fields along $`\gamma (t)`$ of the form $`\lambda (t)𝐰`$, where $`\lambda (t_1)=\lambda (t_2)=0`$ and $`𝐰(x,y)=(x,y;1,f^{}(a))`$.
For example, if $`\varphi (\dot{x},\dot{y})\dot{y}\sqrt{1+\dot{x}^2}=0`$ and $`x(t)`$ is a linear function of time, there are not strongly admissible variations (see ). But there exist admissible variation fields, so there are weakly admissible variations. In fact, the weakly admissible variations $`\mathrm{\Gamma }(\epsilon ,t)=(x(\epsilon ,t),y(\epsilon ,t))`$ have the form $`x(\epsilon ,t)=x(t)+\lambda (t)\epsilon +o(\epsilon )`$, $`y(\epsilon ,t)=y(t)+(a/\sqrt{1+a^2})\lambda (t)\epsilon +o(\epsilon )`$, where $`\lambda (t_1)=\lambda (t_2)=0`$.
### 3.2 The equations of motion of vakonomic mechanics
As we have shown, a critical path of the action with constraints is an admissible path $`\gamma `$ such that $`{\displaystyle _I}_L(\ddot{\gamma }(t)),𝐰(t)dt`$ vanishes for each admissible variation field $`𝐰`$ (theorem 1). To obtain the corresponding Euler-Lagrange equation, we first establish the following proposition.
###### Proposition 2
Given a variational problem $`(Q,L,C,𝒞,𝒲)`$, where $`𝒲`$ are the variation fields satisfying definition 5, let $`\gamma `$ be an admissible path. Then, for any family of functions $`\mu _i(t)`$, $`i=1,\mathrm{},m`$, the first-order variations of the $`_{\dot{\gamma }}Ldt`$ and $`_{\dot{\gamma }}(L+_{i=1}^m\mu _i\varphi ^i)dt`$ with respect to an admissible variation field $`𝐰𝒲`$ coincide.
Proof: In principle, since the variations may not be strongly admissible, it is not clear that the variations of both actions yield the same result. However, using theorem 1 and definition 3, the difference of the first-order variations of the actions is
$$_I_{_{i=1}^m\mu _i\varphi ^i}\ddot{\gamma },𝐰dt=\underset{i=1}{\overset{m}{}}_I\mu _i\mathrm{d}\varphi ^i\dot{\gamma },\text{s}\dot{𝐰}dt\underset{i=1}{\overset{m}{}}_I\mathrm{D}\mu _i(𝒟\varphi ^i\dot{\gamma }),𝐰dt,$$
and both terms vanish whenever $`𝐰`$ is an admissible variation field. Therefore, the variations of the two actions coincide when $`𝐰`$ is an admissible variation field. $`\mathrm{}`$
###### Theorem 2
Given a variational problem $`(Q,L,C,𝒞,𝒲)`$, where $`𝒲`$ are the variation fields satisfying definition 5, let $`\gamma `$ be an admissible path. Then $`\gamma `$ is critical if and only if there exist functions $`\mu _j:I\text{}`$, $`j=1,\mathrm{}m`$, such that
$$_{L+_{i=1}^m\mu _i\varphi ^i}\ddot{\gamma }=0.$$
(3.4)
This is the equation of motion of vakonomic mechanics.
Proof: If equation (3.4) holds then, for each admissible variation field $`𝐰`$, $`_I_{L+_{i=1}^m\mu _i\varphi ^i}\ddot{\gamma },𝐰dt=0`$, which, according to Proposition (2), is equivalent to $`_I_L\ddot{\gamma },𝐰dt=0`$. This shows that $`\gamma `$ is a critical path. So it remains to prove the converse: that equation (3.4) is a necessary condition for the criticity of an admissible path $`\gamma `$. First, notice that the $`\mu _i`$ can be chosen such that
$$_{L+_{i=1}^m\mu _i\varphi ^i}\ddot{\gamma },𝐰_j=0$$
(3.5)
for $`j=1,\mathrm{},m`$. Indeed, by lemma 2 this equation can be written as
$$_L\ddot{\gamma },𝐰_j+\underset{i=1}{\overset{m}{}}\mu _i_{\varphi ^i}\ddot{\gamma },𝐰_j\underset{i=1}{\overset{m}{}}\mathrm{D}\mu _i𝒟\varphi ^i\dot{\gamma },𝐰_j=0$$
for each $`j=1,\mathrm{},m`$. From definition 3 and the choice of $`𝐰_j`$ we have
$$_{\varphi ^i}\ddot{\gamma },𝐰_j=\mathrm{d}\varphi ^i\dot{\gamma },\text{s}\dot{𝐰}_j\mathrm{D}𝒟\varphi ^i\dot{\gamma },𝐰_j=\mathrm{d}\varphi ^i\dot{\gamma },\text{s}\dot{𝐰}_j,$$
for $`j=1,\mathrm{},m`$. That is, we have
$$(\mathrm{D}\mu _i)A_j^i\mu _iB_j^iD_j=0,$$
(3.6)
where $`A`$ and $`B`$ are the matrices we have used previously (3.2). So again we have a linear differential equation that determines the functions $`\mu _i`$ (up to initial conditions). From now on we assume that $`A`$ is the identity matrix; this can be easily done through a linear change of the basis $`(𝐰_i)`$.
If we apply the variational principle for the modified Lagrangian $`L+_{i=1}^m\mu _i\varphi ^i`$, we have
$$\underset{k=1}{\overset{n}{}}_I_{L+_{i=1}^m\mu _i\varphi ^i}\ddot{\gamma },𝐰_k\lambda ^kdt=0$$
for each set of functions $`\lambda ^k`$ yielding an admissible variation field.
If we choose the functions $`\mu _i`$ satisfying (3.5), then the sum is only from $`m+1`$ to $`n`$:
$$\underset{j=m+1}{\overset{n}{}}_I_{L+_{i=1}^m\mu _i\varphi ^i}\ddot{\gamma },𝐰_j\lambda ^jdt=0.$$
(3.7)
This must hold for any choice of the functions $`\lambda ^{m+1},\mathrm{},\lambda ^n`$ giving an admissible variation field. However, as we have shown in the preceding section, the functions $`\lambda ^{m+1},\mathrm{},\lambda ^n`$ are not arbitrary in general, due to the final conditions $`\lambda ^1(t_2)=\mathrm{}=\lambda ^m(t_2)=0`$. Let $`(\overline{\mu }_i)_{1im}`$ be the particular solution of (3.6) satisfying $`\overline{\mu }_i(t_2)=0`$, and let $`(\overline{\nu }_i^j)`$ be the transpose of the fundamental matrix of the homogeneous system of (3.6) satisfying the initial condition $`\overline{\nu }_i^j(t_2)=\delta _i^j`$. Notice that $`\overline{\nu }=(\nu )^1`$. (In general, if $`\nu `$ is a fundamental matrix of $`\dot{x}=Ax`$, then $`(\nu ^t)^1`$ is a fundamental matrix of $`\dot{x}=A^tx`$). Then the general solution of (3.6) is $`\mu _i=\overline{\mu }_i+_{j=1}^m\rho _j\overline{\nu }_i^j`$, where $`\rho _j`$ are arbitrary constants. Suppose for a while that there exist admissible variation fields $`𝐰=_{k=1}^n\lambda ^k𝐰_k`$ with
$$\lambda ^l=_{L+_{i=1}^m(\overline{\mu }_i+_{j=1}^m\rho _j\overline{\nu }_i^j)\varphi ^i}\ddot{\gamma },𝐰_l,$$
(3.8)
for $`l=m+1\mathrm{},n`$. Then (3.7) is a vanishing sum of integrals of squares, which, combined with (3.5), yields the equation of motion (3.4). It remains to prove that the choice of such $`\lambda ^l`$, $`l=m+1,\mathrm{},n`$, gives an admissible variation field. From (3.3) and $`\overline{\nu }=(\nu )^1`$, the variations defined by (3.8) are admissible if and only if
$$_{t_1}^{t_2}\overline{\nu }(s)C(s)[_{L+_{i=1}^m(\overline{\mu }_i+_{j=1}^m\rho _j\overline{\nu }_i^j)\varphi ^i}\ddot{\gamma },𝐰_{(2)}]^tds=0,$$
(3.9)
where $`𝐰_{(2)}`$ denotes the last $`nm`$ vector fields. Now, from lemma 2 and definition 3, we have
$$_{_{j=1}^m\overline{\nu }_j^i\varphi ^j}\ddot{\gamma },𝐰_l=\underset{j=1}{\overset{m}{}}\overline{\nu }_j^i_{\varphi ^j}\ddot{\gamma },𝐰_l\underset{j=1}{\overset{m}{}}(\text{D}\overline{\nu }_j^i)𝒟\varphi ^j\dot{\gamma },𝐰_𝐥=\underset{𝐣=\mathrm{𝟏}}{\overset{𝐦}{}}\overline{\nu }_𝐣^𝐢\mathrm{d}\varphi ^𝐣\gamma ,\text{s}𝐰_𝐥\mathrm{𝐫𝐚𝐧𝐠𝐥𝐞}.$$
(3.10)
Using that $`C_l^i=\mathrm{d}\varphi ^i\dot{\gamma },\text{s}\dot{𝐰}_l`$ and combining (3.9) and (3.10), we obtain the linear system for the $`\rho _j`$
$$\underset{h=1}{\overset{m}{}}\rho _h_{t_1}^{t_2}\underset{l=m+1}{\overset{n}{}}_{_{j=1}^m\overline{\nu }_j^i\varphi ^j}\ddot{\gamma },𝐰_l_{_{k=1}^m\overline{\nu }_k^h\varphi ^k}\ddot{\gamma },𝐰_l=$$
$$=_{t_1}^{t_2}\underset{l=m+1}{\overset{n}{}}_{_{j=1}^m\overline{\nu }_j^i\varphi ^j}\ddot{\gamma },𝐰_l_{L+_{k=1}^m\overline{\mu }_k\varphi ^k}\ddot{\gamma },𝐰_l;.$$
(3.11)
If this system has any solution, then we can find values for $`\rho _j`$, $`j=1,\mathrm{},m`$, such that the functions $`\lambda ^l`$ defined in (3.8) give rise to admissible variations. Now, we prove that this system has always solution. Consider the pre-Hilbert space of continuous vector-valued functions $`𝐟:[t_1,t_2]\text{}^{nm}`$ with the usual scalar product $`𝐟,𝐠=_{l=m+1}^n_{t_1}^{t_2}f_lg_l`$. Let $`V`$ be the finite–dimensional subspace spanned by the $`m`$ vectors
$$𝐞_i=(_{_{j=1}^m\overline{\nu }_j^i\varphi ^j}\ddot{\gamma },𝐰_{m+1},\mathrm{},_{_{j=1}^m\overline{\nu }_j^i\varphi ^j}\ddot{\gamma },𝐰_n),$$
$`i=1,\mathrm{},m`$. Then we can write the system (3.11) as
$$\underset{h=1}{\overset{m}{}}\rho _h𝐞_i,𝐞_h=𝐞_i,𝐯,$$
where $`𝐯=(_{L+_{k=1}^m\overline{\mu }_k\varphi ^k}\ddot{\gamma },𝐰_{m+1},\mathrm{},_{L+_{k=1}^m\overline{\mu }_k\varphi ^k}\ddot{\gamma },𝐰_n)`$.
The solutions of this system are any constants $`\rho _h`$ such that $`_{h=1}^m\rho _h𝐞_h`$ is the orthogonal projection of $`𝐯`$ onto $`V`$. This is well defined, since $`V`$ is finite–dimensional. (Notice that the $`\rho _h`$ may not be unique, since the $`𝐞_h`$ are not necessarily independent). This completes the proof. $`\mathrm{}`$
Remark: Notice that, using lemma 2, the equation of motion may also be written as
$$_L\ddot{\gamma }=\underset{i=1}{\overset{m}{}}\left((\mathrm{D}\mu _i)𝒟\varphi ^i\dot{\gamma }\mu _i_{\varphi ^i}\ddot{\gamma }\right).$$
(3.12)
Remark: In the proofs of the equation of motion of vakonomic mechanics that one can find in the literature, it is usually assumed that the functions $`\lambda ^l`$, $`l=m+1,\mathrm{},n`$, giving the admissible variations are free. Then, the equation of motion is obtained as a straight consequence of (3.7). However, in general, these functions are not absolutely free.
## 4 The variations and equations of motion of nonholonomic mechanics
In this section we are going to show that nonholonomic mechanics may be understood as a variational problem.
Our initial setting is also the submanifold $`C\mathrm{T}Q`$, which may be locally defined by the vanishing of the constraints $`\varphi ^i`$. An admissible path is still a path $`\gamma :IQ`$ such that $`\dot{\gamma }`$ takes its values in $`C`$. Let us define which are the admissible variation fields in nonholonomic mechanics.
###### Definition 6
A variation field $`𝐰`$ of an admissible path $`\gamma `$ is called an admissible variation field (in nonholonomic mechanics) if it is a section of the subbundle $`\text{L}_\gamma ^C\gamma ^{}\mathrm{T}Q`$. That is,
$$\mathrm{vl}(\dot{\gamma }(t),𝐰(t))\mathrm{T}_{\dot{\gamma }(t)}(C).$$
Using the constraints, equivalent statements are
$$\mathrm{d}\varphi \dot{\gamma },\mathrm{vl}(\dot{\gamma },𝐰)=0,$$
(4.1)
or
$$𝒟\varphi \dot{\gamma },𝐰=0,$$
(4.2)
for each constraint $`\varphi `$.
Notice the key difference with respect to vakonomic mechanics: now the admissibility is a $`C^1(I)`$-linear condition on $`𝐰`$. This linearity makes things easier. Next, we obtain the equation of motion of nonholonomic mechanics.
###### Theorem 3
Given a variational problem $`(Q,L,C,𝒞,𝒲)`$, where $`𝒲`$ are the admissible variation fields satisfying definition 6, an admissible path $`\gamma `$ is critical if and only if there exist functions $`\mu _j`$, $`j=1,\mathrm{}m`$, such that
$$_L\ddot{\gamma }=\underset{i=1}{\overset{m}{}}\mu _i𝒟\varphi ^i\dot{\gamma }.$$
(4.3)
Proof: A critical path for nonholonomic mechanics is an admissible path $`\gamma `$ such that the first-order variation of the action, $`{\displaystyle _I}_L(\ddot{\gamma }(t)),𝐰(t)dt,`$ vanishes for each admissible variation field $`𝐰`$. By equation (4.2), this means that $`_L(\ddot{\gamma }(t))`$ is a linear combination of the $`𝒟\varphi ^i\dot{\gamma }`$. Thus, the result follows. $`\mathrm{}`$
Remark: It is obvious that there always exist admissible variation fields in nonholonomic mechanics. For example, if we calculate the admissible variation fields of the example in section 3, we will find that, in this case, they coincide with the admissible variation fields of vakonomic mechanics. But this will not be true in general if the constraints are not integrable.
## 5 The case of integrable constraints
Let us consider the problem of holonomic constraints in the usual sense.
###### Definition 7
Given a differentiable manifold $`Q`$ and a Lagrangian function $`L:\mathrm{T}Q\text{}`$, a holonomic problem is a variational problem where
* The constraint submanifold is given by a submanifold $`PQ`$, thus $`C=\mathrm{T}P`$.
* Admissible paths are paths $`\gamma :IPQ`$.
* Variation fields $`𝐰`$ along $`\gamma `$ are admissible if they are tangent to $`P`$.
Notice that, from any admissible variation field along an admissible path $`\gamma `$, one may construct a variation $`\mathrm{\Gamma }`$ contained in $`P`$. Therefore, it is clear that the problem with holonomic constraints is equivalent to the unconstrained variational problem defined on $`P`$ by taking the restriction of the Lagrangian $`L`$ to $`\mathrm{T}P\mathrm{T}Q`$.
Now, consider the cases of both vakonomic and nonholonomic mechanics when the constraints are defined by an integrable subbundle $`C\mathrm{T}Q`$. In this situation, we have the following equivalence.
###### Theorem 4
If the constraint submanifold $`C`$ is an integrable subbundle of $`\mathrm{T}Q`$, then both vakonomic and nonholonomic mechanics coincide, and they are equivalent to a holonomic constrained problem on each integral submanifold of $`C`$.
Proof: First of all, notice that, in both cases, a path $`\gamma :IQ`$ is admissible (i.e., $`\dot{\gamma }`$ is in $`C`$) if and only if it is contained in an integral submanifold $`PQ`$ of $`C`$. Recall that $`\mathrm{T}_q(P)=C_q`$, for each point in $`P`$.
Let $`𝐰`$ be a variation field of $`\gamma `$. We know that $`𝐰`$ is admissible in the nonholonomic framework if $`\mathrm{vl}(\dot{\gamma }(t),𝐰(t))\mathrm{T}_{\dot{\gamma }(t)}(C)`$ (see definition 6). Since the vertical lift restricts naturally to subbundles, this is equivalent to say that $`\mathrm{vl}(\dot{\gamma }(t),𝐰(t))\text{V}_{\dot{\gamma }(t)}C\text{V}_{\dot{\gamma }(t)}(\mathrm{T}Q)`$, that is, $`𝐰(t)C_{\gamma (t)}`$ or, what is the same, $`𝐰(t)\mathrm{T}_{\gamma (t)}P`$. And this last condition says that $`𝐰`$ is admissible for the holonomic problem. Therefore, we have proved that a nonholonomic problem with integrable constraints is equivalent to a holonomic problem on each integral submanifold of $`C`$.
Now, we show the equivalence with the vakonomic problem. Since $`C`$ is an integrable subbundle of $`\mathrm{T}Q`$, it is known that the integral submanifolds of $`C`$ can be locally described as $`\psi =\text{constant}`$, for some independent functions $`\psi `$ on $`Q`$. This implies that the constraint submanifold $`C`$ can be locally described by $`\varphi =\text{}\psi =0`$. Here, $`\text{}\psi :\mathrm{T}Q\text{}`$ is the differential d$`\psi `$ of $`\psi :Q\text{}`$ considered as a function on the tangent bundle.
Let $`𝐰`$ be a variation field of $`\gamma `$, in the sense of vakonomic mechanics. Then, for any $`\varphi =\text{}\psi `$, we have
$$0=\text{d}\varphi \dot{\gamma },\text{s}\dot{𝐰}=\text{d̃d̃}\psi \text{s}\dot{𝐰}.$$
Using the property $`\text{d̃d̃}\psi =\text{d̃d̃}\psi \text{s}`$, we obtain
$$0=\text{d̃d̃}\psi \dot{𝐰}=\text{dd̃}\psi 𝐰,\dot{𝐰}=\text{d}\varphi 𝐰,\dot{𝐰}=\mathrm{D}(\varphi 𝐰)=\mathrm{D}\text{d}\psi \gamma ,𝐰.$$
Thus, $`𝐰`$ is an admissible variation field in vakonomic mechanics if and only if $`\text{d}\psi \gamma ,𝐰`$ is constant. Since $`𝐰(t_1)=0`$, this constant is zero, $`\text{d}\psi \gamma ,𝐰=0`$, which means that $`𝐰`$ is tangent to $`P`$. Therefore, $`𝐰`$ is admissible for the vakonomic problem if and only if it is admissible for the holonomic problem on $`P`$. $`\mathrm{}`$
## 6 Optimal control and vakonomic mechanics
A problem of optimal control may be given by the following data: a configuration space $`B`$ describing the state variables, a fibre bundle $`\pi :MB`$ whose fibres describe the control variables, a vector field $`Y`$ along the projection of the bundle, $`Y:M\mathrm{T}B`$, and a “Lagrangian function” $`L:M\text{}`$. For a path $`\gamma :IM`$ where $`\pi \gamma `$ (not $`\gamma `$!) has fixed end-points, the problem is to find an extremum of the action
$$_\gamma L(\gamma (t))dt$$
when $`\gamma `$ satisfies the differential equation
$$(\pi \gamma )^.=Y\gamma $$
(6.1)
that rules the evolution of the state variables.
It is easy to show that this is indeed a vakonomic problem on the manifold $`M`$, in which the Lagrangian $`L`$ is very singular, since it does not depend on the velocities. The constraint submanifold $`C\mathrm{T}M`$, given by the differential equation above, is
$$C=\{w_u\mathrm{T}M\mathrm{T}\pi (w_u)=Y(u)\}.$$
(6.2)
In this way, a path $`\gamma `$ is admissible if and only if it is a solution of the differential equation (6.1) or, equivalently, if it takes values in the affine subbundle $`C`$ of $`\mathrm{T}M`$. On the other hand, from the special characteristics of optimal control problems, we can relax the boundary conditions as we have done. Observe that theorems 1 and 2 remain true since $`L`$ does not depend on the velocities and the structure of the constraints (6.1) (they do not depend on the derivatives on the control variables). In coordinates, if $`x^i`$, $`i=1,\mathrm{}m`$, are local coordinates in $`B`$ and $`(x^i,u^\alpha )`$, $`i=1,\mathrm{}m`$, $`\alpha =1,\mathrm{},nm`$, are local coordinates in $`M`$, the action is given by
$$_{t_1}^{t_2}L(x^i,u^\alpha )dt$$
and the constraints are given by a set of first order differential equations
$$\dot{x}^i=Y^i(x,u),i=1,\mathrm{},m,$$
with boundary conditions $`x_1^i=x^i(t_1)`$ to $`x_2^i=x^i(t_2)`$ (there are no boundary conditions on the control variables). Notice that for this vakonomic problem the constraints are very particular: they express the velocities of the state variables in terms of the state and control variables.
Let us identify which variation fields $`𝐰`$ are admissible. Notice that frames for the bundles $`\text{L}_\gamma ^C\gamma ^{}\mathrm{T}(M)`$ are provided with $`(𝐰_\alpha )`$ and $`(𝐰_i,𝐰_\alpha )`$, where for instance
$$𝐰_\alpha =\frac{}{u^\alpha }\gamma ,𝐰_i=\frac{}{x^i}\gamma ,$$
if we have coordinates $`(x^i,u^\alpha )`$ on $`M`$. Writing $`𝐰=\lambda ^i𝐰_i+\lambda ^\alpha 𝐰_\alpha `$, the differential equation (3.1) turns out to be
$$\mathrm{D}\lambda ^i\frac{y^i}{x^j}\lambda ^j=\frac{y^i}{u^\alpha }\lambda ^\alpha .$$
Optimal control theory admits several geometric formulations and expressions of the equation of motion. First, we can give a lagrangian description on the configuration manifold $`\mathrm{T}^{}B\times _BM`$. We only need to define a lagrangian on its tangent space:
$$(x,p,u,\dot{x},\dot{p},\dot{u})=L(x,u)+p,\dot{x}Y(x,u)$$
where we write $`(x;p,u)`$ for the variables of $`\mathrm{T}^{}B\times _BM`$—recall that $`\mathrm{T}(\mathrm{T}^{}B\times _BM)=\mathrm{T}(\mathrm{T}^{}B)\times _{\mathrm{T}B}\mathrm{T}M`$; its elements are pairs of tangent vectors $`(\dot{x},\dot{p},\dot{u})`$ projecting to the same tangent vector $`\dot{x}`$. For a path $`\eta `$ on $`\mathrm{T}^{}B\times _BM`$ the Euler-Lagrange equation $`_{}\ddot{\eta }`$ is readily seen to be equivalent to the vakonomic equation (3.4).
Let us recall that, given a $`l`$-dimensional differentiable manifold $`Q`$ with local coordinates $`(q^A)`$, there is a canonical tensor field on $`\mathrm{T}Q`$, the vertical endomorphism $`S`$ which is a rank-$`l`$ (1,1) tensor field on $`\mathrm{T}Q`$ such that $`KerS=\text{Im}S`$ and whose Nijenhuis tensor $`N_S`$ vanishes. In natural coordinates $`(q^A,\dot{q}^A)`$, the local expression of $`S`$ is given by $`S=\mathrm{d}q^A\frac{}{\dot{q}^A}`$. Also, we have the Liouville vector field $`\mathrm{\Delta }`$ (the infinitesimal generator of the dilations along the fibres on $`\mathrm{T}Q`$), whose local expression is $`\mathrm{\Delta }=\dot{q}^A/\dot{q}^A`$. A vector field $`X`$ in $`\mathrm{T}Q`$ is called a second order differential equation (SODE) if $`S(X)=\mathrm{\Delta }`$. Now, if $`:\mathrm{T}Q\text{}`$ is a Lagrangian function, we can construct the Cartan 1-form associated with $``$, given by $`\theta _{}=S^{}\mathrm{d}`$, the Cartan 2-form $`\omega _{}=\mathrm{d}\theta _{}`$ and the energy function $`E_{}=\mathrm{\Delta }()`$. Then the paths $`\eta `$ solution of the Euler-Lagrange equations are the integral curves of a second order differential equation $`X`$ in $`\mathrm{T}Q`$ satisfying the dynamical equation $`i_X\omega _{}=\mathrm{d}E_{}`$.
Taking $`Q=\mathrm{T}^{}B\times _BM`$ and $`=L(x,u)+p,\dot{x}Y(x,u)`$, we obtain a geometrical expression of the equation of motion of optimal control theory,
$$i_X\omega _{}=\mathrm{d}E_{},$$
where $`E_{}=p,Y(x,u)L(x,u)`$ and $`\omega _{}=\mathrm{d}\theta _{}=\mathrm{d}(p_i\mathrm{d}x^i)=\mathrm{d}x^i\mathrm{d}p_i`$.
Associated to this lagrangian description we can consider the hamiltonian formalism of $``$. This means to consider the manifold $`\mathrm{T}^{}(\mathrm{T}^{}B\times _BM)`$ with its canonical symplectic structure, the Legendre’s transformation of $``$, $`𝒟:\mathrm{T}(\mathrm{T}^{}B\times _BM)\mathrm{T}^{}(\mathrm{T}^{}B\times _BM)`$, and to push forward through it the energy function to a hamiltonian function $`:\mathrm{T}^{}(\mathrm{T}^{}B\times _BM)\text{}`$.
However, the most interesting geometric description of optimal control theory is a presymplectic description which can be constructed on the manifold $`\mathrm{T}^{}B\times _BM`$. Here we consider the 2-form $`\omega `$ obtained by pull-back through $`\mathrm{T}^{}B\times _BM\mathrm{T}^{}B`$ of the canonical 2-form of the last manifold. In local coordinates, $`\omega =\mathrm{d}q\mathrm{d}p`$. Taking the hamiltonian function defined by
$$H(x,u,p)=p,Y(x,u)L(x,u),$$
if $`\eta `$ is a path on $`\mathrm{T}^{}B\times _BM`$, the presymplectic equation
$$i_{\dot{\eta }}\omega =\mathrm{d}H\eta $$
is equivalent to the equation of motion of vakonomic mechanics, in the sense that there is a natural bijection between both sets of solutions. To show this is enough to write the local expressions.
Remark: In optimal control theory the hamiltonian function is usually written as $`H(x,u,p)=p,Y(x,u)\mu _0L(x,u)`$, where $`\mu _0=0,1`$. When $`\mu _0=0`$ we recover the so called abnormal solutions (see ). However, in the vakonomic approach, there are not abnormal solutions. The key issue is that we work with admissible variation fields, not admissible variation curves.
## 7 Constraints defined by a distribution
In this section we present a geometric framework for constrained systems when the constraint submanifold $`C`$ is a distribution (or vector subbundle) of the tangent manifold $`\mathrm{T}Q`$. In local coordinates, this means that the constraints are linear functions on the velocities. The subbundle $`C\mathrm{T}Q`$ can be described in terms of its annihilator, $`C^0\mathrm{T}^{}Q`$. If this is locally described in terms of 1-forms, $`\alpha ^i=\alpha _a^i(q)\mathrm{d}q^a`$ ($`i=1,\mathrm{},m`$, where $`m`$ is the codimension of $`C`$ and $`(q^a)`$ are local coordinates of $`Q`$), then $`C`$ is locally described in terms of the constraints
$$\varphi ^i(v_q)=\alpha ^i(q),v_q=0,i=1,\mathrm{},m.$$
Let us consider the vector bundle $`\mathrm{T}QC^0`$, in which we will set up the dynamics.
On the one hand, given the Lagrangian function $`L`$ on the tangent bundle $`\mathrm{T}Q`$, let $`\theta _L=S^{}\mathrm{d}L`$ be the Lagrange 1-form on $`\mathrm{T}Q`$. Its pull-back along the projection $`\pi _1:\mathrm{T}QC^0\mathrm{T}Q`$ yields the 1-form
$$\theta _1=\pi _1^{}\theta _L$$
on $`\mathrm{T}QC^0`$. Also, using the Liouville vector field $`\mathrm{\Delta }`$ on $`\mathrm{T}Q`$, the energy function associated with $`L`$ in $`\mathrm{T}QC^0`$ is
$$E=\pi _1^{}(\mathrm{\Delta }(L)L).$$
On the other hand, let $`\theta _Q`$ be the canonical 1-form defined on the cotangent bundle $`\mathrm{T}^{}Q`$. If $`j_0:C^0\mathrm{T}^{}Q`$ denotes the canonical inclusion and $`\pi _2:\mathrm{T}QC^0C^0`$ is the projection onto the second factor, then we can take the pull-back of these mappings to construct a 1-form $`\theta _2`$ on $`\mathrm{T}QC^0`$ as
$$\theta _2=(j_0\pi _2)^{}\theta _Q.$$
Using the 1-forms $`\theta _1`$ and $`\theta _2`$, we have a presymplectic form
$$\mathrm{\Omega }=\mathrm{d}(\theta _1+\theta _2).$$
By means of the energy function $`E`$, we obtain a presymplectic dynamics on the extended phase space $`\mathrm{T}QC^0`$ which is equivalent to vakonomic mechanics:
###### Theorem 5
Let $`L:\mathrm{T}Q\text{}`$ be a Lagrangian, and $`C\mathrm{T}Q`$ a vector subbundle. For a path $`\xi `$ in the manifold $`\mathrm{T}QC^0`$, consider the differential equation
$$i_{\dot{\xi }}\mathrm{\Omega }=\mathrm{d}E\xi .$$
(7.1)
This equation is equivalent to the equation of motion of vakonomic mechanics in the following sense:
* If $`\xi =(\xi _1,\xi _2)`$ is a solution of (7.1) and $`\xi _1`$ is the lift of a path in $`Q`$, $`\xi _1=\dot{\gamma }`$, then $`\gamma `$ is an admissible path ($`\dot{\gamma }`$ is in $`C`$) and is a solution of the equation of motion of vakonomic mechanics (3.4).
* Conversely, given an admissible path $`\gamma `$ which is a solution of (3.4), together with the multipliers $`\mu ^i`$, then the path $`\xi (t)=(\dot{\gamma }(t),\mu _i(t)\mathrm{d}\varphi ^i(\dot{\gamma }(t)))`$ is a solution of equation (7.1).
If the Lagrangian is regular then equation (7.1) already implies that $`\xi _1`$ is the lift of a path in $`Q`$.
Proof: It is enough to check the equivalence in local coordinates. We take $`(q^a,v^a,\lambda _i)`$, $`a=1,\mathrm{},n`$, $`i=1,\mathrm{},m`$, as local coordinates in $`\mathrm{T}QC^0`$ (we represent an element of $`C_q^0`$ as $`\lambda _i\alpha ^i(q)`$). Then we have
$`\theta _1`$ $`=`$ $`{\displaystyle \frac{L}{v^a}}\mathrm{d}q^a,`$
$`\theta _2`$ $`=`$ $`\lambda _i\alpha _a^i(q)\mathrm{d}q^a,`$
$`\mathrm{\Omega }`$ $`=`$ $`\left({\displaystyle \frac{^2L}{v^aq^b}}+\lambda _i{\displaystyle \frac{\alpha _a^i}{q^b}}\right)\mathrm{d}q^a\mathrm{d}q^b+{\displaystyle \frac{^2L}{v^av^b}}\mathrm{d}q^a\mathrm{d}v^b+\alpha _a^i\mathrm{d}q^a\mathrm{d}\lambda _i.`$
Since $`E=v^a\left(L/v^a\right)L`$, we also have
$$\mathrm{d}E=\left(v^b\frac{^2L}{q^av^b}\frac{L}{q^a}\right)\mathrm{d}q^a+v^b\frac{^2L}{v^av^b}\mathrm{d}v^a.$$
Now let us consider the path $`\xi (t)=(q^a(t),v^a(t),\lambda ^i(t))`$, with velocity $`\dot{\xi }=(q,v,\lambda ;\dot{q},\dot{v},\dot{\lambda })`$. Then
$`i_{\dot{\xi }}\mathrm{\Omega }`$ $`=`$ $`\left(\dot{q}^b{\displaystyle \frac{^2L}{q^av^b}}+\dot{q}^b{\displaystyle \frac{\alpha _b^i}{q^a}}\lambda _i\dot{q}^b{\displaystyle \frac{^2L}{v^aq^b}}\dot{q}^b{\displaystyle \frac{\alpha _a^i}{q^b}}\lambda _i\dot{v}^b{\displaystyle \frac{^2L}{v^av^b}}\dot{\lambda }_i\alpha _a^i\right)\mathrm{d}q^a`$
$`+\dot{q}^b{\displaystyle \frac{^2L}{v^av^b}}\mathrm{d}v^a+\dot{q}^a\alpha _a^i\mathrm{d}\lambda _i.`$
Therefore, equation $`i_{\dot{\xi }}\mathrm{\Omega }=\mathrm{d}E`$ is equivalent to the three equations
$$(\dot{q}^bv^b)\frac{^2L}{q^av^b}+\dot{q}^b\frac{\alpha _b^i}{q^a}\lambda _i\dot{q}^b\frac{^2L}{v^aq^b}\dot{q}^b\frac{\alpha _a^i}{q^b}\lambda _i\dot{v}^b\frac{^2L}{v^av^b}\dot{\lambda }_i\alpha _a^i=\frac{L}{q^i},$$
(7.2)
$$\dot{q}^b\frac{^2L}{v^av^b}=v^b\frac{^2L}{v^av^b},$$
(7.3)
$$\dot{q}^a\alpha _a^i=0.$$
(7.4)
The fact that $`\xi _1`$ is the lift of a path $`\gamma `$ in $`Q`$ means in coordinates that $`v(t)=\dot{q}(t)`$, so equation (7.3) is an identity. Notice also that if the Lagrangian is regular then the Hessian matrix $`\left(\frac{^2L}{v^av^b}\right)`$ is invertible, therefore in this case equation (7.3) implies that $`v(t)=\dot{q}(t)`$, that is to say, $`\xi _1`$ is the lift of a path in $`Q`$. Then, in equation (7.4) we obtain the constraints $`\varphi ^i(q,\dot{q})=\alpha _a^i(q)\dot{q}^a=0`$, that is, $`\gamma `$ is an admissible path. Finally, we can write equation (7.2) as
$$\ddot{q}^b\frac{^2L}{v^av^b}+\dot{q}^b\frac{^2L}{v^aq^b}+\dot{q}^b\frac{\alpha _a^i}{q^b}\lambda _i+\dot{\lambda }_i\alpha _a^i=\frac{L}{q^a}+\dot{q}^b\frac{\alpha _b^i}{q^a}\lambda _i.$$
But these are the vakonomic equations (3.4) of the extended Lagrangian $`=L+\mu _i\alpha _a^iv^a`$, using the natural identification between the functions $`\mu _i`$ and the coordinates $`\lambda _i`$ of the cotangent vectors. $`\mathrm{}`$
Remark: In a similar way, the vakonomic dynamics can be also defined on the manifold $`CC^0`$. Since $`CC^0`$ is a vector subbundle of $`\mathrm{T}QC^0`$, we can pull-back the $`2`$-form $`\mathrm{\Omega }`$ and the energy function $`E`$ to define a $`2`$-form $`\stackrel{~}{\mathrm{\Omega }}`$ and a new function $`\stackrel{~}{E}`$ in $`CC^0`$. The reader can check that, then the equation of motion of vakonomic mechanics (3.4) is equivalent to find the paths $`\xi =(\xi _1,\xi _2)`$ in $`CC^0`$, where $`\xi _1`$ is the lift of a path in $`Q`$, such that
$$i_{\dot{\xi }}\stackrel{~}{\mathrm{\Omega }}=\mathrm{d}\stackrel{~}{E}\xi .$$
(7.5)
Moreover, if the lagrangian is regular, then $`\stackrel{~}{\mathrm{\Omega }}`$ is a symplectic form. Notice also that this equation, as well as equation (7.1), can also be expressed in terms of vector fields. For instance, when $`\stackrel{~}{\mathrm{\Omega }}`$ is symplectic, the solutions of equation (7.5) are the integral curves of the vector field $`\stackrel{~}{X}`$ such that
$$i_{\stackrel{~}{X}}\stackrel{~}{\mathrm{\Omega }}=\mathrm{d}\stackrel{~}{E}.$$
In the case of nonholonomic mechanics, a similar result can be proved, in the same way as for theorem 5. Let us denote $`\mathrm{\Omega }_1=\mathrm{d}\theta _1`$. Then we have:
###### Theorem 6
Let $`L:\mathrm{T}Q\text{}`$ be a Lagrangian, and $`C\mathrm{T}Q`$ a vector subbundle. For a path $`\xi `$ in the manifold $`\mathrm{T}QC^0`$, consider the differential equation
$$i_{\dot{\xi }}\mathrm{\Omega }_1=\mathrm{d}E\xi +\theta _2\xi .$$
(7.6)
This equation is equivalent to the equation of motion of nonholonomic mechanics in the following sense:
* If $`\xi =(\xi _1,\xi _2)`$ is a solution of (7.6) where $`\xi _1`$ is the lift of a path $`\gamma `$ in $`Q`$, $`\xi _1=\dot{\gamma }`$, then $`\gamma `$ is a solution of the equation of motion of nonholonomic mechanics (4.3).
* Conversely, given a path $`\gamma `$ which is a solution of (4.3), together with the multipliers $`\mu ^i`$, then the path $`\xi (t)=(\dot{\gamma }(t),\mu _i(t)\mathrm{d}\varphi ^i(\dot{\gamma }(t)))`$ is a solution of equation (7.6).
If the Lagrangian is regular then equation (7.6) already implies that $`\xi _1`$ is the lift of a path in $`Q`$. $`\mathrm{}`$
## 8 Conclusions
In this paper we have presented variational calculus (in one dimension) in a geometric framework, aiming to study dynamical systems with non-holonomic constraints (i.e., constraints depending on the positions and the velocities). We have shown that a generalised formulation of variational calculus, in which the admissible paths and the admissible infinitesimal variations are not necessarily related, makes room for the study of dynamical systems subject to non-holonomic constraints from different points of view. This generalized variational calculus encompasses the often-called vakonomic mechanics (which is a strict variational problem with constraints) and the non-holonomic mechanics (based on d’Alembert’s principle).
In the case of vakonomic mechanics, we have provided a geometric procedure to obtain the equation of motion, choosing an appropriate set of admissible infinitesimal variations proving that they always exist.
In the case of non-holonomic mechanics, it is far more simple than in vakonomic mechanics to choose an appropriate set of admissible infinitesimal variations, and the corresponding equation of motion is readily obtained.
Our formulation also provides a neat equivalence between both vakonomic and non-holonomic mechanics when the constraints are integrable (also called holonomic).
We have also found the geometry lying on some particular cases of vakonomic mechanics, namely the case of optimal control and the case where the constraint submanifold is a vector subbundle of the tangent bundle.
All the paper is written for the case of time-independent lagrangian and constraints, but the reader may check that the time-dependent case may be dealt with by adjunction of the time variable in a not too involved way.
## Acknowledgements
The authors thank N. Román-Roy for useful discussions. X.G. and M.C.M.L. acknowledge partial financial support from CICYT TAP 97–0969–C03–01 and PB98–0920. J.M.S. acknowledge partial financial support from CICYT projects PB98–0821 and PB98–0920.
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# Polarised parton densities from the fits to the deep inelastic spin asymmetries on nucleons
## Abstract
We have made next to leading order QCD fit to the deep inelastic spin asymmetries on nucleons and we have determined polarised quark and gluon densities. The functional form for such distributions was inspired by the unpolarised ones given by MRST (Martin, Roberts, Stirling and Thorne) fit. In addition to usually used data points (averaged over $`Q^2`$) we have also considered the sample containing points with the same $`x`$ and different $`Q^2`$. Our fits to both groups of data give very similar results. For the integrated quantities we get rather small gluon polarisation. For the non averaged data the best fit is obtained with vanishing gluon contribution at $`Q^2=1\mathrm{GeV}^2`$. For comparison models with alternative assumptions about quark sea and in particular strange sea behaviour are discussed.
One has quite a lot of data on deep inelastic spin asymmetries for different nucleon targets. The data (recent and old) come from experiments made at SLAC \[1-9\], CERN \[10-15\] and DESY . The newest data on proton and deuteron targets have smaller statistical errors and hence dominate in $`\chi ^2`$ fits.
Since the analysis of the EMC group results an enormous interest were started in studying polarised structure functions. It was suggested that polarised gluons may be responsible for the little spin carried by quarks. Determination of polarised gluon distribution is particularly interesting in this context. After calculation of two loop polarised splitting functions several next to leading (NLO) order QCD fits were performed \[20-24\] and polarised parton distributions (i.e. for quarks, antiquarks and gluons) were determined. Many groups obtained rather highly polarised gluon contribution (however this number is determined with a big error). The aim of this paper is to extend our next to leading order QCD analysis given in by taking into account, in addition to all previously considered data, also proton data from Hermes (DESY) and deuteron data from E155 experiment (SLAC). We will also use more recent MRST fit for partons . In we got negligible gluon contribution at $`Q^2=1\mathrm{GeV}^2`$. We will see that our result about gluon polarisation does not change very much. The method of choosing basic functions for fitting, used in this paper, differs from the other groups. As in we will divide the data into two groups. Let us remind that many experimental groups published data sets for the near values of $`x`$ and different $`Q^2`$ in addition to the ”averaged” data where one averages over $`Q^2`$ (the errors are smaller and $`Q^2`$ dependence is smeared out). In principle when we take into account $`Q^2`$ evolution of polarised and unpolarised functions (in NLO analysis) the first group of data points, i.e. non averaged, should be considered. In most of the fits to experimental data only second group of data, namely with averaged $`Q^2`$ dependence, was used. We will make fits to the both sets of data (the first group contains 417 points and the second 159 points). The results for both groups of data are very similar (the same conclusion was already drawn in ). We will also compare those results with the fits without $`Q^2`$ evolution taken into account (in other words assuming that asymmetries do not depend on $`Q^2`$, as in our previous fits ). One should add that many experimental groups have not succeeded in finding $`Q^2`$ dependence for approximately the same value of $`x`$ and different $`Q^2`$ . In our analysis we limit ourselves, as one usually does, to the data with $`Q^21\text{GeV}^2`$. As was already mentioned in our earlier papers making a fit to spin asymmetries and not directly to $`g_1(x,Q^2)`$ enables to avoid the problem with higher twist contributions which are probably less important in such case ( see for example ) . Experiments on unpolarised targets provide information on the unpolarised quark densities $`q(x,Q^2)`$ and $`G(x,Q^2)`$ inside the nucleon. These densities can be expressed in term of $`q^\pm (x,Q^2)`$ and $`G^\pm (x,Q^2)`$, i.e. densities of quarks and gluons with helicity along or opposite to the helicity of the parent nucleon. The unpolarised quark densities are given by the sum of $`q^+`$, $`q^{}`$ and $`G^+`$ , $`G^{}`$, namely:
$$q=q^++q^{},G=G^++G^{}.$$
(1)
On the other hand the polarised DIS experiments give information about polarised parton densities, i.e. the difference of $`q^+`$, $`q^{}`$ and $`G^+`$, $`G^{}`$:
$$\mathrm{\Delta }q=q^+q^{},\mathrm{\Delta }G=G^+G^{}.$$
(2)
We will try to determine $`q^\pm (x,Q^2)`$ and $`G^\pm (x,Q^2)`$, in other words, we will try to connect unpolarised and polarised data. In principle the asymptotic $`x`$ behaviour of $`q^+`$ and $`q^{}`$ will be taken from the unpolarised data (up to the modifications when some leading order terms vanish). We will use the fit given by MRST . In comparison with their previous fit (called R2) there is different behaviour at small $`x`$ for quark and gluon distributions. It is of course very restrictive assumption that $`\mathrm{\Delta }q`$ and $`\mathrm{\Delta }G`$ have the same behaviour as $`q`$ and $`G`$. On the other hand the small $`x`$ behaviour of unpolarised structure functions is determined from the $`x`$ values of Hera much smaller than in the polarised case. In our further analysis we will consider the integrals over the region measured in the experiments with polarised particles with the hope that in this case the behaviour of $`q^\pm `$ and $`G^\pm `$ could be more plausible. The values of integrals in the whole region ($`0x1`$), involving asymptotic behaviour taken from the unpolarised structure functions, may be not as reliable as in the measured region. As an alternative we also use Regge type behaviour.
It is known that the behaviour of the quark and gluon distributions in small $`x`$ region is extremely important in extrapolation of integrals over whole $`0x1`$ range. It could happen that in $`\mathrm{\Delta }q=q^+q^{}`$ (when we assume that $`q^+`$, $`q^{}`$ and $`q=q^++q^{}`$ have similar $`x`$ dependence) most singular $`x`$ terms cancel (that is especially important in case of valence quarks and sea contributions where the $`x`$ behaviour is relatively singular). We will see later how such description infers the fits and calculated parameters. For the less singular distributions for $`\mathrm{\Delta }u_v`$, $`\mathrm{\Delta }d_v`$ and $`\mathrm{\Delta }M`$ (total sea polarisation) there is no strong dependence of calculated quantities on the extrapolation in an unmeasured region but the fits have higher $`\chi ^2`$. One of the main tasks of considering NLO evolution in $`Q^2`$ is the determination of the gluon contribution $`\mathrm{\Delta }G`$. In $`\overline{MS}`$ scheme $`\mathrm{\Delta }G(x,Q^2)`$ comes in through the higher order corrections. In our fits we obtain $`\mathrm{\Delta }G`$ relatively small contrary to some expectations. When we use the non-averaged sample of data the actually the best fit is with vanishing $`\mathrm{\Delta }G`$ contribution.
Let us start with the formulas for unpolarised quark parton distributions gotten (at $`Q^2=1\mathrm{GeV}^2`$) from the one of recent fits performed by Martin, Roberts, Stirling and Thorne . For the valence quarks one get (in this fit one uses $`\mathrm{\Lambda }_{\overline{MS}}^{n_f=4}=0.3`$ GeV and $`\alpha _s(M_Z^2)=0.120`$):
$`u_v(x)`$ $`=`$ $`0.6051x^{0.5911}(1x)^{3.395}(1+2.078\sqrt{x}+14.56x),`$
$`d_v(x)`$ $`=`$ $`0.0581x^{0.7118}(1x)^{3.874}(1+34.69\sqrt{x}+28.96x),`$ (3)
and for the antiquarks from the sea (the same distribution is for sea quarks):
$`2\overline{u}(x)`$ $`=`$ $`0.4M(x)\delta (x),`$
$`2\overline{d}(x)`$ $`=`$ $`0.4M(x)+\delta (x),`$ (4)
$`2\overline{s}(x)`$ $`=`$ $`0.2M(x).`$
In eq.(4) the singlet contribution $`M=2[\overline{u}+\overline{d}+\overline{s}`$\] is:
$$M(x)=0.2004x^{1.2712}(1x)^{7.808}(1+2.283\sqrt{x}+20.69x),$$
(5)
whereas the isovector part ($`\delta =\overline{d}\overline{u}`$) reads:
$$\delta (x)=1.290x^{0.183}(1x)^{9.808}(1+9.987x33.34x^2).$$
(6)
For the unpolarised gluon distribution one gets:
$$G(x)=64.57x^{0.0829}(1x)^{6.587}(13.168\sqrt{x}+3.251x).$$
(7)
We assume, in an analogy to the unpolarised case, that the polarised quark distributions are of the form: $`x^\alpha (1x)^\beta P(\sqrt{x})`$, where $`P(\sqrt{x})`$ is a polynomial in $`\sqrt{x}`$ and the asymptotic behaviour for $`x0`$ and $`x1`$ (i.e. the values of $`\alpha `$ and $`\beta `$) are the same (except for $`\mathrm{\Delta }M`$, see a discussion below) as in the unpolarised case. Our idea is to split the numerical constants (coefficients of $`P`$ polynomial) in eqs.(3, 5, 6 and 7) in two parts in such a manner that the distributions $`q^\pm (x,Q^2)`$ and $`G^\pm (x,Q^2)`$ remain positive. At the end of the paper we will discuss the consequences of relaxing the positivity conditions. Our expressions for $`\mathrm{\Delta }q(x)=q^+(x)q^{}(x)`$ ($`q(x)=q^+(x)+q^{}(x)`$) are:
$`\mathrm{\Delta }u_v(x)`$ $`=`$ $`x^{0.5911}(1x)^{3.395}(a_1+a_2\sqrt{x}+a_4x),`$
$`\mathrm{\Delta }d_v(x)`$ $`=`$ $`x^{0.7118}(1x)^{3.874}(b_1+b_2\sqrt{x}+b_3x),`$
$`\mathrm{\Delta }M(x)`$ $`=`$ $`x^{0.7712}(1x)^{7.808}(c_1+c_2\sqrt{x}),`$ (8)
$`\mathrm{\Delta }\delta (x)`$ $`=`$ $`x^{0.183}(1x)^{9.808}c_3(1+9.987x33.34x^2),`$
$`\mathrm{\Delta }G(x)`$ $`=`$ $`x^{0.0829}(1x)^{6.587}(d_1+d_2\sqrt{x}+d_3x).`$
It is very important what assumptions one makes about the sea contribution. From the MRST fit for unpolarised structure functions the natural assumption would be: $`\mathrm{\Delta }\overline{s}=\mathrm{\Delta }\overline{d}/2=\mathrm{\Delta }\overline{u}/2`$. If we add the condition that SU(3) combination: $`a_8=\mathrm{\Delta }u+\mathrm{\Delta }d2\mathrm{\Delta }s`$ should be equal to the value determined from the semileptonic hyperon decays, $`\mathrm{\Delta }s`$ is pushed into negative values and so is nonstrange sea. Instead of connecting $`\mathrm{\Delta }s`$ in some way to non-strange sea value we introduce additional free parameters for the strange sea contribution namely
$$\mathrm{\Delta }M_s=x^{0.7712}(1x)^{7.808}(c_{1s}+c_{2s}\sqrt{x}).$$
(9)
For the strange quarks we have additional independent parameters. Hence, in our fits we will start with fourteen parameters. Comparing the expression (5) with (8) and (9) we see that in $`\mathrm{\Delta }M`$ (and $`\mathrm{\Delta }M_s`$) there is no term behaving like $`x^{1.2712}`$ at small $`x`$ (hence, we assume that $`\mathrm{\Delta }M`$ and hence all sea distributions have finite integral) which means that in $`\mathrm{\Delta }M`$ coefficient in front of this term have to be splitted into equal parts in $`\mathrm{\Delta }M^+`$ and $`\mathrm{\Delta }M^{}`$ (the most singular term in the sea contribution drops out). Hence, in the fitting procedure we are using functions that are suggested by the fit to unpolarised data. Maybe not all of them are important in the fit and it could happen that some of the coefficients in eqs.(8,9) taken as free parameters in the fit are small or in some sense superfluous. Putting them to zero or eliminating them increase $`\chi ^2`$ only a little but makes $`\chi ^2/N_{DF}`$ smaller. We will see that that is the case with some parameters introduced in eqs. (7,8). On the other hand we have still relatively strong singular behaviour of $`\mathrm{\Delta }u_v`$ and $`\mathrm{\Delta }M`$ for small $`x`$ values. For comparison we will also consider later the model in which most singular terms are put equal to zero namely $`a_1=c_1=c_{1s}=0`$, which means that plus and minus components have the same coefficients for this power of $`x`$. In the less singular models the dependence of calculated parameters in the unobserved region (below $`x0.003`$) is weak. In the earlier papers we considered the extrapolation of various calculated integrals below $`x=0.003`$ up to 0 assuming Regge type of behaviour for small $`x`$ values. As will be discussed later the less singular models give significantly higher $`\chi ^2`$.
In order to get the unknown parameters in the expressions for polarised quark and gluon distributions (eqs.(8,9)) we calculate the spin asymmetries starting from initial $`Q^2`$ = 1 $`\text{GeV}^2`$ for measured values of $`Q^2`$ and make a fit to the experimental data on spin asymmetries for proton, neutron and deuteron targets. The asymmetry $`A_1(x,Q^2)`$ can be expressed via the polarised structure function $`g_1(x,Q^2)`$ as
$$A_1(x,Q^2)\frac{(1+\gamma ^2)g_1(x,Q^2)}{F_1(x,Q^2)}=\frac{g_1(x,Q^2)}{F_2(x,Q^2)}[2x(1+R(x,Q^2))]$$
(10)
where $`R=[F_2(1+\gamma ^2)2xF_1]/2xF_1`$ whereas $`F_1`$ and $`F_2`$ are the unpolarised structure functions and $`\gamma =2Mx/Q`$. We will take the new determined value of $`R`$ from the . The factor $`(1+\gamma ^2)`$ plays non negligible role for $`x`$ and $`Q^2`$ values measured in SLAC experiments. Polarised structure function $`g_1(x,Q^2)`$ in the next to leading order QCD is related to the polarised quark and gluon distributions $`\mathrm{\Delta }q(x,Q^2)`$, $`\mathrm{\Delta }G(x,Q^2)`$, in the following way:
$`g_1(x,Q^2)`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}e_q^2\{\mathrm{\Delta }q(x,Q^2)+{\displaystyle \frac{\alpha _s}{2\pi }}[\delta c_q\mathrm{\Delta }q(x,Q^2)`$ (11)
$`+{\displaystyle \frac{1}{3}}\delta c_g\mathrm{\Delta }G(x,Q^2)]\}`$
with the convolution * defined by:
$$(Cq)(x,Q^2)=_x^1\frac{dz}{z}C(\frac{x}{z})q(z,Q^2)$$
(12)
The explicit form of the appropriate spin dependent Wilson coefficient $`\delta c_q`$ and $`\delta c_g`$ in the $`\overline{MS}`$ scheme can be found for example in ref. . The NLO expressions for the unpolarised (spin averaged) structure function is similar to the one in eq.(11) with $`\mathrm{\Delta }q(x,Q^2)q(x,Q^2)`$ and the unpolarised Wilson coefficients are given in .
The $`Q^2`$ evolution of the parton densities is governed by the DGLAP equations . For calculating the NLO evolution of the spin dependent parton distributions $`\mathrm{\Delta }q(x,Q^2)`$, $`\mathrm{\Delta }G(x,Q^2)`$ and spin averaged $`q(x,Q^2)`$, $`G(x,Q^2)`$ ones we will follow the method described in . We will calculate Mellin n-th moment of parton distributions $`\mathrm{\Delta }q(x,Q^2)`$ and $`\mathrm{\Delta }G(x,Q^2)`$ according to:
$$\mathrm{\Delta }q^{(n)}(Q^2)=_0^1𝑑xx^{n1}\mathrm{\Delta }q(x,Q^2)$$
(13)
and then use NLO solutions in Mellin n-moment space in order to calculate evolution in $`Q^2`$ for non-singlet and singlet parts.
In calculating evolution of $`\mathrm{\Delta }\mathrm{\Sigma }^{(n)}(Q^2)`$ and $`\mathrm{\Delta }G^{(n)}(Q^2)`$ with $`Q^2`$ we have mixing governed by the anomalous dimension 2x2 matrix . Having evolved moments one can insert them into the n-th moment of eq.(11).
$`g_1^{(n)}(Q^2)`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}e_q^2\{\mathrm{\Delta }q^{(n)}(Q^2)+{\displaystyle \frac{\alpha _s}{2\pi }}[\delta c_q^{(n)}\mathrm{\Delta }q^{(n)}(Q^2)`$ (14)
$`+{\displaystyle \frac{1}{3}}\delta c_g^{(n)}\mathrm{\Delta }G^{(n)}(Q^2)]\}`$
and then numerically Mellin invert the whole expression. In this way we get $`g_1(x,Q^2)`$. The same procedure is applied for the unpolarised structure functions. Having calculated the asymmetries according to equation (10) for the measured in experiments value of $`Q^2`$ we can make a fit to a measured asymmetries on proton, neutron and deuteron targets. We will take into account all together 417 points (193 for proton, 171 for deuteron and 53 for neutron. We will use also the ”experimental” value of $`a_8=0.58\pm 0.1`$ with enhanced (to 3$`\sigma `$) error.
The fit with all fourteen parameters from eqs.(8,9) gives $`\chi ^2=340.4`$. It seems that some of the parameters of the most singular terms are superfluous and we can eliminate them. We will put $`d_1=d_2=0`$ (such assumption gives that $`\delta G/Gx`$ for small $`x`$), $`b_1=0`$ (the most singular term in $`\mathrm{\Delta }d_v`$) and assume $`c_{1s}=c_1`$ (i.e. the most singular terms for strange and nonstrange sea contributions are equal). Fixing these four parameters in the fit practically does not change $`\chi ^2`$ but improves $`\chi ^2/N_{DF}`$. The resulting $`\chi ^2`$ per degree of freedom is better than in the previous fit and one gets $`\chi ^2/N_{DF}`$ =$`\frac{341.1}{41810}`$ =0.84. In this case we get the following values of parameters from the fit to all existing (above mentioned) data for $`Q^21\text{GeV}^2`$ for spin asymmetries:
$$\begin{array}{ccc}a_1=0.61\pm 0.00,\hfill & a_2=6.1\pm 0.19,\hfill & a_4=15.7\pm 0.42,\hfill \\ b_2=1.56\pm 0.20,\hfill & b_3=0.43\pm 0.49,\hfill & \\ c_1=0.40\pm 0.03,\hfill & c_2=4.15\pm 0.00,\hfill & \\ c_{1s}=c_1,\hfill & c_{2s}=0.28\pm 0.83,\hfill & \\ c_3=1.29\pm 2.53,\hfill & & \\ d_3=2.01\pm 11.2.\hfill & & .\hfill \end{array}$$
(15)
Actually the parameter $`d_3`$ could be put equal to zero without increasing $`\chi ^2/N_{DF}`$. We get in this case the smallest $`\chi ^2/N_{DF}`$ =$`\frac{341.1}{4188}`$ =0.83. That means that $`d_3`$ is not well determined in the fit and the best $`\chi ^2/N_{DF}`$ is without gluonic contribution.
The obtained quark and gluon distributions lead for $`Q^2`$ =1 $`GeV^2`$ to the following integrated (over $`x`$) quantities: $`\mathrm{\Delta }u=0.80\pm 0.02,\mathrm{\Delta }d=0.65\pm 0.03,\mathrm{\Delta }s=0.21\pm 0.05,\mathrm{\Delta }u_v=0.67\pm 0.02,\mathrm{\Delta }d_v=0.59\pm 0.02,2\mathrm{\Delta }\overline{u}=0.14\pm 0.03,2\mathrm{\Delta }\overline{d}=0.07\pm 0.03.`$
These numbers yield the following predictions: $`a_0=\mathrm{\Delta }u+\mathrm{\Delta }d+\mathrm{\Delta }s=0.06\pm 0.07,a_3=\mathrm{\Delta }u\mathrm{\Delta }d=1.45\pm 0.02,\mathrm{\Delta }G=0.04\pm 0.19,\mathrm{\Gamma }_1^p=0.111\pm 0.006,\mathrm{\Gamma }_1^n=0.096\pm 0.006,\mathrm{\Gamma }_1^d=0.007\pm 0.005.`$
We have positively polarised sea for up and negatively for down quarks and very strongly negatively polarised sea for strange quarks. Because of the big negative value of $`\mathrm{\Delta }s`$ the quantity $`a_0`$ is negative. The gluon polarisation is small. The value of $`a_3`$ was not assumed as an input in the fit (as is the case in nearly all fits ) and comes out slightly higher than the experimental value. $`\mathrm{\Delta }\delta `$, comes out relatively big from the fit (coefficient in front of $`\mathrm{\Delta }\delta `$ is equal to that in $`\delta `$).
The asymptotic behaviour at small $`x`$ of our polarised quark distributions is determined by the unpolarised ones and these do not have the expected theoretically Regge type behaviour or pQCD which is also used by some groups, to extrapolate results to small values of $`x`$. Some of the quantities in our fit change rapidly for $`x0.003`$.
Hence, we will present quantities integrated over the region from $`x`$=0.003 to $`x`$=1 (it is practically integration over the region which is covered by the experimental data, except of non controversial extrapolation for highest $`x`$). The corresponding quantities are $`\mathrm{\Delta }u=0.85`$ ($`\mathrm{\Delta }u_v=0.56`$, $`2\mathrm{\Delta }\overline{u}=0.29`$), $`\mathrm{\Delta }d=0.48`$ ($`\mathrm{\Delta }d_v=0.57`$, $`2\mathrm{\Delta }\overline{d}=0.09`$), $`\mathrm{\Delta }s=0.12`$, $`a_0=0.25`$, $`\mathrm{\Delta }G=0.04`$, $`\mathrm{\Gamma }_1^p=0.123`$, $`\mathrm{\Gamma }_1^n=0.056`$, $`\mathrm{\Gamma }_1^d=0.036`$, $`a_3=1.32`$. In this region the obtained values of sea contributions are relatively high and those of valence quarks relatively small. Gluon contribution practically vanishes. There is relatively strong dependence of different quantities in the unmeasured region ($`0x0.003`$). May be the unpolarised MRST parton distributions ( with the above mentioned modifications) do not describe quite correctly the small $`x`$ behaviour of polarised parton distributions. On the other hand the fit to the data is very good. So, the values of integrated quantities in the measured region we consider as more reliable then in the whole region. With the value of $`\mathrm{\Delta }s=0.12`$ in the measured region of $`x`$ we have $`a_0=0.25`$ and with $`\mathrm{\Delta }s=0.21`$ in the whole region of $`x`$ $`a_0`$ becomes negative (-0.06).
When we use the quantities calculated in the measured region and extend them to the full $`x`$ region using asymptotic Regge behaviour for small $`x`$ we get $`\mathrm{\Delta }u=0.86`$ ($`\mathrm{\Delta }u_v=0.59`$, $`2\mathrm{\Delta }\overline{u}=0.27`$), $`\mathrm{\Delta }d=0.51`$ ($`\mathrm{\Delta }d_v=0.58`$, $`2\mathrm{\Delta }\overline{d}=0.07`$), $`\mathrm{\Delta }s=0.14`$, $`a_0=0.21`$, $`\mathrm{\Delta }G=0.04`$, $`a_3=1.37`$. We have used $`x^\alpha `$ behaviour for small $`x`$ (with $`0.25\alpha 0.25`$) and the quantities do not depend strongly on a specific value of $`\alpha `$. For the given above values $`\alpha =0`$ was used.
Now, we shall calculate $`\mathrm{\Gamma }^p`$, $`\mathrm{\Gamma }^n`$ and $`\mathrm{\Gamma }^d`$ in the measured region for $`Q^2`$ = 5 $`\text{GeV}^2`$ and compare them with the quantities given by the experimental groups. We get in the region between $`x=0.003`$ and $`x=0.8`$ (covered by the data) $`\mathrm{\Gamma }_1^p=0.132\pm 0.006`$, $`\mathrm{\Gamma }_1^n=0.051\pm 0.007`$ and $`\mathrm{\Gamma }_1^d=0.037\pm 0.006`$. The experimental group SMC present the following values in such region (for $`Q^2`$ = 5 $`\text{GeV}^2`$):
$`\mathrm{\Gamma }_1^p`$ $`=`$ $`0.130\pm 0.007,`$
$`\mathrm{\Gamma }_1^n`$ $`=`$ $`0.054\pm 0.009,`$ (16)
$`\mathrm{\Gamma }_1^d`$ $`=`$ $`0.036\pm 0.005.`$
One can see that our results are in good agreement with experimental values. For comparison we have also made fits using formulas of the simple parton model (as in our papers before ) neglecting evolution of parton densities with $`Q^2`$. More detailed result of these fits (integrated densities and so on) will be given later.
In fig.1 as an example we present our fit to the non averaged data in comparison with measured (averaged over $`Q^2`$) $`g_1/F_1`$ for new proton (Hermes) and deuteron (E155) data. The curves are obtained by joining the calculated values of asymmetries corresponding to actual values of $`x`$ and $`Q^2`$ for measured data points. The curves are not fitted but the difference in fitted asymmetries for averaged and non-averaged data are very small. For asymmetries the curves with $`Q^2`$ evolution taken into account and evolution completely neglected do not differ very much so we do not present them.
In figs. 2 and 3 we show the comparison of our predictions for $`g_1`$ from the basic fit with the measured averaged values for proton, deuteron and neutron data. The values of $`g_1`$ were calculated for the values of $`x`$ and $`Q^2`$ measured for averaged data points in different experiments and then joined together. The agreement is good. On the other hand the spread of experimental points is still substantial. Polarised quark distributions for up and down valence quarks as well as non strange, strange quarks and gluons for $`Q^2`$ = 1 $`\text{GeV}^2`$ are presented in figure 4. Dashed curves represent the $`+`$ and $``$ components for different parton densities. the solid curves corespond to the difference of $`+`$ and $``$ components, the sums of components (not shown) correspond to nonpolarised parton distributions. We see that especially polarised gluon distribution function is really tiny and does not resemble the distribution function for unpolarised case.
This function is also quite different from the gluon distribution (given in ) used to estimate $`\mathrm{\Delta }G/G`$ in COMPASS experiment planned at CERN . For $`x=0.1`$ at $`Q^2=1\mathrm{GeV}^2`$ $`\mathrm{\Delta }G/G=0.01`$ and is below a planned experimental error. In fig.5 we present spin densities for $`u`$ quark. We show $`\mathrm{\Delta }u_v(x,Q^2)`$ obtained in the basic fit (solid line) as well as the same quantity from the fit with $`SU(3)`$ symmetric sea (dashed curve) and from the fit where evolution in $`Q^2`$ was not taken into account (dotted line). In the second plot the same curves for $`\mathrm{\Delta }u(x,Q^2)`$ are presented and one sees from that the values of $`\mathrm{\Delta }u`$ from the basic fit and fit with $`SU(3)`$ symmetric sea are very close in spite of the differences in the valence values. The similar conclusions are also true for the $`d`$ quark densities.
Fixing the value of $`a_8`$ is very important for the fit. When we relax the condition for $`a_8=0.58`$ we get $`\chi ^2=340.8`$, so this number practically does not change. We get the fit with the parameters not very different from our basic fit but with very small $`a_8=0.03`$ and positive $`\mathrm{\Delta }s=0.01`$. It seems that $`\mathrm{\Delta }s`$ is not well determined from the data on spin asymmetries alone but that does not influences strongly the values of $`\mathrm{\Delta }u`$, $`\mathrm{\Delta }d`$ and $`\mathrm{\Delta }G`$.
In order to check how the fit depends on the assumptions made about the sea contribution we have also made fit with $`\mathrm{\Delta }\overline{u}=\mathrm{\Delta }\overline{d}=2\mathrm{\Delta }\overline{s}`$, the assumption that follows directly from MRST unpolarised fit. The $`\chi ^2`$ value increases significantly and per degree of freedom one gets a number $`\chi ^2/N_{DF}`$ =$`\frac{353.2}{4189}`$ =0.86 which is worse than in our basic fit. In this case we have $`\mathrm{\Delta }u=0.80`$ ($`\mathrm{\Delta }u_v=0.87`$, $`2\mathrm{\Delta }\overline{u}=0.07`$), $`\mathrm{\Delta }d=0.61`$ ($`\mathrm{\Delta }d_v=0.40`$, $`2\mathrm{\Delta }\overline{d}=0.21`$), $`\mathrm{\Delta }s=0.07`$, $`a_0=0.11`$, $`\mathrm{\Delta }G=0.07`$ and $`a_8=0.33`$. The quantity $`\mathrm{\Delta }s`$ must be negative in order to get experimental value for $`a_8`$ and because of our assumption $`\mathrm{\Delta }\overline{u}=\mathrm{\Delta }\overline{d}=2\mathrm{\Delta }\overline{s}`$ (the value of $`a_8`$ does not come out correctly in the fit because of our assumption $`c_{1s}=c_1`$) we obtain negative values of non strange sea for up and down quarks. One sees that the values for sea polarisation depend very strongly on assumptions we made(in many papers SU(3) symmetric sea is assumed that also together with fixing of $`a_8`$ value gives negative non strange sea). On the other hand $`\mathrm{\Delta }u=\mathrm{\Delta }u_v+2\mathrm{\Delta }\overline{u}`$ and $`\mathrm{\Delta }d=\mathrm{\Delta }d_v+2\mathrm{\Delta }\overline{d}`$ practically do not change (however, $`\mathrm{\Delta }u_v`$ and $`\mathrm{\Delta }d_v`$ also change). Also $`\mathrm{\Delta }G`$ does not change and is small. We get that the value of $`a_8`$ is not coming out correctly from the fit to spin asymmetries. Fixing $`a_8`$ and making specific assumption about $`\mathrm{\Delta }s`$ introduces shifts in nonstrange sea polarisation (and so in $`\mathrm{\Delta }u_v`$ and $`\mathrm{\Delta }d_v`$) but $`\mathrm{\Delta }u`$, $`\mathrm{\Delta }d`$ do not change. Because the value of $`\mathrm{\Delta }s`$ is needed to determine $`a_0`$ we decided to use additional free parameters for strange sea contribution in order to determine it (with fixed value of $`a_8`$) from the fit to experimental data.
In many papers making fits to experimental data on spin asymmetries the assumption of $`SU(3)`$ symmetric sea was made. We have also for comparison made fit with this assumption. The value of $`\chi ^2=350.6`$ is substantially higher in comparison with our basic fit. In this case we have $`\mathrm{\Delta }u=0.81`$ ($`\mathrm{\Delta }u_v=0.87`$, $`2\mathrm{\Delta }\overline{u}=0.06`$), $`\mathrm{\Delta }d=0.57`$ ($`\mathrm{\Delta }d_v=0.40`$, $`2\mathrm{\Delta }\overline{d}=0.17`$), $`\mathrm{\Delta }s=0.11`$, $`a_0=0.13`$, $`\mathrm{\Delta }G=0.14`$, $`a_3=1.38`$ and $`a_8=0.47`$. In this case we also have shifts in values of valence and sea contributions similar to the case discussed above.
Looking at the dependence of unpolarised quark and gluon densities we see that after elimination of most singular term in $`\mathrm{\Delta }d_v(x)`$ the most singular behaviour for small $`x`$ one has for $`\mathrm{\Delta }u_v(x)`$ and $`\mathrm{\Delta }M(x)`$. For comparison we have investigated the model when in polarised densities these singular contributions are absent. In this case $`\mathrm{\Delta }u_v`$ and $`\mathrm{\Delta }M`$ are $`\sqrt{x}`$ less singular than in our basic fit. For such a fit we get $`\chi ^2/N_{DF}`$ =$`\frac{356.6}{4188}`$ =0.87, i.e significantly higher than in our basic fit. We get in this case: $`\mathrm{\Delta }u=0.77`$ ($`\mathrm{\Delta }u_v=0.57`$, $`2\mathrm{\Delta }\overline{u}=0.20`$), $`\mathrm{\Delta }d=0.38`$ ($`\mathrm{\Delta }d_v=0.63`$, $`2\mathrm{\Delta }\overline{d}=0.25`$), $`\mathrm{\Delta }s=0.10`$, $`a_0=0.28`$, $`\mathrm{\Delta }G=0.22`$. In such fit the integrated quantities taken over the whole range of $`0x1`$ and in the truncated one ($`0.003x1`$) differ very little. The quantity $`\mathrm{\Delta }G`$ is positive and different from zero. So it is possible to get the fit with practically no change of integrated quantities in the region between $`x=0`$ and $`x=0.003`$ but with significantly higher $`\chi ^2`$ value. For $`Q^2`$=1 $`\text{GeV}^2`$ we have $`\mathrm{\Gamma }_1^p=0.122`$ and $`\mathrm{\Gamma }_1^n=0.041`$.
The obtained results can be compared with the fit when instead of 417 points for different $`x`$ and $`Q^2`$ values we take spin asymmetries for only 160 data points with the averaged $`Q^2`$ values for the same $`x`$ (one has smaller errors in this case). In such fit the ratio of number of neutron to number of deuteron and proton data points is increased. It seems that the influence of neutron points is stronger than in basic fit $`\chi ^2/N_{DF}`$ =$`\frac{118.3}{16110}`$ = 0.78 is a little bit better than in our basic fit. The integrated values for quark and gluon densities are: $`\mathrm{\Delta }u=0.79`$ ($`\mathrm{\Delta }u_v=0.65`$, $`2\mathrm{\Delta }\overline{u}=0.14`$), $`\mathrm{\Delta }d=0.66`$ ($`\mathrm{\Delta }d_v=0.60`$, $`2\mathrm{\Delta }\overline{d}=0.06`$), $`\mathrm{\Delta }s=0.22`$, $`a_0=0.09`$, $`\mathrm{\Delta }G=0.31`$ and $`a_3=1.45`$. We see that averaging over $`Q^2`$ and different numbers of data points leads to very similar fit. The values for integrated valence densities and nonstrage sea contribution are only a bit shifted ($`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$ in the whole region of integration do not differ from the basic fit for non averaged data and the same is also true for integrated quantities in the region $`0.003x1`$). Integrated gluon density is relatively small and positive. A little bit higher value for $`\mathrm{\Delta }G=0.31\pm 0.28`$ we do not consider as significant difference. $`\chi ^2/N_{DF}`$ is very good and smaller then in where the same experimental data sample and MRST parton distributions (modified for small values of $`x`$) were used.
As was already mentioned before we have also made for comparison fits neglecting evolution of parton densities with $`Q^2`$ (formulas from the simple parton model). We get for non averaged data sample $`\chi ^2/N_{DF}`$ =$`\frac{349.9}{4189}`$=0.86 (biger than in our basic fit: $`\chi ^2/N_{DF}`$ =0.84): $`\mathrm{\Delta }u=0.66`$ ($`\mathrm{\Delta }u_v=0.56`$, $`2\mathrm{\Delta }\overline{u}=0.10`$), $`\mathrm{\Delta }d=0.49`$ ($`\mathrm{\Delta }d_v=0.49`$, $`2\mathrm{\Delta }\overline{d}=0.0`$), $`\mathrm{\Delta }s=0.20`$, $`a_0=0.03`$, $`a_3=1.14`$, $`\mathrm{\Gamma }_1^p=0.108`$, $`\mathrm{\Gamma }_1^n`$ = -0.082. For averaged data points we get $`\chi ^2/N_{DF}`$ =$`\frac{125.4}{1619}`$=0.83 (this number should be compared with $`\chi ^2/N_{DF}`$ =0.78, the corresponding quantity from the NLO fit) and we have: $`\mathrm{\Delta }u=0.66`$ ($`\mathrm{\Delta }u_v=0.58`$, $`2\mathrm{\Delta }\overline{u}=0.08`$), $`\mathrm{\Delta }d=0.48`$ ($`\mathrm{\Delta }d_v=0.48`$, $`2\mathrm{\Delta }\overline{d}=0.0`$), $`\mathrm{\Delta }s=0.20`$, $`a_0=0.03`$. Hence, $`\chi ^2`$ per degree of freedom is smaller in the case of averaged sample. We see that both fits give very similar results. It means that the averaging of data does not influence the fit when we do not take $`Q^2`$ evolution into account (the differences are also very small in the $`0.003x1`$ region).
It has been pointed out that the positivity conditions could be restrictive and influence the contribution of polarised gluons. We have also made a fit to experimental data without such assumption for polarised partons. The $`\chi ^2`$ value does not changed much $`\chi ^2/N_{DF}`$ =$`\frac{340.7}{41810}`$ =0.84 and we get $`\mathrm{\Delta }u=0.84`$ ($`\mathrm{\Delta }u_v=0.72`$, $`2\mathrm{\Delta }\overline{u}=0.12`$), $`\mathrm{\Delta }d=0.74`$ ($`\mathrm{\Delta }d_v=0.50`$, $`2\mathrm{\Delta }\overline{d}=0.24`$), $`\mathrm{\Delta }s=0.24`$, $`a_0=0.13`$, $`a_3=1.57`$, $`\mathrm{\Delta }G=0.02`$. The results are a little bit different but the value of $`\mathrm{\Delta }G`$ is not influenced by the positivity conditions. The same is also true in the case of averaged data. It seems that our positivity conditions are not very restrictive.
We have made fits for two samples of data with averaged $`Q^2`$ values and non averged ones (adding neutron data from E154 and Hermes experiments) leading to very similar results for calculated parameters (except small difference in $`\mathrm{\Delta }G`$). The value of $`a_3`$ was not fixed in the fit and comes out high in comparison with experimental value. In order to check the influence of different assumptions about strange sea we considered fits without fixing $`a_8`$ value, with $`SU(3)`$ symmetric sea and with modified sea contribution. The models with less singular behaviour for valence $`u`$ quark and sea contribution were also discussed. In most of the modifications the $`\chi ^2`$ value increases significantly. For comparison we have also considered fits to the simple parton model neglecting $`Q^2`$ dependence of parton densities. It seems that splitting of integrated densities $`\mathrm{\Delta }u`$, $`\mathrm{\Delta }d`$ into valence and sea contribution is model dependent ($`\mathrm{\Delta }u`$ and $`\mathrm{\Delta }d`$ do not differ much). The integrated gluon contribution comes out small. The best fits (measured by $`\chi ^2`$ per degree of freedom) we have for zero (for non averaged data points) or rather small (for averaged data) gluon polarisation. It seems that from results of our fits the perspective of measuring in COMPASS $`\mathrm{\Delta }G/G`$ is not very encouraging. The experimental accuracy still must be improved and probabely additional experiments are needed in order to make more precise statements about polarised quark and gluon densities.
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# 1 Introduction
## 1 Introduction
Killing tensors are indispensable tools in the quest for exact solutions in many branches of general relativity as well as classical mechanics . Killing tensors are important for solving the equations of motion in particular space-times.The notable example here is the Kerr metric which admits a second rank Killing tensor .Killing tensors give rise to new exact solutions in perfect fluid Bianchi and Katowski-Sachs cosmologies as well in inflationary models with a scalar field sources . Recently Killing tensors of third rank in $`(1+1)`$ dimensional geometry were investigated and classified .Even more recently the Killing tensors of order two associated with orthogonal separable coordinates for the Klein-Gordon equation in flat 2+1 dimensional space-time were considered as metrics . In a geometrical setting ,symmetries are connected with isometries associated with Killing vectors, and more generally, Killing tensors on the configuration space of the system.An example is the motion of a point particle in a space with isometries ,which is a physicist’s way of studying the geodesic structure of a manifold .We recall that $`K_{\alpha \beta }`$ is a Killing tensor if and only if , for any geodesic motion of a test particle with a world velocity $`p^\alpha `$,the scalar $`K_{\alpha \beta }p^\alpha p^\beta `$ is a constant of motion .The Jacobi’s geometrical model of dynamical systems with a finite number of degrees of freedom was investigated by many authors (see for example Refs.).The essential conclusion was that :the paths of the motions of a dynamical system in the configuration space are identical with the geodesics of the Riemannian manifold obtained by providing the configuration space with the metric given by
$$ds^2=g_{ij}dq_idq_j=2(EV)a_{ij}dq_idq_j.$$
(1)
We mention here that $`T=\frac{1}{2}a_{ij}\dot{q}_i\dot{q}_j`$(the dot signifying derivation with respect to time and $`a_{ij}`$ are functions of the q’s),V is a function of the q’s only and $`T+V=E`$. In it was pointed out that a single Lax tensor may generate an infinite number of tensors of varying ranks.It is very well known that the most general constant on a geodesic motion is of the form
$$K=K_0+\chi _\mu p^\mu +K_{\mu \nu }p^\mu p^\nu +K_{\mu \nu \lambda }p^\mu p^\nu p^\lambda +\mathrm{},$$
(2)
where $`K_0`$ is a constant of motion on the geodesic , $`\chi _\mu `$ is a Killing vector and $`K_{\mu _1\mathrm{}\mu _n}`$ is a Killing tensor of order n. The important point is that if we are using Jacobi’s geometrical model a natural way to produce Killing tensors is to consider the elements of the Lax matrix $`L_{\alpha \beta }`$ as $`L_{\alpha \beta }=L_{\alpha \beta }^\gamma p_\gamma +C_{\alpha \beta }`$ .Here $`C_{\alpha \beta }`$ is a matrix having the elements satisfying the following relations $`tr(C_{\alpha \beta })=K_0`$,$`L_{\alpha \beta }^\alpha =\chi _\beta `$, $`K_{\alpha \beta }=L_{\nu \alpha }^\mu L_{\mu \beta }^\nu `$ and so on. Open three dimensional Toda’s case analyzed in is a special and very interesting case because the Lax tensor generates a Killing tensor of order two which is equal to the metric tensor.We know that in this case Killing tensor of order two is called trivial (see for more details ). Recently the geometric duality between a metric $`g^{\mu \nu }`$ and its non-degenerate Killing tensor $`K^{\mu \nu }`$ and the structural equations of a Killing tensor of order two were analyzed in . An interesting example arises when the manifold admits Killing-Yano tensors because they generate Killing tensors.In addition we know that any manifold having constant curvature admits Killing-Yano tensors and then it admits Killing tensors.
For these reasons the Lax tensor equations on a given manifold and its dual are interesting to investigate.
The plan of this paper is as follows:
In Section 2 the Lax pair tensors are investigated.In Section 3 the geometric duality is presented and the Lax tensors on the dual manifolds are analyzed.In Section 4 the examples are presented.Section 5 contains our comments and remarks.
## 2 Lax pair tensors
Let us consider a Riemannian or pseudo-Riemannian geometry with the metric
$$ds^2=g_{\mu \nu }dq^\mu dq^\nu .$$
(3)
The geodesic equation can be represented by the Hamiltonian
$$H=\frac{1}{2}g^{\mu \nu }p_\mu p_\nu ,$$
(4)
together with the natural Poisson bracket on the cotangent bundle.The geodesic system has the form
$$\dot{q}^\alpha =g^{\alpha \mu }p_\mu ,\dot{p}_\alpha =\mathrm{\Gamma }_\alpha ^{\mu \nu }p_\mu p_\nu .$$
(5)
The complete integrability of this system can be shown with the help of a pair of matrices L and A with entries defined on the phase space and satisfying the Lax pair equation .
$$\dot{L}=\{L,H\}=[L,A].$$
(6)
It follows from (6) that the quantities $`I_k=\frac{1}{k}TrL^k`$ are all constants of motion. If in addition they commute with each other $`\{I_k,I_j\}=0`$ then it is possible to integrate the system completely at least in principle. We know that Lax pair equation is invariant under a transformation of the form
$$\stackrel{~}{U}=ULU^1,\stackrel{~}{A}=UAU^1\dot{U}U^1.$$
(7)
We see that L transforms as a tensor while A transforms as a connection. Typically,the Lax matrices are linear in the momenta and in the geometric setting that may also be assumed to be homogeneous. This motivates the introduction of two third rank geometrical objects $`L_{\beta }^{\alpha }{}_{}{}^{\gamma }`$ and $`A_{\beta }^{\alpha }{}_{}{}^{\gamma }`$ such that the Lax matrices can be written as
$$L=(L_\beta ^\alpha )=(L_{\beta }^{\alpha }{}_{}{}^{\mu }p_\mu ),A=(A_\beta ^\alpha )=(A_{\beta }^{\alpha }{}_{}{}^{\mu }p_\mu ).$$
(8)
We will refer to $`L_{\beta }^{\alpha }{}_{}{}^{\gamma }`$ and $`A_{\beta }^{\alpha }{}_{}{}^{\gamma }`$ as the Lax tensor and the Lax connection , respectively.Defining
$$B=(B_\beta ^\alpha )=(B_{\beta }^{\alpha }{}_{}{}^{\mu }p_\mu )=A\mathrm{\Gamma },$$
(9)
where $`\mathrm{\Gamma }=(\mathrm{\Gamma }_\beta ^\alpha )=(\mathrm{\Gamma }_{\beta }^{\alpha }{}_{}{}^{\mu }\mathrm{p}_\mu )`$ is the Levi-Civita connection with respect to $`g_{\alpha \beta }`$, it then follows that the Lax pair equation takes the covariant form. Let us suppose that a manifold $`g_{\mu \nu }`$ admits a Lax pair tensors $`L_{\alpha \beta \gamma },A_{\alpha \beta \gamma }`$ in such a way that
$$L_{\alpha \beta \gamma ;\delta }+L_{\alpha \beta \delta ;\gamma }=L_{\alpha \mu (\gamma }B_{|\beta |\delta )}^\mu B_{\alpha \mu (\gamma }L_{|\beta |\delta )}^\mu .$$
(10)
Here the parenthesis denotes the full symmetrization. We know that a Killing tensor of order n is a symmetric tensor $`K_{\mu _1\mathrm{}\mu _n}`$ which satisfies the following relation:
$$D_{(\lambda }K_{\mu _1\mathrm{}\mu _n)}=0$$
(11)
where $`D_\mu `$ denote covariant derivative. and Using (10) ,in the case when $`L_{\alpha \beta \gamma }`$ has only symmetric part, we find immediately that $`L_{(\alpha \beta \gamma ;\delta )}=0`$ for $`B_{\beta \gamma }^\alpha =0`$ . Then it is a Killing tensor of order three. Any solution of (10) generates an infinite number of Killing tensors on a given manifold.Of course not all Killing tensors generated by Lax tensors are independent and some of them are trivial Killing tensors .Another important observation is that in the case when we have $`g_{\alpha \beta }=L_{\nu \alpha }^\mu L_{\mu \beta }^\nu `$ we can identify the invariant $`I_2`$ with the geodesic Hamiltonian.
Let us suppose that the manifold admits a Killing tensors $`K_{\alpha \beta }`$ and define a three dimensional tensor as
$$L_{\alpha \beta \gamma }=K_{\beta \gamma ;\alpha }K_{\alpha \gamma ;\beta }.$$
(12)
We conclude immediately that it has the symmetries
$$L_{\alpha \beta \gamma }=L_{[\alpha \beta ]\gamma },L_{[\alpha \beta \gamma ]}=0,$$
(13)
where square brackets denote the anti-symmetrization. After an appropriate grouping of terms and use of the symmetries of the Riemann tensor $`R_{\alpha \beta \gamma \delta }`$ we obtain
$$L_{\alpha \beta (\gamma ;\delta )}=2R_{\alpha \beta \mu (\gamma }K_{\delta )}^\mu 2K_{[\alpha }^\mu R_{\beta ]\mu (\gamma \delta )\mu }.$$
(14)
We are interested now to investigate if (12) satisfies (10).In other words our problem is to find a tensor $`B_{\alpha \beta \gamma }`$ in such a way that (10) is satisfied.Using (12) and (10) we conclude that $`B_{\alpha \beta \gamma }=B_{\beta \alpha \gamma }`$. Let us denote $`V_{\alpha \beta \gamma }`$ as
$`V_{\alpha \beta \gamma }=L_{[\alpha \beta ]\gamma }`$.Taking into account (10) and (14) we find that
$$V_{\alpha \mu (\gamma }B_{|\beta |\delta )}^\mu V_{\beta \mu (\gamma }B_{|\alpha |\delta )}^\mu =R_{\alpha \beta \mu (\gamma }K_{\delta )}^\mu +K_{[\alpha }^\mu R_{\beta ]\mu (\gamma \delta )\mu }.$$
(15)
Solving (15) we can determine $`B_{\beta \gamma }^\alpha `$.
Using definition of a Killing tensor of order two and (12) we get
$$K_{\beta \gamma ;\alpha }=\frac{2}{3}L_{\alpha (\beta \gamma )}.$$
(16)
Conversely , (16) and the conditions $`L_{\alpha \beta \gamma }=L_{\beta \alpha \gamma },L_{[\alpha \beta \gamma ]}=0`$ imply
(12) and that $`K_{\alpha \beta }`$ is a Killing tensor.
Another interesting case is when $`L_{\alpha \beta \gamma }=V_{\alpha \beta \gamma }`$ and in addition we suppose that $`B_{\alpha \beta \gamma }=L_{\alpha \beta \gamma }`$.For this case the Lax equations become
$$V_{\alpha \beta (\gamma ;\delta )}=0,$$
(17)
and we see that (17) looks like Killing-Yano equations.
## 3 Geometric Duality
Let us suppose that the metric $`g_{\mu \nu }`$ admits a Killing tensor field $`K_{\mu \nu }`$.
From the covariant components $`K_{\mu \nu }`$ of the Killing tensor one can construct a constant of motion $`K=\frac{1}{2}K_{\mu \nu }p^\mu p^\nu `$. It can be verified that $`\{H,K\}=0`$. The formal similarity between the constants of motion H and K , and the symmetrical nature of the condition implying the existence of the Killing tensor amount to a reciprocal relation between two different models:the model with Hamiltonian H and constant of motion K, and a model with constant of motion H and Hamiltonian K.The relation between the two models has a geometrical interpretation: it implies that if $`K_{\mu \nu }`$ are the contravariant components of a Killing tensor with respect to the metric $`g_{\mu \nu }`$, then $`g_{\mu \nu }`$ must represent a Killing tensor with respect to the metric defined by $`K_{\mu \nu }`$. When $`K_{\mu \nu }`$ has an inverse we interpret it as the metric of another space and we can define the associated Riemann-Christoffel connection $`\widehat{\mathrm{\Gamma }}_{\mu \nu }^\lambda `$ as usual through the metric postulate $`\widehat{D}_\lambda K_{\mu \nu }=0`$. Here $`\widehat{D}`$ represents the covariant derivative with respect to $`K_{\mu \nu }`$. This reciprocal relation between the metric structure of pairs of spaces constitutes a duality relation: performing the operation of mapping a Killing tensor to a metric twice leads back to the original theory.
The relation between connections $`\widehat{\mathrm{\Gamma }}_{\alpha \beta }^\sigma `$ and $`\mathrm{\Gamma }_{\alpha \beta }^\sigma `$ is
$$\widehat{\mathrm{\Gamma }}_{\alpha \beta }^\mu =\mathrm{\Gamma }_{\alpha \beta }^\mu K^{\mu \delta }D_\delta K_{\alpha \beta }.$$
(18)
In the case when the tensor $`B_{\beta \delta }^\alpha `$ is symmetric in the lower indices and has the form
$$B_{\beta \delta }^\alpha =K^{\alpha \omega }D_\omega K_{\beta \delta }$$
(19)
then (10) becomes
$$\widehat{L}_{\alpha \beta \gamma ;\delta }+\widehat{L}_{\alpha \beta \delta ;\gamma }=0.$$
(20)
Here comma represents the covariant derivative in the dual space. We are interested now to investigate when the original space and its dual admit the same Lax tensors.
Proposition
The manifold and its dual have the same Lax tensors iff
$`2(K^{\sigma \omega }D_\omega K_{\gamma \delta })L_{\alpha \beta \sigma }`$ $`+(K^{\sigma \omega }D_\omega K_{\alpha \delta })L_{\sigma \beta \gamma }+(K^{\sigma \omega }D_\omega K_{\alpha \gamma })L_{\sigma \beta \delta }+`$ (21)
$`(K^{\sigma \omega }D_\omega K_{\beta \delta })L_{\alpha \sigma \gamma }`$ $`+(K^{\sigma \omega }D_\omega K_{\beta \gamma })L_{\alpha \sigma \delta }=0.`$ (22)
Proof.
Let us consider $`L_{\alpha \beta \gamma }`$ the Lax tensor satisfies
$$L_{\alpha \beta (\gamma ;\delta )}=0$$
(23)
and $`\widehat{L}_{\alpha \beta \gamma }`$ be the dual Lax tensor. Using (18) the corresponding dual Lax equations are
$`D_\delta \widehat{L}_{\alpha \beta \gamma }`$ $`+D_\gamma \widehat{L}_{\alpha \beta \delta }+2(K^{\sigma \omega }D_\omega K_{\gamma \delta })\widehat{L}_{\alpha \beta \sigma }+(K^{\sigma \omega }D_\omega K_{\alpha \delta })\widehat{L}_{\sigma \beta \gamma }+`$ (24)
$`(K^{\sigma \omega }D_\omega K_{\alpha \gamma })\widehat{L}_{\sigma \beta \delta }`$ $`+(K^{\sigma \omega }D_\omega K_{\beta \delta })\widehat{L}_{\alpha \sigma \gamma }+(K^{\sigma \omega }D_\omega K_{\beta \gamma })\widehat{L}_{\alpha \sigma \delta }=0.`$ (25)
Let us suppose that $`\widehat{L}_{\alpha \beta \gamma }=L_{\alpha \beta \gamma }`$, then using (23) and (24) we obtain (21) Conversely if we suppose that (21) holds , then from (24) we can deduce immediately that $`L_{\alpha \beta \gamma }=\widehat{L}_{\alpha \beta \gamma }`$. q.e.d.
## 4 Examples
In this section we will present some examples when the equations (10) admit solutions.
A.
Let us consider the n-dimensional Euclidean space and first let us investigate the Lax equations corresponding to $`B_{\beta \gamma }^\alpha =0`$.Then (10) becomes
$$\frac{L_{\alpha \beta \gamma }}{x^\delta }+\frac{L_{\alpha \beta \delta }}{x^\gamma }=0.$$
(26)
The solution of this equation has the form
$$L_{\alpha \beta \gamma }=T_{\alpha \beta \gamma \sigma }x^\sigma +V_{\alpha \beta \gamma },$$
(27)
where $`T_{\alpha \beta \gamma \sigma }`$ and $`V_{\alpha \beta \gamma }`$ are constant tensors and in addition $`T_{\alpha \beta \gamma \sigma }=T_{\alpha \beta \sigma \gamma }`$.
B.
Let us consider now the two dimensional metric
$$ds^2=f(u,v)du^2+g(u,v)dv^2.$$
(28)
We are interested to investigate the Lax tensors when $`L_{\alpha \beta \gamma }`$ is symmetric and $`A_{\alpha \beta \gamma }=\mathrm{\Gamma }_{\alpha \beta \gamma }`$. The non-vanishing Christoffel symbols of (28) are
$`\mathrm{\Gamma }_{11}^1=`$ $`{\displaystyle \frac{\frac{f}{u}}{2f}},\mathrm{\Gamma }_{11}^2={\displaystyle \frac{\frac{f}{v}}{2g}},\mathrm{\Gamma }_{21}^1=\mathrm{\Gamma }_{12}^1={\displaystyle \frac{\frac{f}{v}}{2f}},`$ (29)
$`\mathrm{\Gamma }_{22}^2=`$ $`{\displaystyle \frac{\frac{g}{u}}{2f}},\mathrm{\Gamma }_{22}^2={\displaystyle \frac{\frac{g}{v}}{2g}},\mathrm{\Gamma }_{21}^1=\mathrm{\Gamma }_{12}^2={\displaystyle \frac{\frac{g}{u}}{2g}}.`$ (30)
$`L_{\alpha \beta \gamma }`$ has four independent components $`L_{111},L_{112},L_{122},L_{222}`$ and the independent Lax equations are
$`L_{11(1;u)}=0`$ $`,L_{11(1;v)}=0,L_{11(2;u)}=0,L_{11(2;v)}=0,`$ (31)
$`L_{12(2;u)}=0`$ $`,L_{12(2;v)}=0,L_{22(2;u)}=0,L_{22(2;v)}=0.`$ (32)
We found after some calculations that if the scalar curvature of the manifold corresponding to (28) is 0 then the system (31) is integrable.
Let us consider now the Rindler system. The Rindler system is conventionally denoted by $`\tau `$ and r
$$t=r\mathrm{sinh}\tau ,x=r\mathrm{cos}\tau ,0<r<\mathrm{},\mathrm{}<\tau <\mathrm{},$$
(33)
with coordinates curves (timelike hyperbolas and spacelike straight lines) given by
$$x^2t^2=r^2,\frac{t}{x}=\mathrm{tanh}\tau ,$$
(34)
the metric
$$ds^2=r^2d\tau ^2dr^2,$$
(35)
and the associated Killing tensor
$$k^{ik}=\left(\begin{array}{ccc}& 1\frac{c}{r^2}& 0\\ & 0& c\end{array}\right).$$
(36)
Here c is a constant. The non-zero Christoffel symbols are $`\mathrm{\Gamma }_{11}^2=\mathrm{r}`$, $`\mathrm{\Gamma }_{12}^1=\frac{1}{\mathrm{r}}`$. Solving (31) we found the solution of the Lax equations having the form
$`L_{122}`$ $`=(C_1\mathrm{e}^\tau +C_2\mathrm{e}^{3\tau })r,L_{112}=(C_1\mathrm{e}^\tau +C_2\mathrm{e}^{3\tau })r^2,`$ (37)
$`L_{111}`$ $`=(3C_1\mathrm{e}^\tau C_2\mathrm{e}^{3\tau })r^3,L_{222}=3C_1\mathrm{e}^\tau +C_2\mathrm{e}^{3\tau },`$ (38)
where $`C_1,C_2`$ are constants.
The next step is to find a solution of the form (12) corresponding to the Rindler system. Using (12) and (36) we found immediately the solution having the form
$$L_{121}=\frac{r^3(r^23c)^2}{c(r^2c)^2},L_{122}=0.$$
(39)
Let us consider now the tensor
$$K^{\mu \nu }=g^{\mu \lambda }g^{\nu \delta }K_{\lambda \delta }$$
(40)
and a connection defined as in (18). Using (18) and taking into account (40)and (36) we found a new metric having the non-zero components
$$d\widehat{s}^2=r^2d\tau ^2+\frac{c(r^2c)^2}{(2r^2+r^4c2r^2c^2+c^3)}dr^2.$$
(41)
The scalar curvature corresponding to (41) is $`R=\frac{4(r^2+c)}{c(r^2+c)^3}`$ ,then this metric has no symmetric Lax tensors.
## 5 Concluding remarks
In this paper we investigated the Lax equations on a given manifold and its dual. When a manifold admits a Killing tensor $`K_{\mu \nu }`$ we constructed a tensor $`L_{\alpha \beta \gamma }`$ as $`L_{\alpha \beta \gamma }=K_{\beta \gamma ;\alpha }K_{\alpha \gamma ;\beta }`$ and found the conditions when it is a Lax tensor.In this case $`L_{\alpha \beta \gamma }`$ is antisymmetric in the first two indices and $`B_{\alpha \beta \gamma }`$ should have the same property. If in addition we suppose that $`B_{\beta \alpha \gamma }=L_{\alpha \beta \gamma }`$ we found that (10) has the simple form $`L_{\alpha \beta (\gamma ;\delta )}=0`$. We found the conditions when the manifold and its dual have the same Lax tensors. For the two dimensional manifolds we found that the symmetric Lax tensors exist if the scalar curvature is zero. The solution of the Lax equations for the flat space case , the Rindler system and its dual manifold were found.
Finding the Lax tensors on the manifolds which admits Killing-Yano tensors is an interesting problem and it requires further investigation.
## 6 Acknowledgements
One of the authors (D.B.) would like to thank Ashok Das for helpful discussions. He also would like to thank TUBITAK and NATO for financial support and METU for the hospitality during his working stage at the Department of Physics. This work is partially supported by the Scientific and Technical Research Council of Turkey.
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# Effective Charge and Spin Hamiltonian for the Quarter-Filled Ladder Compound 𝛼'-NaV2O5
## I Introduction
Low-dimensional quantum spin systems have received considerable attention from both theoretical as well as experimental point of view due to their unconventional physical properties. $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub>, which was believed to be a low-dimensional inorganic spin-Peierls (SP) compound has recently been under intense investigation. $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> is an insulator and its magnetic susceptibility data fits very well to the one-dimensional Heisenberg chain model yielding an exchange interaction $`J`$=440 and 560 K for temperatures below and above the spin-Peierls transition temperature T<sub>SP</sub> (T<sub>SP</sub>$``$ 34 K ) respectively . For T$``$ T<sub>SP</sub>, an isotropic drop in the susceptibility corresponding to a singlet-triplet gap of $`\mathrm{\Delta }_{SP}`$=85 K has been observed.
Recent X-ray structure data analysis at room temperature disfavours the previously reported non-centrosymmetric structure ($`C_{2v}^7P2_1mn`$ space group) where V<sup>+4</sup> spin-1/2 ions form a one-dimensional Heisenberg chain, running along the crystallographic b-direction, separated by chains of V<sup>+5</sup> spin-zero ions. But the evidence for the centrosymmetric point group ($`D_{2h}^{13}Pmmn`$) leads to only one type of V-site with a formal valence +4.5 in this compound. The V-sites then form a quarter-filled ladder, running along the b-axis with the rungs along the crystallographic a-axis. In the quarter-filled scenario, the electron spins are not localized at V-ions rather distributed over a V-O-V molecule which has found support by NMR as well as Raman measurements .
The nature of the state below T<sub>SP</sub> is presently under intense investigation. Isobe and Ueda originally proposed a usual spin-Peierls scenario but the detection of two inequivalent V-sites in NMR indicates a more complicated scenario and the possibility of charge ordering. Several types of charge ordering, including ’in-line’ and ’zig-zag’ ordering has been proposed, but only the zig-zag type of ordering has been found to be in agreement with neutron scattering and anomalous X-ray scattering .
Recent determinations of the low-temperature crystal structure found the space group Fmm2 and proposed the existence of three inequivalent V-ions below T<sub>SP</sub> . This scenario was investigated by DMRG (density-matrix renormalization group) and a cluster-operator theory and a strong disagreement with neutron scattering data was found. Ohama et al. recentely observed that the aparent contradiction between cyrstallography (three inequivalent V-ions bewlow T<sub>SP</sub>) and NMR (two inequivalent V-ions bewlow T<sub>SP</sub>) could be resolved when one considers possible subgroups of the originally proposed space group Fmm2.
The crystal structure of $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> at T$`>`$T<sub>SP</sub> is orthorhombic (a=$`11.318\text{Å}`$, b=$`3.611\text{Å}`$, c=$`4.797\text{Å}`$) and consists of double chains of edge-sharing distorted tetragonal VO<sub>5</sub> pyramids running along the orthorhombic b-axis. These double chains are linked together via common corners of the pyramids and form layers. These are stacked along c-direction with no-direct V-O-V links. The Na atoms are located in between these layers. For the orbitals of the d-electrons at V-sites, those with d<sub>xy</sub> symmetry are suggested to be the relevant ones above and below T<sub>SP</sub> . Due to the special orbital structure, the hopping amplitudes $`t_a`$ and $`t_b`$ are much larger than the inter-ladder hopping $`t_{ab}`$, $`t_a`$ and $`t_b`$ being the hopping amplitudes along the rung and the ladder direction respectively ($`t_a0.38\text{eV}`$, $`t_b0.17\text{eV}`$, $`t_{ab}0.012\text{eV}`$). Since $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> is an insulator, it has been assumed that the on-site Coulomb repulsion $`U`$ is sufficiently large in comparision to the hopping amplitudes ($`U2.8\text{eV}`$ from DFT calculation ). Moreover, one has to introduce the inter-site Coulomb repulsions, $`V_a`$, $`V_b`$, $`V_{ab}`$ to obtain the required charge ordering. In fact, it has been shown in a Hartree-Fock calculation that the condition $`U>V_a,V_b,V_{ab}>t_a,t_b,t_{ab}`$ must be fulfilled in order to achieve a complete charge ordering. We consider this and other limits in the present paper.
In the present work, we take into account the charge dynamics to obtain an effective low energy hamiltonian for $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub>. One starts from a pure electronic hamiltonian, which includes electron hopping in and between the ladders as well as the on-site and inter-site Coulomb interactions. The on-site Coulomb interaction $`U`$ is taken to be the largest parameter in our calculation. Since $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> is an insulator and we work at quarter-filling, one can project on a subspace of states which contains one electron per rung. Therefore, it is convenient to use an Ising pseudo-spin variable $`\tau ^z=\pm 1/2`$ corresponding to a rung with an electron on the right/left site of the rung. This is in the same spirit of Kugel and Khomskii’s treatment of the orbital degeneracy problem in Jahn-Teller systems . The spin and the pseudo-spin operators can be written as,
$`S^z={\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}\sigma (R_\sigma ^{}R_\sigma +L_\sigma ^{}L_\sigma ),S^+=R_{}^{}R_{}+L_{}^{}L_{},S^{}=R_{}^{}R_{}+L_{}^{}L_{},`$ (1)
$`\tau ^z={\displaystyle \frac{1}{2}}{\displaystyle \underset{\sigma }{}}(R_\sigma ^{}R_\sigma L_\sigma ^{}L_\sigma ),\tau ^+={\displaystyle \underset{\sigma }{}}R_\sigma ^{}L_\sigma ,\tau ^{}={\displaystyle \underset{\sigma }{}}L_\sigma ^{}R_\sigma ,`$ (2)
which for example yields, $`R_{i,}^{}L_{i,}=\tau _i^+(\frac{1}{2}+S_i^z)`$, $`R_{i,}^{}L_{i,}=\tau _i^+S_i^+`$, $`L_{i,}^{}L_{i,}=(\frac{1}{2}\tau _i^z)S_i^{}`$, etc., where $`\tau ^\pm =\tau ^x+i\tau ^y`$, $`S^\pm =S^x+iS^y`$ and $`R_{i,\sigma }^{}(L_{i,\sigma }^{})`$ are the creation operator of an electron with spin $`\sigma `$ on the right (left) site of the i-th rung of the ladder. In (1) and (2) we have suppressed the site-indices.
## II Intra-ladder Exchange
Let us start with an electronic hamiltonian for the quarter-filled ladder (see Fig. 1), which can be written as, $`H=H_0+H_0^{}+H_I`$, with
$$H_0=t_a\underset{i,\sigma }{}(R_{i,\sigma }^{}L_{i,\sigma }+h.c.)+U\underset{i}{}(n_{i,R}n_{i,R,}+n_{i,L,}n_{i,L,})+V_a\underset{i,\sigma ,\sigma ^{}}{}n_{i,R,\sigma }n_{i,L,\sigma ^{}},$$
(3)
$`H_0^{}=V_b^{}{\displaystyle \underset{i,\sigma ,\sigma ^{}}{}}(n_{i,L,\sigma }n_{i+1,L,\sigma ^{}}+n_{i,R,\sigma }n_{i+1,R,\sigma ^{}})+V_b^{\prime \prime }{\displaystyle \underset{i,\sigma ,\sigma ^{}}{}}(n_{i,L,\sigma }n_{i+1,R,\sigma ^{}}+n_{i,R,\sigma }n_{i+1,L,\sigma ^{}})`$ (4)
$`+V_{ab}{\displaystyle \underset{<m,n>,\sigma ,\sigma ^{}}{}}n_{m,R,\sigma }n_{n,L,\sigma ^{}},`$ (5)
$$H_I=t_b\underset{i,\sigma }{}(R_{i,\sigma }^{}R_{i+1,\sigma }+L_{i,\sigma }^{}L_{i+1,\sigma }+h.c.)+t_{ab}\underset{<m,n>,\sigma }{}(R_{m,\sigma }^{}L_{n,\sigma }+L_{n,\sigma }^{}R_{m,\sigma }),$$
(6)
where $`t_a`$, $`U`$ and $`V_a`$ are the hopping integral, on-site and the inter-site Coulomb repulsion in a rung respectively. $`t_b`$, $`V_b^{}`$ and $`V_b^{\prime \prime }`$ are the hopping integral and the Coulomb interaction in between rungs in a ladder, whereas $`t_{ab}`$ and $`V_{ab}`$ are the inter-ladder hopping and the Coulomb interaction respectively. $`n_{i,R,\sigma }(n_{i,L,\sigma })`$ is the electron density operator with spin $`\sigma `$ in the right (left) site of i-th rung and $`<m,n>`$ denotes the pair of rungs $`m`$ and $`n`$ on adjacent ladders.
We estimate the parameters of the inter-site Coulomb repulsion using a screened Coulomb repulsion, $`V=e^2/(ϵd)`$, where $`ϵ`$ is the dielectric constant and $`d`$ the distance between the respective Vanadium atoms. The distance in between two V-ions along the rung and the leg of the ladder (a- and b-directions) are $`3.502\text{Å}`$ and $`3.611\text{Å}`$. The dielectric constant is $`ϵ=11`$ from microwave and far infrared measurements . One obtains, $`V_a=0.3738\text{eV}`$ and $`V_b^{}=0.3625\text{eV}`$. The diagonal V-V distance in the b-direction in the ladder is $`5.030\text{Å}`$, which implies, $`V_b^{\prime \prime }=0.2643\text{eV}`$. It will be clear from the later discussion that the effective inter-rung Coulomb repulsion $`V_b`$ is given by the difference between $`V_b^{}`$ and $`V_b^{\prime \prime }`$, i.e. $`V_b=V_b^{}V_b^{\prime \prime }`$, which comes out to be small ($`V_b=0.1023\text{eV}`$) compared to $`V_a`$. Next, $`V_{ab}=0.4305\text{eV}`$, as the V-V inter-ladder distance is $`3.0401\text{Å}`$. Note that $`V_{ab}`$ is slightly higher than $`V_a`$ and nearly four times higher than $`V_b`$.
In order to develop a perturbation expansion, we start by considering the case of a single two-leg quarter-filled ladder. The one-electron eigenstates of $`H_0`$ for a single rung consist of bonding and anti-bonding wave functions, which we denote as $`a^{}|0>=\frac{1}{\sqrt{2}}(R^{}L^{})|0>`$ and $`s^{}|0>=\frac{1}{\sqrt{2}}(R^{}+L^{})|0>`$, with eigenenergies $`t_a`$ and $`+t_b`$ respectively. Now let us consider the coupling of the rungs along the legs described by the first term in the hamiltonian $`H_I`$. In order to obtain a coupling between the pseudo-spin and the spin variables, we use here the standard canonical transformation method , which is given by,
$$H_{eff}=e^{iS}He^{iS},$$
(7)
where the operator $`S`$ is determined from the condition
$$H_I+i[S,H_0]=0,$$
(8)
which turns out to be
$$\widehat{S}=\underset{n,n^{}}{}\frac{i}{(E_n^{}E_n)}|n><n|H_I|n^{}><n^{}|.$$
(9)
Thus, the effective hamiltonian can be written as,
$$H_{eff}=H_0+H_0^{}\frac{1}{2}\underset{n,n^{},n^{\prime \prime }}{}\left(\frac{1}{E_n^{}E_n}+\frac{1}{E_n^{}E_n^{\prime \prime }}\right)|n><n|H_I|n^{}><n^{}|H_I|n^{\prime \prime }><n^{\prime \prime }|,$$
(10)
where the initial and the final states $`|n^{\prime \prime }>`$ and $`|n>`$ are the two-rung states, i.e., all possible combinations of the bonding and the anti-bonding states between the nearest neighbour rungs. In the present case, there are sixteen possible such states which are the following:
$$s_{i,\sigma }^{}s_{j,\sigma ^{}}^{}|0>,s_{i,\sigma }^{}a_{j,\sigma ^{}}^{}|0>,a_{i,\sigma }^{}s_{j,\sigma ^{}}^{}|0>,a_{i,\sigma }^{}a_{j,\sigma ^{}}^{}|0>,$$
(11)
with $`\sigma ,\sigma ^{}=,`$. The six intermediate states $`|n^{}>`$ which are the two-particle excited states in a rung, have to be antisymmetric under the exchange of both spin and pseudo-spin coordinates, in accordance with the Pauli principle. Thus, we have two sectors for the excited states depending on the total and the z-component of spin as well as the pseudo-spin quantum numbers which are labeled as, $`|S,S^z;\tau ,\tau ^z>`$. Hence, the states involved are, $`|0,0;1,1>=R_{i,}^{}R_{i,}^{}|0>`$, $`|0,0;1,1>=L_{i,}^{}L_{i,}^{}|0>`$, $`|0,0;1,0>=\frac{1}{\sqrt{2}}(R_{i,}^{}L_{i,}^{}R_{i,}^{}L_{i,}^{})|0>`$ and $`|1,1;0,0>=R_{i,}^{}L_{i,}^{}|0>`$, $`|1,1;0,0>=R_{i,}^{}L_{i,}^{}|0>`$, $`|1,0;0,0>=\frac{1}{\sqrt{2}}(R_{i,}^{}L_{i,}^{}+R_{i,}^{}L_{i,}^{})|0>`$. The eigenenergies of the excited states in the large $`U`$ limit are $`V_a`$ for the spin-triplet states $`|1,S^z;0,0>`$ ($`S^z=1,0+1`$). For the spin-singlets, the eigenenergies are $`U^{}`$ for $`\frac{1}{\sqrt{2}}(|0,0;1,1>+|0,0;1,1>)`$ (symmetric), $`U`$ for $`\frac{1}{\sqrt{2}}(|0,0;1,1>|0,0;1,1>)`$ (antisymmetric) and $`V_a^{}`$ for $`|0,0;1,0>`$, with $`U^{}U+\frac{4t_a^2}{UV_a}`$ and $`V_a^{}V_a\frac{4t_a^2}{UV_a}`$. After some lengthy but straightforward algebra and in the case of large but finite $`U`$, the total effective hamiltonian can be written as,
$$H_{eff}=H_0+H_0^{}+H_{eff}^{intra}+H_{eff}^{inter},$$
(12)
where $`H_{eff}^{intra}`$ is the effective intra-ladder hamiltonian which one can express as, $`H_{eff}^{intra}=H_{eff}^{(t)}+H_{eff}^{(s)}`$, with $`H_{eff}^{(t)}`$ and $`H_{eff}^{(s)}`$ being the contribution due to the intermediate spin-triplet and spin-singlet states respectively. $`H_{eff}^{inter}`$ is the effective inter-ladder hamiltonian which can be derived in a similar way and will be discussed in the next section. In terms of the pseudo-spin and spin variables, the unperturbed hamiltonian $`H_0`$ and $`H_0^{}`$ can be expressed as,
$$H_0=2t_a\underset{i}{}\tau _i^x,$$
(13)
$$H_0^{}=2(V_b^{}V_b^{\prime \prime })\underset{i}{}\left(\frac{1}{4}+\tau _i^z\tau _{i+1}^z\right)+NV_b^{\prime \prime }+V_{ab}\underset{<m,n>}{}\left(\frac{1}{4}\tau _m^z\tau _n^z+\frac{\tau _m^z\tau _n^z}{2}\right),$$
(14)
where $`N`$ is the number of rungs. In a similar way, the effective hamiltonian $`H_{eff}^{(t)}`$ can be written as,
$$H_{eff}^{(t)}=\frac{4t_b^2}{V_a}\underset{i}{}\left(\frac{1}{4}\stackrel{}{\tau }_i\stackrel{}{\tau }_{i+1}\right)\left(\frac{3}{4}+\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right).$$
(15)
It is obvious from the above expression that $`H_{eff}^{(t)}`$ is independent of the Coulomb correlation energy $`U`$. This is due to the fact that while deriving this effective hamiltonian we have used the eigenenergies for the excited states which happen to be $`V_a`$ for this case. Since the effective hamiltonian $`H_{eff}^{(s)}`$ is obtained due to the contribution from the same intermediate states $`|n^{}>`$ and different initial and final states $`|n>`$ and $`|n^{\prime \prime }>`$, it can be written as, $`H_{eff}^{(s)}=H_{eff}^{(s1)}+H_{eff}^{(s2)}+H_{eff}^{(s3)}`$ (which are the contributions due to the antisymmetric, symmetric and $`|0,0;1,0>`$ intermediate states), with
$$H_{eff}^{(s1)}=\frac{4t_b^2}{U}\underset{i}{}\left(\frac{1}{4}2\tau _i^x\tau _{i+1}^x+\stackrel{}{\tau }_i\stackrel{}{\tau }_{i+1}\right)\left(\frac{1}{4}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right),$$
(16)
$`H_{eff}^{(s2)}={\displaystyle \frac{2t_b^2}{U^{}2t_a}}{\displaystyle \underset{i}{}}\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\tau _i^x+\tau _{i+1}^x}{2}}2\tau _i^y\tau _{i+1}^y+\stackrel{}{\tau }_i\stackrel{}{\tau }_{i+1}\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right)`$ (17)
$`{\displaystyle \frac{2t_b^2}{U^{}+2t_a}}{\displaystyle \underset{i}{}}\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{\tau _i^x+\tau _{i+1}^x}{2}}2\tau _i^y\tau _{i+1}^y+\stackrel{}{\tau }_i\stackrel{}{\tau }_{i+1}\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right),`$ (18)
$`H_{eff}^{(s3)}={\displaystyle \frac{2t_b^2}{V_a^{}2t_a}}{\displaystyle \underset{i}{}}\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\tau _i^x+\tau _{i+1}^x}{2}}2\tau _i^z\tau _{i+1}^z+\stackrel{}{\tau }_i\stackrel{}{\tau }_{i+1}\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right)`$ (19)
$`{\displaystyle \frac{2t_b^2}{V_a^{}+2t_a}}{\displaystyle \underset{i}{}}\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{\tau _i^x+\tau _{i+1}^x}{2}}2\tau _i^z\tau _{i+1}^z+\stackrel{}{\tau }_i\stackrel{}{\tau }_{i+1}\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right).`$ (20)
It should be noted here that one gets a non-zero contribution to $`H_{eff}^{(t)}`$ and $`H_{eff}^{(s)}`$ even if $`U=\mathrm{}`$ which will be discussed below. Moreover, it is obvious from the expression for $`H_{eff}^{(t)}`$ (see Eq. (15)) that a finite $`V_a`$ is indeed needed to get a meaningful result otherwise the perturbation becomes ill-defined for $`V_a=0`$.
### Limiting Cases and Discussion
Case I: $`t_a=0`$: This limit implies $`U^{}=U`$ and $`V_a^{}=V_a`$ and thus, $`H_{eff}^{(s1)}`$ and $`H_{eff}^{(s2)}`$ can be combined together to yield,
$$H_{eff}^{(s1)}+H_{eff}^{(s2)}=\frac{8t_b^2}{U}\underset{i}{}\left(\frac{1}{4}+\tau _i^z\tau _{i+1}^z\right)\left(\frac{1}{4}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right),$$
(21)
whereas the effective hamiltonian $`H_{eff}^{(s3)}`$ gets reduced to,
$$H_{eff}^{(s3)}=\frac{4t_b^2}{V_a}\underset{i}{}\left(\frac{1}{4}+\stackrel{}{\tau }_i\stackrel{}{\tau }_{i+1}2\tau _i^z\tau _{i+1}^z\right)\left(\frac{1}{4}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right).$$
(22)
Since $`H_{eff}^{(t)}`$ neither depends on $`t_a`$ nor on $`U`$, it doesn’t get affected in the above limiting case and the same is true for the other cases considered below. The effective hamiltonian derived by Thalmeier and Fulde corresponds to Eq. (21).
Case II: $`t_a=0`$, $`U=\mathrm{}`$: In this limit, which also corresponds to the limit $`V_a2t_a`$, the contribution to the effective hamiltonian from $`H_{eff}^{(s1)}`$ and $`H_{eff}^{(s2)}`$ vanish and thus, the total intra-ladder effective hamiltonian becomes the sum of $`H_{eff}^{(t)}`$ and $`H_{eff}^{(s3)}`$ (see Eq. (15) and (22)). This is what is exactly obtained by Mostovoy and Khomskii but with a different interpretation.
Case III: $`U=\mathrm{}`$, $`V_a2t_a`$: In this case, the effective hamiltonian $`H_{eff}^{(s1)}`$ and $`H_{eff}^{(s2)}`$ vanish but $`H_{eff}^{(s3)}`$ reduces to,
$$H_{eff}^{(s3)}=\frac{t_b^2}{t_a}\underset{i}{}\left(\tau _i^x+\tau _{i+1}^x\right)\left(\frac{1}{4}\stackrel{}{S}_i\stackrel{}{S}_{i+1}\right),$$
(23)
so that the total effective intra-ladder hamiltonian becomes the sum of Eq. (15) and (23). Moreover, since $`V_a2t_a`$, the major contribution will be from Eq. (15).
Case IV: Disordered Phase: In the disordered phase, where the electrons are in the bonding states (see Fig. 2 (a)), we can get an estimate of the effective exchange coupling in the effective hamiltonian by taking the averages over its charge (pseudo-spin) part. One can write down the effective exchange hamiltonian (disregarding the constant factors) as, $`H_{eff}^{exch}=J_{eff}^{exch}_i\stackrel{}{S}_i\stackrel{}{S}_{i+1}`$. In the present case, we have, $`<\tau _i^x>=1/2`$ and $`<\tau _i^y>=0=<\tau _i^z>`$ whereas $`<\tau _i^x\tau _{i+1}^x>=1/4`$ and $`<\tau _i^y\tau _{i+1}^y>=0=<\tau _i^z\tau _{i+1}^z>`$. Thus, the effective exchange coupling due to the hamiltonian $`H_{eff}^{(t)}`$ and $`H_{eff}^{(s1)}`$ vanish, but that of $`H_{eff}^{(s2)}`$ and $`H_{eff}^{(s3)}`$ become, $`\frac{2t_b^2}{U^{}+2t_a}`$ and $`\frac{2t_b^2}{V_a^{}+2t_a}`$ which yields $`J_{eff}^{exch}=2t_b^2(\frac{1}{U^{}+2t_a}+\frac{1}{V_a^{}+2t_a})`$. It is clear that the exchange coupling here is antiferromagnetic ($`>0`$). Using the parameters mentioned in the present work, $`J_{eff}^{exch}`$ is estimated to be 0.08 eV. The expression for $`J_{eff}^{exch}`$ is exactly the same (for the case $`V_a=0`$) as obtained by Horsch and Mack .
Case V: Complete Charge Ordered Phase: We can get an estimate of the effective exchange couplings in the completely charge ordered (zig-zag) phase (see Fig. 2 (b, c)), following the same procedure as that of the disordered phase. Here, one has, $`<\tau _i^x>=0=<\tau _i^y>`$ and $`<\tau _i^z>=1/2`$, $`<\tau _{i+1}^z>=1/2`$ whereas $`<\tau _i^x\tau _{i+1}^x>=0=<\tau _i^y\tau _{i+1}^y>`$ and $`<\tau _i^z\tau _{i+1}^z>=1/4`$. Hence, the effective exchange coupling due to $`H_{eff}^{(s1)}`$ and $`H_{eff}^{(s2)}`$ vanish but that of $`H_{eff}^{(t)}`$ and $`H_{eff}^{(s3)}`$ survive, which ultimately leads to $`J_{eff}^{exch}=2t_b^2[\frac{1}{V_a}\frac{1}{2(V_a^{}2t_a)}\frac{1}{2(V_a^{}+2t_a)}]`$. However, using the parameter values, it is calculated to be -0.17 eV. In the case where $`U=\mathrm{}`$ and for $`0<(V_a/2t_a)<1`$, $`J_{eff}^{exch}`$ becomes ferromagnetic ($`<0`$) whereas for $`(V_a/2t_a)>1`$, it is antiferromagnetic. The variation of $`J_{eff}^{exch}`$ with respect to the parameter $`(V_a/2t_a)`$ is shown in Fig. 4. It is clear from the figure that there exists a minimum in the ferromagnetic region for $`(V_a/2t_a)=1/\sqrt{3}`$, where $`J_{eff}^{exch,min}=2.26(t_b^2/t_a)=0.2\text{eV}`$, which is quite large.
## III Inter-ladder Exchange
Next, let us consider the hopping between the two nearest neighbour ladders, which is described by the second term of the hamiltonian $`H_I`$ (see Eq. (6)), i.e.,
$$H_I^{inter}=t_{ab}\underset{<m,n>,\sigma }{}(R_{m,\sigma }^{}L_{n,\sigma }+L_{n,\sigma }^{}R_{m,\sigma }).$$
(24)
The effective inter-ladder coupling between the charge and the spin degrees of freedom is derived in the same way as has been done for the single ladder case. Since the two-particle excited states in this case are exactly the same as what has been done earlier, the effective inter-ladder hamiltonian can be written as sum of two parts, i.e., $`H_{eff}^{inter}=H_{eff}^{inter(t)}+H_{eff}^{inter(s)}`$, where the superscript ‘$`t`$’ and ‘$`s`$’ have the same meaning discussed in the previous section. The effective hamiltonian $`H_{eff}^{inter(t)}`$ is derived to be,
$`H_{eff}^{inter(t)}={\displaystyle \frac{t_{ab}^2}{2(V_a2t_a)}}{\displaystyle \underset{<m,n>}{}}\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\tau _m^x+\tau _n^x}{2}}2\tau _m^y\tau _n^y+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{3}{4}}+\stackrel{}{S}_m\stackrel{}{S}_n\right)`$ (25)
$`{\displaystyle \frac{t_{ab}^2}{2(V_a+2t_a)}}{\displaystyle \underset{<m,n>}{}}\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{\tau _m^x+\tau _n^x}{2}}2\tau _m^y\tau _n^y+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{3}{4}}+\stackrel{}{S}_m\stackrel{}{S}_n\right)`$ (26)
$`{\displaystyle \frac{t_{ab}^2}{V_a}}{\displaystyle \underset{<m,n>}{}}\left({\displaystyle \frac{1}{4}}2\tau _m^x\tau _n^x+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{3}{4}}+\stackrel{}{S}_m\stackrel{}{S}_n\right),`$ (27)
whereas $`H_{eff}^{inter(s)}`$ can be written as $`H_{eff}^{inter(s)}=H_{eff}^{inter(s1)}+H_{eff}^{inter(s2)}+H_{eff}^{inter(s3)}`$, with
$`H_{eff}^{inter(s1)}={\displaystyle \frac{t_{ab}^2}{2(U2t_a)}}{\displaystyle \underset{<m,n>}{}}`$ (28)
$`\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\tau _m^x+\tau _n^x}{2}}+{\displaystyle \frac{\tau _m^z\tau _n^z}{2}}2\tau _m^z\tau _n^z\tau _m^x\tau _n^z+\tau _m^z\tau _n^x+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_m\stackrel{}{S}_n\right)`$ (29)
$`{\displaystyle \frac{t_{ab}^2}{2(U+2t_a)}}{\displaystyle \underset{<m,n>}{}}`$ (30)
$`\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{\tau _m^x+\tau _n^x}{2}}+{\displaystyle \frac{\tau _m^z\tau _n^z}{2}}2\tau _m^z\tau _n^z+\tau _m^x\tau _n^z\tau _m^z\tau _n^x+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_m\stackrel{}{S}_n\right)`$ (31)
$`{\displaystyle \frac{t_{ab}^2}{U}}{\displaystyle \underset{<m,n>}{}}\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\tau _m^z\tau _n^z}{2}}\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_m\stackrel{}{S}_n\right).`$ (32)
The expression for $`H_{eff}^{inter(s2)}`$ is exactly the same as that of $`H_{eff}^{inter(s1)}`$ except that $`U`$ here is replaced by $`U^{}`$. On the other hand, $`H_{eff}^{inter(s3)}`$ is obtained as,
$`H_{eff}^{inter(s3)}={\displaystyle \frac{t_{ab}^2}{2(V_a^{}2t_a)}}{\displaystyle \underset{<m,n>}{}}\left({\displaystyle \frac{1}{4}}+{\displaystyle \frac{\tau _m^x+\tau _n^x}{2}}2\tau _m^y\tau _n^y+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_m\stackrel{}{S}_n\right)`$ (33)
$`{\displaystyle \frac{t_{ab}^2}{2(V_a^{}+2t_a)}}{\displaystyle \underset{<m,n>}{}}\left({\displaystyle \frac{1}{4}}{\displaystyle \frac{\tau _m^x+\tau _n^x}{2}}2\tau _m^y\tau _n^y+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_m\stackrel{}{S}_n\right)`$ (34)
$`{\displaystyle \frac{t_{ab}^2}{V_a^{}}}{\displaystyle \underset{<m,n>}{}}\left({\displaystyle \frac{1}{4}}2\tau _m^x\tau _n^x+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right)\left({\displaystyle \frac{1}{4}}\stackrel{}{S}_m\stackrel{}{S}_n\right).`$ (35)
Since the derivation of the effective inter-ladder hamiltonian proceeds in the same way as that of the intra-ladder case, the limiting cases will follow the same way as has been done before. Moreover, here also, one needs a finite $`V_a`$ (but $`V_a2t_a`$) in deriving the effective hamiltonian, otherwise the perturbation becomes ill-defined for $`V_a=0`$. Furthermore, one gets a non-zero contribution to the effective inter-ladder hamiltonian even if $`U=\mathrm{}`$ .
### Limiting Cases and Discussion
Case I: $`t_a=0`$: In this limit $`U^{}=U`$ and $`V_a^{}=V_a`$ follow naturally. Thus, $`H_{eff}^{inter(t)}`$ reduces to,
$$H_{eff}^{inter(t)}=\frac{2t_{ab}^2}{V_a}\underset{<m,n>}{}\left(\frac{1}{4}+\tau _m^z\tau _n^z\right)\left(\frac{3}{4}+\stackrel{}{S}_m\stackrel{}{S}_n\right).$$
(36)
Similarly, $`H_{eff}^{inter(s1)}`$, $`H_{eff}^{inter(s2)}`$ and $`H_{eff}^{inter(s3)}`$ become,
$$H_{eff}^{inter(s1)}+H_{eff}^{inter(s2)}=\frac{4t_{ab}^2}{U}\underset{<m,n>}{}\left(\frac{1}{4}+\frac{\tau _m^z\tau _n^z}{2}\tau _m^z\tau _n^z\right)\left(\frac{1}{4}\stackrel{}{S}_m\stackrel{}{S}_n\right),$$
(37)
$$H_{eff}^{inter(s3)}=\frac{2t_{ab}^2}{V_a}\underset{<m,n>}{}\left(\frac{1}{4}+\tau _m^z\tau _n^z\right)\left(1/4\stackrel{}{S}_m\stackrel{}{S}_n\right).$$
(38)
The expression Eq. (37) is the same one as has been obtained by Thalmeier and Fulde .
Case II: $`t_a=0`$, $`U=\mathrm{}`$: In this case, the contributions from $`H_{eff}^{inter(s1)}`$and $`H_{eff}^{inter(s2)}`$ vanish. The contributions from $`H_{eff}^{inter(t)}`$ and $`H_{eff}^{inter(s3)}`$ can be combined together to yield,
$$H_{eff}^{inter}=H_{eff}^{inter(t)}+H_{eff}^{inter(s)}=\frac{2t_{ab}^2}{V_a}\underset{<m,n>}{}\left(\frac{1}{4}+\tau _m^z\tau _n^z\right),$$
(39)
which becomes independent of the spin degrees of freedom.
Case III: $`U=\mathrm{}`$, $`V_a2t_a`$: In this case, the contributions from $`H_{eff}^{inter(s1)}`$ and $`H_{eff}^{inter(s2)}`$ vanish. Combining $`H_{eff}^{inter(t)}`$ and $`H_{eff}^{inter(s3)}`$, the spin-dependence drops out and one obtains,
$$H_{eff}^{inter}=H_{eff}^{inter(t)}+H_{eff}^{inter(s3)}=\frac{t_{ab}^2}{4t_a}\underset{<m,n>}{}\left(\tau _m^x+\tau _n^x\right)\frac{t_{ab}^2}{V_a}\underset{<m,n>}{}\left(\frac{1}{4}2\tau _m^x\tau _n^x+\stackrel{}{\tau }_m\stackrel{}{\tau }_n\right).$$
(40)
Case IV: Disordered Phase: Following the procedure already mentioned for the intra-ladder case, we can get an estimate of the effective exchange coupling between the nearest neighbour ladders (see Fig. 3 (a)) by taking the averages over the charge part in the effective hamiltonian. Here also, one can write down the effective hamiltonian (disregarding the constant factors) as, $`H_{eff}^{inter,exch}=(J_{eff}^{inter,exch(t)}+J_{eff}^{inter,exch(s)})_{<m,n>}\stackrel{}{S}_m\stackrel{}{S}_n=J_{eff}^{inter,exch}_{<m,n>}\stackrel{}{S}_m\stackrel{}{S}_n`$. In the present case, one has, $`<\tau _m^x>=<\tau _n^x>=1/2`$ and $`<\tau _m^y>=<\tau _n^y>=0=<\tau _m^z>=<\tau _n^z>`$ whereas $`<\tau _m^x\tau _n^x>=1/4`$ and $`<\tau _m^y\tau _n^y>=0=<\tau _m^z\tau _n^z>`$. Thus, the effective exchange coupling which is due to $`H_{eff}^{inter(t)}`$, $`H_{eff}^{inter(s1)}`$, $`H_{eff}^{inter(s2)}`$ and $`H_{eff}^{inter(s3)}`$ is given as, $`J_{eff}^{inter,exch}=\frac{t_{ab}^2}{2}[\frac{1}{(V_a+2t_a)}\frac{1}{(U+2t_a)}\frac{1}{(U^{}+2t_a)}\frac{1}{(V_a^{}+2t_a)}]`$. The exchange coupling here is antiferromagnetic and is estimated to be $`(J_{eff}^{inter,exch}/t_{ab}^2)=0.39`$. For $`V_a=0`$, it gives rise to the same expression as obtained by Horsch and Mack .
Case V: Complete Charge Ordered Phase: Here, we have four different completely charge ordered (zig-zag) phase depending on the state through which we compute the averages over the charge part of the effective hamiltonian. In all these cases, one has, $`<\tau _m^x>=<\tau _n^x>=0=<\tau _m^y>=<\tau _n^y>`$ and $`<\tau _m^x\tau _n^x>=0=<\tau _m^y\tau _n^y>`$. In addition to this, one has,
$`<\tau _m^z>=1/2`$, $`<\tau _n^z>=1/2`$ and $`<\tau _m^z\tau _n^z>=1/4`$, where the averages are due to the state $`R_{m,\alpha }^{}L_{n,\beta }^{}|0>`$ (Fig. 3 (b)). The effective exchange coupling for $`H_{eff}^{inter(t)}`$ and $`H_{eff}^{inter(s3)}`$ vanish but that of $`H_{eff}^{inter(s1)}`$ and $`H_{eff}^{inter(s2)}`$ are finite which gives rise to, $`J_{eff}^{inter,exch}=t_{ab}^2[\frac{1}{2(U2t_a)}+\frac{1}{2(U+2t_a)}+\frac{1}{U}+\frac{1}{2(U^{}2t_a)}+\frac{1}{2(U^{}+2t_a)}+\frac{1}{U^{}}]`$. The exchange coupling here is antiferromagnetic and is estimated to be $`(J_{eff}^{inter,exch}/t_{ab}^2)=1.42`$. For $`t_a=0`$, it gives rise to usual super-exchange, i.e. $`J_{eff}^{inter,exch}=4t_{ab}^2/U`$.
$`<\tau _m^z>=1/2`$, $`<\tau _n^z>=1/2`$, and $`<\tau _m^z\tau _n^z>=1/4`$, where the averages are taken over the state $`L_{m,\alpha }^{}R_{n,\beta }^{}|0>`$ (Fig. 3 (c)). In this case, $`J_{eff}^{inter,exch(t)}=0=J_{eff}^{inter,exch(s)}`$ which ultimately corresponds to the case $`J_{eff}^{inter,exch}=0`$.
$`<\tau _m^z>=<\tau _n^z>=1/2`$ and $`<\tau _m^z\tau _n^z>=1/4`$, where the averages here are due to the state $`R_{m,\alpha }^{}R_{n,\beta }^{}|0>`$ (Fig. 3 (d)). The effective exchange couplings for $`H_{eff}^{inter(s1)}`$ and $`H_{eff}^{inter(s2)}`$ vanish. Thus, $`J_{eff}^{inter,exch}`$ is obtained from the contribution due to $`H_{eff}^{inter(t)}`$ and $`H_{eff}^{inter(s3)}`$ which is, $`J_{eff}^{inter,exch}=\frac{t_{ab}^2}{2}[\frac{1}{2(V_a2t_a)}+\frac{1}{2(V_a+2t_a)}+\frac{1}{V_a}\frac{1}{2(V_a^{}2t_a)}\frac{1}{2(V_a^{}+2t_a)}\frac{1}{V_a^{}}]`$. The exchange coupling here is antiferromagnetic and is estimated to be $`(J_{eff}^{inter,exch}/t_{ab}^2)=2.8`$. It vanishes for $`U=\mathrm{}`$ as well as for $`t_a=0`$.
$`<\tau _m^z>=<\tau _n^z>=1/2`$ and $`<\tau _m^z\tau _n^z>=1/4`$, where the averages are taken over the state $`L_{m,\alpha }^{}L_{n,\beta }^{}|0>`$ (Fig. 3 (e)). Here, the effective exchange coupling turns out to be the same as that of (iii).
Case VI: A Phase with one Ladder Completely Charge Ordered and the Nearest Neighbour Disordered: In this case, we have two possibilities, again depending on the state through which one computes the averages over the charge part in the effective hamiltonian. In both the cases, one has, $`<\tau _m^x>=0=<\tau _m^y>`$, $`<\tau _n^y>=0=<\tau _n^z>`$ and $`<\tau _m^x\tau _n^x>=0=<\tau _m^y\tau _n^y>=<\tau _m^z\tau _n^z>`$. Besides, one has,
$`<\tau _m^z>=1/2`$, $`<\tau _n^x>=1/2`$ and $`<\tau _m^z\tau _n^x>=1/4`$, where the averages are taken over the state $`R_{m,\alpha }^{}a_{n,\beta }^{}|0>`$ (Fig. 3 (f)). All the effective exchange couplings in this case turn out to be non-zero and hence $`J_{eff}^{inter,exch}`$ is obtained as, $`J_{eff}^{inter,exch}=\frac{t_{ab}^2}{2}[\frac{1}{2(V_a+2t_a)}+\frac{1}{2V_a}\frac{1}{(U+2t_a)}\frac{1}{U}\frac{1}{(U^{}+2t_a)}\frac{1}{U^{}}\frac{1}{2(V_a^{}+2t_a)}\frac{1}{2V_a^{}}]`$.
$`<\tau _m^z>=1/2`$, $`<\tau _n^x>=1/2`$ and $`<\tau _m^z\tau _n^x>=1/4`$ where the averages are due to the state $`L_{m,\alpha }^{}a_{n,\beta }^{}|0>`$ (Fig. 3 (g)). The effective exchange coupling which is due to $`H_{eff}^{inter(t)}`$ and $`H_{eff}^{inter(s3)}`$ (the contributions due to $`H_{eff}^{inter(s1)}`$ and $`H_{eff}^{inter(s2)}`$ vanish) is given by, $`J_{eff}^{inter,exch}=\frac{t_{ab}^2}{2}[\frac{1}{2(V_a+2t_a)}+\frac{1}{2V_a}\frac{1}{2(V_a^{}+2t_a)}\frac{1}{2V_a^{}}]`$.
In both of the above mentioned cases, the exchange coupling becomes antiferromagnetic and is estimated to be $`(J_{eff}^{inter,exch}/t_{ab}^2)=2.26`$ for case (i) and 1.65 for (ii). However, it vanishes in both cases for $`U=\mathrm{}`$ irrespective of whether $`V_a2t_a`$ or $`V_a2t_a`$.
## IV Discussion and Conclusion
We have derived the effective spin-charge hamiltonian for $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> for both intra-ladder and inter-ladder exchange. We find a rich structure as a function of possible realization of the microscopic parameters. We have found several, in part unexpected, results.
The effective magnetic exchange along the ladder decreases with increasing charge ordering (of zig-zag type). For complete charge ordering, the magnetic exchange becomes ferromagnetic and quite large in magnitude.
The effective inter-ladder magnetic exchange in between two given rungs of a charge-ordered and a charge-disordered ladder changes (only) by a factor of two when the charge-density-wave is shifted by a lattice constant along $`a`$ (compare Fig. 3 (f) and (g)).
There are novel terms of type $`\tau _m^z\tau _n^x`$ in the effective charge-charge inter-ladder interaction. These terms, which are however rather small in magnitude, could in principle stabilize a mixed charge-order configurations like the one illustrated in Fig. 3 (f) and (g).
As a consequence of (i) and (ii), the proposed frustrated spin-cluster model by de Boer et al. seems to be unlikely, since (a) no indications of ferromagnetic couplings have been found experimentally and (b) the coupling in between the one-dimensional spin-cluster chains of Ref. should be, as a consequence of (ii), rather strong. This conclusion is consistent with a recent study of the frustrated spin-cluster model by DMRG and a cluster-operator theory .
Let us note that the change in sign of the effective intra-ladder magnetic exhange shown in Fig. 4 cannot be described accurately by a perturbation expansion in $`t_b`$. Near to the singularity at $`V_a=2t_a`$ the perturbation expansion breaks down and the effective intra-ladder spin-hamiltonian becomes long-ranged.
Ultimately, the reason for the ferromagnetic intra-ladder coupling for the zig-zag ordering found in the above calculation lies in the fact that the charge-ordered state is not the ground-state of $`H_0`$. In fact, the charge-ordered state would be stabilized by $`H_0^{}`$ in the case of a large effective inter-rung Coulomb-coupling $`V_b^{}V_b^{\prime \prime }`$ (see Eq. (14)). For large values of $`V_b^{}V_b^{\prime \prime }2t_a`$, the perturbation expansion would yield (e.g. in mean-field approximation for $`H_0^{}`$) an antiferromagnetic intra-ladder spin-spin coupling. We did not show this calculation here, since our estimated value for the inter-rung Coulomb repulsion $`V_b^{}V_b^{\prime \prime }0.1\text{eV}`$ is about one-order of magnitude too small in order to do the job. We therefore believe that this small value of $`V_b^{}V_b^{\prime \prime }`$ indicates (a) the importance of elastic effects for the stabilization of the observed phase transition at $`T_c=34\text{K}`$ and (b) that the degree of charge ordering is far from complete. This is consistent with the proposal of about 20% charge ordering .
###### Acknowledgements.
The authors would like to thank R. Valentí, J. V. Alvarez, F. Capraro and K. Pozgajcic for several discussions and critical reading of the manuscript.
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# Does Anti-Parallel Spin Always Contain More Information?
## Abstract
We show that the Bloch vectors lying on any great circle is the largest set $`S_L`$ for which the parallel states $`|\stackrel{}{n},\stackrel{}{n}`$ can always be exactly transformed into the anti-parallel states $`|\stackrel{}{n},\stackrel{}{n}`$. Thus more information about $`\stackrel{}{n}`$ is *not* extractable from $`|\stackrel{}{n},\stackrel{}{n}`$ than from $`|\stackrel{}{n},\stackrel{}{n}`$ by any measuring strategy, for $`\stackrel{}{n}S_L`$. Surprisingly, the largest set of Bloch vectors for which the corresponding qubits can be flipped is again $`S_L`$. We then show that probabilistic exact parallel to anti-parallel transformation is not possible if the corresponding anti-parallels span the whole Hilbert space of the two qubits. These considerations allow us to generalise a conjecture of Gisin and Popescu (Phys. Rev. Lett. 83 432 (1999)).
<sup>a</sup>Physics and Applied Mathematics Unit, Indian Statistical Institute, 203 B. T. Road, Calcutta -700035, India
<sup>b</sup>Dept. of Physics, Bose Institute, 93/1 A.P.C. Road, Calcutta - 700009, India
Recently, Gisin and Popescu revealed that for qubits, there exists a measuring strategy (see , ) on the anti-parallel state $`|\stackrel{}{n},\stackrel{}{n}`$ that extracts more information about an arbitrary Bloch vector $`\stackrel{}{n}`$ than that can be extracted from the parallel state $`|\stackrel{}{n},\stackrel{}{n}`$ .
In this paper, we ask whether there exists any proper subset $`S`$ of the unit ball $`B_3=\{\stackrel{}{n}:\stackrel{}{n}R^3,\left|\stackrel{}{n}\right|=1\}`$ for which an anti-parallel state selected at random from the set $`A_S=\{|\stackrel{}{n},\stackrel{}{n}:\stackrel{}{n}S\}`$ carry more information about the Bloch vector $`\stackrel{}{n}`$ than a parallel state from the set $`P_S=\{|\stackrel{}{n},\stackrel{}{n}:\stackrel{}{n}S\}`$. We provide a partial answer to this query. Specifically, we find the largest set $`S_L`$ for which there exists a unitary operator $`U=U(S_L)`$ such that
$$U|\stackrel{}{n},\stackrel{}{n}|M=e^{i\theta (\stackrel{}{n})}|\stackrel{}{n},\stackrel{}{n}|N(\stackrel{}{n})$$
(1)
for all $`\stackrel{}{n}S_L`$, where $`|M`$ and $`|N(\stackrel{}{n})`$ are the states of a possible ancilla. $`|N(\stackrel{}{n})`$ must be independent of $`\stackrel{}{n}`$ except possibly in a phase to satisfy the unitarity of $`U`$ . Consequently an anti-parallel state chosen at random from $`A_{S_L}`$ cannot contain more information about $`\stackrel{}{n}`$ than in the corresponding parallel state from $`P_{S_L}`$.
A related query is how far we can go by just flipping the second state in order to transform $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$. Here one may think that the largest set of Bloch vectors would in this case be a very small subset of $`S_L`$. Surprisingly, as we would show here, this conjecture is not true: the largest set is again $`S_L`$.
We shall show here that universal exact machines for the transformations $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$ and $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$ do not exist. Therefore, in the same vein as one considered deterministic inexact and probabilistic exact cloning when faced with the no-cloning theorem , one can consider deterministic inexact and probabilistic exact parallel to anti-parallel (and antiparallel to parallel) machines. Gisin has considered the deterministic inexact case. In this paper we consider the probabilistic exact machines for parallel to anti-parallel transformations. The motivation behind such consideration is to provide basis for a conjecture on the type of sets of Bloch vectors $`\stackrel{}{n}`$ for which the anti-parallel state $`|\stackrel{}{n},\stackrel{}{n}`$ contains more information about $`\stackrel{}{n}`$ than the corresponding parallel state $`|\stackrel{}{n},\stackrel{}{n}`$. We also provide ground for a conjecture made by Gisin and Popescu in .
Let us first demonstrate the non-existence of a machine that transforms $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$ universally. Suppose that there exists a unitary operator $`U^{}`$ such that
$$U^{}|\stackrel{}{n},\stackrel{}{n}|M=e^{i\vartheta (\stackrel{}{n})}|\stackrel{}{n},\stackrel{}{n}|N$$
(2)
for all $`\stackrel{}{n}`$, where $`|N`$ is independent of $`\stackrel{}{n}`$. First, let us display the action of $`U^{}`$ on the (orthogonal) basis $`|00`$, $`|11`$ and $`\frac{1}{\sqrt{2}}(|01+|10)`$ of the linear space spanned by the parallel states:
$$U^{}|00|M=e^{ia}|01|N$$
$$U^{}|11|M=e^{ib}|10|N$$
(3)
$$U^{}\frac{1}{\sqrt{2}}(|01+|10)|M=e^{ic}(c_1|00+c_2|11)|N$$
For an arbitrary parallel state $`|\stackrel{}{n},\stackrel{}{n}(|\stackrel{}{n}=e^{i\alpha _n}(\mathrm{cos}\frac{\theta _n}{2}|0+e^{i\varphi _n}\mathrm{sin}\frac{\theta _n}{2}|1),\mathrm{\hspace{0.25em}0}\theta _n\pi ,\mathrm{\hspace{0.25em}0}\varphi _n<2\pi )`$, checking for unitarity and linearity, it is easy to show that condition (2) will be satisfied only for those $`\stackrel{}{n}`$’s for which $`\varphi _n=\frac{ab+\pi }{2}+l\pi (l=0,\pm 1,\pm 2,\mathrm{}.)`$.
Thus we have shown that there exists a machine which (unitarily) transforms $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$ for all $`\stackrel{}{n}`$ lying on any great circle of the Bloch sphere. But is this the largest set of Bloch vectors $`\stackrel{}{n}`$ for which the transformation $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$ is possible using a single unitary operator? One would be reluctant to believe it as we have constrained $`U^{}`$ to transform at least two orthogonal states $`|00`$ and $`|11`$ to their corresponding anti-parallels. But we shall show that a great circle is the largest set of Bloch vectors for which the corresponding parallels transform to their anti-parallels by a single unitary operator.
To prove this we first show that for any three parallel states $`|\stackrel{}{n_1},\stackrel{}{n_1}`$, $`|\stackrel{}{n_2},\stackrel{}{n_2}`$ and $`|\stackrel{}{n_3},\stackrel{}{n_3}`$, the Bloch vectors $`\stackrel{}{n_1}`$, $`\stackrel{}{n_2}`$ and $`\stackrel{}{n_3}`$ must lie on the same great circle if the parallel states are to be transformed to their anti-parallels by the same unitary operator. But before that note that any three (distinct) parallel states are linearly independent. The reason is that all non-trivial linear combinations of any two parallel states are entangled and hence no parallel state, which are product states, lie in their linear span. Consider then any three parallel states $`|\stackrel{}{n_i},\stackrel{}{n_i}(i=1,2,3)`$ and suppose that there exists a unitary operator $`U^{\prime \prime }`$ such that
$$U^{\prime \prime }|\stackrel{}{n_i},\stackrel{}{n_i}|M=e^{i\vartheta _i}|\stackrel{}{n_i},\stackrel{}{n_i}|N(i=1,2,3)$$
(4)
The $`\stackrel{}{n_i}`$’s are constrained by the fact that $`U^{\prime \prime }`$ is unitary:
$$\stackrel{}{n_i}|\stackrel{}{n_j}=e^{i(\vartheta _j\vartheta _i)}\stackrel{}{n_i}|\stackrel{}{n_j}(i,j=1,2,3)$$
(5)
where we have assumed that the $`|\stackrel{}{n_i}`$’s are mutually non-orthogonal. Taking
$$|\stackrel{}{n_i}=e^{i\alpha _i}(\mathrm{cos}\frac{\theta _i}{2}|\stackrel{}{n_1}+e^{i\varphi _i}\mathrm{sin}\frac{\theta _i}{2}|\stackrel{}{n_1})(i=2,3)$$
(6)
and using (5), we have $`2\alpha _i=\frac{1}{2}(\vartheta _i\vartheta _1)+(2k+1)\pi `$ (with $`k`$ an integer and $`i=2,3`$) and
$$\varphi _2\varphi _3=n\pi (n=0,\pm 1,\pm 2,\mathrm{}.)$$
(7)
which means that $`\stackrel{}{n_1}`$, $`\stackrel{}{n_2}`$ and $`\stackrel{}{n_3}`$ lie on the same great circle.
In the special case when any two of $`|\stackrel{}{n_1}`$, $`|\stackrel{}{n_2}`$ and $`|\stackrel{}{n_3}`$ are orthogonal, the corresponding Bloch vectors always lie on a great circle of the Bloch sphere and hence the corresponding parallel states can be transformed to their anti-parallels by the unitary operator $`U^{}`$ of (2). Therefore in any case, three parallel states can be transformed to their anti-parallels by the same unitary operator if and only if the corresponding Bloch vectors lie on a great circle of the Bloch sphere. The ‘if’ part follows from the existence of $`U^{}`$ of (2).
An immediate consequence is that there cannot exist a measuring strategy on an anti-parallel state chosen at random from $`A_{S_L}`$ that estimates $`\stackrel{}{n}`$ better than any optimal strategy on a parallel state chosen at random from $`P_{S_L}`$. Being a little imprecise, more information about $`\stackrel{}{n}`$ is not extractable from $`A_{S_L}`$ as compared to $`P_{S_L}`$. Gisin and Popescu had conjectured in that more information about $`\stackrel{}{n}`$ *is* extractable from $`A_T`$ as compared to $`P_T`$ where $`T`$ is the set of the four vertices
$$(0,0,1),(\frac{\sqrt{8}}{3},0,\frac{1}{3}),(\frac{\sqrt{2}}{3},\sqrt{\frac{2}{3}},\frac{1}{3}),(\frac{\sqrt{2}}{3},\sqrt{\frac{2}{3}},\frac{1}{3})$$
(8)
on the Bloch sphere, of a tetrahedron. One can see that the elements of $`T`$ do not lie on any great circle whereby it is not possible to transform $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$ for all $`\stackrel{}{n}T`$ by a single unitary operator. The conjecture thus survives this attempt of defeat. Within a few paragraphs we provide more basis for this conjecture.
Before proceeding further, we parenthetically note that the largest set of $`\stackrel{}{n}`$’s for which $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$ is possible, is again $`S_L`$. The corresponding unitary operator is just the inverse of the unitary operator $`U^{}`$ of (2). Such a transformation would not be possible for a larger set of Bloch vectors as that would enable one to transform a larger set of parallels to their anti-parallels (by the inverse of that (unitary) transformation) which we have seen to be untrue.
We have seen that $`S_L`$ is the largest set of Bloch vectors for which $`|\stackrel{}{n},\stackrel{}{n}`$ goes over to $`|\stackrel{}{n},\stackrel{}{n}`$ by the same machine. We now want to see how far we can go by just flipping the second state. Consider a unitary operator $`U_f`$ such that
$$U_f|0|M=e^{i\vartheta _0}|1|N$$
$$U_f|\stackrel{}{n}|M=e^{i\vartheta _n}|\stackrel{}{n}|N$$
(9)
where
$$|\stackrel{}{n}=e^{i\alpha }(\mathrm{cos}\frac{\theta }{2}|0+e^{i\varphi }\mathrm{sin}\frac{\theta }{2}|1)$$
(10)
Here again $`|N`$ is independent of the input qubits.
As $`U_f`$ is unitary, we must have
$$0|\stackrel{}{n}=e^{i(\vartheta _n\vartheta _0)}1|\stackrel{}{n}$$
(11)
where we have assumed that $`0|\stackrel{}{n}0`$. This gives us
$$2\alpha =\vartheta _n\vartheta _0+(2r+1)\pi $$
(12)
with $`r`$ an integer. This does not give any constraint on $`\stackrel{}{n}`$ which essentially depends on $`\theta `$ and $`\varphi `$. If $`0|\stackrel{}{n}=0`$, then the orthogonal states $`|0`$ and $`|\stackrel{}{n}`$ can obviously be unitarily transformed to the orthogonal states $`|1`$ and $`|\stackrel{}{n}`$. Thus in any case, there exists a unitary operator to flip a qubit chosen at random from two arbitrary but fixed qubits. Next let us assume that
$$U_f|\stackrel{}{m}|M=e^{i\vartheta _m}|\stackrel{}{m}|N$$
(13)
where
$$|\stackrel{}{m}=e^{i\alpha ^{}}(\mathrm{cos}\frac{\theta ^{}}{2}|0+e^{i\varphi ^{}}\mathrm{sin}\frac{\theta ^{}}{2}|1)$$
(14)
Unitarity restricts $`\stackrel{}{m}`$ to lie on the same great circle as of $`\stackrel{}{n}`$ and the Bloch vector for $`|0`$ on the Bloch sphere. As $`|\stackrel{}{m}`$ can be written as a linear combination of $`|0`$ and $`|\stackrel{}{n}`$, we must also check for linearity and constrain $`\stackrel{}{m}`$ accordingly. But surprisingly, as one can easily see, linearity does not give any new constraint. Thus $`S_L`$ is yet again the largest set of Bloch vectors for which the corresponding *qubits* can be flipped. Consequently, the unitary operators for which $`|\stackrel{}{n},\stackrel{}{n}`$ goes over to $`|\stackrel{}{n},\stackrel{}{n}`$ (and also $`|\stackrel{}{n},\stackrel{}{n}`$ to $`|\stackrel{}{n},\stackrel{}{n}`$) for the largest set of $`\stackrel{}{n}`$’s is surprisingly possible for the form $`IU_2`$ where $`I`$ is the identity operator on the Hilbert space of the first qubit and $`U_2`$ is a unitary operator on the Hilbert space of the second qubit. Thus we can actually transform $`|\stackrel{}{m},\stackrel{}{n}`$ to $`|\stackrel{}{m},\stackrel{}{n}`$ (and also $`|\stackrel{}{m},\stackrel{}{n}`$ to $`|\stackrel{}{m},\stackrel{}{n}`$) for all $`\stackrel{}{n}`$ lying on a great circle and for *any* $`\stackrel{}{m}`$.
As we have already transpired, we shall next consider the case of probabilistic exact parallel to anti-parallel transformations. Two arbitrarily chosen parallel states can always be unitarily transformed to the corresponding anti-parallel states. This is because two arbitrarily chosen Bloch vectors always lie on a great circle. The same is not true for *three* arbitrarily chosen parallel states. Hence it is relevant to consider whether a probabilistic (exact) transformation is possible in this case. Precisely, we want to find out whether there exists a unitary operator for which
$$|\stackrel{}{n_i},\stackrel{}{n_i}|M\sqrt{\gamma _i}|\stackrel{}{n_i},\stackrel{}{n_i}|M_i+\sqrt{1\gamma _i}|\mathrm{\Phi }_i$$
(15)
where $`|M`$, $`|M_i`$ are ancilla states, $`|\mathrm{\Phi }_i`$’s belong to the combined Hilbert space of the two qubits and ancilla, $`0<\gamma _i<1`$ and $`I_4|M_iM_i||\mathrm{\Phi }_j=0`$ for $`i,j=1,2,3`$. A matrix theoretical argument following Duan and Guo shows that such an operator always exists and the corresponding optimal success probabilities $`\gamma _i`$ can also be found out.
Let us now try to find whether there exists a unitary operator for which (15) holds for $`i=1,2,\mathrm{},k`$ with $`k4`$. The corresponding parallel states must be linearly dependent as all the parallel states of two qubits span only a three dimensional subspace of the Hilbert space of two qubits. We assume that the set $`S^{}`$ of $`\stackrel{}{n_i}`$’s for $`i=1,2,\mathrm{},k`$ is not a subset of any great circle as in that case one can transform $`|\stackrel{}{n_i},\stackrel{}{n_i}`$ to $`|\stackrel{}{n_i},\stackrel{}{n_i}`$ with unit probabilities. Now in a unitary-reduction process (15), the dimension of the linear span of the domain must always be greater than or equal to the dimension of the linear span of the range. Consequently there is no probabilistic transformation taking the states of $`P_S^{}`$ to $`A_S^{}`$, if the dimension of the linear span of $`A_S^{}`$ is four.
Let $`S`$ be any set of four or more Bloch vectors $`\stackrel{}{n}`$ which do not lie on any great circle and corresponding to which $`A_S`$ spans the whole Hilbert space of the two qubits. We conjecture that there exists a measuring strategy on an anti-parallel state chosen at random from $`A_S`$ that estimates $`\stackrel{}{n}`$ better than the optimal measuring strategy on a parallel state chosen at random from $`P_S`$. The result of Gisin and Popescu , for the universal case, gives support to this conjecture. Again if $`S`$ is any set of four Bloch vectors for which $`A_S`$ is a linearly independent set, there exists an optimal state discriminating strategy for $`A_S`$ which give some kind of estimation of the states in $`S`$, while no such estimation is possible for $`P_S`$ (as $`P_S`$ is a linearly dependent set) which provides further support to our conjecture. The conjecture of Gisin and Popescu , refered to earlier, is a special case of this conjecture. This is because there are four Bloch vectors in $`T`$ and they do not lie on any great circle and the linear span of $`A_T`$ is of dimension four.
For completeness, we note that linearly independent anti-parallel states can always be probabilistically transformed to their parallels with non-zero success probabilities (also non-unit, if the Bloch vectors do not lie on a great circle).
To summarize, we have shown that the Bloch vectors lying on a great circle of the Bloch sphere is the largest set $`S_L`$ for which the corresponding parallels are exactly transformed to their anti-parallels. As a consequence, with the apriori knowledge that an anti-parallel state belongs to $`A_{S_L}`$, one cannot estimate the Bloch vector better than the optimal case with the apriori knowledge that a parallel state belongs to $`P_{S_L}`$. We then found to our surprise that the largest set of Bloch vectors for which the corresponding qubits can be flipped by a single unitary operator is again $`S_L`$. We then showed that probabilistic exact parallel to anti-parallel transformation is not possible if the corresponding anti-parallels span the whole four-dimensional Hilbert space of the two qubits. This allowed us to provide ground for a conjecture made by Gisin and Popescu and also to make a more general conjecture. To justify this general conjecture, we have to find out state estimation strategies for finite sets of parallel and anti-parallel states.
The authors acknowledge useful discussions with Guruprasad Kar, Somshubhro Bandyopadhyay and Debasis Sarkar. U.S. thanks Dipankar Home for encouragement and acknowledges partial support by the Council of Scientific and Industrial Research, Government of India, New Delhi.
References:
N. Gisin and S. Popescu, Phys. Rev. Lett. 83 432 (1999).
S. Massar and S. Popescu, Phys. Rev. Lett. 74, 1259 (1995).
S. Massar, quant-ph/0004035.
The normalized pure product state $`|\psi |\varphi `$ of two qubits, where the (unit) Bloch vectors of $`|\psi `$ and $`|\varphi `$ are $`\stackrel{}{n}`$ and $`\stackrel{}{m}`$ respectively, is denoted here by $`|\stackrel{}{n},\stackrel{}{m}`$.
For any two non-orthogonal states $`|\stackrel{}{m}`$, $`|\stackrel{}{n}`$, from (1) we get $`\left|N(\stackrel{}{m})|N(\stackrel{}{n})\right|=1`$. Then using Cauchy’s inequality we see that $`|N(\stackrel{}{m})`$ and $`|N(\stackrel{}{n})`$ are equal upto a phase.
V. Buzek and M. Hillery, Phys. Rev. A, 54, 1844 (1996); D. Bruss et. al., Phys Rev. A, 57, 2368 (1998).
L. M. Duan and G. C. Guo, Phys. Lett. A, 243, 261 (1998).
L. M. Duan and G. C. Guo, Phys. Rev. Lett., 80, 4999 (1998).
W. K. Wootters and W. H. Zurek, Nature, 299, 802 (1982).
N. Gisin, private communication.
A. Peres, Quantum Theory: Concepts and Methods, Kluwer 1993; A. Chefles and S.M. Barnett, J. Mod. Opt., 45, 1295 (1998).
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# The PSI np data and their effect on the charged 𝜋NN coupling constant
## 1 Introduction
Neutron-proton elastic scattering at backward angles in the medium energy regime has been the subject of several experiments \[1–6\]. The common feature of these investigations is a steep rise of the differential cross section towards the back scattering angle of 180. The slope of the sharp backward peak suggests a connection to the one pion exchange (OPE) amplitude. Several suggestions have been made for a theoretical description of the experiments \[7–16\]. All of these proposals are able to describe the backward spike, but most of them are purely phenomenological and some of them fail to describe other observables.
Precise backward scattering data offer the opportunity to evaluate the pion–nucleon coupling constant $`f_{\pi \mathrm{NN}}^2`$. The medium energy region is particularly suited for this purpose, because the pion pole is not very far from the physical region, so that one could expect a reliable extrapolation to it. However, the values determined so far by this method are not all in accordance with each other and with the value obtained from pion-nucleon scattering .
The results presented in this paper have been obtained from four separate experiments of different angular ranges, labeled I–IV \[4, 19–21\]. Together they span the interval from about 80 to 180 in the centre-of-mass system. The neutron energies range from 200 MeV to 580 MeV in steps of 20 MeV.
## 2 Experiment
The experimental set-ups and techniques of experiments I-IV have been quite similar but with differences in detail of the accelerator performance, the beam arrangement and the detection equipment. Here we outline the common features only and refer to a forthcoming paper for details.
The experiments have been performed at the Paul Scherrer Institute (PSI), the former Swiss Institute for Nuclear Research (SIN). The proton beam of the ring cyclotron of 589 MeV energy consists of bunches with a width of less than 1 ns at a rate of 50.63 MHz for experiments I and III and 16.88 MHz for II and IV, corresponding to bunch spacings of 19.75 ns and 59.25 ns, respectively. The beam current during the data taking was 60–100 $`\mu `$A. The neutrons were produced on a thick target of beryllium (I–III) or carbon (IV). Neutrons escaped through a collimator hole in the beam dump at an angle of 60 mrad with respect to the incident protons. They were shaped by two additional collimators to a beam of about 2 x 2 cm<sup>2</sup> at a distance of 61 m from the production target. Charged particle contaminations were eliminated by cleaning magnets behind the collimators. A lead filter was inserted in order to reduce the photon component of the beam which mainly originates from the decay of neutral pions in the production target.
The continuous neutron energy spectrum consists of the 40 MeV wide quasielastic peak at about 540 MeV and a broad distribution at lower energies, which can be ascribed to inelastic processes . The detailed shape of the spectrum depends on the target material and its thickness.
After a flight path of 61 meters, kept at rough vacuum of about 100 Pa, the neutron beam hits a liquid hydrogen target. For the analysis of the scattering products a magnet spectrometer was installed on a turn table. The spectrometer was equipped with drift chambers and two scintillation counters, a thin one (1 mm thick) in front of the magnet, and a thicker one (1 cm thick) behind it. They allow to measure both, the time-of-flight of the particle detected by the spectrometer, and that of the incoming neutron with respect to the rf-signal of the accelerator. The angular acceptance is almost 20. The average momentum resolution is about 3 % FWHM.
## 3 Data taking and analysis
The data have been taken at different run periods, each extending over several weeks. The contribution of the target surroundings and spectrometer materials was measured with an empty target. It was subtracted after normalization to the neutron intensity of the full target measurement. Corrections have been applied for the event rate dependent dead-time losses.
The long periods of data taking required special attention to the long term variations of the whole system, particularly of the neutron time-of-flight measurement. The stability of the electronics has been checked regularly and close control of the drift chamber gas flow and high voltage supply has been maintained.
The neutron intensity was monitored in three different ways. The integrated primary proton beam intensity was provided from the accelerator control centre. This signal, though convenient, was not sufficiently reliable because of variations of the focusing of the primary beam on the production target. The second monitor consisted of a three stage scintillation counter telescope, which recorded charged particles emerging from a thin polyethylene target placed in the beam at about 44 m from the neutron target. A third monitor was installed behind the magnet spectrometer. It recorded elastic np scattering events from a polyethylene block in coincidence.
These monitors were intended to measure the relative neutron intensity. Above 280 MeV the absolute normalization was performed by the simultaneously recorded $`\mathrm{np}\mathrm{d}\pi ^\mathrm{o}`$ reaction as reference cross section. This is discussed in sect. 5.
### 3.1 Time-of-flight calibration
As mentioned above, the energy spectrum of the incident neutrons is continuous. Therefore, the neutron time–of–flight measurement is the basis of the incident energy determination, and the control of its stability is crucial. It requires careful calibration of the zero point and the conversion gain of the time–to–digital converters. For the control of the time zero point we used the high energy photons in the beam. The lead filter in the neutron beam was removed for these calibration runs, and the hydrogen target was replaced by a lead slab for a higher conversion rate. With the spectrometer field inverted and reduced appropriately, the converted electrons were detected. From the width of the peak the time resolution has been evaluated to be 0.8 ns FWHM including the bunch width of the primary beam and the contribution from electronics. The corresponding neutron energy resolutions vary from 1.2 MeV at 200 MeV to 7.6 MeV at 580 MeV.
These calibrations have been performed in regular intervals. Deviations from the overall mean have been corrected for each run.
### 3.2 Data reduction and event selection
The first step in the off-line analysis was a reduction of the data by setting cuts in order to select good events. These cuts included a unique track in the drift chambers, starting in the target volume and emitted within the full acceptance of the spectrometer in vertical and horizontal direction. Next, the mass of the particle in the spectrometer was determined from the measured time-of-flight through the spectrometer and the measured momentum. A mass selection has been applied by rejecting masses $`m<0.5m_\mathrm{p}`$ and $`m>1.5m_\mathrm{p}`$, where $`m_\mathrm{p}`$ is the proton mass. The upper limit was increased to $`m>2.5m_\mathrm{p}`$ for experiments I and II, in order to keep also deuteron events, which were used for the normalization of the np data (see sect. 5). Finally, low momentum particles have been excluded by cuts at $`p=500`$ MeV/c for protons and $`p=620`$ MeV/c for deuterons.
### 3.3 Neutron energy determination
The emission angle of the recoil proton is determined by the set of drift chambers in front of the magnet, and its momentum by the deflection in the magnetic field. The proton energy is obtained after correction of the energy loss in traversing the material of target and spectrometer. The incident neutron energy is determined from the measured neutron time-of-flight $`t_{\mathrm{}}`$. A complication arises from the fact, that the bunch interval given by the rf-signal (19.75 ns or 59.25 ns) is much shorter than the time-of-flight from the production target to the hydrogen target (more than 200 ns). This introduces ambiguities on the proper multiple of the bunch spacing time which has to be added to the measured value. However, by comparison with the time-of-flight $`t_{\mathrm{calc},\mathrm{}}`$, calculated from the measured recoil proton momentum vector, the correct number of bunch spacings can be determined.
## 4 Experimental results
### 4.1 np differential cross sections
The data were binned in energy intervals of 20 MeV, and angular bins of 0.5 degrees in the centre-of-mass (CM) system. Each bin has been corrected for absorption of the recoil protons in the target and the spectrometer material. This energy dependent correction was below 1 % in all cases. Another correction was necessary due to a small polarization component of the neutron beam in vertical direction of $`.05<|P_\mathrm{n}|<.08`$. This comes about since the neutrons are produced at an angle of 60 mrad (cf. sect. 2). The energy and angle dependent effect of the analyzing power of the np elastic scattering has been taken into account. The correction to the differential cross section was below 2.5 %. The errors of these corrections contribute to the errors of the differential cross sections with less than 0.3 %.
The four sets of data have been combined in an appropriate way. For the experiments I and II a common angular region exists which has been used to evaluate the multiplying factor for experiment I. The uncertainty introduced by this procedure for the different energies has been determined as 0.6 % from the spread of the individual factors from the average. In a similar way the data set from experiment III has been linked to the combined data of set I and set II. Here the uncertainty was 0.7 %. This procedure could not be followed for the normalization of experiment IV because no overlap of data points exists. We therefore fitted for each neutron energy the data points of experiment III for CM-angles smaller than 145 together with the same number of data points of set IV by Legendre polynomials of 4th order with a free scaling factor. Each data point has been given the same statistical weight of 2 % in the fit procedure in order to assign the same weight for the two data sets. The average value of $`\chi ^2`$ for the relative normalization was 1.41, and the mean error of the scaling factor is 2 %.
The statistical errors of the differential cross sections range from 1 % to 2 %. The spread of the data with respect to smooth fits with polynomials of an appropriate order is more like 2–2.5 %. This reflects systematic uncertainties due to
* short term changes in the phase of the rf-signal
* imperfections of the drift chambers and other detectors
* remaining effects of gain shifts of detectors and electronics.
Taken together they are in the order of 1–2 %, depending on the set-up. Systematic errors of 1.4 %, 0.6 % and 1.7 % have been added in quadrature to the data sets II, III and IV, respectively. A systematic error of 0.5 % had been added earlier to data set I .
The relative differential cross sections for the four separate experiments are displayed together in Fig. 1. The angular distributions are normalized to 1.0 at the largest angle of data set II. The numerical values are tabulated in ref. .
### 4.2 Discussion
A common feature of the angular dependence at all energies is the sharp rise towards the backward direction. It is followed by a less steep decrease passing through a wide minimum, which is shifted to smaller angles with decreasing energy. In the transition region of the two different slopes there is an indication of a bump. A similar shape is indicated in the LAMPF data but not as clear.
The energy dependence of the cross section at 180 is shown in Fig. 2. It is roughly constant in our energy range, and is well described by the phase shift predictions. With decreasing angle a deviation of the data from the phase shift predictions develops within the full energy range. This can be clearly seen by a comparison of the cross section ratios $`\sigma (\theta )/\sigma (180^{})`$ as shown in Fig. 2 for $`\theta `$ = 135 and 90. The deviations occur for both phase shift solutions of Arndt et al. and Bystricky et al. . It reflects the fact that for the extreme backward scattering angles several precise and consistent measurements exist which pin down the phase shift solutions, whereas the angular region around $`\theta =90^{}`$ has been covered only scarcely at singular energies so far. The cross section ratios, as displayed in Fig. 2, are decreasing with increasing energy, except for the ratio $`\sigma (90^{})/\sigma (180^{})`$ which is almost energy independent above 250 MeV.
## 5 Absolute normalization
For the absolute normalization of the cross section the incident neutron intensity has to be known as a function of energy. In our case the neutron intensity was obtained for experiment II from a comparison with the simultaneously measured reaction
$$\mathrm{np}\mathrm{d}\pi ^0,$$
(1)
where the deuterons have been recorded by the spectrometer like the recoiling protons of the elastic scattering.
Isospin independence is used to relate the cross section of the process (1) to the cross section of reaction
$$\mathrm{pp}\mathrm{d}\pi ^+,$$
(2)
for which precise data exist. However, isospin symmetry, which relates the two reactions by $`\sigma _{\mathrm{d}\pi ^\mathrm{o}}=\frac{1}{2}\sigma _{\mathrm{d}\pi ^+}`$ is not exact since the masses in the initial and in the final states are different for (1) and (2). Also, the Coulomb interaction for process (2) must be taken into account. This is discussed in detail in section 5.2.
### 5.1 Measurement of $`\frac{𝐝𝝈}{𝐝𝛀}\mathbf{(}\mathrm{𝐧𝐩}\mathbf{}𝐝𝝅^𝐨\mathbf{)}`$
The angular acceptance of the magnet spectrometer is sufficiently large to cover the full angular range of the deuterons, $`\theta _\mathrm{d}^{\mathrm{max}}12^o`$, so that a complete angular distribution of reaction (1) is obtained. The energy of the deuterons emitted in (1) varies strongly with the emission angle. For an incident neutron energy of 400 MeV, as an example, the deuteron energy varies from 130 MeV to 260 MeV. In this energy region the deuteron cross section on nuclei is strongly energy dependent and distorts the angular distribution. Besides the energy loss the absorption by the spectrometer material has to be corrected for. While the energy loss correction is straightforward, the correction for absorption loss is more involved because of the missing knowledge of deuteron cross sections on nuclei.
For the deuteron total cross section on a nucleus with mass number $`A`$ we use
$$\sigma _{\mathrm{d}A}=0.97(\sigma _{\mathrm{n}A}+\sigma _{\mathrm{p}A}).$$
(3)
The material around the target and the spectrometer consists mainly of hydrogen and the self-conjugate nuclei C, N, O and Ar, for which we assume $`\sigma _{\mathrm{n}A}=\sigma _{\mathrm{p}A}`$. The deuteron cross section eq. (3) reduces then to
$$\sigma _{\mathrm{d}A}=1.94\sigma _{\mathrm{n}A}.$$
(4)
Neutron total cross sections on nuclei have been measured in a previous work . For N and Ar, which have not been measured directly, the cross section has been evaluated from a parameterization of the energy and mass number dependence given in . For the deuteron cross section on hydrogen $`\sigma _{\mathrm{dp}}`$ at the energy $`T`$ we used $`\sigma _{\mathrm{pd}}`$ data or the $`\sigma _{\mathrm{np}}`$ data of at $`T/2`$.
The absorption correction is then obtained by summing up the contributions of the different elements according to their thickness in the spectrometer. The resulting correction varies from 5.2 % at $`T{}_{\mathrm{d}}{}^{}=120`$ MeV to 2.5 % at 400 MeV. As in the case of elastic np scattering (cf. sect. 4.1) the effect of a small polarization of the incoming neutrons has been corrected for. This angular and energy dependent correction for the deuteron production is below 0.6 % for $`T_\mathrm{n}<490`$ MeV. At our highest energy it reaches 2.6 %. The errors of both corrections taken together contribute to the errors of the differential production cross sections with less than 0.5 %.
### 5.2 Mass- and Coulomb corrections
Niskanen and Vestama studied the symmetry violating effects for the reactions (1) and (2) in the framework of a coupled-channel method with different pion production mechanisms. Because the final momentum is most relevant in threshold reactions with a large negative Q-value the comparisons are made at the same final pion momenta $`p_\pi `$ in the CM system, commonly expressed by the dimensionless variable $`\eta =p_\pi /m_{\pi ^+}`$. For both reactions, the pion production cross section can be factorized as $`\mathrm{d}\sigma /\mathrm{d}\mathrm{\Omega }=PR`$ with the phase space factor $`P`$ and $`R`$ the sum of the squared matrix elements of the pion production operator. With this ansatz the relative change in the cross sections was determined in good approximation as
$$\frac{\delta \sigma }{\sigma _{\mathrm{av}}}\frac{\delta P}{P_{\mathrm{av}}}+\frac{\delta R}{R_{\mathrm{av}}}$$
(5)
with
$$\begin{array}{ccccc}\hfill \delta \sigma =& 2\sigma _{\mathrm{d}\pi ^o}\sigma _{\mathrm{d}\pi ^+}\hfill & & \hfill \sigma _{\mathrm{av}}=& \frac{1}{2}(2\sigma _{\mathrm{d}\pi ^o}+\sigma _{\mathrm{d}\pi ^+})\hfill \\ \hfill \delta R=& R_{\mathrm{d}\pi ^o}R_{\mathrm{d}\pi ^+}\hfill & & \hfill R_{\mathrm{av}}=& \frac{1}{2}(R_{\mathrm{d}\pi ^o}+R_{\mathrm{d}\pi ^+})\hfill \\ \hfill \delta P=& P_{\mathrm{d}\pi ^o}P_{\mathrm{d}\pi ^+}\hfill & & \hfill P_{\mathrm{av}}=& \frac{1}{2}(P_{\mathrm{d}\pi ^o}+P_{\mathrm{d}\pi ^+}).\hfill \end{array}$$
Niskanen and Vestama calculated the effect due to the mass differences and incorporated the Coulomb correction in the final state only. At higher proton energies, the Coulomb correction for the pp initial state is comparable with the one in the d$`\pi ^+`$ final state. Therefore we use the results of Ref. for the change in the matrix element without Coulomb correction and apply Coulomb penetration factors for a point proton with an extended proton and for a point pion with an extended deuteron. The weighted sum of the Coulomb penetration factors for the different orbital angular momentum states with weighting factors from is used as Coulomb correction $`C_{ec}^2`$ for the squared matrix elements $`R_{0\mathrm{d}\pi ^+}`$ without Coulomb correction,
$$R_{\mathrm{d}\pi ^+}=C_{ec}^2R_{0\mathrm{d}\pi ^+}.$$
(6)
Table I shows some relevant information used to transform the cross section $`\sigma _{\mathrm{d}\pi ^+}`$ of (2) to $`\sigma _{\mathrm{d}\pi ^\mathrm{o}}`$ of the process (1) by the relation
$$\sigma _{\mathrm{d}\pi ^0}=\frac{1}{2}\sigma _{\mathrm{d}\pi ^+}\frac{2+\delta \sigma /\sigma _{\mathrm{av}}}{2\delta \sigma /\sigma _{\mathrm{av}}}.$$
(7)
The uncertainty for $`\delta R/R_{\mathrm{av}}`$ is according to Ref. about 0.01. Together with the Coulomb corrections which we applied to the pp initial state, this increases to 0.015. This error contribution is included in the errors given in col. 9 of Table I.
### 5.3 Parameterization of $`𝝈\mathbf{(}\mathrm{𝐩𝐩}\mathbf{}𝐝𝝅^\mathbf{+}`$)
Cross sections for reaction (2) as well as for the inverse reaction $`\pi ^+\mathrm{d}\mathrm{pp}`$ have been measured with sufficient accuracy in the relevant energy range from threshold up to 640 MeV. The data base \[31–35\] consists of 61 integrated cross sections for the reaction (2) between 288 and 641 MeV incident proton energy, and of 74 cross sections for the inverse process between 1.8 and 174 MeV incident pion energy. The latter ones have been converted to (2) via detailed balance. Statistical and systematic errors of each data point have been added in quadrature (Gaussian). The data set has been fitted by a parameterization similar to Bystricky et al. as a function of $`\eta `$. The fit error of $`\sigma _{\mathrm{d}\pi ^+}`$ ranges from 2.1 % at $`\eta =0.396`$ to 0.7 % at $`\eta =1.582`$.
### 5.4 Absolute normalization above 300 MeV
The absolute normalization factor for the elastic np scattering data is given with respect to the highest measured angle $`\theta _{\mathrm{max}}=179.2^{}`$ of data set II
$$\frac{\mathrm{d}\sigma }{\mathrm{d}\mathrm{\Omega }}(\theta _{\mathrm{max}})_{\mathrm{np}}=\frac{\stackrel{~}{N}_{\mathrm{np}}(\theta _{\mathrm{max}})}{\stackrel{~}{\mathrm{N}}_{\mathrm{d}\pi ^0}d\mathrm{\Omega }_{\mathrm{d}\pi ^\mathrm{o}}}\sigma _{\mathrm{d}\pi ^\mathrm{o}},$$
(8)
where $`\stackrel{~}{N}_{\mathrm{np}}`$ and $`\stackrel{~}{N}_{d\pi ^0}`$ are the solid angle corrected event numbers for the two processes in corresponding energy bins. The results are given in col. 10 of Table I. Close to pion threshold and below, this method can not be applied.
## 6 The $`𝝅\mathrm{𝐍𝐍}`$ coupling constant $`𝒇_{𝝅\mathrm{𝐍𝐍}}^\mathrm{𝟐}`$
### 6.1 The Chew method
The value of the $`\pi NN`$ coupling constant $`f_{\pi \mathrm{NN}}^2`$ has focussed new interest recently in connection with a discussion on a conceivable breaking of charge independence \[37–41\]. The charged coupling constant, $`f_\mathrm{c}^2`$, can be obtained by an extrapolation of the measured backward elastic differential cross section to the pion pole. The standard procedure in the past has been the Chew method , which extrapolates the Chew funcion $`y(x)`$ defined by
$$y(x)=(xx_{\pi pole})^2\frac{\mathrm{d}\sigma }{\mathrm{d}\mathrm{\Omega }}=\underset{j=0}{\overset{j_{\mathrm{max}}}{}}a_jP_j(x)$$
(9)
to the pion pole in the unphysical region. Here $`x=\mathrm{cos}\theta `$ with $`\theta `$ the CM-angle of the scattered neutron and
$`x_{\pi pole}`$ $`=`$ $`\mathrm{cos}\theta _{\pi pole}`$ (10)
$`=`$ $`(m_\pi ^2m_\mathrm{n}^2m_\mathrm{p}^2+2E_\mathrm{n}E_\mathrm{p})/2p_\mathrm{n}^2`$
with $`m_\pi `$, $`m_\mathrm{n}`$ and $`m_\mathrm{p}`$ the masses of the charged pion, the neutron and the proton, $`E_\mathrm{n}`$ and $`E_\mathrm{p}`$ the CM-energies of the neutron and the proton and $`p_\mathrm{n}`$ the CM-momentum of the neutron. The $`P_j`$ are polynomials generated with special recurrence relations, so that they are orthogonal in the range where data points exist . At the pion pole the Chew function gives
$`y(x_{\pi pole})`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{j_{\mathrm{max}}}{}}}a_jP_j(x_{\pi pole})`$ (11)
$`=`$ $`(\mathrm{}\mathrm{c})^2f_\mathrm{c}^4(m_\mathrm{n}+m_\mathrm{p})^4/(4sp_\mathrm{n}^4)`$
where $`s`$ is the total energy squared.
### 6.2 Test with pseudo data
We have tested this extraction method with pseudo data between 240 MeV and 540 MeV neutron laboratory kinetic energy in 100 MeV steps. The data were generated from the regularized OPE model of Gibbs and Loiseau in the angular range from $`80^{}`$ to $`180^{}`$ in $`1^{}`$ steps. The input coupling constant was fixed at $`f_\mathrm{c}^2=0.076`$. Uncertainties of the pseudo data were generated with a
Gaussian random error distribution. If these uncertainties are $`0.1\%`$, the model coupling constant can be reproduced with an error of $`0.25`$ %. For pseudo data with uncertainties of 2 %, corresponding to the present experiment, the error for $`f_\mathrm{c}^2`$ is 2.4 %. Moreover, the extracted coupling constant is systematically smaller than its input value. This is because the extracted $`f_\mathrm{c}^2`$ rises with the maximum polynomial order of the fit as shown in Fig. 3 (open circles). Applying an $`F`$ test to the fit gives a lower maximum order of the polynomials of eq. 9 for data with higher uncertainties. At $`T_\mathrm{n}=440`$ MeV, for instance, and for uncertainties $`0.1\%`$ the $`F`$ test results in the maximum order $`j_{\mathrm{max}}=10`$ and $`f_\mathrm{c}^2=0.0761\pm 0.2\%`$. For uncertainties of $`2\%`$ the $`F`$ test results in $`j_{\mathrm{max}}=7`$ and $`f_\mathrm{c}^2=0.0746\pm 2.4\%`$.
### 6.3 Chew method with conformal mapping
The convergence properties of the polynomial expansion is improved by conformal mapping of the variable $`x`$ before extrapolating to the pole . We used this method similar to Dumbrajs et al. . The position of the neutral pion pole and the onset of the two pion ($`\pi ^o\pi ^+`$) exchange cut are mapped to (+1,0) and (-1,0) in the complex plane, respectively. The other exchange cuts lie on the unit circle. In addition, the data points are symmetrized around zero.
This mapping method was tested as before with pseudo data. Again, the model coupling constant is reproduced for uncertainties $`0.1`$ %. The convergence is faster than without conformal mapping, and the extrapolation error is of the same order or slightly larger. For the example at $`T_\mathrm{n}=440`$ MeV given above, the $`F`$ test results a value of $`f_{\mathrm{c},\mathrm{map}}^2=0.0759\pm 0.3\%`$ with a reduced polynomial order of $`j_{\mathrm{max}}=8`$ for uncertainties of $`0.1\%`$. With uncertainties of 2 %, the $`F`$ test results in the lower value $`j_{\mathrm{max}}=5`$ and $`f_{\mathrm{c},\mathrm{map}}^2=0.0786\pm 1.6\%`$. Thus, the conformal mapping method results in a systematic upward shift for the extracted coupling constant, as shown by the triangles in Fig. 3. Again, this comes about because the extracted $`f_{\mathrm{c},\mathrm{map}}^2`$ depends on the maximum polynomial order, which is lower for less precise pseudo data.
Similar results were obtained with pseudo data from phase shift predictions of Arndt et al. . Again, the input value $`f_\mathrm{c}^2=0.076`$ could be reproduced only with high precision ($`0.1`$ %) pseudo data by both methods. For uncertainties of 2 % the result was systematically too low without, and too high with conformal mapping. Thus, both methods have systematic offsets of opposite sign. In principle, they can be taken into account by a correction factor $`k_{\mathrm{corr}}=f_{\mathrm{input}}^2/f_{\mathrm{extracted}}^2`$.
### 6.4 Determination of $`𝒇_𝒄^\mathrm{𝟐}`$ from our data
Besides the systematic uncertainties caused by the errors of the data, there is also a dependence on the angular range used in the fit. With decreasing range also the extracted coupling constant
is decreasing. This can be seen from the results on $`f_\mathrm{c}^2`$ and $`f_{\mathrm{c},\mathrm{map}}^2`$ obtained with the real data at $`T_\mathrm{n}=440`$ MeV, as shown in Fig. 4. For the fits at this energy all data points were given a statistical error of 2.8 %. With this choice the reduced $`\chi ^2`$ of the fit is of the order of 1 for the whole angular range. The error of $`f_\mathrm{c}^2`$ is about 15 % for an angular range $`\mathrm{\Delta }\theta =15^{}`$ and drops to about 3 % for $`\mathrm{\Delta }\theta =100^{}`$. According to Fig. 4 an angular range $`\mathrm{\Delta }\theta 50^{}`$ is needed. The variation of the extracted $`f^2`$ values is minimal for $`75^{}\mathrm{\Delta }\theta 95^{}`$. For all energies it was found, that the angular range should extend to the region where the differential cross section reaches its minimum. In order to have no different weights from the four data sets, we have chosen for each laboratory kinetic energy a common relative error per data point, which was adjusted to get a reduced $`\chi _\mathrm{R}^2`$ of the order of 1 for the whole angular range. The maximum polynomial order $`j_{\mathrm{max}}`$ was determined either when a minimum of $`\chi _\mathrm{R}^2`$ was reached or when $`F4`$, which corresponds to a significance of at least 95 % for the term with $`j_{\mathrm{max}}`$. The fit error of $`f^2`$ is then below 5 %. The results of the fits at energies where an absolute normalization exists are collected in Table II.
The correction factors $`k_{\mathrm{corr}}`$, which were needed to correct for the systematic offset introduced by the Chew extrapolation were determined for both methods by an extended study with OPE pseudo data. For each energy we have generated 1000 sets of pseudo data, with randomly distributed errors $`s`$ as given by the experimental data (see Table II). For the Chew method without conformal mapping no unique $`j_{\mathrm{max}}`$ was found for a given energy. Therefore we have used the $`k_{\mathrm{corr}}`$ obtained for that value of $`j_{\mathrm{max}}`$ which was obtained for the real data. This complication does not arise with conformal mapping, where the fit criteria (F-test) lead to the same value of $`j_{\mathrm{max}}`$ at all energies. Moreover, the energy variation of $`k_{\mathrm{corr}}`$ is only $`1`$ % as compared with almost 10 % without conformal mapping. Therefore we prefer to extract an energy independent value of $`f^2`$ using the conformal mapping method. The maximum polynomial order $`j_{\mathrm{max}}`$, the resulting value of $`f^2`$, the fit error $`\mathrm{\Delta }f^2`$ and the correction factor $`k_{\mathrm{corr}}`$ is given in Table II for both methods at each energy.
A contribution to the error which has not been taken into account so far is due to the uncertainty of $`2\%`$ in the relative normalization of data set IV with respect to the other sets. It was investigated systematically and a variation of about 1.5 % of the extracted $`f^2`$ was observed. This contribution is not contained in $`\mathrm{\Delta }f^2`$ of Table II. It was, however, added in quadrature to each $`\mathrm{\Delta }f^2`$ before the energy averaged weighted mean $`f_{\pi \mathrm{np}}^2=f^2k_{\mathrm{corr}}`$ has been calculated.
For the prefered method with conformal mapping we obtain as our final result the weighted mean of the charged coupling constant in the energy range $`300\text{ MeV}T_\mathrm{n}560\text{ MeV}`$:
$$f_{\pi \mathrm{np}}^2=0.0760\pm 0.0008.$$
The given error is the propagated error calculated from the individual $`\mathrm{\Delta }f^2`$. It is slightly larger than the standard deviation of the mean.
It has to be kept in mind, that both, the relative error of the data points as well as the absolute normalization enter and limit the determination of $`f^2`$. In our case both error types are of the same order. With more precise data, i.e. smaller errors of the relative cross sections, the fitting procedure would be more stable and the systematic effects be smaller. Changes in the absolute normalization factor and its error, on the other hand, propagate to the $`f^2`$ values only with a factor $`1/2`$ since the Chew function at the pole depends on $`f^4`$.
All method tests and the calculation of $`k_{\mathrm{corr}}`$ have been performed with OPE model data. Therefore, a model contribution should eventually be added to the error margin.
Our value of 0.0760 for the charged coupling constant is remarkably lower than the value obtained by Höhler et al. from pion-nucleon scattering. It is, however, in good agreement with more recent determinations of $`f^2`$ from both, pion-nucleon partial wave analysis and nucleon-nucleon scattering .
## Acknowledgments
This work has been supported by the German Bundesministerium für Bildung und Forschung.
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# 1 Introduction
## 1 Introduction
The recent growing interest in the relativistic nuclear physics, especially at intermediate energies, has stimulated the development of the theory of relativistic kinetic equations (RKE) as well. In particular, the main problems of dynamical RKE derivation were discussed within the framework of different approaches with well grounded traditions in the non-relativistic case, namely the method of the real-time (contour) Green’s functions -, the BBGKY method (method of many particles correlation functions) , the method of the non-equilibrium statistical operator , and several other techniques -.
However, in some approaches, the transition to the relativistic region it is not obvious and can lead to misleading results. This situation was analyzed in Refs.- on a rather simple quantum field model which describes interaction of the fermion and boson subsystems in the mean field approximation. In these works, it was shown that the equations of motion for a two-point correlation function lead to true RKE of the Vlasov type only if one uses some definite simple rules of projection. It is important that the same RKE can be derived also by a direct method based on the Heisenberg equation of motion , -. In the latter approach, any ambiguities in the RKE derivation are thus in general avoided. These features of the direct method represent a suitable expedient to obtain a correct generalized RKE in the frame of a relativistic analogue of the Kadanoff-Baym formalism. This is the main result of our study.
Above–mentioned works , - are devoted to covariant generalizations of the Zubarev method of non-equilibrium statistical operator in terms of the relativistic Wigner functions. However, this formalism presents some technical difficulties to go beyond the Born approximation, while the formalism introduced in this work, based on the real–time Green function, avoids this kind of mathematical troubles.
In order to solve this main problem, we develop a specific relativistic modification of the real-time Green’s functions method in the Wigner representation. This improvement proves to be possible since the transition to the Wigner representation in relativistic theory allows us to introduce a time-like direction which can be connected with the momentum argument of the Wigner function. As a result, the time variable determining the evolution of the system can be defined as a true scalar, which plays the role of the proper time in a given point of the eight-dimensional phase space. This invariant time is used also to define the chronological ordering operation in the real-time Green’s functions . It leads to an obvious modification of this method at all stages of the formalism including the description of the dispersion properties of quantum field systems. Such a modification of the real-time Green’s functions method is convenient also in order to describe the dispersion properties of a quantum-field system at finite temperature and density as it proves to be obviously relativistic invariant for all the steps of the relevant calculations.
The outline of the paper is as follows. In Sec. 2, we discuss shortly the ”direct way” of the derivation of RKE in terms of one-particle relativistic Wigner’s functions. The advantages of this approach are its simplicity and clearness. This fact originates from the use of the Heisenberg picture from the very beginning. However, in this section, we do not discuss any truncation scheme and derivation of closed-form RKE. This choice is convenient for a comparison with the method based on the Kadanoff-Baym motion equations for correlation functions (Sec. 3). A covariant modification of the real-time Green’s functions method is found as an adequate formalism for this aim. As a result, we obtain a generalized RKE of the Kadanoff-Baym type without using any irrelevant assumption, thus allowing us to use effectively the suggested method for a reliable kinetic description of a quite wide class of systems with inner degrees of freedom.
In essence, the only limitation of the suggested method is a request of a polynomial character of interaction of systems under study. It is important that the suggested approach results in a self-consistent, non-contradictory, and unambiguous theory of kinetic equations. To illustrate the possibilities of the theory, two examples are discussed in detail. Firstly, a RKE of the Vlasov type (the mean field approximation, Sec. 2) and then a collision integral of the BUU type (the quasiparticle approximation, Sec. 4). The standard Walecka model, which is a quantum field model of the relativistic nuclear matter consisting of nucleons and two types (scalar and vector) of mesons (quantum hadrodynamics) is used in our work as a toy model.
In Sec. 5, we consider another generalization of relativistic kinetic theory based on a change of the standard averaging procedure under the equilibrium density matrix by an averaging using the non-equilibrium statistical operator. This step leads to a possible kinetic description of strongly non-equilibrium states which are often met in the relativistic nuclear physics. Finally, conclusions are drawn in Sec. 6. Everywhere, we work with natural units, $`\mathrm{}=c=1`$c.
## 2 The motion equation for the Wigner function
To illustrate our approach on a nontrivial example of a system with internal degrees of freedom, let us consider the kinetic description of the Fermi subsystem of a quantum-field system. For a consistent dynamical construction of the relativistic kinetic theory, we will start with the introduction of the one-particle Wigner function of the Fermi subsystem,
$$f_{\alpha \beta }(x,p)=(2\pi )^4𝑑ye^{ipy}<P_{\alpha \beta }(x,y)>\text{ ,}$$
(1)
where the symbol $`<\mathrm{}>=\mathrm{Tr}\mathrm{}\rho `$ denotes the operation of statistical averaging with density matrix $`\rho `$ in the Heisenberg representation and
$$P_{\alpha \beta }(x,y)=\overline{\psi }_\beta (x_+)\psi _\alpha (x_{})\text{,}$$
(2)
$`\psi (x)`$ and $`\overline{\psi }(x)`$ being the usual field operators, and $`x_\pm =x\pm y/2`$.
In the realm of dynamical basis of the theory, we choose the Heisenberg equations of motion, namely,
$$i^\mu A(x)=[A(x),P^\mu ]\text{ ,}$$
(3)
where $`A(x)=A[x,\psi (x),\overline{\psi }(x)]`$ is an arbitrary local operator and $`P^\mu `$ is the total 4-momentum of the system,
$$P^\mu =𝑑\sigma _\nu (xn)T^{\mu \nu }\text{ .}$$
(4)
Here, $`d\sigma _\nu (xn)`$ is a vector element of an arbitrary space-like hyperplane, $`\sigma (n),`$ with a time-like normal vector, $`n^\mu (n^2=1),`$ and $`T^{\mu \nu }=T^{\nu \mu }`$ is the energy-momentum tensor..
The energy and momentum conservation laws lead to the independence of the integrals (4) from the selection of a given hyperplane $`\sigma (n)`$ , and we can fix the time-like direction by means of an external condition (a relevant example will be given below).
In order to describe the dynamics of the system along the time-like direction, $`n^\mu ,`$ and the space translation along the independent space-like directions on the hyperplane, $`\sigma (n)`$, a boost transformation with the <sup>′′</sup>velocity<sup>′′</sup>, $`n^\mu ,`$ can be used for the motion Eqs. (3). Technically, this result is achieved using the projection of Eqs. (3) on the direction $`n^\mu `$ and on the hyperplane, $`\sigma (n)`$. Convolution of Eqs. (3) with the vector, $`n^\mu ,`$ leads to the following dynamical equation ,
$$i\frac{A(x)}{\tau }=[A(x),H(\tau )]\text{ .}$$
(5)
The parameter,
$$\tau =n_\mu x^\mu ,$$
(6)
plays here the role of a proper time in a new coordinate system. In Eq.(5), the derivative along the direction of $`n^\mu `$ is thus introduced,
$$\frac{}{\tau }=n^\mu \frac{}{x^\mu }\text{.}$$
(7)
$`H(\tau )`$ is the scalar Hamiltonian of the system and reads as
$$H(\tau )=n_\mu P^\mu =𝑑\sigma (n)n_\mu T^{\mu \nu }n_\nu \text{,}$$
(8)
because $`d\sigma ^\mu (xn)=n^\mu d\sigma (xn)`$. This procedure of a formal restoration of the relativistic invariance in the description of the time evolution is very popular in the quantum field theory (see, e.g., ).
The projection of Eqs. (3) on the space-like directions is carried out by the convolution of Eqs. (3) with the projection operator,
$$\mathrm{\Delta }^{\mu \nu }=g^{\mu \nu }n^\mu n^\nu ,n_\mu \mathrm{\Delta }^{\mu \nu }=0\text{ .}$$
(9)
The resulting equations provide a description of infinitesimal transforms of the system on the space-like hyperplane, $`\sigma (n)`$. These equations are not interesting for the present aim of developing a kinetic theory -.
Now, we can write the motion equation for the Wigner function (1). Performing the $`\tau `$-differentiation of function (1) and using the Liouville equation in the Heisenberg representation,
$$\frac{d\rho }{d\tau }=0\text{ ,}$$
(10)
after substitution of the motion equation (5) and using the definition (7), we find:
$$\frac{f(x,p)}{\tau }=n^\mu \frac{}{x^\mu }f(x,p)=i(2\pi )^4𝑑ye^{ipy}<[P(x,y),H(\tau )]>\text{ .}$$
(11)
It is worth remarking that, for sake of convenience, the spin indices are dropped out in our notation.
At this stage, we can remove the arbitrariness in the choice of the unit time-like vector, $`n^\mu ,`$ using the specific character of the mixed Wigners representation. Let us assume that the momentum vector, $`p^\mu ,`$ in Eq. (11) is a time-like one (Assumption 1). Then, we can define a unit vector in the direction of $`p^\mu `$ and identify it with the vector, $`n^\mu `$(Assumption 2) , -:
$$n^\mu \stackrel{def}{=}u^\mu =p^\mu /\sqrt{p^2}\text{ ,}u^2=1\text{.}$$
(12)
Assuming a description of non-equilibrium quantum field systems in terms of the Wigner functions (1), often one limits oneself to consider only processes on the mass shell,
$$p^2=m^2\text{ .}$$
(13)
In this respect, introduction of Assumptiones 1 and 2 does not produce any essential restriction even at the off mass shell. In other words, Assumption 1 extracts the physical significant domain, $`𝒫^+,`$ from the full phase space, $`𝒫(x,p)`$. Let us remark that an analogous restriction is usually introduced at the stage of the definition of Wigners functions of the states with positive and negative energies . On the other hand, Assumption 1 can be taken off and go out into the space-like region (see the discussion at the end of Sec. 3). We conserve this restriction here in order to remain into the framework of the ordinary dynamical description.
The effectiveness of the introduction of Assumptiones 1 and 2 is confirmed by the fact that substitution of the relation (12) into Eq. (11) completely removes any arbitrariness in this equation and leads to a <sup>′′</sup>generalized RKE<sup>′′</sup>,
$$\begin{array}{c}p_\mu ^\mu (x)f(x,p)=i(2\pi )^4\sqrt{p^2}𝑑ye^{ipy}<[P(x,y),H(\tau )]>=\\ =i(2\pi )^4\sqrt{p^2}𝑑ye^{ipy}<[\overline{\psi }(x_+),H(\tau _+)]\psi (x_{})+\overline{\psi }(x_+)[\psi (x_{}),H(\tau _{})]>,\end{array}$$
(14)
where we took into account the arbitrariness of the assigned hyperplane into the Hamiltonian (8), $`H(\tau )=H(\tau _+)=H(\tau _{})`$. Essentially, this is the first equation of the BBGKY hierarchy in the Wigner representation. In the mean field approximation, Eq. (14) results to be a non-contradictory RKE of the Vlasov type -, .
In a general case, it is necessary to introduce different truncation schemes to obtain closed-form RKE from Eq. (14). The usual perturbation theory on the coupling constant starting from Eq. (14) was developed in Refs. , -. Within the framework of such a perturbation theory, two types of RKE with collision integrals of the second order have been obtained: RKE of the Bloch type (for the vertices of the <sup>′′</sup>three-tails<sup>′′</sup> type) and RKE of the Boltzmann type (for the effective vertices of the <sup>′′</sup>four-tails<sup>′′</sup> type). The Greens functions method offers an effective way to go beyond the perturbation approach. In the following Section, Eq.(14) will be re-formulated in terms of modified contour Greens functions. In particular, this will allow us to get a unambiguous generalized RKE of the Kadanoff-Baym type, thus opening some perspectives to go beyond the standard perturbation theory without breakdown of the relativistic invariance of kinetic theory.
## 3 Generalized RKE in terms of modified <br>real-time Greens functions
### 3.1 Covariant real-time Greens functions
The modern technique of derivation of the Kadanoff-Baym type kinetic equations is based on the real-time Greens functions method . Now, our aim is to combine together the <sup>′′</sup> generalized RKE<sup>′′</sup>(14) and a system of motion equations of the Kadanoff-Baym type for the real-time Greens functions. As a first step, we will modify this method taking into account the existence of a preferred time-like direction (Assumptiones 1 and 2 in Sec. 2) in the definition of the Greens functions in the Wigner representation. This direction is fixed by relation (12) and leads to scalar time variables (see Eq. (6)),
$$\tau _i=x_i^\mu n_\mu ,n^\mu u^\mu \text{ ,}$$
(15)
where $`x_i^\mu `$ are arguments of the two-points Greens functions. It is implied that identification (12) is fulfilled after a transition to the Wigner representation. The time arguments (15) are convenient for use in the definition of the ordering operation determining the modified real-time Greens functions. Limiting ourselves to the case of a Fermi subsystem, let us introduce the following Greens functions in the Wigner representation,
$$G_{\alpha \beta }^{(k)}(xp)=(2\pi )^4𝑑ye^{ipy}G_{\alpha \beta }^{(k)}(x_+,x_{})\text{ ,}$$
(16)
where index (k) enumerates the elements of the real-time (contour) Greens function,
$$\begin{array}{ccccc}G_{\alpha \beta }^c(x_1,x_2)\hfill & =& i<T^c[\psi _\alpha (x_1)\overline{\psi }_\beta (x_2)]>\hfill & =& G_{12}^{}\text{ ,}\hfill \\ G_{\alpha \beta }^a(x_1,x_2)\hfill & =& i<T^a[\psi _\alpha (x_1)\overline{\psi }_\beta (x_2)]>\hfill & =& G_{12}^{++}\text{ ,}\hfill \\ G_{\alpha \beta }^>(x_1,x_2)\hfill & =& i<\psi _\alpha (x_1)\overline{\psi }_\beta (x_2)>\hfill & =& G_{12}^+\text{ ,}\hfill \\ G_{\alpha \beta }^<(x_1,x_2)\hfill & =& i<\overline{\psi }_\beta (x_2)\psi _\alpha (x_1)>\hfill & =& G_{12}^+\text{ ,}\hfill \end{array}$$
(17)
and $`(T^a)T^c`$ represents (anti)chronological time ordering along the $`\tau `$-axis:
$$\begin{array}{ccc}T^c[A(x_1)B(x_2)]\hfill & \stackrel{def}{=}& \theta (\tau _1\tau _2)A(x_1)B(x_2)\theta (\tau _2\tau _1)B(x_2)A(x_1)\text{ ,}\hfill \\ T^a[A(x_1)B(x_2)]\hfill & \stackrel{def}{=}& \theta (\tau _2\tau _1)A(x_1)B(x_2)\theta (\tau _1\tau _2)B(x_2)A(x_1)\text{ .}\hfill \end{array}$$
(18)
The operators, $`A(x)`$and $`B(x),`$ are selected from the set of the field operators, $`\overline{\psi }(x)`$ and $`\psi (x)`$.
The path-ordered real-time contour Greens functions are defined by the relation
$$G(x_1,x_2)=i<T_c[\psi (x_1)\overline{\psi }(x_2)]>\text{ ,}$$
(19)
where $`T_c`$ is the path-ordering operator on the $`\tau `$-contour which is chosen in the usual way (see Fig. 1).
Fig. 1. The Keldysh-Schwinger contour, $`C,`$ with the upper $`C^+`$ and lower $`C^{}`$ branches.
It is worth stressing that this definition is based on the Wigner representation (16) and Assumptiones 1 and 2. This restriction leads by itself to a break-down of a simple analogy between the Wigner and the Fourier transforms.
Let us return to Eq. (14). From the definitions, (1) and (16,17), it follows that
$$G^<(xp)=if(xp)\text{ .}$$
(20)
Therefore, Eq. (14) can be rewritten in the form
$$p^\mu _\mu (x)G^<(xp)=(2\pi )^4\sqrt{p^2}𝑑ye^{ipy}<[P(x,y),H(\tau )]>\text{ .}$$
(21)
Let us assume now that the Hamiltonian of the fermion subsystem has a polynomial structure and adopt the usual notation,
$$H(\tau )=H_0(\tau )+H_{in}(\tau ),H_{in}(\tau )=H_{MF}(\tau )+H_r(\tau )\text{ .}$$
(22)
Here, $`H_0(\tau )`$ describes the free evolution of the fermion field,
$$H_0(\tau )=𝑑\sigma (xu)\overline{\psi }(x)\{\frac{i}{2}\gamma ^\mu \stackrel{}{_\mu ^{}}(x)m\}\psi (x)\text{ ,}$$
(23)
where $`_\mu ^{}(x)`$ is the space-like derivative,
$$_\mu ^{}(x)=\mathrm{\Delta }_\mu ^\nu _\nu (x)\text{ ,}$$
(24)
and $`\mathrm{\Delta }^{\mu \nu }`$ is the projection operator (9). In the interaction Hamiltonian, $`H_{in}(\tau ),`$ Eq. (22), the $`H_{MF}(\tau )`$ part, corresponding to the mean field approximation, is explicitly separated. Then, $`H_r(\tau )`$ describes the residual part of the interaction.
It is easy to calculate the free motion part in the right-hand side of Eq. (21) for the Hamiltonian (23). For covariant transformations, it is convenient to use the following rule,
$$𝑑\sigma (x^{}u)S(xx^{})y(x^{})=iu_\mu \gamma ^\mu y(x),x\sigma (u)\text{ ,}$$
(25)
where $`y(x)`$ is an arbitrary function of field operators and $`S(x)`$ is the anti-commutation function of the fermion fields, namely $`S(xx^{})=i[\psi (x),\overline{\psi }(x^{})]_+.`$ Relation (25) is based on a covariant generalization of the well-known property of the function $`S(x)`$ (see, e.g., )
$$S(x)_{x^0=0}=i\gamma ^0\delta ^{(3)}(x)\text{ ,}$$
(26)
(a simple proof of Eq. (25) is given in Appendix A). This kind of calculation for Eq. (21) leads to the following result:
$$\begin{array}{c}p^\mu _\mu (x)G^<(xp)+\frac{1}{2}[\widehat{p}\gamma ^\mu ,_\mu (x)G^<(xp)]+im[\widehat{p},G^<(xp)]=\hfill \\ =(2\pi )^4\sqrt{p^2}𝑑ye^{ipy}<[\overline{\psi }(x_+),H_{in}(\tau _+)]\psi (x_{})+\overline{\psi }(x_+)[\psi (x_{}),H_{in}(\tau _{})]>\text{ ,}\hfill \end{array}$$
(27)
where $`\widehat{p}=p^\mu \gamma _\mu `$.
The following step consists in the search of a connection between the right-hand side of Eq. (27) and the mass operators, $`\mathrm{\Sigma }^{(k)}(x_1,x_2),`$ defined on the $`\tau `$-contour. To this aim, let us write the motion equations for the real-time Green’s functions,
$$\begin{array}{ccc}\hfill [i\gamma (x_1)m]G(x_1,x_2)& =& \delta _c(x_1,x_2)+_cd^4x^{}\mathrm{\Sigma }(x_1,x^{})G(x^{},x_2)\text{ ,}\hfill \\ \hfill G(x_1,x_2)[i\gamma \stackrel{}{}(x_2)m]& =& \delta _c(x_1,x_2)+_cd^4x^{}G(x_1,x^{})\mathrm{\Sigma }(x^{},x_2)\text{,}\hfill \end{array}$$
(28)
The integration in the right-hand side of Eqs. (28) is performed according to the rules:
$$_cd^4x\mathrm{}=𝑑\sigma (xu)_c𝑑\tau \mathrm{}.\text{ ,}$$
(29)
and
$$_c𝑑\tau \mathrm{}=_{t_0}^{\mathrm{}}𝑑\tau \mathrm{}_{c^+}_{t_0}^{\mathrm{}}𝑑\tau \mathrm{}_c^{}\text{ ,}$$
(30)
where $`c^+(c^{})`$ is the upper (lower) branch of the $`\tau `$-contour (see Fig. 1). Finally, the delta function on the $`\tau `$-contour is defined as
$$\delta _c(x_1,x_2)=\delta _\sigma (x_1,x_2)\delta _c(\tau _1\tau _2)\text{ ,}$$
(31)
where
$$\delta _c(\tau _1\tau _2)=\{\begin{array}{ccc}\hfill \delta (\tau _1\tau _2)& \text{ if }& \tau _1c^+\text{ and }\tau _2c^+\text{ ,}\\ \hfill \delta (\tau _1\tau _2)& \text{ if }& \tau _1c^{}\text{ and }\tau _2c^{}\text{ ,}\\ \hfill 0& & \text{ otherwise }.\end{array}$$
(32)
$`\delta _\sigma (x)`$ is the $`d=3`$ delta function of the space arguments on the hyperplane $`\sigma (xu)`$.
For purpose of comparison with Eq. (27), let us write the motion equations for the correlation function, $`G^<(x_1,x_2)`$. From Eqs. (28), the Kadanoff-Baym type equations follow
$$[i\gamma (x_1)m\mathrm{\Sigma }_{MF}(x_1)]G^<(x_1,x_2)=𝑑x^{}\{\mathrm{\Sigma }^<(x_1,x^{})G^A(x^{},x_2)+\mathrm{\Sigma }^R(x_1,x^{})G^<(x^{},x_2)\}\text{ ,}$$
(33)
$$G^<(x_1,x_2)[i\gamma \stackrel{}{}(x_2)m\mathrm{\Sigma }_{MF}(x_2)]=𝑑x^{}\{G^<(x_1,x^{})\mathrm{\Sigma }^A(x^{},x_2)+G^R(x_1,x^{})\mathrm{\Sigma }^<(x^{},x_2)\}\text{ .}$$
(34)
Here, the mean field part of the mass operator was extracted according to the following prescription:
$$\begin{array}{ccc}\hfill \mathrm{\Sigma }(x_1,x_2)& =& \mathrm{\Sigma }_{MF}(x_1)\delta (x_1x_2)+\hfill \\ & +& \theta (\tau _1\tau _2)\mathrm{\Sigma }^>(x_1,x_2)+\theta (\tau _2\tau _1)\mathrm{\Sigma }^<(x_1,x_2)\text{ .}\hfill \end{array}$$
(35)
Moreover, the retarded and advanced Green’s functions were introduced,
$$\begin{array}{ccc}G^{R/A}(x_1,x_2)\hfill & =& \pm \theta [\pm (\tau _1\tau _2)]\{G^>(x_1,x_2)G^<(x_1,x_2)\}=\hfill \\ & =& G^c(x_1,x_2)G^{</>}(x_1,x_2)=G^{>/<}(x_1,x_2)G^a(x_1,x_2)\text{ .}\hfill \end{array}$$
(36)
In a similar way, the retarded and advanced components of the mass operator are defined. A mutual exchange of the symbols $`>\text{ and }<`$ transforms Eqs. (33) and (34) into the corresponding equations for the correlations functions, $`G^>(x_1,x_2)`$. By this way, the Kadanoff-Baym equations (33) and (34) have the usual form however they are defined using covariant modified constituents (Green’s functions and mass operators). However, these small corrections are essential for the search of the agreement with RKE (27).
The transition to the mixed Wigner’s representation for the motion Eqs. (33) and (34) is usually done with a limitation to the lowest orders of the gradient decomposition (the first order into the drift part of RKE and order zero for the collision part) , approximation that is justified only for rather slow kinetic processes. This allows us to write the motion equations in the simplest local form. This restriction casts some doubt upon the description of non-equilibrium processes in nuclear matter at extreme conditions (see, e.g., ). However, for simplicity’s sake, we keep here this approximation only. Then, the motion Eqs. (33) and (34) in the Wigner representation are transformed into the following equations:
$`[\widehat{p}+{\displaystyle \frac{i}{2}}\gamma (x)m\mathrm{\Sigma }_{MF}(x)+{\displaystyle \frac{i}{2}}_\mu (x)\mathrm{\Sigma }_{MF}(x)^\mu (p)]G^<(xp)=`$
$`=\mathrm{\Sigma }^<(x)G^A(xp)+\mathrm{\Sigma }^R(x)G^<(xp)\text{ ,}`$ (37)
$`G^<(xp)[\widehat{p}{\displaystyle \frac{i}{2}}\gamma \stackrel{}{}(x)m\mathrm{\Sigma }_{MF}(x){\displaystyle \frac{i}{2}}\underset{\mu }{\overset{}{}}(p)^\mu (x)\mathrm{\Sigma }_{MF}(x)]=`$
$`=G^<(xp)\mathrm{\Sigma }^A(xp)+G^R(xp)\mathrm{\Sigma }^<(xp)\text{ .}`$ (38)
We have neglected gradient terms in the right-hand sides of these equations .
### 3.2 Generalized RKE
Let us fulfill now the concordance of the ”generalized RKE” (14) and the system of equations of the Kadanoff-Baym type (3.1) and (3.1). Obviously, for this aim, it is necessary to find such a combination of these equations which provides a coincidence with the left-hand side of Eq. (27) describing the free evolution of the fermion subsystem. This assumption will permit us to identify (in the corresponding orders of the gradient expansion) also the other parts of Eq. (27) and Eqs. (3.1), (3.1).
In order to perform such a construction, it is sufficient to multiply Eq. (3.1) by $`i\widehat{p}`$ from the right and Eq. (3.1) from the left and to subtract the second result from the first one, i.e., it is necessary to consider the following algebraic combination:
$$Eq.(\text{3.1})i\widehat{p}i\widehat{p}Eq.(\text{3.1})\text{ .}$$
(39)
Then, we get the following equation
$`p^\mu _\mu (x)G^<(xp)+{\displaystyle \frac{1}{2}}[\widehat{p}\gamma _\mu ,^\mu (x)G^<(xp)]+im[\widehat{p},G^<(xp)]`$
$`i\{G^<(xp)\mathrm{\Sigma }_{MF}(x)\widehat{p}\widehat{p}\mathrm{\Sigma }_{MF}(x)G^<(xp)\}+`$ (40)
$`+{\displaystyle \frac{1}{2}}\{_\mu (p)G^<(xp)^\mu (x)\mathrm{\Sigma }_{MF}(x)\widehat{p}+\widehat{p}_\mu (x)\mathrm{\Sigma }_{MF}(x)^\mu (p)G^<(xp)\}=S(xp)\text{ ,}`$
where $`S(xp)`$ is the collision integral,
$`S(xp)=i\widehat{p}\{\mathrm{\Sigma }^<(xp)G^A(xp)+\mathrm{\Sigma }^R(xp)G^<(xp)\}`$
$`i\{G^<(xp)\mathrm{\Sigma }^A(xp)+G^R(xp)\mathrm{\Sigma }^<(xp)\}\widehat{p}\text{ .}`$ (41)
Now, from a comparison of Eqs. (27) and (3.2), it follows an equivalence of the left and right parts of those relations:
$`(2\pi )^4\sqrt{p^2}{\displaystyle 𝑑ye^{ipy}<[\overline{\psi }(x_+),H_{MF}(\tau _+)]\psi (x_{})+\overline{\psi }(x_+)[\psi (x_{}),H_{MF}(\tau _{})]>}=`$
$`=i\{G^<(xp)\mathrm{\Sigma }_{MF}(x)\widehat{p}\widehat{p}\mathrm{\Sigma }_{MF}(x)G^<(xp)\}`$ (42)
$`{\displaystyle \frac{1}{2}}\{_\mu (p)G^<(xp)^\mu (x)\mathrm{\Sigma }_{MF}(x)\widehat{p}+\widehat{p}_\mu (x)\mathrm{\Sigma }_{MF}(x)^\mu (p)G^<(xp)\}\text{ ,}`$
$$S(xp)=(2\pi )^4\sqrt{p^2}𝑑ye^{ipy}<[\overline{\psi }(x_+),H_r(\tau _+)]\psi (x_{})+\overline{\psi }(x_+)[\psi (x_{}),H_r(\tau _{})]>\text{.}$$
(43)
Let us remind that these equalities are correct only in the lowest orders of the gradient expansions.
Both parts of the equality (3.2) describe the force contribution (with the inverse signs) of the mean fields into the convective part of RKE (27) or (3.2). Let us illustrate this aspect by means of an actual model. We thus refer to the standard Walecka model, where the Hamiltonian for the interaction of nucleons with massive vector and scalar meson fields reads as :
$$H_{in}(\tau )=𝑑\sigma (xu)\overline{\psi }(x)\{g_v\omega ^\mu (x)\gamma _\mu g_s\varphi (x)\}\psi (x)\text{ .}$$
(44)
In the mean field approximation, the relevant Hamiltonian, $`H_{MF}(\tau ),`$ can be deduced from this one by replacing the meson fields operators with their mean values, $`\varphi <\varphi >,`$ and $`\omega ^\mu <\omega ^\mu >`$. Then, the corresponding contribution to the mass operator is equal to
$$\mathrm{\Sigma }_{MF}(x)=g_v<\omega ^\mu (x)>\gamma _\mu g_s<\varphi (x)>\text{.}$$
(45)
In the first order of the gradient expansion, both RKE (27) and (3.2) lead to identical results in the frame of the mean field approximation:
$$\begin{array}{ccc}P^\mu _\mu (x)G^<(xP)+\frac{1}{2}^\mu (x)M(x)\{\widehat{P},_\mu (P)G^<(xP)\}\hfill & +& \\ +g_vP^\mu F_{\mu \nu }(x)^\nu (P)G^<(xP)+\frac{1}{2}[\widehat{P}\gamma ^\mu ,_\mu (x)G^<(xP)]\hfill & +& \\ +iM(x)[\widehat{P},G^<(xP)]\frac{1}{2}g_vF_{\mu \nu }(x)[\widehat{P}\gamma ^\mu ,^\nu (P)G^<(xP)]\hfill & =& S(xP)\text{ ,}\hfill \end{array}$$
(46)
where $`M(x)=mg_s<\varphi (x)>`$ is the effective nucleon mass in the mean field approximation, $`F^{\mu \nu }(x)=^\mu <\omega ^\nu >^\nu <\omega ^\mu >`$ and $`P^\mu =p^\mu g_v<\omega ^\mu (x)>`$ is the kinetic momentum.
RKE (46) converts in RKE of the Vlasov type when conditions of the collision effects are neglected $`(S(xP)=0)`$. The Vlasov RKE in this form was obtained and analyzed for the first time in Refs. -. These equations are rather complicated because the spin and meson degrees of freedom as well as states with positive and negative energies are taken into account. On the basis of this RKE, a series of more simple, particular cases can be analyzed. The simplest situation corresponds to the case of a spin saturated system without antinucleon states, where all the spin-dependent effects in the RKE can be neglected (the spin effects were investigated separately in Ref. ). Using the Clifford decomposition of the Wigner function, we get – in the quasi-classical limit – the well-known RKE of the Vlasov type in the quantum hadrodynamics,
$`P^\mu _\mu (x)f^S(xP)+M^\mu (x)M_\mu (P)f^S(xP)+g_vP_\mu F^{\mu \nu }_\nu (P)f^S(xP)=0\text{ ,}`$ (47)
where $`f^S(xP)`$ corresponds to the scalar part of the Clifford decomposition of function $`G^<(xP)`$ (see Eq. (20)). It is important that the basic RKE (46) leads to a correct RKE (47) of the Vlasov type for a spin saturated system without any additional assumption, like that of the rescaling procedure type (see the detailed discussion in ).
Technical details of derivation of Eqs.(46)-(47) are given in Refs. -. It is not difficult to obtain analogous RKE also without use of the gradient expansion. Such a RKE has a non-local character in the momentum space.
Calculation of collision integrals can be performed using two alternative schemes based either on formula (3.2) or (43). In the latter case, a standard perturbation theory was developed in Refs. . Within the framework of the Walecka model (44), this approach allows us to derive collision integrals of the Bloch type (when the meson subsystem is considered on equal terms with the nucleon one) or of the Boltzmann type (if the role of mesons is confined only to the formation of the effective nucleon-nucleon interaction).
Let us return to the representation of the collision integral (3.2) based on the method of modified real-time Green’s functions. Our subsequent transformations are based on the analogy with the nonrelativistic theory . From the definition (36), it follows that the Green’s functions and mass operator satisfy the identities:
$$\begin{array}{c}G^R(xp)G^A(xp)=G^>(xp)G^<(xp)\text{ ,}\\ \mathrm{\Sigma }^R(xp)\mathrm{\Sigma }^A(xp)=\mathrm{\Sigma }^>(xp)\mathrm{\Sigma }^<(xp)\text{ .}\end{array}$$
(48)
Using these relations, the collision integral (3.2) can be transformed into the form
$$\begin{array}{ccc}S(xp)\hfill & =& \frac{i}{2}\{\widehat{p}[\mathrm{\Sigma }^>(xp)G^<(xp)\mathrm{\Sigma }^<(xp)G^>(xp)]+\hfill \\ & +& [G^<(xp)\mathrm{\Sigma }^>(xp)G^>(xp)\mathrm{\Sigma }^<(xp)]\widehat{p}\}+S^{off}(xp)\text{ ,}\hfill \end{array}$$
(49)
where $`S^{off}(xp)`$ is the off-mass-shell part of the collision integral,
$$\begin{array}{ccc}S^{off}(xp)\hfill & =& \frac{i}{2}\{\widehat{p}\mathrm{\Sigma }^<(xp)[G^R(xp)+G^A(xp)]+\widehat{p}[\mathrm{\Sigma }^R(xp)+\mathrm{\Sigma }^A(xp)]G^<(xP)\hfill \\ & & G^<(xp)[\mathrm{\Sigma }^R(xp)+\mathrm{\Sigma }^A(xp)]\widehat{p}[G^R(xp)+G^A(xp)]\mathrm{\Sigma }^<(xp)\widehat{p}\}\text{ .}\hfill \end{array}$$
(50)
Of course, the off–shell effects enter in the whole collision operator (49), but function (50) refers to them exclusively for the present aims. The separation into two parts, (49) and (50), is ascertained by means of spectral properties of the Green’s functions and the mass operator. In our formalism, the corresponding spectral decompositions can be fulfilled in a covariant form (see Appendix B). This circumstance is a very desirable property of the relativistic kinetic theory which is intended for a description of macroscopic motions of a medium. Let us remark that the off-shell processes play the most important role in the relativistic nuclear physics of intermediate and high energies.
The generalized RKE (3.2) with the collision integral (49), (50) are the main results of this Section. It is worth remarking that a few other versions of generalized RKE have been previously obtained in the literature (see, e.g., and ). We have in mind the cases of derivation of RKE for the Fermi subsystem based either on the Dirac type equations of motion for the field operator or on the equations of the Kadanoff-Baym type (33), (34). However, it is clear that this procedure contains some elements of an indetermination and arbitrariness on the stage of the construction of the motion equation for the Wigner function. In considered our case, the selection rule (39) has a reliable dynamical basis (see Sec. 2) which was tested in Refs. -. Of course, the generalized RKE (3.2) can be obtained also by a direct way on the basis of Eq. (27), without using the rule (39), with a help of the Dyson equation and definitions of the type of Eqs. (17).
RKE (3.2), (49), and (50) have a more complicated matrix structure. This characteristic behavior arises by taking into account inner degrees of freedom which are intrinsic for the considered Fermi subsystem and, in particular, states with positive and negative energies. Indeed, the presence of $`\gamma ^0`$ matrices before the time derivative in equations of the Kadanoff-Baym type (33) and (34) ensures the ”direct” time movement for states of the spinor fields with positive energies and the ”inverse” time direction for states with negative energies. But for the kinetic theory, it is characteristic the existence of a single time arrow. The rule (39) just ensures a covariant way of projection of Eqs. (33) and (34) on such a single time direction.
Finally, let us return once more to the role of the assumption $`p^2>0`$ (Assumption 1). A consecutive realization of this limitation leads to appreciable complication of the considered integrals. Similar restriction is absent in the traditional approach where an integration is carried out in the whole momentum space. In order to get traditional results in the frame of our formalism, it is necessary to implement the extension of obtained RKE into space-like region. There are three aspects of such a continuation:
1) into the integrals with respect to the momentum space,
2) into the time differential operations,
3) into the time integral operations.
The extension of the momentum integrals into the region, $`p^2<0,`$ is trivial in the quasiparticle approximation (Sec. 4).
As for the derivation of the Vlasov RKE, the considered continuation does not lead to some changes of differential operations with respect to the time (e.g., to a change $`/\tau `$ by $`i/\tau `$). Indeed, in the mean field approximation, the calculation of the force part of the Vlasov RKE (the right-hand side in Eq.(27)) results in the cancellation of $`\sqrt{p^2}`$ once rule (25) has been used. In the other words, RKE conserves its physical content.
Finally, in the Markovian approximation, integral time operations do not lead to any modification of results at the continuation into the region $`p^2<0`$.
The possibility of continuation beyond the above-mentioned cases needs further researches.
## 4 Quasiparticle approximation
In this Section, we derive a covariant generalization of the quasiparticle approximation (QPA), widely used in the kinetic theory, which allows us to leave, by the simplest way, the framework of the usual perturbation theory. After that, we check possibilities of the suggested method on an example of derivation of collision integral of the Boltzmann-Uehling-Uhlenbeck (BUU) type for the standard Walecka model of relativistic nuclear matter consisting of nucleons, scalar and vector mesons with the interaction Hamiltonian (44). Finally, we discuss briefly a feasible generalization of the QPA in the model (44).
Let us introduce a non-equilibrium spectral function in the Wigner representation (see the first formula in Eq.$`(48)`$)
$$a(xp)=i[G^R(xp)G^A(xp)]=i[G^>(xp)G^<(xp)]\text{ ,}$$
(51)
This function satisfies the following sum rule
$$(2\pi )^3𝑑Ea(x;p_{},E)=\gamma u\text{ ,}$$
(52)
which is a simple consequence of the single-time anti-commutation relations for the field operators and of the rule (25). In Eq. (52), the integration is carried out over the longitudinal energy, $`E=pn;p_{}^\mu =\mathrm{\Delta }^{\mu \nu }p_\nu `$ is the transverse momentum ($`\mathrm{\Delta }^{\mu \nu }`$ is the projection operator (9)).
$`G^>(xp)`$ and $`G^<(xp),`$ can be represented as:
$$\begin{array}{ccc}G^<(xp)\hfill & =& ia(xp)(xp)\text{ ,}\hfill \\ G^>(xp)\hfill & =& ia(xp)[1(xp)]\text{ ,}\hfill \end{array}$$
(53)
where $`(xp)`$ is a unknown function.
To determine the non-equilibrium spectral function, we consider relation (51) and the corresponding motion equations for the retarded and advanced Green’s functions. We write at once these equations in the minimal order of the gradient expansion, namely,
$$\begin{array}{c}\hfill [\widehat{p}+\frac{i}{2}\gamma (x)m\mathrm{\Sigma }_{MF}(x)+\frac{i}{2}_\mu (x)\mathrm{\Sigma }_{MF}(x)^\mu (p)]G^{R/A}(xp)=\\ \hfill =1+\mathrm{\Sigma }^{R/A}(xp)G^{R/A}(xp),\\ \hfill G^{R/A}(xp)[\widehat{p}\frac{i}{2}\gamma \stackrel{}{}(x)m\mathrm{\Sigma }_{MF}(x)\frac{i}{2}\underset{\mu }{\overset{}{}}(p)^\mu (x)\mathrm{\Sigma }_{MF}(x)]=\\ \hfill =1+G^{R/A}(xp)\mathrm{\Sigma }^{R/A}(xp).\end{array}$$
(54)
Since the gradient expansion orders in the drift and collisions parts of RKE (3.2) are fixed, we can restrict ourselves to order zero of this decomposition in the estimation of correlation functions (53), so that ,
$$[\widehat{p}m\mathrm{\Sigma }_{MF}(x)\mathrm{\Sigma }^{R/A}(xp)]G^{R/A}(xp)=1.$$
(55)
In our toy model, $`\mathrm{\Sigma }_{MF}(x)`$ is defined by Eq. (45) and hence,
$$[\widehat{P}M(x)\mathrm{\Sigma }^{R/A}(xp)]G^{R/A}(xp)=1.$$
(56)
Let us assume now — as a working approximation — that the mass operator, $`\mathrm{\Sigma }(xp),`$ has the same matrix structure as the mean field part (45) (in the Walecka model (44), this implies a possibility to neglect the tensor part of Clifford’s decomposition of the mass operator),
$$\mathrm{\Sigma }(xp)\mathrm{\Sigma }_s(xp)+\gamma ^\mu \mathrm{\Sigma }_\mu (xp)$$
(57)
(a more general form is given in Ref.). In this expression as well as the following ones, the marks $`R/A`$ are omitted for simplicity’s sake. Eq. (56) has then a matrix structure like the Dirac equation of noninteracting fields (quasiparticle structure)
$$[\widehat{𝒫}(xp)(xp)]G^{R/A}(xp)=1,$$
(58)
where
$$𝒫_\mu (xp)=P_\mu (xp)\mathrm{\Sigma }_\mu (xp),$$
(59)
$$(xp)=M(xp)\mathrm{\Sigma }_s(xp).$$
(60)
Therefore, Eq. (58) can be easily solved:
$$G^{R/A}(xp)=\frac{\widehat{𝒫}+}{𝒫^2^2\pm iϵE}.$$
(61)
Then, from Eqs. (51) and (61), we get
$$a(x𝒫)=2\pi (\widehat{𝒫}+)\delta [Re(𝒫^2^2)]\{\theta (E)\theta (E)\}.$$
(62)
This result is correct, provided that the quasiparticle excitations are weakly damped, i.e.,
$$|Re(𝒫^2^2)|>>|Jm(𝒫^2^2)|.$$
(63)
It is easy to verify that the spectral function (62) satisfies the sum rule (52).
The substitution of the non-equilibrium spectral function (62) into Eqs. (53) leads to the following expressions for the correlation functions,
$$\begin{array}{ccc}G^<(x𝒫)\hfill & =& 2\pi i(\widehat{𝒫}+)\delta (𝒫_{}^{}{}_{}{}^{2}_{}^{}{}_{}{}^{2})(x𝒫^{}),\hfill \\ G^>(x𝒫)\hfill & =& 2\pi i(\widehat{𝒫}+)\delta (𝒫_{}^{}{}_{}{}^{2}_{}^{}{}_{}{}^{2})[1(x𝒫^{})],\hfill \end{array}$$
(64)
where $`𝒫^{}=Re𝒫`$ and $`^{}=Re`$.
As an illustration, we present a derivation of a collision integral of the BUU type in the Walecka model (44). This example is well known in the literature, therefore it is convenient for testing our method. We confined ourselves to the Born approximation only in the calculation of the self energy parts, $`\mathrm{\Sigma }^{>/<},`$ in the collision integral (49). The corresponding diagrams are shown in Fig. 2.
Fig. 2. Born diagrams for the self-energy part, $`i\mathrm{\Sigma }^<=i\mathrm{\Sigma }^+`$. The solid, dashed and wavy lines denote the propagators of nucleons, scalar and vector mesons, respectively.
In a general case, this collision integral has a rather complicated form since we take into account simultaneously the spin and meson degrees of freedom and states with positive and negative energies. The simplest situation corresponds to a spin saturated nuclear matter without antinuclear component. In order to get the corresponding collision integral, we have to perform a transform from the spinor representation to the spin one for states with positive energies and to take into account only the diagonal part of the Wigner function in the spin representation. As a result, for function $`(x𝒫)`$ in Eq. (64), we obtain (for simplicity, we omit here and below all primes in functions $`^{}\text{ and }𝒫^{}`$)
$$_{\alpha \beta }(x𝒫)\underset{r,s=1,2}{}u_\alpha ^s\overline{u}_\beta ^r\mathrm{\Phi }_{sr}(x𝒫)\frac{(\widehat{𝒫}+)_{\alpha \beta }}{2𝒫^0}\mathrm{\Phi }(x𝒫),$$
(65)
where $`\mathrm{\Phi }_{sr}(x𝒫)`$ is the Wigner function of states with positive energy in the spin representation. For a spin-saturated nucleon subsystem we have $`\mathrm{\Phi }_{rs}\mathrm{\Phi }\delta _{rs}`$, where $`\mathrm{\Phi }(x𝒫)`$ is the scalar Wigner function. In Eq.(65), we used the completeness relation for the Dirac spinor
$$\underset{r=1,2}{}u^r\overline{u}^r=\frac{\widehat{𝒫}+}{2\sqrt{^2+\overline{𝒫}^2}}.$$
Substitution of Eq. (65) into the first formula of Eqs. $`(64)`$ leads to the relation
$$G^{<(+)}(x𝒫)=i\pi \frac{\widehat{𝒫}+}{𝒫^0}\delta \left(𝒫^0\sqrt{^2+\overline{𝒫}^2}\right)F(x𝒫),$$
(66)
where a new auxiliary function is introduced:
$$F(x𝒫)=\frac{}{𝒫^0}\mathrm{\Phi }(x𝒫).$$
(67)
This expression differs from the true Wigner function, $`\mathrm{\Phi }(x𝒫),`$ by the scaling factor, $`/𝒫^0`$ . Introduction of such a rescaled Wigner function (67) is convenient in order to write the collision integral.
A tedious but straightforward evaluation of the on-mass-shell part of the collision integral (49) leads to the following result:
$`S^{(+)}(x𝒫)=(2\pi )^5{\displaystyle \underset{i=2}{\overset{4}{}}\frac{d^3𝒫_i}{\sqrt{^2+\overline{𝒫}_i^2}}W(𝒫,𝒫_2,𝒫_3,𝒫_4)\delta (𝒫𝒫_2𝒫_3𝒫_4)}`$
$`\{F(x𝒫)F(x𝒫_2)\overline{F}(x𝒫_3)\overline{F}(x𝒫_4)\overline{F}(x𝒫)\overline{F}(x𝒫_2)F(x𝒫_3)F(x𝒫_4)\},`$ (68)
where $`\overline{F}=1F`$ is the Pauli blocking factor. Here, the transition rates are defined as follows:
$$W=W_{ss}+W_{vv}+W_{sv}.$$
(69)
The partial transition rates are given by exchange of the corresponding pairs of the scalar or vector mesons:
$`W_{ss}(𝒫𝒫_2𝒫_3𝒫_4)`$ $`=`$ $`4g_s^4\{4D_s^2(𝒫𝒫_3)[^4+^2(𝒫𝒫_3+𝒫_2𝒫_4)+`$
$`+(𝒫𝒫_3)(𝒫_2𝒫_4)]D_s(𝒫𝒫_3)D_s(𝒫𝒫_4)[^4+`$
$`+^2(𝒫𝒫_2+𝒫𝒫_3+𝒫𝒫_4+𝒫_2𝒫_3+𝒫_2𝒫_4+𝒫_3𝒫_4)+`$
$`+(𝒫_2𝒫_4)(𝒫𝒫_3)+(𝒫_2𝒫_3)(𝒫𝒫_4)(𝒫𝒫_2)(𝒫_3𝒫_4)]\}\text{ ,}`$
$`W_{vv}(𝒫𝒫_2𝒫_3𝒫_4)`$ $`=`$ $`8g_v^4\{4D_v^2(𝒫𝒫_3)[2^4^2(𝒫_2𝒫_4+𝒫𝒫_3)+`$
$`+(𝒫𝒫_2)(𝒫_3𝒫_4)+(𝒫𝒫_4)(𝒫_2𝒫_3)]`$
$`D_v(𝒫𝒫_3)D_v(𝒫𝒫_4)[2^4+^2(𝒫_2𝒫_3+𝒫𝒫_4+`$
$`+𝒫𝒫_3+𝒫𝒫_2+𝒫_3𝒫_4+𝒫_2𝒫_4)2(𝒫𝒫_2)(𝒫_3𝒫_4)]\}\text{ ,}`$
$`W_{sv}(𝒫𝒫_2𝒫_3𝒫_4)`$ $`=`$ $`8g_s^2g_v^2\{4D_s(𝒫𝒫_3)D_v(𝒫𝒫_3)^2[𝒫𝒫_2+`$
$`+𝒫𝒫_4+𝒫_2𝒫_3+𝒫_3𝒫_4]+D_s(𝒫𝒫_3)D_v(𝒫𝒫_4)[2^4+`$
$`+^2(𝒫𝒫_2+2𝒫𝒫_3𝒫𝒫_4+2𝒫_2𝒫_4𝒫_2𝒫_3`$
$`𝒫_3𝒫_4)+2(𝒫𝒫_3)(𝒫_2𝒫_4)]\}\text{ .}`$
Finally, $`D_{s,v}(𝒫)`$ are the Fourier transforms of the mesons Green’s functions, $`D_{s,v}(𝒫)=(𝒫^2m_{s,v}^2)^1\text{ ,}`$ and $`m_s`$ and $`m_v`$ are masses of the scalar and vector mesons, respectively.
It is worth noticing that similar results have been recently obtained in Ref. within the framework of non-covariant real-time Green’s functions method. A comparison of our results with those of Ref. shows an overall good agreement (up to a replacement of the fermionic Green function calculated in the mean–field approximation by that obtained in the more general approximation (64)). However, it needs to keep in mind that the standard Walecka model is too primitive for a proper description of real properties of nuclear matter at intermediate energies. Already at an equilibrium state, it leads to an equation of state stiffer than the expected one and, at moderately high density and temperature, the effective masses of nucleons become very small or even negative . At present, two directions of the Walecka model improvement exist: either an increase of the number of its constituents or a transition to a nonlinear generalization of the model . The first way leads to technical complications of the kinetic theory (by virtue of an essential increase of the number of constituents for real system), the second one remains still unexplored for a description of non-equilibrium states.
## 5 Connection with the non-equilibrium statistical operator method
As previously outlined, in the definition of Wigner’s (1) and Green’s functions (17), the statistical averaging was performed by means of the equilibrium density matrix, $`\rho ,`$ in the Heisenberg picture corresponding ordinarily to a system state at an infinite past. This fact limits the applicability of the theory to the case of weakly non-equilibrium states which can be described defining only slowly varying thermodynamic functions of the system such as temperature, chemical potential, etc. Now, let us discuss briefly a possibility of a generalization of the formalism suggested in Sec. 3 which is based on the change of the statistical averaging procedure, using the non-equilibrium statistical operator, $`\rho (t)`$. It is expected that such a generalization will allow us to describe strongly non-equilibrium states of a system . The non-equilibrium statistical operator method was adapted to fulfill requirements of the relativistic kinetic theory in our previous works . Here, we discuss only a problem where this method is combined with the covariant formalism of real-time Green’s functions method (Sec. 3).
Let us write the motion equation for the non-equilibrium statistical operator in a differential form
$$d\rho (\tau )/d\tau =ϵ\{\rho (\tau )\rho _q(\tau )\}\text{ ,}$$
(71)
and, in an integral form ($`ϵ`$ is an infinitely small value, $`ϵ>0`$, and will approach to zero after the execution of the thermodynamic limiting transition),
$`\rho (\tau )=\rho _q(\tau )+i{\displaystyle \underset{\mathrm{}}{\overset{\tau }{}}}𝑑\tau ^{}e^{ϵ(\tau ^{}\tau )}\{[\rho _q(\tau ^{}),H(\tau ^{})]+id\rho _q(\tau ^{})/d\tau ^{}\}\text{ .}`$ (72)
Here, the time, $`\tau ,`$ is defined by Eq. (6) with a subsequent introduction of Assumptiones 1 and 2 and relation (12). An infinitely small source in the right-hand side of Eq. (71) is introduced to break the symmetry of this equation with respect to the time reflection (in the Wigner sense). The form of Eq. (71) allows us to select the retarded solution of the Liouville equation on the basis of an analogy with the formal scattering theory. This implies that after completing our calculations the thermodynamical limit should be taken and, then, the limit $`ϵ0`$ assumed.
The quasi-equilibrium statistical operator, $`\rho _q(\tau ),`$ in Eq. (71) is the asymptotic form of the non-equilibrium statistical operator, $`\rho (\tau ),`$ at the kinetic stage of evolution, when $`\tau \mathrm{}`$. An explicit form of operator $`\rho _q(\tau )`$ can be derived from the principle of the maximum of the information entropy under a supplementary condition of the given averaging value of $`<P_{\alpha \beta }(xy)>_{q\tau }`$, where $`<\mathrm{}>_{q\tau }=Tr\mathrm{}\rho _q(\tau )`$ and the operator, $`P_{\alpha \beta }(xy),`$ is defined by Eq. (2). For our task, it is convenient to associate the proper time, $`\tau ,`$ with the slow variable $`x^\mu `$ according to the expression (6).
As a result, we have
$$\rho _q(\tau )=exp\{S(\tau )\}\text{ ,}$$
(73)
where $`S(\tau )`$ is the entropy operator of the system at the kinetic stage,
$$S(\tau )=\mathrm{\Phi }(\tau )+𝑑\sigma (x|u)d^4yP_{\alpha \beta }(x,y)F_{\alpha \beta }(x,y),$$
(74)
and $`\mathrm{\Phi }(\tau )`$ is the normalizing functional defined by condition $`Tr\rho _q(\tau )=1`$. The Lagrange factor, $`F_{\alpha \beta }(x,y),`$ is determined by the self-consistency condition
$$<P_{\alpha \beta }(xy)>_\tau =<P_{\alpha \beta }(xy)>_{q\tau }\text{ .}$$
(75)
In the left-hand side of this equality, the averaging is performed by means of the non-equilibrium statistical operator $`\rho (\tau )`$, i.e., $`<\mathrm{}>_\tau =Tr\mathrm{}\rho (\tau )`$. Relation (75) shows a full equivalence of both statistical operators, $`\rho (\tau )`$ and $`\rho _q(\tau ),`$ at the kinetic stage of evolution. Hence, both variants of the Wigner function,
$$f_{\alpha \beta }(x,p)=(2\pi )^4𝑑ye^{ipy}<P_{\alpha \beta }(x,y)>_\tau =(2\pi )^4𝑑ye^{ipy}<P_{\alpha \beta }(x,y)>_{q\tau }\text{ ,}$$
(76)
can be used for a kinetic description of both quasi-equilibrium and non-equilibrium states.
The Green’s functions, (16)-(19) and (36), also permit an analogous generalization for the case of strong non-equilibrium states.
$$G_{(noneq)}(xp)=i(2\pi )^4𝑑ye^{ipy}<T_c[\psi (x_+)\overline{\psi }(x_{})]>_\tau .$$
(77)
The form of the Kadanoff-Baym Eqs. (33), (34) remains also the same, taking into account that the $`ϵ`$-terms generated by the right-hand of Eq. (72) vanish after a realization of the thermodynamical limit. However, a restriction to the lowest terms of the gradient decomposition is incompatible with a description of strong non-equilibrium states. Therefore, the corresponding results of Sec. 4 require a further generalization. This field of relativistic kinetic theory is still scarcely investigated.
Let us remark that Eq. (72) can be rewritten in a form convenient for the construction of an iteration scheme under the energy of particle-particle interactions. In full analogy with the non-relativistic case , we obtain
$`\rho (\tau )=\rho _q(\tau )+i{\displaystyle \underset{\mathrm{}}{\overset{\tau }{}}}d\tau ^{}e^{ϵ(\tau ^{}\tau )}\{[\rho _q(\tau ^{}),H_{in}(\tau ^{})]+`$
$`+{\displaystyle }d\sigma (x^{}|u){\displaystyle }d^4y^{}<[P_{\alpha \beta }(x^{},y^{}),H_{in}(\tau ^{})]>_{q\tau ^{}}{\displaystyle \frac{\delta \rho _q(\tau ^{})}{\delta <P_{\alpha \beta }(x^{},y^{})>_\tau ^{}}}\}\text{ .}`$ (78)
As shows our experience, the second addend in the right-hand side of Eq. (5) removes usually all disconnected diagrams of a given order from the set of diagrams which are supplied by the first addend.
Some problems of this approach arise on the level of the non-equilibrium thermodynamical Wick’s theorem. However, there is a set of methods which permits to overcome these difficulties . c
## 6 Conclusions
In the present work, we have demonstrated an interesting possibility for a covariant generalization of the real-time Green’s functions method which opens new ways after a transition to the Wigner representation. In fact, this generalization is based on the admissibility of introduction in every point of the Minkowski space a unit time-like vector constructed with the momentum vector of the corresponding point of the phase space. Therefore, for a definition of the preferred time-like direction in the Minkowski space, we used a possibility contained in the Wigner function itself. The relevant Assumptiones (Sec. 2) allow us to write, starting from the Heisenberg representation, a motion equation in a covariant form for the one-particle Wigner function (Sec. 2) and, after that, to obtain a generalized RKE in the frame of the Kadanoff-Baym technique (Sec. 3). Such a way is found especially effective for a kinetic description of subsystems with inner degrees of freedom and allows us to eliminate some ambiguities of relativistic kinetic theory based on the Kadanoff-Baym equations.
As a test of the suggested approach, the Fermi sector of the well known Walecka model of the relativistic nuclear matter was chosen (Sec. 3 and 4). For this model, RKE of the Vlasov type (Sec. 3) and of the BUU type (Sec. 4) were obtained. In the latter case, the collision integral was derived only for a spin saturated system without antinucleon states. The calculations were performed within the framework of the quasiparticle approximation (QPA). We discuss also a possibility of an extension of this approximation. An agreement with a set of well-known simple results in the literature denotes a noncontradictory character of the proposed generalizations. It is safe to assume that methods suggested here provide also a correct description of spin degrees of freedom, states with negative energies (e.g., for the annihilation channel in the theory of non-equilibrium electron-positron plasma), and so on.
Finally, the covariant modification of the real-time Green’s functions method suggested in Sec. 3 was generalized to a kinetic description of strongly non-equilibrium states (Sec. 5). For this aim, we suggest to use the non-equilibrium statistical operator method.
The non-equilibrium relativistic nuclear matter represents a good example of very complicated object for theoretical and experimental investigations. Of course, the simplest toy model and our present results discussed in Sec. 4 are only a rough approximation of the reality. At present, a lot of actual problems still remain in the relativistic kinetic theory, e.g., going beyond the quasiparticle approximation, taking into account nonlocal (and, in particular, non-Markovian) effects in kinetic processes, cluster decomposition, etc. (see, e.g., ). An important feature of the examined approach is a lack of some dynamical constraints into a corresponding quantum field theory (these constraints are introduced in order to eliminate ”unnecessary”, non-physical degrees of freedom). The dynamical introduction of these constraints into the Kadanoff-Baym formalism is a rather non-trivial problem. A necessity to consider such constraints creates serious problems for the kinetic theory (e.g., in order to describe the $`\mathrm{\Delta }`$-isobar subsystem of nuclear matter). From our point of view, the approach suggested in the present work can serve as a reliable basis for such kind of further researches in the relativistic kinetic theory.
## Acknowledgments
We wish to thank V. G. Morozov for stimulating discussions and S. Mrówczyński for addressing our attention to some questions discussed here.
This work was partly supported by the State Committee of Russian Federation for Higher Education under grant N 97-0-6.1-4 and was completed under the auspices of the U.S. Department of Energy by the Los Alamos National Laboratory under contract no. W-7405-ENG-36.
One from authors (S.A.S.) thanks the Soros Education Program for support.
## Appendices
## Appendix A Proof of Eq. (25)
Let us derive a proof of the Eq. (25) using a boost transformation.
$$d^3z^{}S(𝐳𝐳^{},0)f(z^{})=i\gamma ^0f(z)\text{ ,}$$
$`(A1)`$
for an arbitrary function of the field operator $`f(z)`$. The integration here is fulfilled over hyperplane of the constant time, $`z^0=z_{}^{}{}_{}{}^{0}`$. Let us perform a transition to an arbitrary frame of reference using a transformation from the homogeneous Lorentz group
$$x^\mu =\mathrm{\Lambda }_\nu ^\mu (u)z^\nu ,x_{}^{}{}_{}{}^{\mu }=\mathrm{\Lambda }_\nu ^\mu (u)z_{}^{}{}_{}{}^{\nu }.$$
$`(A2)`$
The velocity vector, $`u^\mu ,`$ defines simultaneously the orientation of the hyperplane, $`\sigma (u)`$. The transformation of the commutator function, $`S(𝐳𝐳^{},0),`$ in Eq. (A1) is obtained using the corresponding unitary operator, $`U(\mathrm{\Lambda })`$ ,
$$S(𝐳𝐳^{},0)=S[\mathrm{\Lambda }(xx^{})]=U(\mathrm{\Lambda })S(xx^{})U^1(\mathrm{\Lambda })\text{ ,}$$
$`(A3)`$
where $`x,x^{}\sigma (u)`$. Similarly, we have:
$$f(z)=U(\mathrm{\Lambda })f(x)U^1(\mathrm{\Lambda })\text{ .}$$
$`(A4)`$
Let us take into account also the relation
$$U^1(\mathrm{\Lambda })\gamma ^\mu U(\mathrm{\Lambda })=\mathrm{\Lambda }_\nu ^\mu (u)\gamma ^\nu \text{ .}$$
$`(A5)`$
This equality is a consequence of the invariance under the Lorentz transformation (A2) of the scalar product of $`\psi (x)`$ and $`\varphi (x)`$ ,
$$(\psi ,\varphi )_\sigma =𝑑\sigma ^\mu (x|u)\overline{\psi }(x)\gamma _\mu \varphi (x)\text{ .}$$
Finally, it is necessary to consider the equalities $`\mathrm{\Lambda }_{\mathrm{\hspace{0.33em}0}}^0(u)=u^0`$ and $`\mathrm{\Lambda }_k^0(u)=u^k(k=1,2,3)`$. Formulas (A2)-(A3) allow us to realize the transformation from Eq. (A1) to Eq. (25). Another proof of Eq. (25) can be found in Ref. .
## Appendix B Spectral properties of the modified Green’s functions
The spectral representations of the covariant contour Green’s functions introduced in Sec. 3 can be also considered in a covariant form (we follow the work were a similar approach was fulfilled in the framework of a non-covariant formalism). The main feature of these representations is the selection of the time arguments in the form (15). As an example, let us consider the spectral decomposition of the retarded and advanced Green’s functions. According to the definitions, (16) and (36), we have (see the last remark in Sec. 3.2)
$$G^{R/A}(xp)=\pm (2\pi )^4𝑑ye^{ipy}\{G^>(x_+,x_{})G^<(x_+,x_{})\}\theta [\pm y_\mu u^\mu ]\text{ ,}$$
where $`u^\mu `$ is the unit vector (12). From here, it follows
$$G^{R/A}(xp)=\pm \frac{1}{2\pi i}d^4p^{}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\eta \delta ^{(4)}(p^\mu p^\mu +u^\mu \eta )\frac{G^>(xp^{})G^<(xp^{})}{\pm \eta i\epsilon }\text{ ,}$$
or
$$G^{R/A}(xp)=\pm \frac{1}{2}[G^>(xp)G^<(xp)]+\stackrel{}{G}(xp)\text{ ,}$$
$`(B1)`$
where
$$\stackrel{}{G}(xp)=\frac{1}{2\pi i}d^4p^{}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\eta \delta ^{(4)}(p^\mu p^\mu +u^\mu \eta )[G^>(xp^{})G^<(xp^{})]𝒫\left(\frac{1}{\eta }\right)\text{ .}$$
$`(B2)`$
Let us transform the last relation using the decomposition of an arbitrary vector, $`a^\mu ,`$ on the time-line direction, $`u_\mu `$ (12), and orthogonal to it space-like component,
$$a^\mu =u^\mu (au)+\mathrm{\Delta }_\nu ^\mu a^\nu a_{}^\mu +a_{}^\mu \text{ ,}$$
where the projection operator, $`\mathrm{\Delta }^{\mu \nu },`$ is defined by Eq. (9). Then, we can write the following representation
$$\delta ^{(4)}(p^\mu p^\mu +u^\mu \eta )=\delta [n_\mu (p^\mu p^\mu +u^\mu \eta )]\delta _\sigma (p_{}^\mu p_{}^\mu )\text{ ,}$$
$`(B3)`$
where function $`\delta _\sigma (p_{}^\mu )`$ is defined on the space-like hyperplane $`\sigma (u)`$ . Using Eq. (B3), we can rewrite Eq. (B2) in the final form:
$$\stackrel{}{G}(xp)=\frac{1}{2\pi i}\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\omega 𝒫\left(\frac{1}{\omega \sqrt{p^2}}\right)\{G^>(x;p_{},\omega )G^<(x;p_{},\omega )\}\text{ .}$$
$`(B4)`$
The denominator, $`(\omega \sqrt{p^2}),`$ in Eq. (B4) shows clearly that function $`\stackrel{}{G}(xp)`$ is just the off-mass-shell part of the Green’s functions, $`G^{R/A}(xp)`$. As far as we have
$$G^R(xp)+G^A(xp)=2\stackrel{}{G}(xp),$$
and an analogous relation for the mass operator, $`\mathrm{\Sigma }^{R/A}`$, the last term in Eq. (49) represents indeed the off-mass-shell part of the collision integral (48).
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# Approximate equation relevant to axial oscillations on slowly rotating relativistic stars
## I Introduction
The surprising discovery of the r-mode instability in rotating stars has inspired the study of the axial oscillations . The physical mechanism of the instability is the same as that in the polar modes, so-called the radiation reaction instability found by Chandrasekhar , and Friedman and Schutz . The r-mode oscillations seem to be more important since they are unstable on an inviscid rotating fluid even for small angular velocity. This instability can explain the spin-down process of newly born neutron stars which rotate with nearly Kepler frequency. The gravitational waves associated with the unstable r-mode oscillations may be promising detectable sources on the ground based laser interferometric detectors. It was proposed that the unstable mode might also play a key role on the spin of the accreting white dwarfs . However, Lindblom showed that the possibility is unlikely realized. It was also proposed that the r-mode instability could provide the loss of angular momentum from the accretion disk by the gravitational waves to halt the spin-up in the low mass X-ray binaries . The proposed conclusion crucially depends on the poorly understood dissipation mechanism . In this way, the r-mode oscillations enrich the astrophysical implications. Recent review of this subject is given by Friedman and Lockitch .
Most of the studies are however limited to idealized situation. Some effects are added to the simplified models to examine the validity. For example, Rezzolla, Lamb and Shapiro suggested that the magnetic field of a neutron star is wound up during the non-linear growth of the unstable mode, and that the energy is not transfered to the gravitational radiation so much. However, their calculation is not self-consistent non-linear one, so that the magnetic effect on the r-mode is not conclusive at moment. The relativistic effects are also important for the oscillations in neutron stars. The relativistic factor is of order 0.2, so that the frequency could slightly shift from the Newtonian calculation. More important effects of the general relativity are gravitational wave and frame dragging. Each of them leads to qualitatively a different result. Some authors already calculated the r-modes in general relativity. The frame dragging effect is regarded as a kind of differential rotation . The rotational effects are however limited to the lowest order. Mathematically, the treatment is insufficient, since the modes are degenerate at the order. It is necessary to include the higher order rotational corrections. Such a task will be significantly complicated when the rotation and relativistic effects are simultaneously considered.
One simplification for the non-rotating case is realized as the decomposition of the spherical harmonic function. Each oscillation mode can be specified for each index $`l,m.`$ Furthermore, the polar and axial modes are definitely determined by an appropriate combination of the harmonic functions. However, the functions should be mixed in the presence of the rotation. Considering the slow rotation, the coupling is weak, so that the entangled range can be restricted. It is not known a priori to specify the axial oscillations on the rotating stars by a few spherical harmonic functions. Indeed, Lockitch and Friedman calculated the normal mode by the sum of infinite number of the spherical harmonics indices. We will consider a different approach in this paper. Our treatment is suitable for the initial-value problem. We do not consider a single Fourier mode $`e^{i\sigma t}`$ with respect to time. Suppose that the initial disturbance at $`t=0`$ is produced with a certain symmetry, which can be assumed to be specified by a few number of spherical harmonic functions. What happens in the subsequent evolution? If the oscillation possesses the symmetry, the oscillation preserves the symmetry. Otherwise, new patterns with different spherical harmonics indices are induced. It is clear that the truncated approximation to the finite number of the spherical harmonics becomes worse for large $`t`$ in general. Therefore, our method is constructing the pulsation equation adequate for small $`t.`$ It is difficult to address the valid domain of $`t`$ beforehand. For example, the r-mode oscillation can be well described by a few number of the spherical harmonics indices for the oscillation in the uniform density with the Cowling approximation .
Present authors applied the method to the axial oscillation in a rotating relativistic star with the Cowling approximation , where the metric perturbations were neglected. They took account of the rotational correction up to third order to examine the oscillation equation. They pointed out the importance of the frame-dragging effect, which causes different property unlike the Newtonian case. Thus, the perturbative approach proves useful in providing a physical understanding of many processes. In this paper, we extend the approximation scheme to include the metric perturbations. The remainder of the paper is organized as follows. In Sec. II, we formulate the perturbation scheme to solve the linearized Einstein equations. We employ the slow rotation approximation, i.e., angular velocity is assumed to be small expansion parameter. We look for the solution in which axial-led functions are dominant, where the axial-led functions mean the functions describing axial modes in the absence of the rotation. We use the terminology polar-led in the same way. In our scheme, the lowest order equations are determined only by the axial-led functions. They were already derived elsewhere but are reviewed in Sec. III. In Sec. IV, first-order corrections to the polar-led functions are shown. In Sec. V, second-order correction terms are added to the same equations in the form as the lowest ones. In Sec. VI, concluding remarks are given. Throughout this paper, we work in the geometrical units of $`c=G=1`$.
## II Perturbation scheme
We consider a rotating star with uniform angular velocity $`\mathrm{\Omega }O(\epsilon ),`$ where $`\epsilon `$ is a small rotational parameter. The metric and fluid quantities describing the equilibrium state can be calculated by the slow rotation formalism . They are summarized in the Appendix A. We next investigate the perturbations from the state. The metric perturbations can be described by six functions. Working in the Regge-Wheeler gauge , the perturbations are expressed as
$$h_{\mu \nu }=\underset{l,m}{}\left(\begin{array}{cccc}e^\nu H_{0lm}(t,r)Y_{lm}& H_{1lm}(t,r)Y_{lm}& h_{0lm}(t,r)\frac{_\varphi Y_{lm}}{\mathrm{sin}\theta }& h_{0lm}(t,r)\mathrm{sin}\theta _\theta Y_{lm}\\ \text{sym}& e^\lambda H_{2lm}(t,r)Y_{lm}& h_{1lm}(t,r)\frac{_\varphi Y_{lm}}{\mathrm{sin}\theta }& h_{1lm}(t,r)\mathrm{sin}\theta _\theta Y_{lm}\\ \text{sym}& \text{sym}& r^2K_{lm}(t,r)Y_{lm}& 0\\ \text{sym}& \text{sym}& \text{sym}& r^2\mathrm{sin}^2\theta K_{lm}(t,r)Y_{lm}\end{array}\right),$$
(1)
where $`Y_{lm}=Y_{lm}(\theta ,\varphi )`$ represents spherical harmonics. The symbol “sym” indicates that the missing components of $`h_{\mu \nu }`$ are to be found from the symmetry $`h_{\mu \nu }=h_{\nu \mu }`$. The angular part is expanded with an appropriate combination of the harmonics. In a same way, the fluid perturbations are described by five functions. They are the pressure perturbation $`\delta p`$, density perturbation $`\delta \rho `$ and three components of the 4-velocity ($`\delta u_r,\delta u_\theta ,\delta u_\varphi `$). The component $`\delta u_t`$ can be determined by the condition, $`\delta u_\mu u^\mu =0`$. These perturbed quantities are also expanded as
$`\delta p`$ $`=`$ $`{\displaystyle \underset{l,m}{}}\delta p_{lm}(t,r)Y_{lm},`$ (2)
$`\delta \rho `$ $`=`$ $`{\displaystyle \underset{l,m}{}}\delta \rho _{lm}(t,r)Y_{lm},`$ (3)
$`(\rho +p)\delta u_r`$ $`=`$ $`e^{\nu /2}{\displaystyle \underset{l,m}{}}R_{lm}(t,r)Y_{lm},`$ (4)
$`(\rho +p)\delta u_\theta `$ $`=`$ $`e^{\nu /2}{\displaystyle \underset{l,m}{}}\left[V_{lm}(t,r)_\theta Y_{lm}{\displaystyle \frac{U_{lm}(t,r)}{\mathrm{sin}\theta }}_\varphi Y_{lm}\right],`$ (5)
$`(\rho +p)\delta u_\varphi `$ $`=`$ $`e^{\nu /2}{\displaystyle \underset{l,m}{}}\left[V_{lm}(t,r)_\varphi Y_{lm}+U_{lm}(t,r)\mathrm{sin}\theta _\theta Y_{lm}\right].`$ (6)
These eleven functions are determined by ten components of the linearized Einstein field equations,
$$\delta G_{\mu \nu }=8\pi \delta T_{\mu \nu },$$
(7)
and the adiabatic condition for the perturbations,
$$\delta p+\xi p=\frac{\mathrm{\Gamma }p}{p+\rho }(\delta \rho +\xi \rho ),$$
(8)
where $`\mathrm{\Gamma }`$ is the adiabatic index and $`\xi `$ is the Lagrange displacement.
Now we will solve the pulsation equations by the expansion of the spherical harmonics. In the spherically symmetric star, the equations are decoupled for each harmonic index $`(l,m).`$ The perturbations can also be decoupled into the axial and polar perturbations. They are respectively described by the axial functions $`𝒜_{lm}`$ $`(U_{lm},h_{0,lm},h_{1,lm}),`$ and the polar functions $`𝒫_{lm}`$ $`(\delta p_{lm},\delta \rho _{lm},R_{lm},V_{lm},`$ $`H_{0,lm},H_{1,lm},H_{2,lm},K_{lm}).`$ In the presence of rotation, the perturbations are described by the mixed state of them. From now on, we call these functions as the axial-led ones for $`𝒜_{lm}`$ and the polar-led ones for $`𝒫_{lm}.`$ Since the slow rotation is associated with the perturbation with $`l=1,`$ the formal relation between $`𝒜_{lm}`$ and $`𝒫_{lm}`$ in the Eqs.(7)-(8) can schematically be expressed as
$`0=[𝒜_{lm}]+\times [𝒫_{l\pm 1m}]+^2\times [𝒜_{lm},𝒜_{l\pm 2m}]+\mathrm{},`$ (9)
$`0=[𝒫_{lm}]+\times [𝒜_{l\pm 1m}]+^2\times [𝒫_{lm},𝒫_{l\pm 2m}]+\mathrm{},`$ (10)
where the symbol $``$ denotes some functions of order $`\epsilon ,`$ and the square bracket formally represents the relation among the perturbation functions therein. These selection rules follow from the addition of angular momenta. We moreover assume that the axial-led functions are dominant in the slowly rotating star, i.e., $`𝒜_{lm}𝒫_{lm}.`$ This assumption is not valid for the some cases. If the star and its perturbations obey the same one-parameter equation of state, r-modes and g-modes are coupled at the lowest order in general. They are respectively described by $`𝒫_{lm}`$ and $`𝒜_{lm},`$ which have zero frequency in the spherical isentropic star. They form hybrid modes in the rotating star . However, the r-modes are discriminated from the g-modes for non-isentropic stars, since the g-mode frequencies are non-zero in the spherical limit. In that case, we may use the assumption $`𝒜_{lm}𝒫_{lm}`$ to solve the r-mode oscillations, and expand as
$$𝒜_{lm}=𝒜_{lm}^{(1)}+\epsilon ^2𝒜_{lm}^{(2)}+\mathrm{},𝒫_{lm}=\epsilon (𝒫_{lm}^{(1)}+\epsilon ^2𝒫_{lm}^{(2)}+\mathrm{}).$$
(11)
Substituting these functions into Eqs.(9)-(10), and comparing each order of $`\epsilon `$, we have the following equations of $`\epsilon ^n(n=0,1,2)`$,
$`0`$ $`=`$ $`[𝒜_{lm}^{(1)}],`$ (12)
$`0`$ $`=`$ $`[\epsilon 𝒫_{l\pm 1m}^{(1)}+\times 𝒜_{lm}^{(1)}],`$ (13)
$`0`$ $`=`$ $`[\epsilon ^2𝒜_{lm}^{(2)}]+\times [\epsilon 𝒫_{l\pm 1m}^{(1)}]+^2\times [𝒜_{lm}^{(1)},𝒜_{l\pm 2m}^{(1)}]`$ (14)
$`=`$ $`[\epsilon ^2𝒜_{lm}^{(2)}+^2\times 𝒜_{lm}^{(1)}].`$ (15)
We have here assumed that the perturbation in the lowest order is described by a single component of spherical harmonic, that is, $`𝒜_{l^{}m}^{(1)}=0`$, for $`l^{}l`$. The polar-led functions in Eq.(14) are eliminated by Eq.(13). Equation (12) represents the axial oscillation in the lowest order. Equation (15) is the second-order form of it, and the term $`^2\times 𝒜_{lm}^{(1)}`$ can be regarded as the rotational corrections. The method to solve the equations is straightforward. The first-order equations are solved by the axial-led functions. The polar-led functions are expressed using them. We have the second-order equations with the corrections expressed by the axial-led functions in the lowest order. These equations are successively solved in the following sections. In the actual calculations, we also assume that the time variation of the oscillation is slow and proportional to $`\mathrm{\Omega },`$ i.e., $`_t\mathrm{\Omega }O(\epsilon ).`$ This is true in the r-mode frequency, $`_t(12/l(l+1))m\mathrm{\Omega }.`$
## III Lowest-order calculation
In this section, we review the equations governing the axial oscillations at the lowest order. The radial functions are decoupled for each spherical harmonic index. We denote it as $`L(l,m)`$ for the abbreviation. The relevant functions are $`h_{0,L},h_{1,L}`$ and $`U_L.`$ They are calculated from three components, essentially $`(t\varphi )`$ $`(r\varphi )`$ and $`(\theta \varphi )`$ components, of the Einstein equations. We define a function $`\mathrm{\Phi }_L^{(1)}`$ as
$$\mathrm{\Phi }_L^{(1)}=\frac{h_{0,L}^{(1)}}{r^2},$$
(16)
where the superscript <sup>(1)</sup> denotes the lowest order term in Eq.(11). The relation between the metric functions is given by
$$h_{1,L}^{(1)}=\frac{r^4e^\nu }{[l(l+1)2]}\left[(_Tim\varpi )\mathrm{\Phi }_L^{(1)}+\frac{2im\varpi ^{}}{l(l+1)}\mathrm{\Phi }_L^{(1)}\right],$$
(17)
where $`_T=_t+im\mathrm{\Omega }`$ denotes a time-derivative in a co-rotating frame and a prime denotes a derivation with respect to $`r.`$ The axial velocity function is expressed by two ways:
$`(_Tim\chi )U_L^{(1)}`$ $`=`$ $`4\pi (p_0+\rho _0)r^2e^\nu _T\mathrm{\Phi }_L^{(1)},`$ (18)
$`U_L^{(1)}`$ $`=`$ $`{\displaystyle \frac{j^2r^2}{4}}\left[{\displaystyle \frac{1}{jr^4}}(jr^4\mathrm{\Phi }_L^{(1)})^{}(v+16\pi (p_0+\rho _0)e^\lambda )\mathrm{\Phi }_L^{(1)}\right],`$ (19)
where
$`\chi `$ $`=`$ $`{\displaystyle \frac{2}{l(l+1)}}\varpi ={\displaystyle \frac{2}{l(l+1)}}(\mathrm{\Omega }\omega ),`$ (20)
$`v`$ $`=`$ $`{\displaystyle \frac{e^\lambda }{r^2}}\left[l(l+1)2\right],`$ (21)
$`j`$ $`=`$ $`e^{(\lambda +\nu )/2}.`$ (22)
Eliminating $`U_L^{(1)}`$ in Eqs.(18)-(19), we have a second-order differential equation for $`\mathrm{\Phi }_L^{(1)}`$,
$`0`$ $`=`$ $`[\mathrm{\Phi }_L^{(1)}]`$ (23)
$`=`$ $`\left(_Tim\chi \right)\left[{\displaystyle \frac{1}{jr^4}}\left(jr^4\mathrm{\Phi }_L^{(1)}\right)^{}v\mathrm{\Phi }_L^{(1)}\right]+16\pi im\chi (p_0+\rho _0)e^\lambda \mathrm{\Phi }_L^{(1)}.`$ (24)
The linear operator $``$ is defined above. This equation is reduced to a singular eigen-value problem , by assuming the time dependence for the mode as the form $`e^{i\sigma t}.`$ The equation can be written as
$$(\varpi \mu )\left[\frac{1}{jr^4}(jr^4\mathrm{\Phi }_L^{(1)})^{}v\mathrm{\Phi }_L^{(1)}\right]=q\mathrm{\Phi }_L^{(1)},$$
(25)
where
$`\mu `$ $`=`$ $`{\displaystyle \frac{l(l+1)}{2m}}(\sigma m\mathrm{\Omega }),`$ (26)
$`q`$ $`=`$ $`{\displaystyle \frac{1}{jr^4}}\left(jr^4\varpi ^{}\right)^{}=16\pi (p_0+\rho _0)e^\lambda \varpi 0.`$ (27)
There is a singular point $`r_0`$ in Eq.(25) unless $`q(r_0)=0,`$ corresponding to the real value of $`\mu =\varpi (r_0).`$ It is evident that the singularity originates from the mismatch in Eq.(24), i.e., the first term vanishes whereas the second never does. When the first term, which is formally of first order, is small enough, then higher order corrections become important. This situation is very similar to the inviscid shear flows. When the viscosity is small enough, the stability is almost determined by the Rayleigh equation, the perturbation equation for the inviscid theory. The Rayleigh equation has critical points for some mean fluid. The viscous corrections should be included to determine the behavior near the neighborhood of the critical points. Therefore, the equation should be replaced by the Orr-Sommerfeld equation derived from the Navier-Stokes equation.
As we will show in the subsequent sections, the function $`\mathrm{\Phi }_L^{(1)}`$ can also affect the polar-led functions. When one considers the equations of the next order, $`\mathrm{\Phi }_L^{(2)}`$, additional terms appear in the form (24). The terms depend on different aspect of the background flow, and play an important role as the viscosity terms.
## IV First-order corrections in polar-led functions
In this section, we will derive the equations governing the polar-led functions, $`H_0^{(1)}`$, $`H_1^{(1)}`$, $`H_2^{(1)}`$, $`K^{(1)}`$, $`\delta p^{(1)}`$, $`\delta \rho ^{(1)}`$, $`R^{(1)}`$ and $`V^{(1)}`$. As shown in Sec. II, the functions with $`(l\pm 1,m)`$ are coupled with the axial-led functions with $`(l,m).`$ We will shorten the subscript of the spherical harmonic index as e.g., $`\delta p_\pm ^{(1)}\delta p_{l\pm 1m}^{(1)}`$. These eight functions are determined by seven components of the linearized Einstein field equations and one thermodynamical relation. The calculations are straightforward, but results are sometimes messy. The pressure $`\delta p_\pm ^{(1)}`$ and density perturbations $`\delta \rho _\pm ^{(1)}`$ are expressed by
$`4\pi \delta p_\pm ^{(1)}`$ $`=`$ $`2\pi (p_0+\rho _0)(H_{0,\pm }^{(1)}+2T_\pm \mathrm{\Omega }r^2e^\nu \mathrm{\Phi }_L^{(1)})+2S_\pm \varpi U_L^{(1)},`$ (28)
$`4\pi \delta \rho _\pm ^{(1)}`$ $`=`$ $`4\pi {\displaystyle \frac{\rho _0^{}}{\nu ^{}}}(H_{0,\pm }^{(1)}+2T_\pm \mathrm{\Omega }r^2e^\nu \mathrm{\Phi }_L^{(1)})`$ (30)
$`4S_\pm {\displaystyle \frac{e^{\nu /2}}{\nu ^{}}}\left(e^{\nu /2}\varpi U_L^{(1)}\right)^{}+2T_\pm {\displaystyle \frac{e^\nu }{\nu ^{}r^2}}(\varpi r^2e^\nu )^{}U_L^{(1)},`$
where
$$S_+=\frac{l}{l+1}Q_+,S_{}=\frac{l+1}{l}Q_{},T_+=lQ_+,T_{}=(l+1)Q_{},$$
(31)
$$Q_+=\sqrt{\frac{(l+1)^2m^2}{(2l+1)(2l+3)}},Q_{}=\sqrt{\frac{l^2m^2}{(2l1)(2l+1)}}.$$
(32)
These quantities are expressed by the axial-led functions ($`\mathrm{\Phi }_L^{(1)},U_L^{(1)})`$ and $`H_{0,\pm }^{(1)}.`$ There is another relation among $`H_{0\pm }^{(1)},`$ $`K_\pm ^{(1)}`$ and $`\delta \rho _\pm ^{(1)}`$ in the field equations. Eliminating $`\delta \rho _\pm ^{(1)}`$ by Eq.(30), we have a second-order differential equation for the metric perturbations. The equation can be regarded as relativistic version of the Poisson equation, $`^2\delta \varphi =4\pi \delta \rho ,`$ for the gravitational potential $`\delta \varphi `$ and the density perturbation $`\delta \rho .`$ In the relativistic case, the Newtonian potential is replaced by $`H_0^{(1)}`$ or $`K^{(1)}.`$ The second-order differential equation is explicitly given by
$`K_\pm ^{(1)}{\displaystyle \frac{e^\lambda }{\nu ^{}r^2}}\{n_\pm (K_\pm ^{(1)}H_{0,\pm }^{(1)})+2(4\pi (p_0+\rho _0)r^2+e^\lambda 1)H_{0,\pm }^{(1)}\}=s_1,`$ (33)
$`(K_\pm ^{(1)}H_{0,\pm }^{(1)})^{}\nu ^{}H_{0,\pm }^{(1)}=\left(1+{\displaystyle \frac{\nu ^{}r}{2}}\right)s_1+s_2,`$ (34)
where
$`n_\pm `$ $`=`$ $`2+l^{}(l^{}+1)|_{l^{}=l\pm 1},`$ (35)
$`s_1`$ $`=`$ $`Q_\pm \left[8\varpi rU_L^{(1)}+32\pi (p_0+\rho _0)\varpi r^3e^\nu \mathrm{\Phi }_L^{(1)}+2\varpi ^{}j^2r^3\mathrm{\Phi }_L^{(1)}\right]`$ (38)
$`+S_\pm \left[8\left\{r+{\displaystyle \frac{e^\lambda }{\nu ^{}}}\right\}\varpi U_L^{(1)}{\displaystyle \frac{4e^\nu }{\nu ^{}}}\left\{l(l+1)\omega e^\lambda +16\pi p_0^{}\varpi r^3\varpi ^{}r\right\}\mathrm{\Phi }_L^{(1)}+2\varpi ^{}j^2r^3\mathrm{\Phi }_L^{(1)}\right]`$
$`+T_\pm \left[{\displaystyle \frac{2e^\nu }{\nu ^{}}}\left\{(22e^\lambda l(l+1)e^\lambda )\omega +8\pi \mathrm{\Omega }(p_0+\rho _0)r^2e^\lambda \varpi ^{}r\right\}\mathrm{\Phi }_L^{(1)}{\displaystyle \frac{\varpi ^{}r^2e^\nu }{\nu ^{}}}\mathrm{\Phi }_L^{(1)^{}}\right],`$
$`s_2`$ $`=`$ $`(r+{\displaystyle \frac{2}{\nu ^{}}})[S_\pm [4\varpi e^\lambda U_L^{(1)}2\{l(l+1)\omega e^\lambda 2\varpi ^{}r\}e^\nu \mathrm{\Phi }_L^{(1)}]`$ (42)
$`+T_\pm \{{\displaystyle \frac{1}{2}}\varpi ^{}r^2e^\nu \mathrm{\Phi }_L^{(1)}+[\{(l(l+1)+2)\omega 8\pi (p_0+\rho _0)\mathrm{\Omega }r^2\}e^{\lambda \nu }+2\omega e^\nu ]\mathrm{\Phi }_L^{(1)}\}]`$
$`+\left[4S_\pm {\displaystyle \frac{\varpi ^{}re^\nu }{\nu ^{}}}\mathrm{\Phi }_L^{(1)}+T_\pm \left(2\omega r^2e^\nu \mathrm{\Phi }_L^{(1)}+{\displaystyle \frac{(2\varpi ^{}r8\omega )e^\nu }{\nu ^{}}}\mathrm{\Phi }_L^{(1)}\right)\right].`$
When the axial-led function at the lowest order $`\mathrm{\Phi }_L^{(1)}`$ is given, the functions $`H_{0,\pm }^{(1)}`$ and $`K_\pm ^{(1)}`$ are solved with appropriate boundary conditions. In a similar way, we can solve the other polar-led functions, $`H_{1,\pm }^{(1)}`$, $`H_{2,\pm }^{(1)}`$, $`R_\pm ^{(1)}`$ and $`V_\pm ^{(1)}`$ by ($`\mathrm{\Phi }_L^{(1)},U_L^{(1)})`$ and $`(H_{0,\pm }^{(1)},K_\pm ^{(1)}).`$ The expressions for these four polar-led functions are omitted here, since they are eliminated in the following calculations, and never appear in the final results. However, here is a comment on using the adiabatic condition Eq.(8). The time derivative of it can be written as
$`4\pi _T\left(\delta p_\pm ^{(1)}{\displaystyle \frac{\mathrm{\Gamma }p_0}{p_0+\rho _0}}\delta \rho _\pm ^{(1)}\right)`$ $`=`$ $`{\displaystyle \frac{A\mathrm{\Gamma }p_0e^\nu }{p_0+\rho _0}}\left(e^\lambda R_\pm ^{(1)}{\displaystyle \frac{3im\xi _2}{r^2}}Q_\pm U_L^{(1)}\right),`$ (43)
where
$`\xi _2`$ $`=`$ $`{\displaystyle \frac{2}{\nu ^{}}}\left(h_2+{\displaystyle \frac{1}{3}}\varpi ^2r^2e^\nu \right),`$ (44)
$`A`$ $`=`$ $`{\displaystyle \frac{\rho _0^{}}{p_0+\rho _0}}{\displaystyle \frac{p_0^{}}{\mathrm{\Gamma }p_0}}.`$ (45)
This thermodynamical relation determines the function $`R_\pm ^{(1)}`$ unless the Schwarzschild discriminant $`A`$ vanishes. Otherwise, we have one constraint for the function $`U_L^{(1)}`$ through Eqs.(28)-(30), and the function $`R_\pm ^{(1)}`$ should be specified in another way. (For example, see the method by in the Cowling approximation.) The mathematical drawback for the isentropic case $`A=0`$ is related with the coupling of the g-modes. Both r-modes and g-modes are degenerate to zero frequency in the non-rotating star, and hence a particular treatment is necessary . From now on, we will consider the case $`A0`$ only.
## V Including second-order corrections
So far we have considered ten components of the Einstein equations and one thermodynamical relation. They were limited to the lowest-order form with respect to the rotational parameter. In this section, we consider how the next order terms modify the equation (24) derived in Sec. III. The relevant equations for this purpose are three components of the Einstein equations, i.e., $`(t\varphi )`$ $`(r\varphi )`$ and $`(\theta \varphi )`$ components, which are used in the leading order equation in Sec. III. These equations contain the relations among the axial-led functions of second-order $`h_{0,L}^{(2)}`$, $`h_{1,L}^{(2)}`$, $`U_L^{(2)}`$ and the polar-led functions calculated in the previous section. We follow the same procedure as done in Sec. III. Defining the function $`\mathrm{\Phi }_L^{(2)}=h_{0,L}^{(2)}/r^2`$ and eliminating $`h_{1,L}^{(2)}`$ and $`U_L^{(2)},`$ we eventually have the equation governing the axial oscillations. It can be written in the following form
$`[\mathrm{\Phi }_L^{(2)}]=𝒟[\mathrm{\Phi }_L^{(1)},h_{1,L}^{(1)},U_L^{(1)}]+𝒢[H_{0,\pm }^{(1)},K_\pm ^{(1)}],`$ (46)
where the operator $``$ is defined in Eq.(24) and the right hand side means the second-order rotational corrections. They consist of several terms as
$`𝒟[\mathrm{\Phi }_L^{(1)},h_{1,L}^{(1)},U_L^{(1)}]`$ $`=`$ $`𝒟_0+\alpha _1h_{1,L}^{(1)}+\alpha _2_t\mathrm{\Phi }_L^{(1)}+im(\beta _1U_L^{(1)}+\beta _2\mathrm{\Phi }_L^{(1)}+\beta _3\mathrm{\Phi }_L^{(1)}),`$ (47)
$`𝒢[H_{0,\pm }^{(1)},K_\pm ^{(1)}]`$ $`=`$ $`(A_1_T+imA_2)K_\pm ^{(1)}+(B_1_T+imB_2)H_{0,\pm }^{(1)}.`$ (48)
The explicit forms of the coefficients $`\alpha _i`$, $`\beta _i`$, $`A_i`$ and $`B_i`$ are given in Appendix B. They are expressed by the quantities determined by the stellar model in the equilibrium. Since the term $`𝒟_0`$ contains higher order derivatives, we explicitly show below:
$`𝒟_0`$ $`=`$ $`32c_3{\displaystyle \frac{\varpi e^\nu }{jr^2}}\left\{{\displaystyle \frac{(\rho _0+p_0)r^2}{jA\nu ^{}}}\left({\displaystyle \frac{\varpi _TU_L^{(1)}}{\rho _0+p_0}}\right)^{}\right\}^{}`$ (52)
$`+\left\{{\displaystyle \frac{16c_2\varpi ^2e^\nu }{(\rho _0+p_0)jr^2}}\left[{\displaystyle \frac{(\rho _0+p_0)e^\nu }{A\varpi j\nu ^{}}}(\varpi r^2e^\nu )^{}\right]^{}\left({\displaystyle \frac{8c_1e^\nu }{Aj^2\nu ^{}r^4}}\right)[(\varpi r^2e^\nu )^{}]^2\right\}_TU_L^{(1)}`$
$`+(_Tim\chi )\left\{\left({\displaystyle \frac{\nu ^{}}{r^2}}{\displaystyle \frac{2}{r^3}}\right)(_Tim\varpi )+2im{\displaystyle \frac{\varpi ^{}}{r^2}}\right\}h_{1,L}^{(1)}`$
$`e^{\lambda \nu }(_Tim\chi )(_Tim\varpi )^2\mathrm{\Phi }_L^{(1)},`$
where
$$c_n=\frac{l+1}{l^n}Q_{}^2+(1)^{n1}\frac{l}{(l+1)^n}Q_+^2.$$
(53)
From Eqs.(24) and (46), a function $`\mathrm{\Phi }_L=\mathrm{\Phi }_L^{(1)}+\epsilon ^2\mathrm{\Phi }_L^{(2)}`$ satisfies the following equation, which is correct up to $`O(\epsilon ^2),`$
$`[\mathrm{\Phi }_L]=𝒟[\mathrm{\Phi }_L,h_{1,L},U_L]+𝒢[H_{0,\pm },K_\pm ].`$ (54)
The quantities without the superscript satisfy the same relations as in the leading order, i.e., Eqs.(17), (19), (33) and (34), which are adequate approximation to this order. This equation is of course reduced to the leading order equation (24), when the second-order rotational effects and the coupling to the polar modes are neglected.
The first term in $`𝒟_0`$ contains fourth derivative of $`\mathrm{\Phi }_L`$ with respect to $`r`$, since $`U_L`$ can be expressed by the second derivative of $`\mathrm{\Phi }_L`$ as shown in Eq.(19). The highest derivative term of $`\mathrm{\Phi }_L`$ with respect to $`r`$ is therefore given by
$$𝒟_{}[\mathrm{\Phi }_L]=8c_3\frac{\varpi e^\nu }{jr^2}\left\{\frac{(\rho _0+p_0)r^2}{jA\nu ^{}}\left(\frac{j\varpi }{(\rho _0+p_0)r^2}(jr^4_T\mathrm{\Phi }_L^{})^{}\right)^{}\right\}^{}.$$
(55)
Neglecting $`𝒢`$ and $`𝒟`$ except $`𝒟_{}`$ in Eq.(54) leads to $`[\mathrm{\Phi }_L]=𝒟_{}[\mathrm{\Phi }_L],`$ which is analogous to the Orr-Sommerfeld equation in the incompressible shear flow. (See Ref. .) The term (55) effectively gives the ’viscosity’ in the viscous fluid. The viscosity is important for the stability of the flows. For the small Reynolds number, the laminar flow is realized, whereas the flow becomes turbulence above a critical Reynolds number. The effective Reynolds number $`R_e`$ in Eq.(55) is estimated from dimensional argument as
$$R_e\frac{A\nu ^{}}{\varpi ^2}.$$
(56)
This is roughly the square of the ratio of the g-mode frequency to rotational one. The viscosity term will play a key role on the singular point of the first-order equation, but the consequence is not clear at moment. It is necessary to explore further how the effective Reynolds number should operate in the stability and so on.
## VI Concluding remarks
In this paper, we have explored an effective theory to describe the axial oscillation on a slowly rotating stars. The approximate equation governing the oscillation is constructed from the Einstein equations. The equation is derived by assuming that the angular dependence of the oscillation is dominated by a single component of spherical harmonics. The assumption in general breaks down as the evolution of oscillation. There are coupling terms of order $`\epsilon `$ in the rotating fluids. Other oscillation patterns with different spherical harmonics will gradually be produced through the rotational coupling. For this reason, the equation is valid for small $`t`$, and can be used to examine the time evolution as the initial-value problem. The rotational effects up to third order are involved in this paper, so that the regime of application is enlarged. The equation derived here is also irrelevant to the singular point found in the first order one.
The equation also shows a remarkable property. It is evident that the axial oscillation strongly couples to g-mode oscillations. Viscosity-like term arises from the polar pieces related to the g-mode oscillations. The ’viscosity’ term originated from considering the sub-system only, i.e., a single component of the spherical harmonics. There are no production and extinction in the whole system, but, e.g., ’energy’ of a component is partially transfered to the others. This transportation is regarded as dissipative effect so far as a particular sub-system is concerned. The term also has a significant implication. The condition $`A=0`$ is a good approximation for cold neutron stars, so that the coupling may be neglected. On the other hand, it is not clear that the condition holds, in particular for newly born hot neutron stars, in which the r-mode instability sets in.
In this paper, we concentrated the equations only inside a star, to be more precise, the equations for the region $`A0.`$ The pulsation equation derived here should be solved with appropriate boundary conditions. The boundary conditions are determined from matching with the equations outside, or regularity conditions. For example, the regularity condition of a function $`\mathrm{\Phi }`$ is given by $`\mathrm{\Phi }r^{l1}`$ near the center. Depending on details of the stellar structure, the solution for the $`A0`$ region may be matched to the solution for isentropic region, $`A=0`$. Furthermore, the interior solution should be matched to the exterior one at the stellar surface. The exterior perturbation equation in vacuum is not derived here, but the form should be reduced to the wave equation describing gravitational wave. The perturbation equation should be solved by out-going boundary condition at infinity. One question may arise. Is it possible to calculate the radiation reaction at this order? Newtonian estimate indicates that the back-reaction is of order $`\epsilon ^{2m+2}`$ for $`m2.`$ Our expansion of the rotational parameter is limited to the third-order, and higher order corrections are necessary to examine the effect in a consistent way. The radiation-reaction effect is also a kind of dissipative one, so that the accurate evolution for a long period is necessary.
## Acknowledgements
We would like to thank John Friedman and Nils Andersson for helpful discussions. This was supported in part by the Grant-in-Aid for Scientific Research Fund of the Ministry of Education, Science and Culture of Japan(08NP0801).
## A Equilibrium configuration of a slowly rotating perfect fluid
We here summarize the equilibrium of a slowly rotating star to explain our notation. The equilibrium state with uniform angular velocity $`\mathrm{\Omega }O(\epsilon )`$ can be described by stationary and axisymmetric metric $`g_{\mu \nu }`$, 4-velocity $`u^\mu =(u^t,0,0,u^\varphi )`$, pressure $`p`$ and energy density $`\rho `$ of the fluid. The rotational corrections up to $`O(\epsilon ^3)`$ is needed to assure the consistency in our analysis. The metric is given by
$`ds^2`$ $`=`$ $`e^\nu [1+2(h_0+h_2P_2)]dt^2+e^\lambda [1+{\displaystyle \frac{2e^\lambda }{r}}(m_0+m_2P_2)]dr^2`$ (A2)
$`+r^2(1+2k_2P_2)\left\{d\theta ^2+\mathrm{sin}^2\theta \left[d\varphi \left(\omega +W_1W_3{\displaystyle \frac{1}{\mathrm{sin}\theta }}{\displaystyle \frac{dP_3}{d\theta }}\right)dt\right]^2\right\},`$
where $`P_l=P_l(\mathrm{cos}\theta )(l=2,3)`$ is the Legendre polynomial of order $`l`$. The metric functions introduced above obey the following ordering in $`\epsilon `$: $`\omega O(\epsilon )`$, $`h_0,h_2,m_0,m_2,k_2O(\epsilon ^2)`$, $`W_1,W_3O(\epsilon ^3)`$. These are functions of radial coordinate $`r`$ only. The components of the 4-velocity are
$$u^t=(g_{tt}2g_{t\varphi }\mathrm{\Omega }g_{\varphi \varphi }\mathrm{\Omega }^2)^{1/2},u^\varphi =\mathrm{\Omega }u^t.$$
(A3)
The pressure and energy density are respectively given by
$$p=p_0+\{p_{20}+p_{22}P_2\},\rho =\rho _0+\{\rho _{20}+\rho _{22}P_2\},$$
(A4)
where $`p_0`$ and $`\rho _0`$ are the pressure and energy density of non-rotating fluid. The centrifugal force of order $`\epsilon ^2`$ alters the configuration shape, which corresponds to the quantities in the braces. The functions $`p_{20},p_{22},\rho _{20}`$ and $`\rho _{22}`$ are related with the metric functions of order $`\epsilon ^2`$. We rather use the metric functions to eliminate the pressure and density of order $`\epsilon ^2`$ in the oscillation equations.
## B The second-order terms
### 1 Coefficient of $`h_{1,L}^{(1)}`$, $`\alpha _1`$
$`\alpha _1`$ $`=`$ $`c_2\left[32\pi (\rho _0+p_0)\left({\displaystyle \frac{2}{\nu ^2j^2r^2}}\varpi (\omega e^\nu )^{}+{\displaystyle \frac{2}{r}}\varpi \omega (\omega +3\varpi )\varpi ^{}\right){\displaystyle \frac{8}{r^2e^\lambda }}(\omega r)^{}\varpi ^{}\right]`$ (B1)
$`+`$ $`c_1[5\pi (\rho _0^{}+3p_0^{})\varpi ^2+({\displaystyle \frac{28}{3r^2e^\lambda }}{\displaystyle \frac{5}{12r^2}}+{\displaystyle \frac{5}{2r^3\nu ^{}}})\omega \varpi ^{}{\displaystyle \frac{97}{12re^\lambda }}\varpi ^2`$ (B3)
$`+4\pi (p_0+\rho _0)({\displaystyle \frac{5\varpi ^2}{r}}+23\varpi \varpi ^{}+{\displaystyle \frac{28}{3}}\omega \varpi ^{}{\displaystyle \frac{56}{3r}}\omega \varpi {\displaystyle \frac{12}{j^2r^2\nu ^2}}\varpi (\omega e^\nu )^{})]`$
$`+`$ $`{\displaystyle \frac{m^2}{l(l+1)}}[2\pi (l(l+1){\displaystyle \frac{11}{2}})(\rho _0^{}+3p_0^{})\varpi ^2+{\displaystyle \frac{4\pi }{r^2}}(l(l+1){\displaystyle \frac{15}{2}})(\rho _0+p_0)(\varpi ^2r^2)^{}`$ (B6)
$`{\displaystyle \frac{e^\nu }{2\nu ^{}}}\left(l(l+1)+{\displaystyle \frac{3}{2}}\right)(e^\nu r^2)^{}\omega \varpi ^{}{\displaystyle \frac{1}{2re^\lambda }}\left(l(l+1){\displaystyle \frac{1}{6}}\right)\varpi ^2+{\displaystyle \frac{1}{3r^2}}(8e^\lambda 7)\omega \varpi ^{}`$
$`+16\pi (\rho _0+p_0)({\displaystyle \frac{\varpi ^2}{r}}{\displaystyle \frac{e^\lambda }{r^2\nu ^2}}\varpi \varpi ^{}{\displaystyle \frac{e^\lambda }{\nu ^{}r^2}}\omega \varpi +{\displaystyle \frac{2}{3}}\omega r^2(\varpi r^2)^{}+3\varpi \varpi ^{})]`$
$`+`$ $`{\displaystyle \frac{4}{r^3}}e^{(6\nu \lambda )/4}(e^{(2\nu 3\lambda )/4}h_0^{})^{}{\displaystyle \frac{e^\nu }{r^4}}(l(l+1)4((\nu ^{}r)^26)e^\lambda )h_0^{}`$ (B13)
$`{\displaystyle \frac{1}{j^2r^9}}\left(\nu ^{}r^7e^{2\lambda }\left({\displaystyle \frac{e^\lambda }{r}}m_0\right)^{}\right)^{}{\displaystyle \frac{e^\nu }{r^4}}(l(l+1)2+4e^\lambda )\left({\displaystyle \frac{e^\lambda }{r}}m_0\right)^{}`$
$`{\displaystyle \frac{2\pi }{3}}(\rho _0^{}+3p_0^{})\left\{\left[l(l+1)4{\displaystyle \frac{(e^\nu r^2)^{}}{re^\nu }}\right]\varpi ^2{\displaystyle \frac{12\nu ^{}}{j^2r^2}}m_0\right\}`$
$`+[4\pi (\rho _0+p_0)r^2(2(\nu ^{}r)^2+\nu ^{}\lambda ^{}r^2+12\nu ^{}r+2\lambda ^{}r+4)+82\lambda ^{}r]{\displaystyle \frac{m_0}{j^2r^6}}`$
$`+{\displaystyle \frac{32\pi ^2}{3}}(\rho _0+p_0)^2\varpi ^2e^{\lambda \nu }(r^2e^\nu )^{}{\displaystyle \frac{16\pi }{3r^5}}e^{\nu /2}(r^6\varpi ^2e^{\nu /2})^{}p_0^{}`$
$`{\displaystyle \frac{8\pi }{3r}}(\rho _0+p_0)[l(l+1)\varpi (\varpi r)^{}8\omega \varpi +4r\mathrm{\Omega }\varpi ^{}2(3e^\lambda )\varpi ^2]`$
$`\left\{{\displaystyle \frac{l(l+1)}{2r^2}}\left(1+{\displaystyle \frac{2}{\nu ^{}r}}\right)+{\displaystyle \frac{4}{3r^2}}(2e^\lambda 1)\right\}\omega \varpi ^{}+{\displaystyle \frac{1}{6re^\lambda }}\left\{l(l+1)43\nu ^{}r\right\}\varpi ^2.`$
### 2 Coefficient of $`_t\mathrm{\Phi }_L^{(1)}`$, $`\alpha _2`$
$`\alpha _2`$ $`=`$ $`c_2[{\displaystyle \frac{32\pi e^{\lambda \nu }}{\nu ^{}}}(\varpi +\mathrm{\Omega })\varpi r^2(\rho _0^{}+3p_0^{})32\pi (p_0+\rho _0){\displaystyle \frac{e^{\lambda \nu }}{(r\nu ^{})^2}}\{(3\varpi +\chi )(4r^2e^\lambda \varpi `$ (B15)
$`2r\varpi e^\nu (r^2e^\nu )^{}r^2(r^2e^\nu )^{}(\varpi e^\nu )^{}+4r^2e^\lambda \omega \varpi \nu ^{}(r^4\omega \varpi ^{}+2(r^4\varpi ^2)^{})\}{\displaystyle \frac{4}{r\nu ^{}}}(r^2e^\nu )^{}\omega \varpi ^{}]`$
$`+`$ $`c_1[16\pi (\rho _0^{}+3p_0^{})\mathrm{\Omega }(\varpi +\mathrm{\Omega }){\displaystyle \frac{r^2e^{\lambda \nu }}{\nu ^{}}}`$ (B18)
$`16\pi (p_0+\rho _0){\displaystyle \frac{e^{\lambda \nu }}{\nu ^2}}\{4e^\lambda \omega \varpi 2r\nu ^{}(\varpi +\mathrm{\Omega })^2r^2\nu ^{}\varpi ^{}(\varpi +3\mathrm{\Omega })\}`$
$`+{\displaystyle \frac{8e^\nu }{r\nu ^{}}}\omega ^22(r+{\displaystyle \frac{2}{\nu ^{}}})e^\nu \omega \varpi ^{}{\displaystyle \frac{49}{24}}r^2e^\nu \varpi ^22f_1+f_2+f_4]`$
$`+`$ $`{\displaystyle \frac{m^2}{l(l+1)}}[l(l+1)({\displaystyle \frac{2}{r\nu ^{}}}1)e^{\lambda \nu }\omega ^216\pi (p_0+\rho _0)r^2e^{\lambda \nu }\varpi ^2{\displaystyle \frac{71}{24}}r^2e^\nu \varpi ^2`$ (B20)
$`f_1+f_23f_3f_4]`$
$`+`$ $`l(l+1)\left({\displaystyle \frac{2e^{2\lambda }}{r^3}}m_0+\left(1{\displaystyle \frac{2}{r\nu ^{}}}\right)e^{\lambda \nu }\omega ^2\right)\nu ^{}\left(\left({\displaystyle \frac{e^\lambda m_0}{r}}\right)^{}+h_0^{}\right)+{\displaystyle \frac{4e^{2\lambda }m_0}{r^3}}`$ (B22)
$`16\pi p_0^{}e^\lambda \left(e^\lambda m_0+{\displaystyle \frac{r^3}{3}}e^\nu \varpi ^2\right)+2r^2e^\nu \varpi ^2+f_1{\displaystyle \frac{2}{3}}f_2+f_3,`$
where
$`f_1`$ $`=`$ $`{\displaystyle \frac{8\pi e^\lambda }{3}}(p_0+\rho _0)(\{8e^\nu r^212r^2j^2+{\displaystyle \frac{64e^\nu }{(\nu ^{})^2}}(4\pi r^2e^\lambda p_0r\nu ^{})\}\varpi ^2`$ (B24)
$`+{\displaystyle \frac{8r^2}{\nu ^{}}}e^\nu (1+{\displaystyle \frac{2}{\nu ^{}r}})\varpi \varpi ^{}+r^{12}j^2\left({\displaystyle \frac{\varpi }{r^8}}\right)^{}\varpi ^{})+{\displaystyle \frac{4r}{3}}e^\nu \omega \varpi ^{},`$
$`f_2`$ $`=`$ $`4\pi (\rho _0^{}+3p_0^{})(3e^\lambda \xi _2+2r^3e^\nu \varpi ^2)16\pi (p_0+\rho _0)^2r^4e^{\lambda \nu }\varpi ^212\left(1{\displaystyle \frac{2}{r\nu ^{}}}\right){\displaystyle \frac{e^\lambda }{r^2}}k_2`$ (B26)
$`+{\displaystyle \frac{6}{r^3}}\left(13e^\lambda +4\pi (\rho _0+p_0)r^2e^\lambda \right)\xi _2+{\displaystyle \frac{4}{r\nu ^{}}}(13e^\lambda )e^\nu \varpi ^2+\left({\displaystyle \frac{r}{\nu ^{}}}r^2e^\lambda \right)e^\nu \varpi ^2,`$
$`f_3`$ $`=`$ $`l(l+1)\left({\displaystyle \frac{e^\lambda }{2r^2}}\left(k_2{\displaystyle \frac{\nu ^{}}{2}}\xi _2\right){\displaystyle \frac{4\pi }{3}}(p_0+\rho _0)r^2e^{\lambda \nu }\varpi ^2{\displaystyle \frac{e^{\lambda \nu }}{6}}\varpi ^2{\displaystyle \frac{r^2}{12}}e^\nu \varpi ^2\right)`$ (B29)
$`+{\displaystyle \frac{2\nu ^{}e^\lambda }{r^2}}\xi _2+{\displaystyle \frac{4}{3}}e^{\lambda \nu }\varpi ^2,`$
$`f_4`$ $`=`$ $`{\displaystyle \frac{3e^\lambda }{4r^2}}\left(k_2{\displaystyle \frac{\nu ^{}}{2}}\xi _2\right)4\pi (p_0+\rho _0)\left(\left({\displaystyle \frac{e^\lambda }{6}}+4\right)\varpi ^2r^2e^\nu {\displaystyle \frac{8}{\nu ^2}}e^{2\lambda \nu }\omega \varpi {\displaystyle \frac{r^{12}}{3}}e^\nu \left({\displaystyle \frac{\varpi }{r^8}}\right)^{}\varpi ^{}\right)`$ (B31)
$`{\displaystyle \frac{e^{\lambda \nu }}{4}}\varpi ^2+\left({\displaystyle \frac{1}{r\nu ^{}}}{\displaystyle \frac{1}{2}}\right)e^{\lambda \nu }\omega ^2.`$
### 3 Coefficient of $`U_L^{(1)}`$, $`\beta _1`$
$`\beta _1`$ $`=`$ $`c_3\left[32\chi \varpi ^2\left\{{\displaystyle \frac{e^{\lambda /2}}{r^2}}\left({\displaystyle \frac{e^{\lambda /2}r^2}{\nu ^{}}}\right)^{}+\left({\displaystyle \frac{8\pi e^{2\lambda }}{(\nu ^{})^2}}\right)(\rho _0+p_0)\right\}\right]`$ (B32)
$``$ $`c_2\left[{\displaystyle \frac{16}{\nu ^{}j^2r^4}}(\chi \varpi ^2r^4e^\nu )^{}+4\chi \varpi ^{}\left\{(\omega r^2)^{}+{\displaystyle \frac{2e^\lambda }{\nu ^{}}}\omega \right\}\right]`$ (B33)
$`+`$ $`c_1[{\displaystyle \frac{12\chi }{jr}}\left({\displaystyle \frac{\xi _2}{jr}}\right)^{}18{\displaystyle \frac{\xi _2\chi ^{}}{j^2r^2}}12{\displaystyle \frac{\chi k_2}{j^2r^2}}`$ (B35)
$`+{\displaystyle \frac{120W_3}{l(l+1)j^2r^2}}+4\chi \varpi ^2e^\lambda \{1+4\pi (\rho _0+p_0)r^2\}+{\displaystyle \frac{2}{3r^6}}\chi \varpi ^{}(\omega r^8)^{}]`$
$`+`$ $`{\displaystyle \frac{m^2}{l(l+1)}}\left[6{\displaystyle \frac{(\chi r^2)^{}}{j^2r^4}}\xi _2+12{\displaystyle \frac{\chi k_2}{j^2r^2}}{\displaystyle \frac{120W_3}{l(l+1)j^2r^2}}+{\displaystyle \frac{4}{3}}\chi \varpi ^{}(\omega r^2)^{}\right]`$ (B36)
$`+`$ $`2\chi (e^\nu )^{}\left\{h_0^{}+{\displaystyle \frac{1}{j^2}}\left({\displaystyle \frac{m_0}{re^\nu }}\right)^{}\right\}+{\displaystyle \frac{8W_1+48W_3}{l(l+1)j^2r^2}}+{\displaystyle \frac{32\pi }{3}}\chi \varpi ^2p_0^{}r^3e^\lambda {\displaystyle \frac{4}{3}}\chi \varpi ^{}(\omega r^2)^{}.`$ (B37)
### 4 Coefficient of $`\mathrm{\Phi }_L^{(1)}`$, $`\beta _2`$
$`\beta _2`$ $`=`$ $`4c_2\chi \left[4\pi (\rho _0+p_0)\left\{4\varpi (\omega r)^{}\varpi ^{}\omega r\right\}r^3e^\nu \varpi ^{}(\omega r)^{}j^2r^2\right]`$ (B39)
$`+`$ $`c_1\chi r^2e^\nu [{\displaystyle \frac{5\pi }{2}}(\rho _0^{}+3p_0^{})\varpi ^2r^2+\varpi ^{}\{({\displaystyle \frac{5}{4\nu ^{}r}}{\displaystyle \frac{5}{24}})\omega +({\displaystyle \frac{14}{3}}\omega {\displaystyle \frac{97}{24}}\varpi ^{}r)e^\lambda \}`$ (B41)
$`4\pi (\rho _0+p_0)\{{\displaystyle \frac{5}{2}}\varpi ^2r+{\displaystyle \frac{44}{3}}\varpi \omega r{\displaystyle \frac{85}{6}}\varpi ^{}\varpi r^2{\displaystyle \frac{14}{3}}\varpi ^{}\omega r^2+{\displaystyle \frac{6}{\nu ^2j^2}}\varpi (\omega e^\nu )^{}\}]`$
$`+`$ $`m^2[4\pi (\rho _0^{}+3p_0^{})\chi \varpi r^4e^\nu ({\displaystyle \frac{1}{4}}\varpi {\displaystyle \frac{11}{16}}\chi )`$ (B44)
$`+\chi ^{}\omega e^\nu \left\{{\displaystyle \frac{(r^2e^\nu )^{}}{4\nu ^{}}}\varpi e^\nu +{\displaystyle \frac{3}{8\nu ^{}}}\chi r{\displaystyle \frac{19}{48}}\chi r^2\right\}\chi j^2r^2\left\{{\displaystyle \frac{1}{4}}\varpi ^2r{\displaystyle \frac{2}{3}}\chi ^{}\omega {\displaystyle \frac{1}{48}}\chi ^{}\varpi ^{}r\right\}`$
$`+4\pi (\rho _0+p_0)\chi r^2e^\nu \{(\varpi r)^{}\varpi r{\displaystyle \frac{11}{4}}\chi \varpi r{\displaystyle \frac{8}{3}}\chi \omega r+{\displaystyle \frac{1}{j^2\nu ^2}}\chi (\omega e^\nu )^{}{\displaystyle \frac{1}{12}}\chi \varpi ^{}r^2+{\displaystyle \frac{2}{3}}\chi ^{}\omega r^2\}]`$
$`+`$ $`\left(\chi \varpi +{\displaystyle \frac{\nu ^{}}{2}}\chi ^{}r^2e^\lambda \right)\left(h_0^{}+\left({\displaystyle \frac{e^\lambda m_0}{r}}\right)^{}\right)+{\displaystyle \frac{2\pi }{3}}(\rho _0^{}+3p_0^{})(4\chi \varpi )\varpi ^2r^4e^\nu `$ (B47)
$`(\lambda ^{}+\nu ^{})j^2\left[{\displaystyle \frac{1}{9}}((\varpi r)^3)^{}+{\displaystyle \frac{4}{3}}\left({\displaystyle \frac{\varpi }{r}}\right)^{}\chi \varpi r^4+{\displaystyle \frac{2}{3}}\left({\displaystyle \frac{\varpi }{r^4}}\right)^{}\chi \omega r^7+{\displaystyle \frac{\nu ^{}}{6}}\varpi ^2\chi ^{}r^4+{\displaystyle \frac{\nu ^{}r}{2j^2}}\chi ^{}m_0\right]`$
$`+\left[{\displaystyle \frac{2}{3}}\chi {\displaystyle \frac{1}{2}}\varpi {\displaystyle \frac{1}{\nu ^{}r}}\varpi \right]\varpi ^{}\omega r^2e^\nu +{\displaystyle \frac{1}{6}}(\varpi ^2r8\chi ^{}\omega +4\chi ^{}\varpi ^{}r)\varpi j^2r^2.`$
### 5 Coefficient of $`\mathrm{\Phi }_L^{(1)}`$, $`\beta _3`$
$`\beta _3`$ $`=`$ $`c_2[16\pi r^2e^\nu \chi \varpi (r^2\omega )^{}(\rho _0^{}+3p_0^{})16\pi (\rho _0+p_0)e^\nu \chi (16\pi (\rho _0+p_0)r^4e^\lambda \varpi (\varpi +\mathrm{\Omega })`$ (B50)
$`+2r^2e^\lambda \mathrm{\Omega }\varpi \left(1+{\displaystyle \frac{8}{r\nu ^{}}}{\displaystyle \frac{32\pi }{\nu ^2}}e^\lambda p_0\right)6r\omega (r^2\varpi )^{}{\displaystyle \frac{2re^\lambda }{\nu ^2}}\varpi ^{}(r\nu ^{}\varpi +2\mathrm{\Omega }r\varpi ^{})`$
$`+r^3\varpi ^{}(4\varpi +5r\varpi ^{}))+{\displaystyle \frac{4}{r\nu ^{}}}(r^2e^\nu )^{}\chi \omega \varpi ^{}4e^\nu r^2\chi ^{}\varpi ^{}(4\pi (\rho _0+p_0)r^2\omega +e^\lambda (r\omega )^{})]`$
$`+`$ $`c_1[{\displaystyle \frac{16\pi e^{\lambda \nu }}{\nu ^2}}(\rho _0+p_0)\mathrm{\Omega }\{112\pi r^2e^\lambda \varpi \chi p_04e^\lambda \chi \omega 2r\nu ^{}(9\varpi +\omega )\chi +6r\chi \varpi ^{}+7r^2\nu ^{}\chi \varpi ^{}`$ (B52)
$`+r^2\nu ^{}\omega \chi ^{}\}+2re^\nu \omega \chi ^{}(2\varpi +(1{\displaystyle \frac{2}{r\nu ^{}}})\omega )r^2e^\nu \varpi ^{}\chi ^{}(\varpi (1+{\displaystyle \frac{2}{r\nu ^{}}})\omega )+g_1]`$
$`+`$ $`{\displaystyle \frac{m^2}{l(l+1)}}[4\pi (\rho _0+p_0)e^{\lambda \nu }\varpi ({\displaystyle \frac{23}{4}}r^2\varpi ^2+{\displaystyle \frac{15}{4}}r^2\varpi \omega +{\displaystyle \frac{4e^\lambda }{\nu ^2}}\mathrm{\Omega }\omega )+{\displaystyle \frac{e^{\lambda \nu }}{2}}(5\varpi +\omega )(1{\displaystyle \frac{2}{r\nu ^{}}})\omega ^2`$ (B55)
$`+{\displaystyle \frac{16\pi }{\nu ^{}}}(\rho _0^{}+3p_0^{})r^2e^{\lambda \nu }\mathrm{\Omega }^2(\varpi +\mathrm{\Omega })32\pi (\rho _0+p_0)\varpi \chi e^{\lambda \nu }\left((\mathrm{\Omega }+\omega )r^2+{\displaystyle \frac{\mathrm{\Omega }}{\nu ^2}}(1e^\lambda )\right)`$
$`+r^2\mathrm{\Omega }\varpi ^2e^\nu +4\chi e^\nu (e^\lambda \varpi ^2{\displaystyle \frac{2}{r\nu ^{}}}\omega ^2)g_1+3g_2g_3]m^2(1{\displaystyle \frac{2}{r\nu ^{}}})e^{\lambda \nu }\mathrm{\Omega }\omega ^2`$
$`+`$ $`[l(l+1)(\mathrm{\Omega }\chi )\{(1{\displaystyle \frac{2}{\nu ^{}r}})e^{\lambda \nu }\omega ^2{\displaystyle \frac{2e^{2\lambda }}{r^3}}m_0\}+{\displaystyle \frac{8\pi }{3}}(\rho _0+p_0)r^2e^{\lambda \nu }\varpi \{\varpi ^24(\varpi \omega )\chi \}`$ (B59)
$`+\mathrm{\Omega }\left\{4\pi e^\lambda (\rho _0+p_0)\varpi \left({\displaystyle \frac{r^2}{12}}e^\nu \varpi +{\displaystyle \frac{4}{\nu ^2}}e^{\lambda \nu }\omega \right)+{\displaystyle \frac{4}{r^3}}e^{2\lambda }m_0\nu ^{}\left(\left({\displaystyle \frac{e^\lambda }{r}}m_0\right)^{}+h_0^{}\right)\right\}`$
$`+{\displaystyle \frac{4}{r^3}}e^{2\lambda }\{4\pi (8\pi (\rho _0+p_0)r^2\chi +\chi \mathrm{\Omega })p_0^{}r^3\chi \}\left(m_0+{\displaystyle \frac{r^3}{3}}j^2\varpi ^2\right)`$
$`+{\displaystyle \frac{1}{r}}\{(4+\nu ^{}r16\pi p_0^{}r^3)\chi +\chi ^{}r\}(\left({\displaystyle \frac{e^\lambda }{r}}m_0\right)^{}+h_0^{})g_2+g_3],`$
where
$`l(l+1)g_1`$ $`=`$ $`4\pi (\rho _0^{}+3p_0^{})\{{\displaystyle \frac{3}{4}}(8\varpi +3\mathrm{\Omega })e^\lambda \xi _2+{\displaystyle \frac{r^2\mathrm{\Omega }}{\nu ^{}}}e^{\lambda \nu }(\varpi +\mathrm{\Omega })(4\varpi +3\mathrm{\Omega })`$ (B67)
$`+{\displaystyle \frac{5}{6}}e^\nu r^2\varpi ^2(4(r^2\varpi )^{}11\mathrm{\Omega }r)\}+{\displaystyle \frac{20\pi }{3}}(\rho _0+p_0)^2e^{\lambda \nu }r^4\varpi ^2(8\varpi +11\mathrm{\Omega })`$
$`+4\pi (\rho _0+p_0)e^\lambda \{{\displaystyle \frac{3}{2r}}(5\varpi 3\omega )\xi _2+{\displaystyle \frac{5}{12}}r^2j^2(8\varpi 11\mathrm{\Omega })(12\varpi ^2(r\varpi ^{})^2)`$
$`+{\displaystyle \frac{44}{3}}r^2j^2\varpi \varpi ^{}(r^2\varpi )^{}26r^3j^2\mathrm{\Omega }\varpi \varpi ^{}+{\displaystyle \frac{17}{3}}e^\nu r^2\varpi ^3{\displaystyle \frac{1123}{24}}e^\nu r^2\mathrm{\Omega }\varpi ^2`$
$`+{\displaystyle \frac{\mathrm{\Omega }e^\nu }{\nu ^2}}(22e^\lambda \varpi \omega +2r\nu ^{}(20\varpi ^2+12\varpi \omega +3\omega ^2)+8r\varpi \varpi ^{}+(32\varpi +13\omega )r^2\nu ^{}\varpi ^{})\}`$
$`+(5\varpi 3\omega )\{({\displaystyle \frac{45}{16r^2}}{\displaystyle \frac{6}{\nu ^{}r^3}})e^\lambda k_2{\displaystyle \frac{3}{2r^3}}(13e^\lambda {\displaystyle \frac{\nu ^{}re^\lambda }{16}})\xi _2+{\displaystyle \frac{e^{\lambda \nu }}{16}}\varpi ^2`$
$`+(1{\displaystyle \frac{2}{r\nu ^{}}}){\displaystyle \frac{e^{\lambda \nu }}{8}}\omega ^2{\displaystyle \frac{2e^\nu }{r\nu ^{}}}\omega ^2+(3e^\lambda 1){\displaystyle \frac{\varpi ^2e^\nu }{r\nu ^{}}}{\displaystyle \frac{r}{4\nu ^{}}}e^\nu (\varpi ^{})^2\}{\displaystyle \frac{7}{\nu ^{}}}e^\nu \mathrm{\Omega }\omega \varpi ^{}`$
$`+{\displaystyle \frac{re^\nu }{6}}(23\varpi 9\omega )\omega \varpi ^{}+{\displaystyle \frac{r^2e^\nu }{96}}(437\varpi 403\omega )\varpi ^2+{\displaystyle \frac{5}{12}}r^2j^2\varpi ^2(11\mathrm{\Omega }8(r\varpi )^{}),`$
$`g_2`$ $`=`$ $`l(l+1)e^\lambda \mathrm{\Omega }\left\{{\displaystyle \frac{4\pi }{3}}(\rho _0+p_0)r^2e^\nu \varpi ^2{\displaystyle \frac{1}{2r^2}}\left(k_2{\displaystyle \frac{\nu ^{}}{2}}\xi _2\right)+{\displaystyle \frac{e^\nu }{6}}\varpi ^2+{\displaystyle \frac{r^2j^2}{12}}\varpi ^2\right\}`$ (B77)
$`+4\pi (\rho _0^{}+3p_0^{})e^\lambda \{(2\chi +{\displaystyle \frac{\mathrm{\Omega }}{2}})\xi _2+{\displaystyle \frac{j^2}{3}}\varpi (r^3\varpi (\mathrm{\Omega }4\chi )2\chi (r^4\omega )^{})`$
$`{\displaystyle \frac{2}{\nu ^{}}}r^2e^\nu \mathrm{\Omega }^2(\varpi +\mathrm{\Omega }))\}{\displaystyle \frac{32\pi ^2}{3}}(\rho _0+p_0)^2r^4e^{\lambda \nu }\varpi (\mathrm{\Omega }\varpi 4\chi (\varpi +2\mathrm{\Omega }))`$
$`+4\pi (\rho _0+p_0)\{e^\lambda \mathrm{\Omega }({\displaystyle \frac{\xi _2}{r}}+{\displaystyle \frac{r^2j^2}{6}}(12\varpi ^2+8r\varpi ^{}\varpi (r\varpi ^{})^2))`$
$`{\displaystyle \frac{4}{r}}e^\lambda \chi \xi _2{\displaystyle \frac{2}{3}}r^2e^\nu (12(\varpi +2\omega )\varpi \chi +4r(\varpi +3\omega )\chi \varpi ^{}(5\varpi +2\omega )r^2\varpi ^{}\chi ^{})\}`$
$`{\displaystyle \frac{1}{r^3}}\left(r\nu ^{}e^\lambda \left({\displaystyle \frac{\varpi }{2}}+{\displaystyle \frac{29}{16}}\mathrm{\Omega }2\chi \right)(13e^\lambda )(\mathrm{\Omega }4\chi )\right)\xi _2`$
$`+{\displaystyle \frac{e^\lambda }{r^2}}\left(\varpi {\displaystyle \frac{19}{8}}\mathrm{\Omega }+8\chi +{\displaystyle \frac{4}{r\nu ^{}}}(\mathrm{\Omega }4\chi )\right)k_2{\displaystyle \frac{r^2}{6}}\varpi ^2e^\nu \left(e^\lambda \mathrm{\Omega }+{\displaystyle \frac{47}{8}}\mathrm{\Omega }+\varpi {\displaystyle \frac{\mathrm{\Omega }}{r\nu ^{}}}\right)`$
$`e^{\lambda \nu }\varpi ^2\left({\displaystyle \frac{\varpi }{3}}+{\displaystyle \frac{29}{24}}\mathrm{\Omega }+{\displaystyle \frac{2}{3\nu ^{}r}}(3e^\lambda )(\mathrm{\Omega }4\chi )\right)`$
$`+{\displaystyle \frac{2r}{3\nu ^{}}}e^\nu \chi ^{}\left\{e^\lambda \nu ^{}\varpi ^{}(r\varpi +(r^2\omega )^{})\varpi \varpi ^{}+2\nu ^{}\varpi \omega +r\nu ^{}\varpi ^{}(3\varpi \omega )\right\},`$
$`g_3`$ $`=`$ $`\mathrm{\Omega }e^\nu \{{\displaystyle \frac{8\pi re^\lambda }{\nu ^{}}}(\rho _0+p_0)(2(\mathrm{\Omega }+\varpi )^2+r\varpi ^{}(3\mathrm{\Omega }+\varpi ))`$ (B79)
$`{\displaystyle \frac{e^\lambda }{4}}(1{\displaystyle \frac{2}{r\nu ^{}}})\omega ^2+{\displaystyle \frac{4}{\nu ^{}r}}\omega ^2(r+{\displaystyle \frac{2}{\nu ^{}}})\omega \varpi ^{}\}.`$
### 6 $`𝒢[H_{0,\pm }^{(1)},K_\pm ^{(1)}]`$
$`𝒢[H_{0,\pm }^{(1)},K_\pm ^{(1)}]`$ $`=`$ $`A_1_TK_\pm ^{(1)}+B_1_TH_\pm ^{(1)}+imA_2K_\pm ^{(1)}+imB_2H_\pm ^{(1)},`$ (B80)
where
$`A_1`$ $`=`$ $`Q_\pm \left[{\displaystyle \frac{2e^\lambda }{r^2}}\omega +64\pi (\rho _0+p_0)\left({\displaystyle \frac{e^\nu }{j^2r\nu ^{}}}\right)^2\varpi \right]`$ (B82)
$`+S_\pm \left[16\pi e^\lambda (\rho _0+p_0)\left\{\chi {\displaystyle \frac{2e^\lambda }{r^2\nu ^2}}(\varpi +\chi )\right\}{\displaystyle \frac{2e^\lambda }{r^2}}\left(1{\displaystyle \frac{2}{r\nu ^{}}}\right)\omega \right]+T_\pm \left[{\displaystyle \frac{e^{\lambda \nu }}{\nu ^{}}}\left({\displaystyle \frac{e^\nu }{r^2}}\right)^{}\omega \right],`$
$`B_1`$ $`=`$ $`Q_\pm B_3+S_\pm [{\displaystyle \frac{16\pi e^\lambda }{r^2\nu ^2}}(\rho _0+p_0)\{2\chi (e^\lambda +r\nu ^{})+2e^\lambda \varpi r^2\nu ^{}\varpi ^{}(2\varpi +\omega )(1{\displaystyle \frac{3}{4}}r\nu ^{})r\nu ^{}\}`$ (B85)
$`(2\varpi +\omega )\left({\displaystyle \frac{8\pi e^\nu \rho _0^{}}{j^2\nu ^{}}}\right){\displaystyle \frac{4}{r^3\nu ^{}}}\omega ]+T_\pm [{\displaystyle \frac{2e^\lambda }{r^3\nu ^{}}}(1{\displaystyle \frac{r\nu ^{}}{2}})\omega ],`$
$`A_2`$ $`=`$ $`Q_\pm \left[+2\chi e^\lambda \left(8\pi \mathrm{\Omega }(\rho _0+p_0){\displaystyle \frac{\omega }{r^2}}\right)\right]+S_\pm \left[2\chi e^\lambda \left\{8\pi \mathrm{\Omega }(\rho _0+p_0)\left(1{\displaystyle \frac{2}{\nu ^{}r}}\right)\omega \right\}\right]`$ (B87)
$`+T_\pm \left[+\chi e^\lambda \left({\displaystyle \frac{1}{r^2}}{\displaystyle \frac{2}{\nu ^{}r^3}}\right)\omega \right],`$
$`B_2`$ $`=`$ $`Q_\pm \left[\chi B_3+8\pi \chi e^\lambda (\rho _0+p_0)\varpi \right]+S_\pm \left[\chi B_32\chi e^\lambda \left\{{\displaystyle \frac{e^\nu \omega }{\nu ^{}}}\left({\displaystyle \frac{e^\nu }{r^2}}\right)^{}+4\pi (\rho _0+p_0)\varpi \right\}\right]`$ (B89)
$`+T_\pm \left[{\displaystyle \frac{e^{\lambda \nu }}{\nu ^{}}}\left({\displaystyle \frac{e^\nu }{r^2}}\right)^{}\chi \omega +4\pi (\rho _0+p_0)\left\{{\displaystyle \frac{(r^4\chi ^2e^\nu )^{}}{j^2r^4\nu ^{}}}8\left({\displaystyle \frac{\chi e^\nu }{r\nu ^{}}}\right)^2\right\}\right],`$
$`B_3`$ $`=`$ $`(2\varpi +\omega )\left\{{\displaystyle \frac{8\pi e^\nu \rho _0^{}}{j^2\nu ^{}}}+4\pi e^\lambda (\rho _0+p_0)\left({\displaystyle \frac{4}{r\nu ^{}}}3\right)\right\}+\left({\displaystyle \frac{4}{r^3\nu ^{}}}(1e^\lambda )+{\displaystyle \frac{2e^\lambda }{r^2}}\right)\omega `$ (B91)
$`+{\displaystyle \frac{16\pi e^\lambda }{\nu ^{}}}(\rho _0+p_0)\left(\varpi ^{}{\displaystyle \frac{4e^\nu \varpi }{j^2r^2\nu ^{}}}\right).`$
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# Intertwining operator superalgebras and vertex tensor categories for superconformal algebras, II
## 0 Introduction
It has been known that the $`N=2`$ Neveu-Schwarz superalgebra is one of the most important algebraic objects realized in superstring theory. The $`N=2`$ superconformal field theories constructed from its discrete unitary representations of central charge $`c<3`$ are among the so-called “minimal models.” In the physics literature, there have been many conjectural connections among Calabi-Yau manifolds, Landau-Ginzburg models and these $`N=2`$ unitary minimal models. In fact, the physical construction of mirror manifolds \[GP\] used the conjectured relations \[Ge1\] \[Ge2\] between certain particular Calabi-Yau manifolds and certain $`N=2`$ superconformal field theories (Gepner models) constructed from unitary minimal models (see \[Gr\] for a survey). To establish these conjectures as mathematical theorems, it is necessary to construct the $`N=2`$ unitary minimal models mathematically and to study their structures in detail.
In the present paper, we apply the theory of intertwining operator algebras developed by the first author in \[H3\], \[H4\] and \[H5\] and the tensor product theory for modules for a vertex operator algebra developed by Lepowsky and the first author in \[HL1\]\[HL6\], \[HL8\] and \[H1\] to construct the intertwining operator algebras and vertex tensor categories for $`N=2`$ superconformal unitary minimal models. The main work in this paper is to verify that the conditions to use the general theories are satisfied for these models. The main technique used is the representation theory of the $`N=2`$ Neveu-Schwarz algebras, which has been studied by many physicists and mathematicians, especially by Eholzer and Gaberdiel \[EG\], Feigin, Semikhatov, Sirota and Tipunin \[FSST\] \[FST\], and Adamović \[A1\] \[A2\].
The present paper is organized as follows: In Section 1, we recall the notion of $`N=2`$ superconformal vertex operator superalgebra. In Section 2, we recall and prove some basic results on representations of unitary minimal $`N=2`$ superconformal vertex operator superalgebras and on representations of $`N=2`$ superconformal vertex operator superalgebras in a much more general class. Section 3 is devoted to the proof of the convergence and extension properties for products of intertwining operators for unitary minimal $`N=2`$ superconformal vertex operator superalgebras and for vertex operator superalgebras in the general class. Our main results on the intertwining operator superalgebra structure and vertex tensor category structure are given in Section 4.
Acknowledgment: The research of Y.-Z. H. is supported in part by NSF grant DMS-9622961.
## 1 $`N=2`$ superconformal vertex operator superalgebras
In this section we recall the notion of $`N=2`$ superconformal vertex operator algebra and basic properties of such an algebra. These algebras have been studied extensively by physicists. The following precise version of the definition is from \[A1\] and \[HZ\]:
###### Definition 1.1
An $`N=2`$ superconformal vertex operator superalgebra is a vertex operator superalgebra $`(V,Y,\mathrm{𝟏},\omega )`$ together with odd elements $`\tau ^+`$, $`\tau ^{}`$ and an even element $`\mu `$ satisfying the following axioms: Let
$`Y(\tau ^+,x)`$ $`=`$ $`{\displaystyle \underset{r+\frac{1}{2}}{}}G^+(r)x^{r3/2},`$
$`Y(\tau ^{},x)`$ $`=`$ $`{\displaystyle \underset{r+\frac{1}{2}}{}}G^{}(r)x^{r3/2},`$
$`Y(\mu ,x)`$ $`=`$ $`{\displaystyle \underset{n}{}}J(n)x^{n1}.`$
Then $`V`$ is a direct sum of eigenspaces of $`J(0)`$ with integral eigenvalues which modulo $`2`$ give the $`_2`$ grading for the vertex operator superalgebra structure, and the following $`N=2`$ Neveu-Schwarz relations hold: For $`m,n`$, $`r,s+\frac{1}{2}`$
$`[L(m),L(n)]`$ $`=`$ $`(mn)L(m+n)+{\displaystyle \frac{c}{12}}(m^3m)\delta _{m+n,0},`$
$`[J(m),J(n)]`$ $`=`$ $`{\displaystyle \frac{c}{3}}m\delta _{m+n,0},`$
$`[L(m),J(n)]`$ $`=`$ $`nJ_{m+n},`$
$`[L(m),G^\pm (r)]`$ $`=`$ $`\left({\displaystyle \frac{m}{2}}r\right)G^\pm (m+r),`$
$`[J(m),G^\pm (r)]`$ $`=`$ $`\pm G^\pm (m+r),`$
$`[G^+(r),G^{}(s)]`$ $`=`$ $`2L(r+s)+(rs)J(r+s)+{\displaystyle \frac{c}{3}}(r^2{\displaystyle \frac{1}{4}})\delta _{r+s,0},`$
$`[G^\pm (r),G^\pm (s)]`$ $`=`$ $`0`$
where $`L(m)`$, $`m`$, are the Virasoro operators on $`V`$ and $`c`$ is the central charge of $`V`$.
Modules and intertwining operators for an $`N=1`$ superconformal vertex operator superalgebra are modules and intertwining operators for the underlying vertex operator superalgebra.
The $`N=2`$ superconformal vertex operator superalgebra defined above is denoted by $`(V,Y,\mathrm{𝟏},\tau ^+,\tau ^{},\mu )`$ (without $`\omega `$ since
$$\omega =L(2)\mathrm{𝟏}=\frac{1}{2}[G^+(\frac{3}{2}),G^{}(\frac{1}{2})]\mathrm{𝟏}+\frac{1}{2}J(2)\mathrm{𝟏})$$
or simply by $`V`$. Note that a module $`W`$ for a vertex operator superalgebra (in particular the algebra itself) has a $`_2`$-grading called sign in addition to the $``$-grading by weights. We shall always use $`W^0`$ and $`W^1`$ to denote the even and odd subspaces of $`W`$. If $`W`$ is irreducible, there exists $`h`$ such that $`W=W^0W^1`$ where $`W^0=_{nh+}W_{(n)}`$ and $`W^1=_{nh++1/2}W_{(n)}`$ are the even and odd parts of $`W`$, respectively. We shall always use the notation $`||`$ to denote the map from the union of the even and odd subspaces of a vertex operator superalgebra or of a module for such an algebra to $`_2`$ by taking the signs of elements in the union.
The notion of $`N=2`$ superconformal vertex operator superalgebra can be reformulated using odd formal variables. (In the $`N=1`$ case, this reformulation was given by Barron \[Ba1\] \[Ba2\].)
As in Part I (\[HM\]), for $`l`$ symbols $`\phi _1,\mathrm{},\phi _l`$, consider the exterior algebra of the vector space over $``$ spanned by these symbols and denote this exterior algebra by $`[\phi _1,\mathrm{},\phi _l]`$. For any vector space $`E`$, we also have the vector space
$`E[\phi _1,\mathrm{},\phi _l],`$
$`E[x_1,\mathrm{},x_k][\phi _1,\mathrm{},\phi _l],`$
$`E[x_1,x_1^1,\mathrm{},x_k,x_k^1][\phi _1,\mathrm{},\phi _l],`$
$`E[[x_1,\mathrm{},x_k]][\phi _1,\mathrm{},\phi _l],`$
$`E[[x_1,x_1^1,\mathrm{},x_k,x_k^1]][\phi _1,\mathrm{},\phi _l],`$
$`E\{x_1,\mathrm{},x_k\}[\phi _1,\mathrm{},\phi _l]`$
and
$$E((x_1,\mathrm{},x_k))[\phi _1,\mathrm{},\phi _l].$$
If $`E`$ is a $`_2`$-graded vector space, then there are natural structures of modules over the ring $`[x_1,\mathrm{},x_k][\phi _1,\mathrm{},\phi _l]`$ on these spaces.
Let $`(V,Y,\mathrm{𝟏},\tau ^+,\tau ^{},h)`$ be an $`N=2`$ superconformal vertex operator superalgebra. Let
$`\tau _1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\tau ^++\tau ^{}),`$
$`\tau _2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\tau ^+\tau ^{})`$
and
$`Y(\tau _1,x)`$ $`=`$ $`{\displaystyle \underset{r+\frac{1}{2}}{}}G_1(r)x^{r3/2},`$
$`Y(\tau _2,x)`$ $`=`$ $`{\displaystyle \underset{r+\frac{1}{2}}{}}G_2(r)x^{r3/2}.`$
We define the vertex operator map with odd variables
$`Y:VV`$ $``$ $`V((x))[\phi _1,\phi _2]`$
$`uv`$ $``$ $`Y(u,(x,\phi _1,\phi _2))v`$
by
$`Y(u,(x,\phi _1,\phi _2))v`$ $`=`$ $`Y(u,x)v+\phi _1Y(G_1(1/2)u,x)v`$
$`+\phi _2Y(G_2(1/2)u,x)v`$
$`+\phi _1\phi _2Y(G_1(1/2)G_2(1/2)u,x)v`$
for $`u,vV`$. (We use the same notation $`Y`$ to denote the vertex operator map and the vertex operator map with odd variables.) In particular, we have:
$`Y(h,(x,\phi _1,\phi _2))`$ $`=`$ $`Y(h,x)\sqrt{1}\phi _1Y(\tau _2,x)`$
$`\sqrt{1}\phi _2Y(\tau _1,x)2\sqrt{1}\phi _1\phi _2Y(\omega ,x).`$
Also, if we introduce
$`\phi ^+`$ $`=`$ $`{\displaystyle \frac{\phi _1+\sqrt{1}\phi _2}{\sqrt{2}}}`$
$`\phi ^{}`$ $`=`$ $`{\displaystyle \frac{\phi _1+\sqrt{1}\phi _2}{\sqrt{2}}},`$
then we can write
$`Y(h,(x,\phi _1,\phi _2))`$ $`=`$ $`Y(h,x)+\phi ^+Y(\tau ^+,x)`$
$`+\phi ^{}Y(\tau ^{},x)+2\phi ^+\phi ^{}Y(\omega ,x).`$
We have:
###### Proposition 1.2
The vertex operator map with odd variables satisfies the following properties:
1. The vacuum property:
$$Y(\mathrm{𝟏},(x,\phi _1,\phi _2))=1$$
where $`1`$ on the right-hand side is the identity map on $`V`$.
2. The creation property: For any $`vV`$,
$$Y(v,(x,\phi _1,\phi _2))\mathrm{𝟏}V[[x]][\phi _1,\phi _2],$$
$$\underset{(x,\phi _1,\phi _2)(0,0,0)}{lim}Y(v,(x,\phi _1,\phi _2))\mathrm{𝟏}=v.$$
3. The Jacobi identity: In $`(\text{End}V)[[x_0,x_0^1,x_1,x_1^1,x_2,x_2^1]][\phi _1,\phi _2,\psi _1,\psi _2]`$, we have
$`x_0^1\delta \left({\displaystyle \frac{x_1x_2\phi _1\psi _1\phi _2\psi _2}{x_0}}\right)Y(u,(x_1,\phi _1,\phi _2))Y(v,(x_2,\psi _1,\psi _2))`$ (1.1)
$`(1)^{|u||v|}x_0^1\delta \left({\displaystyle \frac{x_2x_1+\phi _1\psi _1+\phi _2\psi _2}{x_0}}\right)`$
$`Y(v,(x_2,\psi _1,\psi _2))Y(u,(x_1,\phi _1,\phi _2))`$
$`=x_2^1\delta \left({\displaystyle \frac{x_1x_0\phi _1\psi _1\phi _2\psi _2}{x_2}}\right)`$
$`Y(Y(u,(x_0,\phi _1\psi _1,\phi _2\psi _2))v,(x_2,\psi _1,\psi _2)),`$
for $`u,vV`$ which are either even or odd.
4. The $`G_i(1/2)`$-derivative property: For any $`vV`$, $`i=1,2`$,
$$Y(G_i(1/2)v,(x,\phi _1,\phi _2))=\left(\frac{}{\phi _i}+\phi _i\frac{}{x}\right)Y(v,(x,\phi _1,\phi _2)).$$
5. The $`L(1)`$-derivative property: For any $`vV`$,
$$Y(L(1)v,(x,\phi _1,\phi _2))=\frac{}{x}Y(v,(x,\phi _1,\phi _2)).$$
6. The skew-symmetry: For any $`u,vV`$ which are either even or odd,
$`Y(u,(x,\phi _1,\phi _2))v`$
$`=(1)^{|u||v|}e^{xL(1)+\phi _1G_1(1/2)+\phi _2G_2(1/2)}Y(v,(x,\phi _1,\phi _2))u.\text{ }`$
The proof of this result is straightforward and we omit it.
As in the $`N=1`$ case \[HM\], we can also reformulate the data and axioms for modules and intertwining operators for an $`N=2`$ superconformal vertex operator superalgebra using odd variables. Here we give the details of the corresponding reformulation of the data and axioms for intertwining operators.
Let $`W_1`$, $`W_2`$ and $`W_3`$ be modules for an $`N=2`$ superconformal vertex operator superalgebra $`V`$ and $`𝒴`$ an intertwining operator of type $`\left(\genfrac{}{}{0pt}{}{W_3}{W_1W_2}\right)`$. We define the corresponding intertwining operator map with odd variable
$`𝒴:W_1W_2`$ $``$ $`W_3\{x\}[\phi _1,\phi _2]`$
$`w_{(1)}w_{(2)}`$ $``$ $`𝒴(w_{(1)},(x,\phi _1,\phi _2))w_{(2)}`$
by
$`𝒴(w_{(1)},(x,\phi _1,\phi _2))w_{(2)}`$ $`=`$ $`𝒴(w_{(1)},x)w_{(2)}+\phi _1𝒴(G_1(1/2)w_{(1)},x)w_{(2)}`$
$`+\phi _2𝒴(G_2(1/2)w_{(1)},x)w_{(2)}`$
$`+\phi _1\phi _2𝒴(G_1(1/2)G_2(1/2)w_{(1)},x)w_{(2)}`$
for $`u,vV`$. Then we have:
###### Proposition 1.3
The intertwining operator map with odd variable satisfies the following properties:
1. The Jacobi identity: In $`\text{Hom}(W_1W_2,W_3)\{x_0,x_1,x_2\}[\phi _1,\phi _2,\psi _1,\psi _2]`$, we have
$`x_0^1\delta \left({\displaystyle \frac{x_1x_2\phi _1\psi _1\phi _2\psi _2}{x_0}}\right)Y(u,(x_1,\phi _1,\phi _2))𝒴(w_{(1)},(x_2,\psi _1,\psi _2))`$
$`(1)^{|u||w_{(1)}|}x_0^1\delta \left({\displaystyle \frac{x_2x_1+\phi _1\psi _1+\phi _2\psi _2}{x_0}}\right)`$
$`𝒴(w_{(1)},(x_2,\psi _1,\psi _2))Y(u,(x_1,\phi _1,\phi _2))`$
$`=x_2^1\delta \left({\displaystyle \frac{x_1x_0\phi _1\psi _1\phi _2\psi _2}{x_2}}\right)`$
$`𝒴(Y(u,(x_0,\phi _1\psi _1,\phi _2\psi _2))w_{(1)},(x_2,\psi _1,\psi _2)),`$
for $`uV`$ and $`w_{(1)}W_1`$ which are either even or odd.
2. The $`G_i(1/2)`$-derivative property: For any $`vV`$, $`i=1,2`$,
$$𝒴(G_i(1/2)v,(x,\phi _1,\phi _2))=\left(\frac{}{\phi _i}+\phi _i\frac{}{x}\right)𝒴(v,(x,\phi _1,\phi _2)).$$
3. The $`L(1)`$-derivative property: For any $`vV`$,
$$𝒴(L(1)w_{(1)},(x,\phi _1,\phi _2))=\frac{}{x}𝒴(w_{(1)},(x,\phi _1,\phi _2)).$$
4. The skew-symmetry: There is a linear isomorphism
$$\mathrm{\Omega }:𝒱_{W_1W_2}^{W_3}𝒱_{W_2W_1}^{W_3}$$
such that
$`\mathrm{\Omega }(𝒴)(w_{(1)},(x,\phi _1,\phi _2))w_{(2)}`$
$`=(1)^{|w_{(1)}||w_{(2)}|}e^{xL(1)+\phi _1G_1(1/2)+\phi _2G_2(1/2)}`$
$`𝒴(w_{(2)},(e^{\pi i}x,\phi _1,\phi _2)))w_{(1)}`$
for $`w_{(1)}W_1`$ and $`w_{(2)}W_2`$ which are either even or odd.
The proof of this result is similar to the proof of Proposition 1.2 and is omitted.
## 2 Unitary minimal $`N=2`$ superconformal vertex operator superalgebras
In this section, we recall the constructions and results on unitary minimal $`N=2`$ superconformal vertex operator superalgebras and their representations. New results needed in later sections are also proved. We then introduce in this section a class of $`N=2`$ superconformal vertex operator superalgebras and generalize most of the results for unitary minimal $`N=2`$ superconformal vertex operator superalgebras to algebras in this class.
The $`N=2`$ Neveu-Schwarz Lie superalgebra is the vector space
$$_nL_n_{r+\frac{1}{2}}G_r^+_{r+\frac{1}{2}}G_r^{}_nJ_nC$$
equipped with the following $`N=2`$ Neveu-Schwarz relations:
$`[L_m,L_n]`$ $`=`$ $`(mn)L_{m+n}+{\displaystyle \frac{C}{12}}(m^3m)\delta _{m+n,0},`$
$`[J_m,J_n]`$ $`=`$ $`{\displaystyle \frac{c}{3}}m\delta _{m+n,0},`$
$`[L_m,J_n]`$ $`=`$ $`nJ_{m+n},`$
$`[L_m,G_r^\pm ]`$ $`=`$ $`\left({\displaystyle \frac{m}{2}}r\right)G_{m+r}^\pm ,`$
$`[J_m,G_r^\pm ]`$ $`=`$ $`\pm G_{m+r}^\pm ,`$
$`[G_r^+,G_s^{}]`$ $`=`$ $`2L_{r+s}+(rs)J_{r+s}+{\displaystyle \frac{C}{3}}(r^2{\displaystyle \frac{1}{4}})\delta _{r+s,0},`$
$`[G_r^\pm ,G_s^\pm ]`$ $`=`$ $`0`$
for $`m,n`$, $`r,s+\frac{1}{2}`$. For simplicity, we shall simply denote the $`N=2`$ Neveu-Schwarz Lie superalgebra by $`𝔫𝔰(2)`$ in this paper.
We now construct representations of the $`N=2`$ Neveu-Schwarz Lie superalgebra. Consider the two subalgebras
$`𝔫𝔰^+(2)`$ $`=`$ $`_{n>0}L_n_{r>0}G_r^+_{r>0}G_r^{}_{n>0}J_n,`$
$`𝔫𝔰^{}(2)`$ $`=`$ $`_{n<0}L_n_{r<0}G_r^+_{r<0}G_r^{}_{n<0}J_n`$
of $`𝔫𝔰(2)`$. Let $`U()`$ be the functor from the category of Lie superalgebras to the category of associative algebras obtained by taking the universal enveloping algebras of Lie superalgebras. For any representation of $`𝔫𝔰(2)`$, we shall use $`L(n)`$, $`n`$, $`G^\pm (r)`$, $`r+\frac{1}{2}`$, and $`J(n)`$, $`n`$, to denote the representation images of $`L_n`$, $`G^\pm (r)`$ and $`J_n`$.
For any $`c,h,q`$, the Verma module $`M_{𝔫𝔰(2)}(c,h,q)`$ for $`𝔫𝔰(2)`$ is a free $`U(𝔫𝔰^{}(2))`$-module generated by $`\mathrm{𝟏}_{c,h,q}`$ such that
$`𝔫𝔰^+(2)\mathrm{𝟏}_{c,h,q}`$ $`=`$ $`0,`$
$`L(0)\mathrm{𝟏}_{c,h,q}`$ $`=`$ $`h\mathrm{𝟏}_{c,h,q},`$
$`C\mathrm{𝟏}_{c,h,q}`$ $`=`$ $`c\mathrm{𝟏}_{c,h,q},`$
$`J(0)\mathrm{𝟏}_{c,h,q}`$ $`=`$ $`q\mathrm{𝟏}_{c,h,q}.`$
There exists a unique maximal proper submodule $`J_{𝔫𝔰(2)}(c,h,q)`$ of $`M_{𝔫𝔰(2)}(c,h,q)`$. It is easy to see that when $`c0`$, $`\mathrm{𝟏}_{c,0,0}`$, $`G^\pm (3/2)\mathrm{𝟏}_{c,0,0}`$ and $`L(2)\mathrm{𝟏}_{c,0,0}`$ are not in $`J_{𝔫𝔰(2)}(c,0,0)`$. Let
$$L_{𝔫𝔰(2)}(c,h,q)=M_{𝔫𝔰(2)}(c,h,q)/J_{𝔫𝔰(2)}(c,h,q)$$
and
$$V_{𝔫𝔰(2)}(c,0,0)=M_{𝔫𝔰(2)}(c,0,0)/G^+(1/2)\mathrm{𝟏}_{c,0,0},G^{}(1/2)\mathrm{𝟏}_{c,0,0}$$
where $`G^+(1/2)\mathrm{𝟏}_{c,0,0},G^{}(1/2)\mathrm{𝟏}_{c,0,0}`$ is the submodule of $`M_{𝔫𝔰(2)}(c,0,0)`$ generated by $`G^\pm (1/2)\mathrm{𝟏}_{c,0,0}`$. Then $`L_{𝔫𝔰(2)}(c,0,0)`$ and $`V_{𝔫𝔰(2)}(c,0,0)`$ have the structures of vertex operator superalgebras with the vacuum $`\mathrm{𝟏}_{c,0,0}`$, the Neveu-Schwarz elements $`G^\pm (3/2)\mathrm{𝟏}_{c,0,0}`$ and the Virasoro element $`L(2)\mathrm{𝟏}_{c,0,0}`$ (see \[A1\]).
In \[EG\], Eholzer and Gaberdiel showed, among other things, that among vertex operator superalgebras of the form $`L_{𝔫𝔰(2)}(c,0,0)`$, the only ones having finitely many irreducible modules are the “unitary” ones $`L_{𝔫𝔰(2)}(c_m,0,0)`$ for nonnegative integers $`m`$, where $`c_m=\frac{3m}{m+2}`$. The following result was proved by Adamović in \[A1\] and \[A2\] using the results obtained Adamović and Milas \[AM\], Feigin, Semikhatov and Tipunin \[FST\] and Doerrzapf \[D\]:
###### Theorem 2.1
The vertex operator superalgebra $`L_{𝔫𝔰(2)}(c,0,0)`$ has finitely many irreducible modules and every module for $`L_{𝔫𝔰(2)}(c,0,0)`$ is completely reducible if and only if
$$c=c_m=\frac{3m}{m+2}$$
where $`m`$ is a nonnegative integer. A set of representatives of the equivalence classes of irreducible modules for $`L_{𝔫𝔰(2)}(c_m,0,0)`$ is
$$\{L_{𝔫𝔰(2)}(c_m,h_m^{j,k},q_m^{j,k})\}_{j,k_{\frac{1}{2}},\mathrm{\hspace{0.33em}0}j,k,j+k<m}$$
where $`_{\frac{1}{2}}=\{\frac{1}{2},\frac{3}{2},\frac{5}{2},\mathrm{}\}`$ and
$`h_m^{j,k}`$ $`=`$ $`{\displaystyle \frac{jk\frac{1}{4}}{m+2}},`$
$`q_m^{j,k}`$ $`=`$ $`{\displaystyle \frac{jk}{m+2}}\text{ }`$
For any $`m0`$, we call the vertex operator algebra $`L_{𝔫𝔰(2)}(c_m,0,0)`$ a unitary minimal $`N=2`$ superconformal vertex operator superalgebra.
###### Proposition 2.2
Let $`j_i,k_i_{\frac{1}{2}}`$, $`i=1,2,3`$, satisfying $`0j_i,k_i,j_i+k_i<m`$ and $`𝒴`$ an intertwining operator of type
$$\left(\genfrac{}{}{0pt}{}{L_{𝔫𝔰(2)}(c_m,h_m^{j_3,k_3},q_m^{j_3,k_3})}{L_{𝔫𝔰(2)}(c_m,h_m^{j_1,k_1},q_m^{j_1,k_1})L_{𝔫𝔰(2)}(c_m,h_m^{j_2,k_2},q_m^{j_2,k_2})}\right).$$
(2.1)
Then we have:
1. For any $`w_{(1)}L_{𝔫𝔰(2)}(c_m,h_m^{j_1,k_1},q_m^{j_1,k_1})`$ and $`w_{(2)}L_{𝔫𝔰(2)}(c_m,h_m^{j_2,k_2},q_m^{j_2,k_2})`$,
$$𝒴(w_{(1)},x)w_{(2)}x^{h_m^{j_3,k_3}h_m^{j_1,k_1}h_m^{j_2,k_2}}L_{𝔫𝔰(2)}(c_m,h_m^{j_3,k_3},q_m^{j_3,k_3})((x^{1/2})).$$
2. Let $`\mathrm{\Delta }=h_m^{j_3,k_3}h_m^{j_1,k_1}h_m^{j_2,k_2}`$ and $`w_{(i)}=\mathrm{𝟏}_{c_m,h_m^{j_i,k_i},q_m^{j_i,k_i}}`$, $`i=1,2,3`$, the lowest weight vectors in $`L_{𝔫𝔰(2)}(c_m,h_m^{j_i,k_i},q_m^{j_i,k_i})`$. Then the map $`𝒴`$ is uniquely determined by the maps
$`(w_{(1)})_{\mathrm{\Delta }1},`$
$`(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2},`$
$`(G^{}(1/2)w_{(1)})_{\mathrm{\Delta }1/2},`$
$`(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }`$
from the one-dimensional subspace of $`W_2`$ spanned by $`w_{(2)}`$ to the one-dimensional subspace of $`W_3`$ spanned by $`w_{(3)}`$. That is, if these maps are $`0`$, then $`𝒴=0`$.
3. If $`q_m^{j_3,k_3}`$ is not equal to one of the numbers $`q_m^{j_1,k_1}+q_m^{j_2,k_2}`$, $`q_m^{j_1,k_1}+q_m^{j_2,k_2}1`$ and $`q_m^{j_1,k_1}+q_m^{j_2,k_2}+1`$, then the space
$$𝒱_{L_{𝔫𝔰(2)}(c_m,h_m^{j_1,k_1},q_m^{j_1,k_1})L_{𝔫𝔰(2)}(c_m,h_m^{j_2,k_2},q_m^{j_2,k_2})}^{L_{𝔫𝔰(2)}(c_m,h_m^{j_3,k_3},q_m^{j_3,k_3})}$$
of intertwining operators of type (2.1) is $`0`$. If $`q_m^{j_3,k_3}=q_m^{j_1,k_1}+q_m^{j_2,k_2}\pm 1`$, it is at most $`1`$-dimensional. If $`q_m^{j_3,k_3}=q_m^{j_1,k_1}+q_m^{j_2,k_2}`$, it is at most $`2`$-dimensional.
Proof. Conclusion 1 is clear since the three modules are irreducible.
Conclusion 2 can be proved similarly to the proof of the similar statement in the $`N=1`$ case in \[HM\]. Here we give a different proof.
Suppose that
$`(w_{(1)})_{\mathrm{\Delta }1}w_{(2)},`$ (2.2)
$`(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)},`$ (2.3)
$`(G^{}(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)},`$ (2.4)
$`(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }w_{(2)}`$ (2.5)
are all equal to $`0`$ but $`𝒴0`$.
Using the associator formula (obtained by taking residue in $`x_1`$ in the Jacobi identity defining intertwining operators)
$`𝒴(Y(u,x_0)w,x_2)`$ $`=`$ $`Y(u,x_0+x_2)𝒴(w,x_2)`$
$`(1)^{|u||w|}\text{Res}_{x_1}x_0^1\delta \left({\displaystyle \frac{x_2x_1}{x_0}}\right)𝒴(w,x_2)Y(u,x_1)`$
repeatedly, we see that $`𝒴=0`$ if $`𝒴(w_{(1)},x)=0`$. Thus $`𝒴(w_{(1)},x)0`$.
Using the commutator formula (obtained by taking residue in $`x_0`$ in the Jacobi identity defining intertwining operators)
$`Y(u,x_1)𝒴(w,x_2)(1)^{|u||w|}𝒴(w,x_2)Y(w,x_1)`$
$`=\text{Res}_{x_0}x_2^1\delta \left({\displaystyle \frac{x_1x_0}{x_2}}\right)𝒴(Y(u,x_0)w,x_2)`$
together with the $`N=2`$ Neveu-Schwarz algebra relations, the $`L(1)`$-derivative property and the definition of the lowest weight vector $`w_{(1)}`$ repeatedly, we see that $`𝒴(w_{(1)},x)=0`$ if (2.2)–(2.5) are all equal to $`0`$. Thus these four vectors cannot be all $`0`$. We have a contradiction.
We prove Conclusion 3 now. By Conclusion 2, we need only estimate the number of nonzero vectors in the set of four vectors (2.2)–(2.5).
We need the following:
###### Lemma 2.3
The following equality hold:
$$q_m^{j_3,k_3}(w_{(1)})_{\mathrm{\Delta }1}w_{(2)}=(q_m^{j_1,k_1}+q_m^{j_2,k_2})(w_{(1)})_{\mathrm{\Delta }1}w_{(2)},$$
(2.6)
$`q_m^{j_3,k_3}(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)}`$ (2.7)
$`=(q_m^{j_1,k_1}+q_m^{j_2,k_2}+1)(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)},`$
$`q_m^{j_3,k_3}(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)}`$ (2.8)
$`=(q_m^{j_1,k_1}+q_m^{j_2,k_2}1)(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)},`$
$`q_m^{j_3,k_3}(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }w_{(2)}`$ (2.9)
$`=(q_m^{j_1,k_1}+q_m^{j_2,k_2})(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }w_{(2)}`$
$`+(2h_m^{j_1,k_1}+q_m^{j_1,k_1})(w_{(1)})_{\mathrm{\Delta }1}w_{(2)}.`$
Proof. A straightforward calculation gives
$$J(0)(w_{(1)})_{\mathrm{\Delta }1}w_{(2)}=(q_m^{j_1,k_1}+q_m^{j_2,k_2})(w_{(1)})_{\mathrm{\Delta }1}w_{(2)},$$
$`J(0)(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)}`$
$`=(q_m^{j_1,k_1}+q_m^{j_2,k_2}+1)(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)},`$
$`J(0)(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)}`$
$`=(q_m^{j_1,k_1}+q_m^{j_2,k_2}1)(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_{(2)},`$
$`J(0)(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }w_{(2)}`$
$`=(q_m^{j_1,k_1}+q_m^{j_2,k_2})(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }w_{(2)}`$
$`+(2h_m^{j_1,k_1}+q_m^{j_1,k_1})(w_{(1)})_{\mathrm{\Delta }1}w_{(2)}.`$
But on the other hand, note that (2.2)–(2.5) are all (zero or nonzero) constant multiple of $`w_{(3)}`$ and thus all have $`U(1)`$ charge $`q_m^{j_3,k_3}`$. So we have (2.6)–(2.9).
We prove Conclusion 3 using this lemma now. If $`q_m^{j_3,k_3}`$ is not equal to one of the numbers $`q_m^{j_1,k_1}+q_m^{j_2,k_2}`$ $`q_m^{j_1,k_1}+q_m^{j_2,k_2}1`$ and $`q_m^{j_1,k_1}+q_m^{j_2,k_2}+1`$, then from (2.6)–(2.9), we conclude that (2.2)–(2.5) are all equal to $`0`$. Thus the space of intertwining operators is $`0`$.
If $`q_m^{j_3,k_3}=q_m^{j_1,k_1}+q_m^{j_2,k_2}+1`$, then by (2.6), (2.8) and (2.9), (2.2), (2.4) and (2.5) must be $`0`$. Thus there is at most one nonzero vector (2.3). So the dimension is at most $`1`$.
Similarly in the case of $`q_m^{j_3,k_3}=q_m^{j_1,k_1}+q_m^{j_2,k_2}1`$, we can show that $`(w_{(1)})_{\mathrm{\Delta }1}w_{(2)}`$, (2.2), (2.3) and (2.5) must be $`0`$ and thus the dimension is at most $`1`$.
If $`q_m^{j_3,k_3}=q_m^{j_1,k_1}+q_m^{j_2,k_2}`$, then by (2.7) and (2.8), (2.3) and (2.4) must both be $`0`$. Thus we have at most two nonzero vecotrs (2.2) and (2.5) and the dimension is at most $`2`$.
###### Remark 2.4
In the proofs above we do not use the particular properties, except the irreducibility, of $`L(c_m,h_m^{j_i,k_i},q_m^{j_i,k_i})`$, $`i=1,2,3`$. Thus the conclusions of Proposition 2.2 remain true if we replace $`c_m`$ by an arbitrary $`c`$ and $`L(c_m,h_m^{j_i,k_i},q_m^{j_i,k_i})`$, $`i=1,2,3`$ by $`L(c,0,0)`$-modules $`L(c,h_i,q_i)`$, $`i=1,2,3`$, if $`L(c,h_i,q_i)`$ are irreducible.
###### Definition 2.5
An irreducible module for $`L_{𝔫𝔰(2)}(c,0,0)`$ is said to be chiral (anti-chiral) if
$$G^+(1/2)w=0,(G^{}(1/2)w=0)$$
where $`w`$ is a nonzero lowest weight vector of the module.
Note that in the case $`c=c_m`$ we have only finitely many chiral (anti-chiral) modules.
###### Corollary 2.6
Assume that $`L_{𝔫𝔰(2)}(c_m,h_m^{j_1,k_1},q_m^{j_1,k_1})`$ is chiral or anti-chiral. Then the dimension of the space
$$𝒱_{L_{𝔫𝔰(2)}(c_m,h_m^{j_1,k_1})L_{𝔫𝔰(2)}(c_m,h_m^{j_2,k_2})}^{L_{𝔫𝔰(2)}(c_m,h_m^{j_3,k_3})}$$
is at most $`1`$.
Proof. Assume that $`L_{𝔫𝔰(2)}(c_m,h_m^{j_1,k_1},q_m^{j_1,k_1})`$ is chiral. Then
$$(G^+(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_2=0.$$
We claim that in this case
$$(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }w_{(2)}=2(w_{(1)})_{\mathrm{\Delta }1}w_{(2)}.$$
(2.10)
In fact, the commutator formula for $`G^+(1/2)`$ and $`G^{}(1/2)`$ gives
$$G^+(1/2)G^{}(1/2)=G^{}(1/2)G^+(1/2)+2L(1).$$
Thus
$`(G^+(1/2)G^{}(1/2)w_{(1)})_\mathrm{\Delta }w_2`$ $`=`$ $`(G^{}(1/2)G^+(1/2)w_{(1)})_\mathrm{\Delta }w_2`$ (2.11)
$`+2(L(1)w_{(1)})_\mathrm{\Delta }w_2`$
$`=`$ $`2(L(1)w_{(1)})_\mathrm{\Delta }w_2.`$
From the $`L(1)`$-derivative property for intertwining operators, we obtain
$$(L(1)w_{(1)})_\mathrm{\Delta }=(w_{(1)})_{\mathrm{\Delta }1}.$$
Thus the right-hand side of (2.11) becomes $`2(w_{(1)})_{\mathrm{\Delta }1}w_{(2)}`$, proving (2.10). On the other hand, from (2.6) and (2.8), we see that the vectors $`(w_{(1)})_{\mathrm{\Delta }1}w_2`$ and $`(G^{}(1/2)w_{(1)})_{\mathrm{\Delta }1/2}w_2`$ cannot be nonzero at the same time. Thus in this case, the dimension of the space spanned by the four vectors (2.2)–(2.5) is at most $`1`$. Equivalently, the corollary is proved.
###### Remark 2.7
Note that the operator $`J(0)`$ plays an essential role in the proofs of Conclusion 3 in Proposition 2.2 and Corollary 2.6.
###### Remark 2.8
After the first version of the present paper was finished, we received a preprint \[A3\] from Adamović in which the fusion rules for $`L_{𝔫𝔰(2)}(c_m,0,0)`$ are calculated explicitly and a stronger complete reducibility theorem is proved. But for the purpose of the present paper, we shall not need these stronger results.
Combining Theorem 2.1 and the second or third conclusion of Proposition 2.2, we obtain:
###### Corollary 2.9
The unitary minimal $`N=2`$ superconformal vertex operator superalgebras are rational in the sense of \[HL1\], that is, the following three conditions are satisfied:
1. Every module for such an algebra is completely reducible.
2. There are only finitely many inequivalent irreducible modules for such an algebra.
3. The fusion rules among any three (irreducible) modules are finite.
We also have:
###### Proposition 2.10
Any finitely-generated lower truncated generalized module $`W`$ for $`L_{𝔫𝔰(2)}(c_m,0,0)`$ is an ordinary module.
Proof. The proof is the same as the corresponding result in \[HM\]. We repeat it here since it is simple. Suppose that $`W`$ is generated by a single vector $`wW`$. Then by the Poincaré-Birkhoff-Witt theorem and the lower truncation condition, every homogeneous subspace of $`U(𝔫𝔰(2))w`$ is finite-dimensional, proving the result.
Let $`m_i`$, $`i=1,\mathrm{},n`$, be positive integers and let $`V=L_{𝔫𝔰(2)}(c_{m_1},0,0)\mathrm{}L_{𝔫𝔰(2)}(c_{m_n},0,0)`$. From the trivial generalizations of the results proved in \[FHL\] and \[DMZ\] to vertex operator superalgebras, $`V`$ is a rational $`N=2`$ superconformal vertex operator superalgebra, a set of representatives of equivalence classes of irreducible modules for $`V`$ can be listed explicitly and the fusion rules for $`V`$ and can be calculated easily.
We introduce a class of $`N=2`$ superconformal vertex operator vertex operator superalgebras:
###### Definition 2.11
Let $`m_i`$, $`i=1,\mathrm{},n`$, be positive integers. An $`N=2`$ superconformal vertex operator superalgebra $`V`$ is said to be in the class $`𝒞_{m_1;\mathrm{};m_n}`$ if $`V`$ has a vertex operator subalgebra isomorphic to $`L_{𝔫𝔰(2)}(c_{m_1},0,0)\mathrm{}L_{𝔫𝔰(2)}(c_{m_n},0,0)`$.
###### Proposition 2.12
Let $`V`$ be an $`N=2`$ superconformal vertex operator superalgebra in the class $`𝒞_{m_1;\mathrm{};m_n}`$. Then any finitely-generated lower truncated generalized $`V`$-module $`W`$ is an ordinary module.
Proof. The proof is similar to the proofs of Proposition 3.7 in \[H2\] and Proposition 2.7 in \[HM\]. Here we only point out the main difference. As in \[H2\] and \[HM\], we discuss only the case $`n=2`$. Similar to the proofs of Proposition 3.7 in \[H2\] and Proposition 2.7 in \[HM\], using the Jacobi identity, the $`N=2`$ Neveu-Schwarz relations, in particular, the formulas $`[G^+(1/2),G^{}(1/2)]=2L(1)`$, $`(G^+(1/2))^2=(G^{}(1/2))^2=0`$, and Theorem 4.7.4 of \[FHL\], we can reduce our proof in the case of $`n=2`$ to the finite-dimensionality of the space spanned by the elements of the form
$`A(L(1)^{l_1}G^+(1/2)^{k_1}G^{}(1/2)^{k_2}u_{(j)}^{(1)})_{j_1}Bw_{(t)}^{(1)}`$
$`C(L(1)^{l_2}G^+(1/2)^{k_3}G^{}(1/2)^{k_4}u_{(j)}^{(2)})_{j_2}Dw_{(t)}^{(2)},`$
$`l_1,l_2`$, $`k_1,k_2,k_3,k_4=0,1`$, $`j_1,j_2`$, $`t=1,\mathrm{},c`$, $`j=1,\mathrm{},d`$, where $`A`$ (or $`C`$) are products of operators of the forms $`L(a)_1)`$ (or $`L(a)_2`$), $`J(a)_1`$ (or $`J(a)_2`$), $`a_+`$, and $`G^\pm (b)_1`$ (or $`G^\pm (b)_2`$), $`b/2`$, $`B`$ (or $`D`$) are products of operators of the forms $`L(a)_1`$ (or $`L(a)_2`$), $`J(a)_1`$ (or $`J(a)_2`$), $`a_+`$, and $`G^\pm (b)_1`$ (or $`G^\pm (b)_2`$), $`b+\frac{1}{2}`$, where $`u_{(j)}^{(i)}`$, $`j=1,\mathrm{},d`$, $`i=1,2`$, are elements of $`V`$ such that the $`L_{𝔫𝔰(2)}(c_{m_i},0,0)`$-submodules generated by them are isomorphic to $`L(c_{m_i},h_{m_i}^{r_j,s_j},q_{m_i}^{r_j,s_j})`$ for some $`r_j,s_j_{\frac{1}{2}}`$ satisfying $`0r_j,s_j,r_j+s_j<m_i`$ with the images of $`u_{(j)}^{(i)}`$, $`j=1,\mathrm{},d`$, $`i=1,2`$, as the lowest weight vectors and such that $`V`$ is isomorphic to the direct sum of these submodules, and where $`w_{(t)}^{(i)}`$, $`t=1,\mathrm{},c`$, $`i=1,2`$, are homogeneous elements of some irreducible $`L_{𝔫𝔰(2)}(c_{m_i},0,0)`$-modules. The remaining argument in the proof is the same as the corresponding parts in \[H2\] and \[HM\] and is omitted here.
## 3 The convergence and the extension property for products of intertwining operators
In this section, we study products of intertwining operators for the unitary minimal $`N=2`$ superconformal vertex operator superalgebras. The main result is the folllwing:
###### Theorem 3.1
Let $`m`$ be a positive integer. Then intertwining operators for the vertex operator superalgebra $`L(c_m,0,0)`$ satisfy the convergence and extension property for products of intertwining operators introduced in \[H1\].
An immediate consequence is a similar result for the vertex operator algebras in the class $`𝒞_{m_1;\mathrm{};m_n}`$. See Theorem 3.6.
Instead of proving Theorem 3.1 by deriving differential equations with regular singularities satisfied by the matrix elemenets of products of intertwining operators of lowest weight vectors, as in \[H1\], \[HL7\] and \[HM\], we use the so-called anti-Kazama-Suzuki mapping \[FST\], which reduces the problem to the study of intertwining operators for a vertex operator algebra constructed from an $`\widehat{𝔰𝔩}_2`$ integrable lowest weight representation.
First we give some auxiliary constructions and discuss some results obtained using the so-called “anti-Kazama-Suzuki mapping” (introduced in \[FST\]).
Fix a positive integer $`m`$. As before, we let $`c_m=\frac{3m}{m+2}`$ and denote by $`L_{𝔫𝔰(2)}(c_m,0,0)`$ the unitary minimal $`N=2`$ superconformal vertex operator superalgebra.
Let $`L`$ be a rank one lattice generated by $`\alpha `$ with the bilinear form $`,`$ given by
$$\alpha ,\alpha =1,$$
and let $`l=L_{}`$. As in \[FLM\], we have a vertex superalgebra
$$V_LS(\widehat{l}_{})[L].$$
Note that $`V_L`$ is super since $`L`$ is odd and $`V_L`$ does not satisfy the grading-restriction conditions for the grading obtained from the usual Virasoro element for lattice vertex algebras since the bilinear form is not positive definite.
Let
$$\stackrel{~}{\omega }_{V_L}=\frac{\alpha (1)^2}{2}+\frac{\alpha (2)}{2}.$$
It can be verified easily that the component operators of the vertex operator associated to $`\stackrel{~}{\omega }_{V_L}`$ satisfy the Virasoro relations with central charge $`4`$. In particular, the component operator for the $`2`$-nd power of $`x`$ gives a $`/2`$-grading for $`V_L`$. With this grading, $`V_L`$ is a $`/2`$-graded vertex superalgebra. However, it is easy to see that this grading is not truncated from below. Thus $`V_L`$ with $`\stackrel{~}{\omega }_{V_L}`$ as its Virasoro element fails to be a vertex operator superalgebra.
We also need the following construction of the so-called “Liouville scalar model:” Let $`\widehat{}`$ be the Lie algebra with a basis $`a(n)`$, $`n`$, and $`d`$, satisfying the bracket relations
$`[a(m),a(n)]`$ $`=`$ $`m\delta _{m+n,0}d,`$
$`[a(m),d]`$ $`=`$ $`0`$
for $`m,n`$. Let $`M(1,s)`$ be the corresponding irreducible highest weight module with central charge $`1`$ and highest weight $`s`$. It is well-known that $`M(1,0)=S(\widehat{})`$ has a vertex operator algebra structure with the Virasoro element $`\omega _{M(1,0)}=\frac{a(1)^2}{2}`$. Consider the vertex algebra structure on $`M(1,0)`$ together with a different Virasoro element
$$\stackrel{~}{\omega }_{M(1,0)}=\frac{a(1)^2}{2}+\frac{ia(2)}{2}.$$
A straightforward calculation shows that the vertex algebra structure on $`M(1,0)`$ together with the Virasoro element $`\stackrel{~}{\omega }_{M(1,0)}`$ is a vertex operator algebra with the central charge $`4`$ (see, for example, \[L\] and \[FF\]). We shall denote this vertex operator algebra by $`V_{\mathrm{Liou}}`$ (the vertex operator algebra associated to the Liouville scalar model). In addition, every $`M(1,s)`$, $`s`$, is an irreducible $`V_{\mathrm{Liou}}`$-module and any $`V_{\mathrm{Liou}}`$-module on which $`a(0)`$ acts semisimply is completely reducible. We will work only with such modules $`V_{\mathrm{Liou}}`$-module, which are enough for our purposes.
The anti-Kazama-Suzuki mapping gives us a structure of an $`\widehat{𝔰𝔩}_2`$-module on $`L_{𝔫𝔰(2)}(c_m,h_m^{j,k},q_m^{j,k})V_L`$ for $`j,k_{\frac{1}{2}}`$, $`0<j,k,j+k<m`$. Consider the vectors (as in \[FST\] and \[A1\])
$`𝔢`$ $`=`$ $`G^+(3/2)\mathrm{𝟏}_{c_m,0,0}e^\alpha ,`$
$`𝔣`$ $`=`$ $`{\displaystyle \frac{m+2}{2}}G^{}(3/2)\mathrm{𝟏}_{c_m,0,0}e^\alpha ,`$
$`𝔥`$ $`=`$ $`m\mathrm{𝟏}_{c_m,0,0}\alpha (1)+(m+2)J(1)\mathrm{𝟏}_{c_m,0,0}e^0`$
in $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$. Then the vertex operators $`Y(𝔢,x)`$, $`Y(𝔣,x)`$ and $`Y(𝔥,x)`$ for the $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$-module $`L_{𝔫𝔰(2)}(c_m,h_m^{j,k},q_m^{j,k})V_L`$ give a representation of $`\widehat{𝔰𝔩}_2`$ of level $`m`$ on $`L_{𝔫𝔰(2)}(c_m,h_m^{j,k},q_m^{j,k})V_L`$. The main observation in \[FST\] is that $`L_{𝔫𝔰(2)}(c_m,h_m^{j,k},q_m^{j,k})V_L`$ is completely reducible as an $`\widehat{𝔰𝔩}_2`$-module. In the special case $`h_m^{j,k}=q_m^{j,k}=0`$, we obtain a vertex subalgebra of $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$ which is isomorphic as a vertex algebra to the underlying vertex algebra of the vertex operator algebra $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ on the integrable highest-weight $`𝔰𝔩(2)`$-module of level $`m`$ and highest weight $`0`$ (as in \[A1\]). The Virasoro elelement for $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ is given by the Sugawara-Segal construction, and if we identify this vertex subalgebra with $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$, then
$`\omega _{L_{\widehat{𝔰𝔩}(2)}(m,0)}`$ $`=`$ $`\omega _{L_{𝔫𝔰(2)}(c_m,0,0)}e^0+{\displaystyle \frac{m+2}{4}}J(1)^2\mathrm{𝟏}_{c_m,0,0}e^0`$
$`{\displaystyle \frac{m+2}{2}}J(1)\mathrm{𝟏}_{c_m,0,0}\alpha (1)+{\displaystyle \frac{m}{4}}\mathrm{𝟏}_{c_m,0,0}\alpha (1)^2,`$
where $`\omega _{L_{\widehat{𝔰𝔩}(2)}(m,0)}`$ and $`\omega _{L_{𝔫𝔰(2)}(c_m,0,0)}`$ are the Virasoro elements for $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ and $`L_{𝔫𝔰(2)}(c_m,0,0)`$, respectively. We see that under the isomorphism from this vertex subalgebra of $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$ to $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$, the Virasoro element in $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ is not the image of the Virasoro element of $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$.
To get the correct Virasoro element (as in \[FST\]), we consider the vertex subalgebra of $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$ generated by the element
$$\rho =\sqrt{\frac{m+2}{2}}(J(1)\mathrm{𝟏}_{c_m,0,0}e^0\mathrm{𝟏}_{c_m,0,0}\alpha (1)).$$
It is straightforward to verify that this vertex subalgebra is actually isomorphic to $`V_{\mathrm{Liou}}`$. Straightforward calculations also show that $`Y(\rho ,x)`$ commutes with $`\widehat{𝔰𝔩}_2`$ generators. So $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$ is isomorphic to a vertex subalgebra of $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$ as well and we shall, for convenience, identify $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$ with this vertex subalgebra. It is easy to see that the Virasoro element of $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$ is identified with the Virasoro element of $`L_{𝔫𝔰(2)}(c_m,0,0)V_L`$. Thus $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$ is a vertex operator subalgebra.
Now we recall a key result from \[FST\] and \[FSST\] (in the case of unitary modules), slightly reformulated in the language of vertex operator algebras:
###### Theorem 3.2
As a generalized $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-module, $`L_{𝔫𝔰(2)}(c_m,h,q)V_L`$ decomposes as
$$_{k\{0,1,\mathrm{},m\}}_{sI_s}L_{\widehat{𝔰𝔩}(2)}(m,k)M(1,s).$$
where $`s`$ runs through certain infinite index set $`I_s`$.
We know that $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ is rational (see \[FZ\]). Although $`V_{\mathrm{Liou}}`$ is not rational, any module on which $`a(0)`$ acts semisimply is completely reducible, as we mentioned above. Any irreducible module for $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ is isomorphic to $`L_{\widehat{𝔰𝔩}(2)}(m,i)`$ for some $`i\{1,\mathrm{},m\}`$ and any irreducible module for $`V_{\mathrm{Liou}}`$ is isomorphic to $`M(1,s)`$ for some $`s`$. Thus by the result in \[FHL\] on modules for a tensor product of vertex operator algebras, any irreducible module for $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$ is isomorphic to $`L_{\widehat{𝔰𝔩}(2)}(m,i)M(1,s)`$ for some $`i\{0,\mathrm{},m\}`$, $`s`$.
Proposition 2.7 in \[DMZ\] and its proof can be generalized trivially to the case that one of the vertex operator algebra is an irrational vertex operator algebra like $`V_{\mathrm{Liou}}`$, such that in particular, any $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-module with $`\mathrm{𝟏}_{L_{\widehat{𝔰𝔩}(2)}(m,0)}a(0)`$ ($`\mathrm{𝟏}_{L_{\widehat{𝔰𝔩}(2)}(m,0)}`$ being the vacuum vector of $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$) acting semi-simply is completely reducible. Now suppose that $`M`$ is such an $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-module. Then it follows that $`M`$ is completely reducible. So it has a decomposition
$$_{\beta B}M_\beta $$
(3.1)
where $`B`$ is an index set. Note that here the sum might be infinite (comparing with the rational case where this sum is always finite). Since any irreducible $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-module is isomorphic to $`L_{\widehat{𝔰𝔩}(2)}(m,i)M(1,s)`$ for some $`i\{0,\mathrm{},m\}`$, $`s`$, $`M_\beta `$ for any $`\beta B`$ is isomorphic to such a module.
We need the following:
###### Lemma 3.3
Let $`𝒴`$ be an intertwining operator of type
$$\left(\genfrac{}{}{0pt}{}{L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_3)}{L_{\widehat{𝔰𝔩}(2)}(m,i_1)M(1,s_1)L_{\widehat{𝔰𝔩}(2)}(m,i_2)M(1,s_2)}\right).$$
Then
$$𝒴=𝒴^{}𝒴^{\prime \prime }$$
where $`𝒴^{}`$ and $`𝒴^{\prime \prime }`$ are intertwining operators of types $`\left(\genfrac{}{}{0pt}{}{L_{\widehat{𝔰𝔩}(2)}(m,i_3)}{L_{\widehat{𝔰𝔩}(2)}(m,i_1)L_{\widehat{𝔰𝔩}(2)}(m,i_2)}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{M(1,s_3)}{M(1,s_1)M(1,s_2)}\right)`$, respectively. In particular, all fusion rules for irreducible modules for $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$ are $`0`$ or $`1`$.
Proof. It is enough to show that there is a linear injective map from
$$𝒱_{(L_{\widehat{𝔰𝔩}(2)}(m,i_1)M(1,s_1))(L_{\widehat{𝔰𝔩}(2)}(m,i_2)M(1,s_2))}^{L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_3)},$$
the space of intertwining operators of type
$$\left(\genfrac{}{}{0pt}{}{L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_3)}{L_{\widehat{𝔰𝔩}(2)}(m,i_1)M(1,s_1)L_{\widehat{𝔰𝔩}(2)}(m,i_2)M(1,s_2)}\right),$$
to
$$𝒱_{L_{\widehat{𝔰𝔩}(2)}(m,i_1)L_{\widehat{𝔰𝔩}(2)}(m,i_2)}^{L_{\widehat{𝔰𝔩}(2)}(m,i_3)}𝒱_{M(1,s_1)M(1,s_2)}^{M(1,s_3)},$$
where $`𝒱_{L_{\widehat{𝔰𝔩}(2)}(m,i_1)L_{\widehat{𝔰𝔩}(2)}(m,i_2)}^{L_{\widehat{𝔰𝔩}(2)}(m,i_3)}`$ and $`𝒱_{M(1,s_1)M(1,s_2)}^{M(1,s_3)}`$ are the space of intertwining operators of type $`\left(\genfrac{}{}{0pt}{}{L_{\widehat{𝔰𝔩}(2)}(m,i_3)}{L_{\widehat{𝔰𝔩}(2)}(m,i_1)L_{\widehat{𝔰𝔩}(2)}(m,i_2)}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{M(1,s_3)}{M(1,s_1)M(1,s_2)}\right)`$, respectively. But this follows from Proposition 2.10 in \[DMZ\] which in turn is a consequence of a result in \[FHL\] on irreducible modules for a tensor product vertex operator algebra and a result in \[FZ\] giving an isomorphism between a space of intertwining operators and a certain vector space. Since the fusion rules for irreducible $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$-modules are $`0`$ and $`1`$ (see \[FZ\]) and the same is true for irreducible $`V_{\mathrm{Liou}}`$-modules, the lemma is proved.
###### Remark 3.4
Note that this result should be very usefull in calculating the fusion algbera for $`L_{𝔫𝔰(2)}(c_m,0,0)`$ since every such intertwining operator for $`L_{𝔫𝔰(2)}(c_m,0,0)`$ factors as a tensor product of an intertwining operator for vertex operator algebra $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ and an intertwining operator for the vertex operator algebra associated to the Heisenberg algebra. In the present paper, the exact values of the fusion rules are not what we are interested and thus we shall not calculate them here. We leave this calculation to the interested readers. (After the first version of the present paper was finished, we received a preprint \[A3\] from Adamović in which the fusion rules for $`L_{𝔫𝔰(2)}(c_m,0,0)`$ are calculated explicitly.)
Using Lemmas 3.3, we obtain:
###### Proposition 3.5
For fixed $`i_1,i_2\{1,\mathrm{},m\}`$ and $`s_1,s_2`$, if $`s_3s_1+s_2`$, the space
$$𝒱_{(L_{\widehat{𝔰𝔩}(2)}(m,i_1)M(1,s_1))(L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_2))}^{L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_3)}$$
is $`0`$. In particular, there are only finitely many pairs $`(i_3,s_3)\{1,\mathrm{},m\}\times `$ such that the space
$$𝒱_{(L_{\widehat{𝔰𝔩}(2)}(m,i_1)M(1,s_1))(L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_2))}^{L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_3)}$$
are not $`0`$.
Proof. Let
$$𝒴𝒱_{(L_{\widehat{𝔰𝔩}(2)}(m,i_1)M(1,s_1))(L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_2))}^{L_{\widehat{𝔰𝔩}(2)}(m,i_3)M(1,s_3)}.$$
By Lemma 3.3,
$$𝒴=𝒴^{}𝒴^{\prime \prime }$$
where $`𝒴^{}`$ and $`𝒴^{\prime \prime }`$ are intertwining operators of types $`\left(\genfrac{}{}{0pt}{}{L_{\widehat{𝔰𝔩}(2)}(m,i_3)}{L_{\widehat{𝔰𝔩}(2)}(m,i_1)L_{\widehat{𝔰𝔩}(2)}(m,i_2)}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{M(1,s_3)}{M(1,s_1)M(1,s_2)}\right)`$, respectively. It is clear that if $`s_3s_1+s_2`$, $`𝒴^{\prime \prime }=0`$, proving the result.
Proof of Theorem 3.1: Let $`W_l`$, $`l=1,\mathrm{},5`$, be irreducible $`L_{𝔫𝔰(2)}(c_m,0,0)`$-modules and $`𝒴_1`$ and $`𝒴_2`$ intertwining operators of types $`\left(\genfrac{}{}{0pt}{}{W_4}{W_1W_5}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{W_5}{W_2W_3}\right)`$, respectively. Consider the (formal) matrix coefficients
$$w_{(4)}^{},𝒴_1(w_{(1)},x_1)𝒴_2(w_{(2)},x_2)w_{(3)},$$
(3.2)
where $`w_{(l)}W_l`$, $`l=1,2,3`$ and $`w_{(4)}^{}W_4^{}`$. We shall identify $`W_l`$, $`l=1,2,3`$, with $`W_le^0`$ in $`W_lV_L`$ and $`W_4^{}`$ with $`W_4^{}e^0`$ in $`W_4^{}V_L`$. In particular, we use the same notations $`w_{(l)}`$, $`l=1,2,3`$, to denote $`w_{(l)}e^0`$, and $`w_{(4)}^{}`$ to denote $`w_{(4)}^{}e^0`$.
We extend intertwining operators $`𝒴_1`$ and $`𝒴_2`$ uniquely to intertwining operators (denoted by the same notations $`𝒴_1`$ and $`𝒴_2`$) of type
$$\left(\genfrac{}{}{0pt}{}{W_4V_L}{W_1V_LW_5V_L}\right),$$
(3.3)
and
$$\left(\genfrac{}{}{0pt}{}{W_5V_L}{W_2V_LW_3V_L}\right),$$
(3.4)
respectivly. By Theorem 3.2, $`W_lV_L`$, $`l=1,2,3`$, and $`W_4^{}W_4^{}`$ are generalized modules for $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$ and are completely reducible. So $`w_{(l)}=_{k=1}^{p_l}w_{(l)}^{(k)}`$, $`l=1,2,3`$, and $`w_{(4)}^{}=_{k=1}^{p_4}w_{(4)}^{(k)}`$ where $`w_{(l)}^{(k)}`$, $`k=1,\mathrm{},p_i`$, $`l=1,2,3,4`$, are elements of direct summands $`M_l^{(k)}`$ (irreducible $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-modules) in $`W_i`$ for $`i=1,2,3`$ or $`W_4^{}`$ for $`i=4`$. Thus (3.2) is equal to
$$\underset{k_1=1}{\overset{p_1}{}}\underset{k_2=1}{\overset{p_2}{}}\underset{k_2=1}{\overset{p_3}{}}\underset{k_4=1}{\overset{p_4}{}}w_{(4)}^{(k_4)},𝒴_{k_1W_5}^{k_4}(w_{(1)}^{(k_1)},x_1)𝒴_{k_2k_3}^{W_5}(w_{(2)}^{(k_2)},x_2)w_{(3)}^{(k_3)},$$
(3.5)
where $`𝒴_{k_1W_5}^{k_4}`$ and $`𝒴_{k_2k_3}^{W_5}`$ are intertwining operators of types $`\left(\genfrac{}{}{0pt}{}{M_4^{(k_4)}}{M_1^{(k_1)}W_5}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{W_5}{M_2^{(k_2)}M_3^{(k_3)}}\right)`$, respectively.
By Theorem 3.2, $`W_5`$ is a completely reducible generalized $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-module. By Proposition 3.5, (3.5) is equal to
$$\underset{k_1=1}{\overset{p_1}{}}\underset{k_2=1}{\overset{p_2}{}}\underset{k_2=1}{\overset{p_3}{}}\underset{k_4=1}{\overset{p_4}{}}\underset{k_5=1}{\overset{p_5}{}}w_{(4)}^{(k_4)},𝒴_{k_1k_5}^{k_4}(w_{(1)}^{(k_1)},x_1)𝒴_{k_2k_3}^{k_5}(w_{(2)}^{(k_2)},x_2)w_{(3)}^{(k_3)},$$
(3.6)
where $`𝒴_{k_1k_5}^{k_4}`$ and $`𝒴_{k_2k_3}^{k_5}`$ are intertwining operators of types $`\left(\genfrac{}{}{0pt}{}{M_4^{(k_4)}}{M_1^{(k_1)}M_5^{(k_5)}}\right)`$ and $`\left(\genfrac{}{}{0pt}{}{M_5^{(k_5)}}{M_2^{(k_2)}M_3^{(k_3)}}\right)`$, respectively, and $`M_5^{(k_5)}`$, $`k_5=1,\mathrm{},p_5`$, are irreducible irreducible $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-submodules of $`W_5`$ as an $`L_{\widehat{𝔰𝔩}(2)}(m,0)V_{\mathrm{Liou}}`$-module.
By Proposition 3.3,
$$𝒴_{k_1k_5}^{k_4}=(𝒴_{k_1k_5}^{k_4})^{}(𝒴_{k_1k_5}^{k_4})^{\prime \prime }$$
and
$$𝒴_{k_2k_3}^{k_5}=(𝒴_{k_2k_3}^{k_5})^{}(𝒴_{k_2k_3}^{k_5})^{\prime \prime },$$
where $`(𝒴_{k_1k_5}^{k_4})^{}`$ and $`(𝒴_{k_2k_3}^{k_5})^{}`$ are intertwining operators for the vertex operator algebras $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ and $`(𝒴_{k_1k_5}^{k_4})^{\prime \prime }`$ and $`(𝒴_{k_2k_3}^{k_5})^{\prime \prime }`$ are intertwining operators for the vertex operator algebras $`V_{\mathrm{Liou}}`$. Thus (3.6) is equal to a finite sum of series of the form
$`\stackrel{~}{w}_{(4)},\stackrel{~}{𝒴}_1(\stackrel{~}{w}_{(1)},x_1)\stackrel{~}{𝒴}_2(\stackrel{~}{w}_{(2)},x_2)\stackrel{~}{w}_{(3)}\stackrel{~}{w}_{(4)},\stackrel{~}{\stackrel{~}{𝒴}}_1(\stackrel{~}{\stackrel{~}{w}}_{(1)},x_1)\stackrel{~}{\stackrel{~}{𝒴}}_2(\stackrel{~}{\stackrel{~}{w}}_{(2)},x_2)\stackrel{~}{\stackrel{~}{w}}_{(3)},`$ (3.7)
where $`\stackrel{~}{𝒴}_1`$ and $`\stackrel{~}{𝒴}_2`$ are intertwining operators among irreducible modules for $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ and $`\stackrel{~}{\stackrel{~}{𝒴}}_1`$ and $`\stackrel{~}{\stackrel{~}{𝒴}}_2`$ are intertwining operators among irreducible modules for $`V_{\mathrm{Liou}}`$.
In \[HL7\], it was proved that intertwining operators for the vertex operator algebra $`L_{\widehat{𝔰𝔩}(2)}(m,0)`$ satisify the convergence and extension property for products using the Knizhnik-Zamolodchikov eqautions. The convergence and extension property for products of intertwining operators for the vertex operator algebra $`V_{\mathrm{Liou}}`$ can be proved trivially by a straightford calculation. Using the convergence and extension properties for products of intertwining operators for these vertex operator algebras and using the fact proved above that (3.2) is a finite sum of series of the form (3.7), we conclude that (3.2) is convergent when we substitute $`e^{n_i\mathrm{log}z_i}`$ for $`x_i^n`$ with $`z_1,z_2`$ satisfying $`|z_1|>|z_2|>0`$ and it can be analytically extended to an analytic function in the region $`|z_2|>|z_1z_2|>0`$ of the form
$$\underset{i=1}{\overset{j}{}}z_2^{r_i}(z_1z_2)^{s_i}f_i\left(\frac{z_1z_2}{z_2}\right).$$
(3.8)
We still need to prove the following: There exists $`N`$ (which does not depend on $`w_{(1)}`$ and $`w_{(2)}`$) such that
$$\mathrm{wt}(w_{(1)})+\mathrm{wt}(w_{(2)})+s_i>N,$$
(3.9)
for $`i=1,\mathrm{},j`$. The existence of $`N`$ follows (as in the $`N=0`$ and $`N=1`$ cases in \[H2\] and \[HM\], respectively) from an induction argument for the $`N=2`$ superconformal algebra. Since any $`L_{𝔫𝔰(2)}(c_m,0,0)`$-module is completely reducible and since there are only finitely many irreducible $`L_{𝔫𝔰(2)}(c_m,0,0)`$-modules, we need only prove this existence in the case that $`W_1`$ and $`W_2`$ are irreducible. When $`w_{(1)}`$ and $`w_{(2)}`$ are lowest weight vectors, (3.2) is absolutely convergent in the region $`|z_1|>|z_2|>0`$ and can be analytically extended to an analytic function in in the region $`|z_2|>|z_1z_2|>0`$ of the form (3.8). We choose an $`N`$ such that for these lowest weight vectors, (3.9) holds. For general $`w_{(1})`$ and $`w_{(2)}`$, we use induction instead of the proof above to show that (3.2) in the region $`|z_1|>|z_2|>0`$ and can be analytically extended to an analytic function in in the region $`|z_2|>|z_1z_2|>0`$ of the form (3.8). In addition, the induction also shows that (3.9) holds for the $`N`$ we choose.
An immediate consequence of Theorem 3.1 is the following:
###### Theorem 3.6
Let $`m_i`$, $`i=1,\mathrm{},n`$, be positive integers and $`V`$ an $`N=2`$ superconformal vertex operator superalgebra in the class $`𝒞_{m_1;\mathrm{},m_n}`$. Then intertwining operators for $`V`$ satisfy the convergence and extension property for products of intertwining operators introduced in \[H1\].
We omit the proof since it is the same as the corresponding result in \[H2\] and \[HL7\].
## 4 Intertwining operator superalgebras and vertex tensor categories for $`N=2`$ unitary minimal models
Let $`m_i`$, $`i=1,\mathrm{},n`$, be $`n`$ positive integers and $`V`$ a vertex operator superalgebra in the class $`𝒞_{m_1;\mathrm{};m_n}`$. Using Corollary 2.9, Proposition 2.12 and Theorem 3.6 above, and Theorems 3.1 and 3.2 in \[H2\], which are proved in \[H2\] using results in \[HL1\]\[HL6\] and \[H1\], we obtain the following:
###### Theorem 4.1
The conclusions of Theorems 4.1, 4.2, and 4.6 and the first conclusion of Corollary 4.7 in Section 4 of \[HM\] holds for $`V`$.
The notions of intertwining operator algebra in \[H3\] (see also \[H4\] and \[H5\]) and $`N=1`$ superconformal intertwining operator superalgebra in \[HM\] can be generalized easily to the following notion:
###### Definition 4.2
An $`N=2`$ superconformal intertwining operator superalgebra is an intertwining operator superalgebra $`W`$ together with three elements $`\tau ^+`$, $`\tau ^{}`$ and $`\mu `$ such that $`(W^e,Y,\mathrm{𝟏},\tau ^+,\tau ^{})`$ is an $`N=2`$ superconformal vertex operator algebra.
Then we have:
###### Theorem 4.3
The conclusion of Theorem 4.4 in \[HM\] with “$`N=1`$” repalced by “$`N=2`$” holds for $`V`$.
Department of Mathematics, Kerchof Hall, University of Virginia, Charlottesville, VA 22904-4137
and
Department of Mathematics, Rutgers University, 110 Frelinghuysen Rd., Piscataway, NJ 08854-8019 (permanent address)
E-mail address: yzhuang@math.rutgers.edu
Department of Mathematics, Rutgers University, 110 Frelinghuysen Rd., Piscataway, NJ 08854-8019
E-mail address: amilas@math.rutgers.edu
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# Fermi–liquid theory of superfluid asymmetric nuclear matter
## Abstract
Influence of asymmetry on superfluidity of nuclear matter with triplet–singlet pairing of nucleons (in spin and isospin spaces) is considered within the framework of a Fermi–liquid theory. Solutions of self–consistent equations for the critical temperature and the energy gap at $`T=0`$ are obtained with the use of Skyrme effective nucleon interaction. It is shown, that if the chemical potentials of protons and neutrons are determined in the approximation of ideal Fermi–gas, then the energy gap for some values of density and asymmetry parameter of nuclear matter demonstrates double–valued behavior. However, with account for the feedback of pairing correlations through the normal distribution functions of nucleons two–valued behavior of the energy gap turns into universal one–valued behavior. At $`T=0`$ superfluidity arises and disappears as a result of a first order phase transition in density.
It is well established, that neutron–proton ($`np`$) pairing plays an essential role in the description of superfluidity of finite nuclei with $`N=Z`$ (see Refs. and references therein) and symmetric nuclear matter . In astrophysical context $`np`$ pairing correlations can be important for the description of $`r`$–process and cooling of neutron stars, which permit pion or kaon condensation . In this Rapid Communication we shall investigate the influence of asymmetry on $`np`$ superfluidity of nuclear matter. Previously this problem was treated with the use of various approaches and potentials of NN interaction. In particular, the cases of $`{}_{}{}^{3}S_{1}^{}`$$`{}_{}{}^{3}D_{1}^{}`$ and $`{}_{}{}^{3}D_{2}^{}`$ pairing were considered in Refs. on the basis of the Thouless criterion for the thermodynamic $`T`$ matrix. As a potential of NN interaction, the Graz II and Paris potentials were chosen, respectively. Superfluidity in $`{}_{}{}^{3}S_{1}^{}`$$`{}_{}{}^{3}D_{1}^{}`$ pairing channel was studied also in Ref. within BCS theory of superconductivity with the use of the Paris potential in the separable form. Investigations, based on the Thouless criterion, deduce the suppression of $`np`$ pairing correlations with increase of isospin asymmetry. However, the Thouless criterion can be exploited for finding the critical temperature only, but does not permit one to draw any conclusions about superfluidity with a finite gap. The studies in Ref. , based on the BCS theory, were carried out with the use of the bare interaction and the single particle spectrum of a free Fermi gas and give, thus, overestimated values of the energy gap. The effect of ladder–renormalized single particle spectrum on the magnitude of the energy gap in $`{}_{}{}^{3}S_{1}^{}`$$`{}_{}{}^{3}D_{1}^{}`$ pairing channel was investigated in Ref. . The Argonne $`V_{14}`$ potential was explored as input for determination of the single particle energy and the bare interaction in the form of the Paris potential was used to evaluate the energy gap. The use of the bare interaction in the gap equation seems to be a very strong simplification, because medium polarization strongly reduces the magnitude of the gap (see Ref. for the influence of the polarization effects on the pairing force in the $`{}_{}{}^{1}S_{0}^{}`$ channel). In principle, the effective pairing interaction should be obtained by means of Brueckner renormalization, which gives the correct interaction after modifying the bare interaction for the effect of nuclear medium. However, the issue of a microscopic many–body calculation of the effective pairing potential is a complex one and still is not solved. For this reason, it is quite a natural step to develop some kind of a phenomenological theory, where instead of microscopical calculation of the pairing interaction one exploits the phenomenological effective interaction. We shall investigate the influence of asymmetry on superfluid properties of nuclear matter, using Landau’s semiphenomenological theory of a Fermi–liquid (FL). In the Fermi–liquid model the normal and anomalous FL interaction amplitudes are taken into account on an equal footing. This will allow us to consider consistently the influence of the FL amplitudes on superfluid properties of nuclear matter. Besides, as a potential of NN interaction we choose the Skyrme effective forces, describing the interaction of two nucleons in the presence of nucleon medium. The Skyrme forces are widely used in the description of nuclear system properties and, in particular, they were exploited for the description of superfluid properties of finite nuclei as well as infinite symmetric nuclear matter -.
The basic formalism is laid out in more detail in Ref. , where superfluidity of symmetric nuclear matter was studied. As shown there, superfluidity with triplet–singlet (TS) pairing of nucleons (total spin $`S`$ and isospin $`T`$ of a pair are equal $`S=1`$, $`T=0`$) is realized near the saturation density in symmetric nuclear matter with the Skyrme interaction. Therefore, we shall study further the influence of asymmetry on superfluid properties of TS phase of nuclear matter. For the states with the projections of total spin and isospin $`S_z=T_z=0`$ the normal distribution function $`f`$ and the anomalous distribution function $`g`$ have the form
$$f(𝐩)=f_{00}(𝐩)\sigma _0\tau _0+f_{03}(𝐩)\sigma _0\tau _3,g(𝐩)=g_{30}(𝐩)\sigma _3\sigma _2\tau _2$$
(1)
where $`\sigma _i,\tau _k`$ are the Pauli matrices in spin and isospin spaces. The operator of quasiparticle energy $`\epsilon `$ and the matrix order parameter $`\mathrm{\Delta }`$ of the system for the energy functional, being invariant with respect to rotations in spin and isospin spaces, have the analogous structure
$$\epsilon (𝐩)=\epsilon _{00}(𝐩)\sigma _0\tau _0+\epsilon _{03}(𝐩)\sigma _0\tau _3,\mathrm{\Delta }(𝐩)=\mathrm{\Delta }_{30}(𝐩)\sigma _3\sigma _2\tau _2.$$
(2)
Using the minimum principle of the thermodynamic potential and procedure of block diagonalization , one can express evidently the distribution functions $`f_{00},f_{03}`$ and $`g_{30}`$ in terms of the quantities $`\epsilon `$ and $`\mathrm{\Delta }`$:
$`f_{00}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\xi _{00}}{4E}}\left(\mathrm{tanh}{\displaystyle \frac{E+\xi _{03}}{2T}}+\mathrm{tanh}{\displaystyle \frac{E\xi _{03}}{2T}}\right),`$ (3)
$`f_{03}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(\mathrm{tanh}{\displaystyle \frac{E+\xi _{03}}{2T}}\mathrm{tanh}{\displaystyle \frac{E\xi _{03}}{2T}}\right),`$ (4)
$`g_{30}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_{30}}{4E}}\left(\mathrm{tanh}{\displaystyle \frac{E+\xi _{03}}{2T}}+\mathrm{tanh}{\displaystyle \frac{E\xi _{03}}{2T}}\right).`$ (5)
Here
$`E=\sqrt{\xi _{00}^2+\mathrm{\Delta }_{30}^2},`$ $`\xi _{00}=\epsilon _{00}\mu _{00}^0,`$ $`\xi _{03}=\epsilon _{03}\mu _{03}^0,`$
$`\mu _{00}^0=`$ $`{\displaystyle \frac{\mu _p^0+\mu _n^0}{2}},\mu _{03}^0=`$ $`{\displaystyle \frac{\mu _p^0\mu _n^0}{2}},`$
$`T`$ is temperature, $`\mu _p^0`$ and $`\mu _n^0`$ are chemical potentials of protons and neutrons. To obtain the closed system of equations for the quasiparticle energy $`\epsilon `$ and the energy gap $`\mathrm{\Delta }`$, it is necessary to express the quantities $`\epsilon ,\mathrm{\Delta }`$ through the distribution functions $`f`$ and $`g`$. For this purpose one has to set the energy functional $`(f,g)`$ of the system. In the case of asymmetric nuclear matter with TS pairing of nucleons the energy functional is characterized by two normal $`U_0,U_2`$ and one anomalous $`V_1`$ FL amplitudes . Differentiating the functional $`(f,g)`$ with respect to $`g`$ and using Eq. (5), one can obtain the gap equation in the form
$`\mathrm{\Delta }_{30}(𝐩)`$ $`=`$ $`{\displaystyle \frac{1}{4𝒱}}{\displaystyle \underset{𝐪}{}}V_1(𝐩,𝐪){\displaystyle \frac{\mathrm{\Delta }_{30}(𝐪)}{E(𝐪)}}`$ (7)
$`\times \left\{\mathrm{tanh}{\displaystyle \frac{E(𝐪)+\xi _{03}(𝐪)}{2T}}+\mathrm{tanh}{\displaystyle \frac{E(𝐪)\xi _{03}(𝐪)}{2T}}\right\}.`$
The anomalous interaction amplitude $`V_1`$ in the Skyrme model reads
$`V_1(𝐩,𝐪)`$ $`=`$ $`t_0(1+x_0)+{\displaystyle \frac{1}{6}}t_3\varrho ^\beta (1+x_3)`$ (9)
$`+{\displaystyle \frac{1}{2\mathrm{}^2}}t_1(1+x_1)(𝐩^{\mathrm{\hspace{0.17em}2}}+𝐪^{\mathrm{\hspace{0.17em}2}}),`$
where $`\varrho `$ is density of nuclear matter, $`t_i,x_i,\beta `$ are some phenomenological parameters, which differ for various versions of the Skyrme forces (later we shall use the SkP potential ). Equation (7) should be solved jointly with equations
$`{\displaystyle \frac{1}{𝒱}}{\displaystyle \underset{𝐩}{}}\{2{\displaystyle \frac{\xi _{00}(𝐩)}{E(𝐩)}}`$ (10)
$`\times (\mathrm{tanh}{\displaystyle \frac{E(𝐩)+\xi _{03}(𝐩)}{2T}}`$ $`+\mathrm{tanh}{\displaystyle \frac{E(𝐩)\xi _{03}(𝐩)}{2T}})\}=\varrho ,`$ (11)
$$\frac{1}{𝒱}\underset{𝐩}{}\left\{\mathrm{tanh}\frac{E(𝐩)+\xi _{03}(𝐩)}{2T}\mathrm{tanh}\frac{E(𝐩)\xi _{03}(𝐩)}{2T}\right\}=\alpha \varrho ,$$
(12)
being the normalization conditions for the normal distribution functions $`f_{00},f_{03}`$. In Eq. (12) the quantity $`\alpha =(\varrho _n\varrho _p)/\varrho `$ is the asymmetry parameter of nuclear matter, $`\varrho _p,\varrho _n`$ are the partial number densities of protons and neutrons. Note that the account of the normal FL amplitudes in the case of the effective Skyrme interaction, being quadratic in momenta, is reduced to the renormalization of free nucleon masses and chemical potentials. Expressions for the quantities $`\xi _{00},\xi _{03}`$, which enter in Eqs. (7),(11),(12), with regard for the explicit form of the amplitudes $`U_0,U_2`$ read
$$\xi _{00}=\frac{p^2}{2m_{00}}\mu _{00},\xi _{03}=\frac{p^2}{2m_{03}}\mu _{03},$$
where the effective nucleon mass $`m_{00}`$ and effective isovector mass $`m_{03}`$ are defined by the formulas:
$`{\displaystyle \frac{\mathrm{}^2}{2m_{00}}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{}^2}{2m_{00}^0}}+{\displaystyle \frac{\varrho }{16}}[3t_1+t_2(5+4x_2)],`$ (13)
$`{\displaystyle \frac{\mathrm{}^2}{2m_{03}}}`$ $`=`$ $`{\displaystyle \frac{\alpha \varrho }{16}}[t_1(1+2x_1)t_2(1+2x_2)],`$ (14)
$`m_{00}^0`$ being the bare mass of a nucleon. The renormalized chemical potentials $`\mu _{00},\mu _{03}`$ should be determined from Eqs. (11), (12) and in the leading approximation on the ratios $`T/\epsilon _F,\mathrm{\Delta }/\epsilon _F`$ have the form
$$\mu _{00}=\frac{1}{2}(\mu _p+\mu _n),\mu _{03}=\frac{1}{2}(\mu _p\mu _n);\mu _{p,n}=\frac{\mathrm{}^2k_{F_{p,n}}^2}{2m_{p,n}},$$
(15)
where $`k_{F_{p,n}}=(3\pi ^2\varrho _{p,n})^{1/3}`$, $`m_p`$ and $`m_n`$ are the proton and neutron effective masses, defined as
$$\frac{2}{m_{00}}=\frac{1}{m_p}+\frac{1}{m_n},\frac{2}{m_{03}}=\frac{1}{m_p}\frac{1}{m_n}.$$
The critical temperature of transition to TS superfluid phase is found from Eq. (7), determining the energy gap, in the linear on $`\mathrm{\Delta }`$ approximation. Considering, that the interaction amplitude $`V_1`$ is not equal to zero only in a narrow layer near the Fermi–surface, $`|\xi _{00}|\theta `$ (we shall set $`\theta =0.1\mu _{00}`$) and entering new variables $`x=\xi _{00}/\mu _{00}`$, $`\theta _0=\theta /\mu _{00}`$, $`T_\mu =T/\mu _{00}`$, we present this equation in the form
$`1`$ $`=`$ $`{\displaystyle \frac{g}{4}}{\displaystyle _{\theta _0}^{\theta _0}}{\displaystyle \frac{dx}{x}}`$ (17)
$`\times \left\{\mathrm{tanh}{\displaystyle \frac{x[1+\frac{m_{00}}{m_{03}}]+\psi }{2T_\mu }}+\mathrm{tanh}{\displaystyle \frac{x[1\frac{m_{00}}{m_{03}}]\psi }{2T_\mu }}\right\},`$
$`g=\nu _FV_1(p=p_F,q`$ $`=`$ $`p_F),\psi ={\displaystyle \frac{m_{00}}{m_{03}}}{\displaystyle \frac{\mu _{03}}{\mu _{00}}},`$
$`\nu _F={\displaystyle \frac{m_{00}p_F}{2\pi ^2\mathrm{}^3}},`$ $`p_F`$ $`=\sqrt{2m_{00}\mu _{00}}`$
The results of numerical integration of Eq. (17) are shown in Fig. 1. For small values of asymmetry $`\alpha `$ there exist such regions of large and low densities of nuclear matter, for which we have two critical temperatures. When $`\alpha `$ increases, these regions begin to approach and at some value $`\alpha =\alpha _c`$ ($`\alpha _c0.071`$) it takes place contiguity of the regions, so that we have always only two critical temperatures (for the regions, where solutions exist). For $`\alpha >\alpha _c`$ the phase curves are separated from the density axis and turn into the closed oval curves. Under further increase of $`\alpha `$ the dimension of the oval curves is reduced and at some $`\alpha =\alpha _m`$ the oval curves contract to a point ($`\alpha _m0.092`$). For the values $`\alpha >\alpha _m`$ the triplet–singlet superfluidity fails. Note, that our results concerning two–valued behavior of the critical temperature qualitatively agree with the results of Ref. , where $`{}_{}{}^{3}D_{2}^{}`$ pairing of nucleons with the Paris NN potential was considered.
As follows from Eq. (7), an equation determining the energy gap at $`T=0`$ has the form
$`\mathrm{\Delta }_{30}(𝐩)`$ $`=`$ $`{\displaystyle \frac{1}{4𝒱}}{\displaystyle \underset{𝐪}{}}V_1(𝐩,𝐪){\displaystyle \frac{\mathrm{\Delta }_{30}(𝐪)}{E(𝐪)}}`$ (19)
$`\times \left\{\text{sgn}(E(𝐪)+\xi _{03}(𝐪))+\text{sgn}(E(𝐪)\xi _{03}(𝐪))\right\}`$
Passing in Eq. (19) to integration on a layer, we arrive at the equation for determining the dimensionless gap $`y=\mathrm{\Delta }_{30}/\mu _{00}`$:
$$1=\frac{g}{4}_{\theta _0}^{\theta _0}\frac{dx}{\sqrt{x^2+y^2}}\left\{1+\text{sgn}\left(\sqrt{x^2+y^2}\frac{m_{00}}{m_{03}}x\psi \right)\right\}$$
(20)
(for $`\alpha >0`$ it holds true $`m_{03}>0,\mu _{03}<0`$). The contribution to the integral gives the domain on $`x`$, for which the function, standing as an argument of the function sgn, is positive. In particular, such values of the gap, density and asymmetry parameter of nuclear matter are possible, that this function has no roots with respect to $`x`$. In this case the whole domain from $`\theta _0`$ to $`\theta _0`$ gives the contribution to the integral in Eq. (20) and we arrive at the equation of the BCS type at $`T=0`$ with the solution $`y=\theta _0/(\mathrm{sinh}(1/g))`$.
Let us first find the solutions of Eq. (20) in the case when the chemical potentials $`\mu _{00},\mu _{03}`$ are given in the main approximation on $`\mathrm{\Delta }/\epsilon _F`$. The results of numerical integration of Eq. (20) are presented in Fig. 2. In the case of symmetric nuclear matter ($`\alpha =0`$) we obtain the phase curve with one–valued behavior of the gap. For small values of asymmetry $`\alpha `$ there exist such regions of large and low densities of nuclear matter (excluding some vicinity of the point $`\varrho =0`$), for which we have two values of the energy gap, where one of these values is the solution of the BCS type and practically coincides with the corresponding value of the gap for the case $`\alpha =0`$. For $`\alpha `$ less than some $`\alpha _c`$, the energy gap demonstrates double–valued behavior in the intervals $`(\varrho _{min},\varrho _{min}^{})`$ and $`(\varrho _{max}^{},\varrho _{max})`$, while in the interval $`(\varrho _{min}^{},\varrho _{max}^{})`$ it has one–valued behavior (see Table I for the values of various boundary points $`\varrho `$). When $`\alpha `$ increases, the regions with double–valued behavior of the gap begin to approach and at $`\alpha =\alpha _c`$ (the same $`\alpha _c`$ as for the phase curves $`T_c(\varrho )`$) it takes place contiguity of the regions with two solutions. For $`\alpha >\alpha _c`$ two branches of the phase curves are separated from the density axis and combine to one curve, beginning and ending in some points of the phase curve with $`\alpha =0`$. In this case the energy gap differs from zero only in the interval $`(\varrho _{min},\varrho _{max})`$, where it has double–valued behavior. When $`\alpha `$ increases further, the boundary points of the phase curves move towards and at some $`\alpha =\alpha _m^{}`$ (not equal to $`\alpha _m`$ for the phase curves $`T_c(\varrho )`$) the branches of solutions contract to a point. The value $`\alpha _m^{}`$ determines the maximum value of the asymmetry parameter, when TS superfluidity exists at $`T=0`$ ($`\alpha _m^{}0.179`$).
Let us now find the solutions of Eq. (20) while accounting for the influence of the finite size of the gap on the chemical potentials $`\mu _{00},\mu _{03}`$. The results of integration of Eq. (20) in this case are presented in Fig. 3. Here the solutions of Eqs. (20), (11), and (12) obtained for different $`\alpha `$, correspond to the different parts of the dome–shaped curve, contained by the dashed lines of the same type for a given $`\alpha `$. One can see, that taking into account the feedback of the finite size of the gap through the normal distribution functions $`f_{00},f_{03}`$ in Eqs. (3), (4) leads to the qualitative change: instead of two–valued behavior of the gap we have universal one-valued behavior. The first solution of the BCS type, obtained in uncoupled calculation, remains practically unchanged in self–consistent treatment of the gap equation (20) and with sufficiently high accuracy equals to its value at $`\alpha =0`$ in the self–consistent determination. The second solution in uncoupled scheme, to which corresponds the smaller gap width, under simultaneous iterations of Eqs. (11), (12), (20) tends to the first solution of the BCS type. Taking into account the finiteness of the gap results in the reduction of the threshold asymmetry, at which superfluidity disappears, to the value $`\alpha _m^{}0.029`$. Thus, in spite of the smallness of the ratio $`\mathrm{\Delta }/\epsilon _F`$ ($`\mathrm{\Delta }/\epsilon _F0.12`$ for all densities $`\varrho `$), the backward influence of pairing correlations is significant. This is explained by the fact, that if the quantity $`\mathrm{\Delta }`$ in Eqs. (11), (12) differs from zero, then the absolute value of the chemical potential $`\mu _{03}`$ increases a few times as against its value at $`\mathrm{\Delta }=0`$. The increase of $`|\mu _{03}|`$ is equivalent to the increase of the effective shift between neutron and proton Fermi surfaces, that leads to significant reduction of the threshold asymmetry. As at $`\alpha 0`$ the gap is finite everywhere, superfluidity arises and disappears under changing density by means of a first order phase transition. In principle, this phase transition can be observed in laboratory conditions under the study of intermediate–energy heavy ion reactions. If we assume that the final stage of the reaction can be described as an expansion of a compound nucleon system , formed in a heavy ion collision, then under lowering density this disassembling phase can undergo a first order phase transition in density from the normal to superfluid state.
The temperature dependence of the energy gap was studied in Refs. with the use of the bare interaction in the gap equation in the form of the Paris and Argonne V<sub>14</sub> potentials, respectively. Our results agree qualitatively with theirs at $`T=0`$. However, in our calculations with the effective density–dependent NN interaction we obtain the gap as a function of density (not at fixed density) and this allows us to find new important features in behavior of the energy gap.
In summary, we studied TS superfluidity of asymmetric nuclear matter in the FL model with density–dependent Skyrme effective interaction (SkP force). In the FL approach the normal and anomalous FL amplitudes are taken into account on an equal footing and this allows us to consider consistently within the framework of a phenomenological theory the influence of medium effects on superfluid properties of nuclear matter. It is shown, that for some values of density and asymmetry parameter of nuclear matter the critical temperature of a second order phase transition demonstrates double–valued behavior, that agrees with the results of previous studies. In the case when the chemical potentials $`\mu _{00},\mu _{03}`$ (half of a sum and half of a difference of the proton and neutron chemical potentials, respectively) are determined in the approximation of ideal Fermi–gas, the energy gap demonstrates for some values of density and asymmetry parameter the double–valued behavior. If we consider the feedback of pairing correlations through the dependence of the normal distribution functions of nucleons from the energy gap, then the energy gap drastically changes its behavior from two–valued to universal one–valued character. In spite of relative smallness of the ratio $`\mathrm{\Delta }/\epsilon _F`$, taking into account of the finite size of the gap in chemical potentials leads to the significant increase of absolute value of $`\mu _{03}`$ and, hence, to considerable reduction of the threshold asymmetry, at which superfluidity at $`T=0`$ disappears. In self–consistent determination the energy gap at $`T=0`$ as a function of density has a finite width and normal–to–superfluid and superfluid–to–normal phase transitions should appear as a first order phase transitions in density. Among the other problems we note here the study of multi–gap superfluidity in asymmetric nuclear matter.
Acknowledgement. The authors thank A. Sedrakian for reading the preliminary version of the manuscript and valuable comment. Financial support of BMBF and STCU (grant $`\mathrm{\#}1480`$) is acknowledged.
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# Dilaton-driven brane inflation in type IIB string theory
## I Introduction
The idea that our universe might be a domain wall embedded in a higher dimensional space has attracted much interest recently. It is possible that the fundamental scale of gravity can be lowered to the electroweak scale by introducing large extra dimensions . Randall and Sundrum proposed scenarios that our observed universe is embedded in a five-dimensional bulk, in which the background metric is curved along the extra dimension due to the negative bulk cosmological constant. These models have been studied extensively because they might provide the solution to the gauge hierarchy problem and cosmological constant problem . The cosmology can be different from the conventional four dimensional one.
One can naturally attempt to justify the scenario within a well defined framework for higher dimensional quantum theory of gravity such as string theory. Many attempts have been made to apply this idea to string theory in the context with D-branes , where the standard model gauge bosons as well as charged matter arise as fluctuations of the D-branes. An early example of this is the Horava-Witten picture for the nonperturbative heterotic $`E_8\times E_8`$ string . The spacetime includes a compact dimension with an orbifold structure. Matter is confined to the hypersurface which forms the boundaries of the spacetime. Within the string theory context, it is natural that our observable four-dimensional world is a three-brane embedded in ten dimensional string. In such theories, one of the important issue is the cosmological evolution of our universe. Many cosmological models associated to brane universe have been suggested. The models can be classified into two categories. The first is that the domain walls (branes) are static solution of the underlying theory and the cosmological evolution of our universe is due the time evolution of energy density on the domain wall (brane) . The second is that the cosmological evolution of our universe is due to the motion of our brane-world in the background of gravitational field of the bulk . We will focus on the second approach in this paper.
It is shown that the motion of the brane in ambient space induces cosmological expansion (or contraction) on our universe simulating various kinds of matter or a cosmological constant. In other words the cosmological expansion is not due to energy density on our universe but somewhere else. This is the idea by mirage cosmology . Friedman-like equations were derived for various bulk background field solutions. In , the motion of a three-brane, in a background of type 0 theory was examined.
In this paper, employing the formalism of , we will study how the presence of matter field on the background geometry affects the cosmological evolution of the brane universe. More specifically we will consider cosmological evolution of type IIB theory with two different background geometries, one without axion and the other with axion. We will compare the two and see the difference.
The organization of the paper is as follows. In Sec. II we will briefly review the formalism of mirage cosmology in ref. and set up some preliminaries for our calculation. In Sec. III we consider the type IIB theory and its background solution with and without axion field. In Sec. IV, using the background solutions of Sec. III, we find the cosmological evolution of the three-brane under the background. Finally section V is devoted to conclusions and discussion.
## II Formalism
In this section, we consider a probe brane moving in a generic static spherically symmetric background. We ignore its back reaction to the ambient space. As the brane moves in a geodesic, the induced world-volume mertic becomes a function of time. The cosmological evolution is possible from the brane resident point of view. We will focus on a D3-brane case. For this purpose we parametrize the metric of a D3-brane as
$$ds_{10}^2=g_{00}(r)dt^2+g(r)(d\stackrel{}{x})^2+g_{rr}(r)dr^2+g_S(r)d\mathrm{\Omega }_5^2,$$
(1)
and there are dilaton field $`\varphi `$ as well as RR (Ramond-Ramond) background $`C(r)=C_{0\mathrm{}3}(r)`$. The probe brane will in general move in this background and its dynamics is governed by the Dirac-Born-Infeld (DBI) action. In maximally supersymmetric case, ignoring the fermions, it is given by
$`S`$ $`=`$ $`T_3{\displaystyle d^4\xi e^\varphi \sqrt{det(\widehat{G}_{\alpha \beta }+(2\pi \alpha ^{})F_{\alpha \beta }B_{\alpha \beta })}}`$ (2)
$`+`$ $`T_3{\displaystyle d^4\xi \widehat{C}_4}+\mathrm{anomaly}\mathrm{terms},`$ (3)
where the induced metric on the brane is
$$\widehat{G}_{\alpha \beta }=G_{\mu \nu }\frac{x^\mu }{\xi ^\alpha }\frac{x^\nu }{\xi ^\beta }$$
(4)
with similar expressions for other fields. Generally the motion of a probe D3-brane have a nonzero angular momentum in the transverse directions. We can write the relevant part of the Lagrangian, in the static gauge $`x^\alpha =\xi ^\alpha (\alpha =0,1,2,3)`$, as
$$L=\sqrt{A(r)B(r)\dot{r}^2D(r)h_{ij}\dot{\phi }^i\dot{\phi }^j}C(r),$$
(5)
where $`h_{ij}\phi ^i\phi ^j`$ is the line element on the unit five-sphere ($`i,j=5,\mathrm{},9`$),
$$A(r)=g^3(r)|g_{00}(r)|e^{2\varphi },B(r)=g^3(r)g_{rr}(r)e^{2\varphi },D(r)=g^3(r)g_S(r)e^{2\varphi },$$
(6)
and $`C(r)`$ is the RR background. The momenta of the system are given by
$`p_r={\displaystyle \frac{B(r)\dot{r}}{\sqrt{A(r)B(r)\dot{r}^2D(r)h_{ij}\dot{\phi }^i\dot{\phi }^j}}},`$ (7)
$`p_i={\displaystyle \frac{D(r)h_{ij}\dot{\phi }^j}{\sqrt{A(r)B(r)\dot{r}^2D(r)h_{ij}\dot{\phi }^i\dot{\phi }^j}}}.`$ (8)
Calculating the Hamiltonian and demanding the conservation of energy, we have
$$H=C(r)\frac{A(r)}{\sqrt{A(r)B(r)\dot{r}^2D(r)h_{ij}\dot{\phi }^i\dot{\phi }^j}}=E,$$
(9)
where $`E`$ is the total energy of the brane. Also from the conservation of the total angular momentum $`h^{ij}p_ip_j=\mathrm{}^2`$, we have
$$h_{ij}\dot{\phi }^i\dot{\phi }^j=\frac{\mathrm{}^2(A(r)B(r)\dot{r}^2)}{D(r)(D(r)+\mathrm{}^2)}.$$
(10)
Substituting equation (10) into (9) and solving with respect to $`\dot{r}^2`$, we have the equation for the radial variable as
$$\dot{r}^2=\frac{A}{B}\left\{1\frac{A}{(C+E)^2}\frac{D+\mathrm{}^2}{D}\right\}.$$
(11)
Plugging equation (11) back into (10), we have the equation for the angular variable
$$h_{ij}\dot{\phi }^i\dot{\phi }^j=\frac{A^2\mathrm{}^2}{D^2(C+E)^2}.$$
(12)
The induced four-dimensional metric on the three-brane universe is
$$d\widehat{s}_{4d}^2=(g_{00}+g_{rr}\dot{r}^2+g_Sh_{ij}\dot{\phi }^i\dot{\phi }^j)dt^2+g(d\stackrel{}{x})^2.$$
(13)
Using equation (11) and (12), this reduces to
$$d\widehat{s}_{4d}^2=\frac{g_{00}^2g^3e^{2\varphi }}{(C+E)^2}dt^2+g(d\stackrel{}{x})^2=d\eta ^2+g(r(\eta ))(d\stackrel{}{x})^2,$$
(14)
where we defined, for the standard form of a flat expanding universe, the cosmic time $`\eta `$ as
$$d\eta =\frac{|g_{00}|g^{3/2}e^\varphi }{|C+E|}dt.$$
(15)
If we define the scale factor as $`a^2=g`$, we can calculate, from the analogue of the four-dimensional Friedman equation, the Hubble constant $`H=\dot{a}/a`$
$$\left(\frac{\dot{a}}{a}\right)^2=\frac{(C+E)^2g_Se^{2\varphi }|g_{00}|(g_Sg^3+\mathrm{}^2e^{2\varphi })}{4|g_{00}|g_{rr}g_Sg^3}\left(\frac{g^{}}{g}\right)^2,$$
(16)
where the dot denotes the derivative with respect to cosmic time and the prime denotes the derivative with respect to $`r`$. The right hand side of (16) can be interpreted as the effective matter density on the probe brane
$$\frac{8\pi }{3}\rho _{\mathrm{eff}}=\frac{(C+E)^2g_Se^{2\varphi }|g_{00}|(g_Sg^3+\mathrm{}^2e^{2\varphi })}{4|g_{00}|g_{rr}g_Sg^3}\left(\frac{g^{}}{g}\right)^2.$$
(17)
We have also
$`{\displaystyle \frac{\ddot{a}}{a}}`$ $`=`$ $`\left(1+{\displaystyle \frac{g}{g^{}}}{\displaystyle \frac{}{r}}\right){\displaystyle \frac{(C+E)^2g_Se^{2\varphi }|g_{00}|(g_Sg^3+\mathrm{}^2e^{2\varphi })}{4|g_{00}|g_{rr}g_Sg^3}}\left({\displaystyle \frac{g^{}}{g}}\right)^2`$ (18)
$`=`$ $`\left(1+{\displaystyle \frac{1}{2}}a{\displaystyle \frac{}{a}}\right){\displaystyle \frac{8\pi }{3}}\rho _{\mathrm{eff}}.`$ (19)
Equating the above to $`(4\pi /3)(\rho _{\mathrm{eff}}+3p_{\mathrm{eff}})`$, we can find the effective pressure
$$p_{\mathrm{eff}}=\rho _{\mathrm{eff}}\frac{1}{3}a\frac{}{a}\rho _{\mathrm{eff}}.$$
(20)
The apparent scalar curvature of the four-dimensional universe is
$$R_{4d}=6\left\{\frac{\ddot{a}}{a}+\left(\frac{\dot{a}}{a}\right)^2\right\}=8\pi \left(4+a_a\right)\rho _{\mathrm{eff}}.$$
(21)
We have given the formalism for simple D3-brane case. The geodesic motion of Dp-brane in the background of the $`\mathrm{Dp}^{}`$-brane with $`p^{}>p`$ can be generalized easily. In the case $`p=p^{}`$, there exists the additional Wess-Zumino term $`T_p\widehat{C}_{p+1}`$ in the DBI action which modifies the equation of the probe brane as well as the induced metric. This modification turn out to be the shift $`EE+C`$ where $`C=C_{0\mathrm{}p}`$ .
## III The type IIB background solution
Here we will consider the background geometry of type IIB theory with five-form flux through an $`S^5`$. We will also assume, for the metric, (3+1)-dimensional Poincaré invariance $`ISO(1,3)`$ since we need the theory defined on the Minkowski space-time. In addition we will preserve the $`SO(6)`$ symmetry of the $`AdS_5\times S^5`$. As a result, the $`ISO(1,3)\times SO(6)`$ invariant ten-dimensional metric with $`N`$ units of five-form flux through an $`S^5`$, in the Einstein frame, can be written as
$`ds_{10}^2`$ $`=`$ $`\widehat{g}_{MN}dx^Mdx^N=e^{\frac{10}{3}\chi +2\sigma }(dt^2+d\stackrel{}{x}^2+dr^2)+L^2e^{2\chi }d\mathrm{\Omega }_5^2,`$ (22)
$`F_5`$ $`=`$ $`{\displaystyle \frac{N\sqrt{\pi }}{2\mathrm{V}\mathrm{o}\mathrm{l}S^5}}(\mathrm{vol}_{S^5}+\mathrm{vol}_{S^5}),`$ (23)
where $`\chi `$, $`\sigma `$, and also the dilaton $`\varphi `$ and the axion $`\eta `$ are allowed to depend only on the radial coordinate $`r`$.
The equations of motion in type IIB supergravity, truncated to the fields of our interests, are given by
$`\widehat{}^2\varphi `$ $`=`$ $`e^{2\varphi }(\eta )^2,`$ (24)
$`\widehat{}^2\eta `$ $`=`$ $`2(_M\varphi )(^M\eta ),`$ (25)
$`\widehat{R}_{MN}`$ $`=`$ $`{\displaystyle \frac{1}{2}}_M\varphi _N\varphi {\displaystyle \frac{1}{2}}e^{2\varphi }_M\eta _N\eta +{\displaystyle \frac{\kappa ^2}{6}}F_{MKLPQ}F_N{}_{}{}^{KLPQ},`$ (26)
where hat means that the operators are expressed in ten-dimensional terms and $`M,N,\mathrm{}=0,\mathrm{},9`$. The equation of motion for the five-form field is
$$\widehat{}_MF^{MKLPQ}=0,$$
(27)
which is satisfied with the self-duality condition (23). The Einstein equation in $`S^5`$ direction is
$$\widehat{R}_{ij}=\left(\frac{\kappa N}{2\pi ^{5/2}}\right)^2g_{ij},i,j=5,\mathrm{},9$$
(28)
which is automatically satisfied if
$$L^4=\frac{\kappa N}{2\pi ^{5/2}}.$$
(29)
The remaining equations can be expressed in purely five-dimensional terms
$`^2\varphi `$ $`=`$ $`e^{2\varphi }(_\mu \eta )^2,`$ (30)
$`^2\eta `$ $`=`$ $`2(_\mu \varphi )(^\mu \eta ),`$ (31)
$`^2\chi `$ $`=`$ $`{\displaystyle \frac{4}{L^2}}\left(e^{\frac{16}{3}\chi }e^{\frac{40}{3}\chi }\right),`$ (32)
$`R_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \varphi _\nu \varphi {\displaystyle \frac{1}{2}}e^{2\varphi }_\mu \eta _\nu \eta +{\displaystyle \frac{40}{3}}_\mu \chi _\nu \chi {\displaystyle \frac{g_{\mu \nu }}{L^2}}\left({\displaystyle \frac{20}{3}}e^{\frac{16}{3}\chi }{\displaystyle \frac{8}{3}}e^{\frac{40}{3}\chi }\right),`$ (33)
where the new metric
$$ds_5^2=g_{\mu \nu }dx^\mu dx^\nu =e^{2\sigma }(dt^2+d\stackrel{}{x}^2+dr^2)$$
(34)
should be used to compute $`R_{\mu \nu }`$ and to contract indices. The equations in (32) can be derived from the five-dimensional action
$$S=\frac{1}{2\kappa _5^2}d^5x\sqrt{g}\left\{R\frac{1}{2}(\varphi )^2+\frac{1}{2}e^{2\varphi }(\eta )^2\frac{40}{3}(\chi )^2+\frac{1}{L^2}\left(20e^{\frac{16}{3}\chi }8e^{\frac{40}{3}\chi }\right)\right\},$$
(35)
where the gravitational couplings $`\kappa _5`$ and $`\kappa `$ in five and ten dimensions are related by
$$\frac{1}{2\kappa _5^2}=\frac{\pi ^3L^5}{2\kappa ^2}=\frac{N^2}{8\pi ^2}\frac{1}{L^3}.$$
(36)
Note the minus sign in front of the axion kinetic term, which is the result of the Hodge-duality rotation of the type-IIB nine-form .
### A Solution without axion
In general one can reduce the equations of motion (32) to a set of coupled non-linear second order ordinary differential equations in $`\varphi `$, $`\eta `$, $`\chi `$, and $`\sigma `$. These equations are too complicated to solve in general, but there is an obvious simplification with $`\chi =\eta =0`$ . Then the equations in (32) become much simpler
$`e^{5\sigma }_re^{3\sigma }_r\varphi =0,`$ (37)
$`_r^2\sigma +3(_r\sigma )^2={\displaystyle \frac{4}{L^2}}e^{2\sigma },`$ (38)
$`4_r^2\sigma ={\displaystyle \frac{1}{2}}(_r\varphi )^2+{\displaystyle \frac{4}{L^2}}e^{2\sigma }.`$ (39)
Equations (38) and (39) are obtained from $`(tt)`$ and $`(rr)`$ components of the Einstein equation. Because $`\chi =0`$ there is no distinction between the five-dimensional Einstein metric and ten-dimensional metric restricted to the five-dimensional noncompact subspace. The equation (37) can be integrated to give
$$\varphi (r)=\varphi _{\mathrm{}}+\frac{B}{L}_0^r𝑑\stackrel{~}{r}e^{3\sigma (\stackrel{~}{r})},$$
(40)
where $`\varphi _{\mathrm{}}`$ is the value of the dilaton at the boundary of the asymptotically $`AdS_5`$ geometry and $`B`$ is an integration constant. Substituting (40) into (38) and (39), defining new variable $`ur/L`$, we obtain
$`(_u\sigma )^2=e^{2\sigma }+{\displaystyle \frac{B^2}{24}}e^{6\sigma },`$ (41)
$`_u^2\sigma =e^{2\sigma }{\displaystyle \frac{B^2}{8}}e^{6\sigma }.`$ (42)
The second equation follows from differentiating the first, so we see that (40), (41) and (42) are consistent system of equations despite being overdetermined.
One can understand (41) as a mechanical analog of a classical particle with unit mass moving in the potential
$$V(\sigma )=\frac{1}{2}e^{2\sigma }\frac{B^2}{48}e^{6\sigma },$$
(43)
with zero enegry. If $`B=0`$, the solution is pure $`AdS_5`$ with constant dilaton. To have a solution with nonconstant dilaton, we take $`B>0`$. However, the $`B0`$ geometry is geodesically incomplete and singular at some point $`u=u_0`$. To find $`u_0`$ explicitly, we integrate equation (41)
$$u=_\sigma ^{\mathrm{}}\frac{d\stackrel{~}{\sigma }}{\sqrt{e^{2\stackrel{~}{\sigma }}+\frac{B^2}{24}e^{6\stackrel{~}{\sigma }}}}=\frac{3^{1/8}\mathrm{\Gamma }(3/8)\mathrm{\Gamma }(1/8)}{8^{7/8}\sqrt{\pi }B^{1/4}}\sqrt{\frac{8}{3}}\frac{e^{3\sigma }}{B}F(\frac{3}{8},\frac{1}{2};\frac{11}{8};\frac{24e^{8\sigma }}{B^2}),$$
(44)
where $`F(\alpha ,\beta ;\gamma ;z)`$ is the usual hypergeometric function. The second term vanishes as $`\sigma \mathrm{}`$, so we find
$$u_0=\frac{3^{1/8}\mathrm{\Gamma }(3/8)\mathrm{\Gamma }(1/8)}{8^{7/8}\sqrt{\pi }B^{1/4}}.$$
(45)
Also we find the dilaton in terms of $`\sigma `$ by solving the equation (40)
$`\varphi `$ $`=\varphi _{\mathrm{}}+\sqrt{{\displaystyle \frac{3}{2}}}\mathrm{coth}^1\sqrt{1+(24/B^2)e^{8\sigma }}`$ (47)
$`=\varphi _{\mathrm{}}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{3}{2}}}\mathrm{ln}{\displaystyle \frac{\sqrt{1+(24/B^2)e^{8\sigma }}+1}{\sqrt{1+(24/B^2)e^{8\sigma }}1}}.`$
We can write the ten-dimensional Einstein metric explicitly if we use $`\sigma `$ as the radial variable
$$ds_{10}^2=e^{2\sigma }(dt^2+d\stackrel{}{x}^2)+\frac{L^2d\sigma ^2}{1+(B^2/24)e^{8\sigma }}+L^2d\mathrm{\Omega }_5^2.$$
(48)
We can cancel the factors of $`B^2/24`$ from (47) and (48) by replacing by $`\sigma \sigma +(1/8)\mathrm{ln}(B^2/24)`$, $`t(B^2/24)^{1/8}t`$ and $`x_i(B^2/24)^{1/8}x_i`$ for $`i=1,2,3`$. If we also rescale $`r(B^2/24)^{1/8}r`$ then the net result is the same as if we had set $`B^2=24`$. Choice of the radial coordinate in $`AdS_5`$ corresponds to choice of one out of a given class of conformally equivalent boundary metrics. Thus the freedom to change $`B`$ in the solution corresponds to the asymptotic scale invariance of the boundary theory.
### B Solution with axion
In the previous subsection we described the simplest case in which one can have a solution with nontrivial dilaton. Here we will consider the case with only $`\chi =0`$ to find the solution with nontrivial axion field . The equations in (32) can be written
$`{\displaystyle \frac{1}{\sqrt{g}}}_\mu \left(\sqrt{g}g^{\mu \nu }_\nu \varphi \right)=e^{2\varphi }_\nu \eta _\mu \eta g^{\mu \nu },`$ (49)
$`{\displaystyle \frac{1}{\sqrt{g}}}_\mu \left(\sqrt{g}g^{\mu \nu }e^{2\varphi }_\nu \eta \right)=0,`$ (50)
$`R_{\mu \nu }={\displaystyle \frac{4}{L^2}}g_{\mu \nu }+{\displaystyle \frac{1}{2}}_\mu \varphi _\nu \varphi {\displaystyle \frac{1}{2}}e^{2\varphi }_\mu \eta _\nu \eta .`$ (51)
Integral of the axion in equation (50) is
$$\eta ^{}=\eta _0e^{3\sigma }e^{2\varphi },$$
(52)
where the prime denotes derivative with respect to $`u=r/L`$ and $`\eta _0`$ is an integration constant. Inserting this expression for $`\eta `$ into equation (49) we obtain the differential equation for dilaton
$$e^{3\sigma }(e^{3\sigma }\varphi ^{})^{}=\eta _0^2e^{2\varphi },$$
(53)
and integrating once we have
$$e^{6\sigma }\varphi ^2=\eta _0^2e^{2\varphi }+\stackrel{~}{B}^2,$$
(54)
where $`\stackrel{~}{B}`$ is another arbitrary constant. Equations (52) and (54) are sufficient to proceed and solve for the function $`\sigma (u)`$ that appears in the metric (51). The Einstein equation (51) becomes
$`_u^2\sigma `$ $`=e^{2\sigma }{\displaystyle \frac{1}{8}}(_u\varphi )^2+{\displaystyle \frac{1}{8}}(_u\eta )^2`$ (55)
$`_u^2\sigma +3(_u\sigma )^2`$ $`=4e^{2\sigma }.`$ (56)
Inserting (52) and (54) into (55) and (56), we obtain
$`_u^2\sigma `$ $`=e^{2\sigma }{\displaystyle \frac{\stackrel{~}{B}^2}{8}}e^{6\sigma },`$ (57)
$`(_u\sigma )^2`$ $`=e^{2\sigma }{\displaystyle \frac{\stackrel{~}{B}^2}{24}}e^{6\sigma }.`$ (58)
The above equations are exactly the same as those without axion (see (41) and (42) ). So the expressions for $`u`$ and $`u_0`$ in (44) and (45) are still valid in the presence of axion field. Also we find the dilaton in terms of $`\sigma `$ by solving the equation (54)
$`e^\varphi `$ $`={\displaystyle \frac{\stackrel{~}{B}}{\eta _0}}\mathrm{sinh}\left\{\mathrm{ln}\left({\displaystyle \frac{\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}+1}{\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}1}}\right)^{(1/2)\sqrt{3/2}}\right\}`$ (60)
$`={\displaystyle \frac{\stackrel{~}{B}}{2\eta _0}}\left\{\left({\displaystyle \frac{\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}+1}{\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}1}}\right)^{(1/2)\sqrt{3/2}}\left({\displaystyle \frac{\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}+1}{\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}1}}\right)^{(1/2)\sqrt{3/2}}\right\}.`$
Finally from equation (52), we can find the solution for axion
$$\eta =\frac{\eta _0}{\stackrel{~}{B}}\left\{\frac{(\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}+1)^{\sqrt{3/2}}+(\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}1)^{\sqrt{3/2}}}{(\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}+1)^{\sqrt{3/2}}(\sqrt{1+(24/\stackrel{~}{B}^2)e^{8\sigma }}1)^{\sqrt{3/2}}}\right\}.$$
(61)
## IV Brane cosmology
In this section we will consider the cosmology probe D3-brane when it is moving along a geodesic in the background type IIB solutions of the previous section.
### A Without axion field
The metric of D3-brane (1) using the background solution (48) is
$$|g_{00}(r)|=e^{2\sigma },g(r)=e^{2\sigma },g_{rr}(r)\frac{L^2d\sigma ^2}{1+(B^2/24)e^{8\sigma }},g_S(r)=L^2.$$
(62)
To apply the formalism of Sec. II we also need to express RR field in terms of $`\sigma `$. From the ansatz for the RR field
$$C_{0123}=C(r),F_{0123r}=\frac{dC(r)}{dr},$$
(63)
equation (27) becomes
$$\frac{dC(r)}{dr}=2Qg^2g_S^{5/2}\sqrt{g_{rr}}$$
(64)
where $`Q`$ is a constant. Using the solution of the metric in (62), the $`RR`$ field can be integrated with appropriate normalization,
$$C=\sqrt{1+(24/B^2)e^{8\sigma }}.$$
(65)
Now we can calculate the effective density on the brane using equations (47), (62) and (65)
$`{\displaystyle \frac{8\pi }{3}}\rho _{\mathrm{eff}}={\displaystyle \frac{1+(24/B^2)e^{8\sigma }}{L^2}}`$ $`[\{E+\sqrt{1+(24/B^2)e^{8\sigma }}\}^2e^{8\sigma }\left({\displaystyle \frac{\sqrt{1+(24/B^2)e^{8\sigma }}+1}{\sqrt{1+(24/B^2)e^{8\sigma }}1}}\right)^{\sqrt{3/2}}`$ (67)
$`\{1+{\displaystyle \frac{\mathrm{}^2}{L^2}}e^{6\sigma }\left({\displaystyle \frac{\sqrt{1+(24/B^2)e^{8\sigma }}+1}{\sqrt{1+(24/B^2)e^{8\sigma }}1}}\right)^{\sqrt{3/2}}\}].`$
If we rescale with $`B^2=24`$ and using $`a=e^\sigma `$, we get
$`{\displaystyle \frac{8\pi }{3}}\rho _{\mathrm{eff}}={\displaystyle \frac{1+a^8}{L^2}}`$ $`[(\sqrt{1+a^8}+{\displaystyle \frac{E}{a^4}})^2\left({\displaystyle \frac{\sqrt{1+a^8}+1}{\sqrt{1+a^8}1}}\right)^{\sqrt{3/2}}`$ (69)
$`\{1+{\displaystyle \frac{\mathrm{}^2}{L^2}}a^6\left({\displaystyle \frac{\sqrt{1+a^8}+1}{\sqrt{1+a^8}1}}\right)^{\sqrt{3/2}}\}],`$
where the range of $`a`$ is $`0<a<\mathrm{}`$, while the range of $`\sigma `$ is $`\mathrm{}<\sigma <\mathrm{}`$. When the universe brane is moving towards the singularity ($`\sigma \mathrm{},a0`$) the universe is contracting while it is moving outward ($`\sigma \mathrm{},a\mathrm{}`$) it is expanding. We also can calculate the scalar curvature of the four-dimensional universe from (21).
Far from the black brane, one can see that $`\rho _{\mathrm{eff}}a^4`$. The cosmological expansion due to the brane motion is indistinguishable from the one by radiation on the brane. This is the idea of the mirage cosmology. If we use the effective density (69) it blows up $`\rho _{\mathrm{eff}}a^{8(2+\sqrt{3/2})}`$ as $`a0`$. Also if we move $`a0`$, the ten-dimensional metric becomes
$$ds_{10}^2=a^2(dt^2+(d\stackrel{}{x})^2)+\frac{L^2da^2}{a^2}+L^2d\mathrm{\Omega }_5,$$
(70)
which is an $`AdS_5\times S^5`$ space. Thus the brane develops an initial singularity as it reaches $`a=0`$ where the description of our formalism breaks down.
### B With axion field
In this case, the only difference on the effective density comes from the form of the dilaton. Still with $`\stackrel{~}{B}^2=24,e^\sigma =a`$, we have
$`{\displaystyle \frac{8\pi }{3}}\rho _{\mathrm{eff}}`$ $`={\displaystyle \frac{1+a^8}{L^2}}`$ (73)
$`\times [\left({\displaystyle \frac{12}{\eta _0}}\right)^2(\sqrt{1+a^8}+{\displaystyle \frac{E}{a^4}})^2\{\left({\displaystyle \frac{\sqrt{1+a^8}+1}{\sqrt{1+a^8}1}}\right)^{\sqrt{3/2}}+\left({\displaystyle \frac{\sqrt{1+a^8}+1}{\sqrt{1+a^8}1}}\right)^{\sqrt{3/2}}2\}`$
$`1{\displaystyle \frac{\mathrm{}^2}{L^2}}\left({\displaystyle \frac{12}{\eta _0}}\right)^2a^6\{\left({\displaystyle \frac{\sqrt{1+a^8}+1}{\sqrt{1+a^8}1}}\right)^{\sqrt{3/2}}+\left({\displaystyle \frac{\sqrt{1+a^8}+1}{\sqrt{1+a^8}1}}\right)^{\sqrt{3/2}}2\}].`$
Near the brane the effective density blows up $`\rho _{\mathrm{eff}}a^{8(2+\sqrt{3/2})}`$ as $`a0`$ which has the same functional dependence as in the case without axion. However, far from the brane ($`a\mathrm{}`$), it gives $`\rho _{\mathrm{eff}}1/L^2`$. This negative cosmological constant means that the expansion of the universe stops at somewhere and eventually recollapses. Comparing equation (73) with (69), we see that the effective density becomes negative faster if there is axion field. The presence of the axion field do not play any important role in the early stage of the evolution but its coupling to other field gives different evolution at late stage.
## V Discussion
We considered the motion of a brane universe moving in a background bulk space of type IIB string theory. For two different backgrounds which give nontrivial dilaton profile, one without axion field and the other with axion, we have derived the Friedman-like equations. These give the cosmological evolution which is similar to the one by matter density on the universe brane. As the brane moves towards the singularity (smaller values of radial coordinate) it contracts and while if it moves away from the black brane it expands. So an observer on the three-brane will see that the universe is expanding. The presence of axion field in the background changes the dilaton profile but it does not change the induced metric. Since dilaton, as well as the induced metric, plays an important role in the effective density, the cosmological evolutions are different for two different backgrounds. For both cases, the effective density blows up as we move toward the singularity showing the initial singularity problem and becomes negative due to the angular momenta $`\mathrm{}^2`$ on the brane meaning the recollapse of the universe. The functional dependence on the radial coordinate shows that when there is axion field in the ambient space the recollapsing of the universe occurs faster compared with the case without axion field. It seems that this phenomenon is true if we do the same calculation with field other than axion.
## Acknowledgement
I would like to thank N. Kaloper, S. P. Kim and E. Kiritsis for discussions and suggestions.
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# 𝒞𝒫ODD ANOMALOUS INTERACTIONS OF HIGGS BOSON IN ITS PRODUCTION AT PHOTON COLLIDERS
## 1 Introduction
In paper we studied potentialities to discover $`𝒞𝒫`$even anomalous interactions of the Higgs boson via its production at $`\gamma \gamma `$ and $`\gamma e`$colliders. Below we bring under analysis effects of $`𝒞𝒫`$-parity violating anomalies. They result in the polarization and azimuthal asymmetries in the Higgs boson production. With new opportunities for variation of photon polarization at Photon Colliders , the Higgs boson production at $`\gamma \gamma `$ and $`\gamma e`$ colliders has an exceptional potential in the extraction of these anomalies. To some extent, some similar issues have been considered in Refs. . However, in the analysis there the polarization potential was not used in its complete form and some natural degrees of freedom in the parameter space were not considered. Besides, the authors cited consider either $`\gamma \gamma `$ or $`\gamma e`$ collisions separately. In the present paper we have in mind that experiments in $`\gamma \gamma `$ and $`\gamma e`$ collider modes will supplement each other and provide complementary opportunities in investigating Higgs boson anomalous interactions.
In our analysis we assume the Higgs boson to be discovered by the time the photon collider starts operating, so that its basic properties will be known by that time. For the definiteness, we assume that the Higgs boson coupling constants will be found experimentally to lie close to their SM values. Our substantial idea is the necessity of step-by-step strategy in studying anomalous effects. Namely, the first step is the study of $`H\gamma \gamma `$ anomalies in $`\gamma \gamma `$collisions and the second step is using $`\gamma e`$ collisions for the study of $`HZ\gamma `$ anomalies assuming $`H\gamma \gamma `$ anomalies (both $`𝒞𝒫`$even and $`𝒞𝒫`$odd) to be studied at the first stage (with higher accuracy than it is possible at the second stage).
The $`\gamma \gamma `$ and $`\gamma e`$ colliders will be the specific modes of the future Linear Colliders (in addition to the $`e^+e^{}`$ mode) with the following typical parameters ($`E`$ and $`_{ee}`$ are the electron energy and luminosity of the basic $`e^+e^{}`$ collider).
* Characteristic photon energy $`E_\gamma 0.8E`$.
* Annual luminosity is typically $`_{\gamma \gamma }200`$ fb<sup>-1</sup>.
* Mean energy spread $`\mathrm{\Delta }E_\gamma 0.07E_\gamma `$.
* Mean photon helicity $`\lambda _\gamma 0.95`$ with variable sign .
* Circular polarization of photons can be transformed into the linear one .
The effective photon spectra for these colliders are given in Ref. . With the above properties, considering photon beams at the Photon Collider as roughly monochromatic is good approximation for our purposes.
The value of effects which can be observed in experiment is given by the expected accuracy in the measuring of the cross sections under interest. For the $`\gamma \gamma `$colliders the expected accuracy in the measuring of Higgs boson decay width will be 2% or better . For $`e\gamma eH`$process we assume the achievable accuracy to be $`5÷10`$%.
Throughout the paper we denote by $`\lambda `$ and $`\zeta /2`$ the average helicities of photons and electrons and by $`\mathrm{}`$ the average degree of the photon linear polarization. We use some $`𝒮`$notations: $`s_W=\mathrm{sin}\theta _W`$, $`c_W=\mathrm{cos}\theta _W`$, $`v_e=14\mathrm{sin}^2\theta _W`$ and $`v=246`$ GeV (Higgs field v.e.v.). In the numerical calculations we assume the Higgs boson to lie in the most expected mass interval 110–250 GeV. Some further notation is borrowed from ref. .
## 2 Sources of $`𝒞𝒫`$violation. Parameterization
We consider below triple Higgs boson anomalous interactions $`H\gamma \gamma `$ and $`HZ\gamma `$ in the processes $`\gamma \gamma H`$ and $`\gamma eeH`$. The quartic interactions lie beyond our scope.
One can imagine two possible mechanisms of $`𝒞𝒫`$violation in the interactions of the Higgs boson. First, the observed Higgs boson can be a mixture of purely scalar and pseudoscalar fields, as it can happen in the multi-doublet Higgs models or in $`𝒮𝒮`$, see for details, e.g. and Sec. 5 as an example. In this case $`𝒞𝒫`$violating effects could be either weak or strong.
The second possibility is that the Higgs boson itself is $`𝒞𝒫`$even fundamentally but underlying interactions can break the $`𝒞𝒫`$parity conservation law. In this case we expect small $`𝒞𝒫`$violating effects in the interactions of Higgs boson with known particles. In turn, this type of $`𝒞𝒫`$violation can be caused either by effects in the underlying theory, similar to the aforementioned mixing, or by fundamental effects related, for example, to the breaking of unitarity of $`S`$-matrix at very small distances. (In the latter case the $`𝒞𝒫`$breaking can originate, in principle, from the possibility that the $`S`$-matrix is unitary only when written it in terms of observable asymptotic states and the unitarity appears broken if the space of states is expanded by adding the unobservable unstable $`H`$ final states.)
A natural question then arises, namely, whether we can distinguish between these two possible causes of $`𝒞𝒫`$violation: i.e. whether the energy scale of $`𝒞𝒫`$violation $`\mathrm{\Lambda }`$ is low or high. In order to answer this question, one should study how corresponding amplitudes depend on additional kinematical variables, such as total energy $`\sqrt{s}`$, photon virtualities $`Q^2`$ etc., i.e. on $`Q^2/\mathrm{\Lambda }^2`$, $`s/\mathrm{\Lambda }^2`$, etc. Indeed, in the first case the dependence on these parameters could be observable, while in the second case the above dimensionless parameters are small and the corresponding amplitudes appear independent from these kinematical variables. (The latter case is described usually with the aid of effective lagrangians.) However, a specific feature of reaction $`\gamma \gamma H`$ is that its kinematics is fixed. This makes it impossible to observe any additional dependence on $`\mathrm{\Lambda }`$. As we turn to process $`e\gamma eH`$, one kinematical degree of freedom appears, namely, the virtuality $`Q^2`$ of the exchanged photon or $`Z`$. However, as shown in , the bulk of the cross section comes from region $`Q^2/M_H^21`$, which again leaves us unable to learn about the source of $`𝒞𝒫`$violation.
The outcome of this discussion can be summarized as follows: when considering real Higgs boson production in the two processes discussed, the above two sources of $`𝒞𝒫`$violation are indistinguishable in the discussed experiments.
Given this, we follow a natural procedure to describe the deviation of discussed production amplitudes from their $`𝒮`$values in a universal manner. We parameterize the $`H\gamma \gamma `$ and $`HZ\gamma `$ amplitudes (which will be also referred to as effective $`H\gamma \gamma `$ and $`HZ\gamma `$ vertices) in the operator form, similar to that for the effective lagrangian:
$$\begin{array}{c}_{\gamma \gamma H}=\frac{1}{v}\left[G_\gamma HF^{\mu \nu }F_{\mu \nu }+i\stackrel{~}{G}_\gamma HF^{\mu \nu }\stackrel{~}{F}_{\mu \nu }\right],\\ _{\gamma ZH}=\frac{1}{v}\left[G_ZHZ^{\mu \nu }F_{\mu \nu }+i\stackrel{~}{G}_ZHZ^{\mu \nu }\stackrel{~}{F}_{\mu \nu }\right].\end{array}$$
(1)
Here $`F^{\mu \nu }`$ and $`Z^{\mu \nu }`$ are the standard field strengths for the electromagnetic and $`Z`$ field and $`\stackrel{~}{F}^{\mu \nu }=\epsilon ^{\mu \nu \alpha \beta }F_{\alpha \beta }/2`$. Dimensionless parameters $`G_i`$ are effective coupling constants. They are sums of well known $`𝒮`$contributions (see e.g. for normalization)<sup>1</sup><sup>1</sup>1 With the proposed experimental accuracy, when doing the final numerical calculation, one should, of course, use $`H\gamma \gamma `$ coupling with radiative corrections . and anomalous parts $`g_i`$ (”anomalies”), describing the strength of interactions beyond $`𝒮`$, which are generally complex:
$$G_i=G_i^{SM}+g_i,\stackrel{~}{G}_i=g_{Pi};g_a=|g_a|e^{i\xi _a}.$$
(2)
The complex values of ”couplings” $`g_a`$ are quite natural. Indeed, recall that even $`G_i^{SM}`$ are complex due to contributions, for example, of $`b`$-quark loop in the amplitude. The same is valid in various versions of the first variant of $`𝒞𝒫`$violation. One particular example of this is discussed in Sec. 5, where the anomaly can be defined simply as the difference between the minimal $`𝒮`$and Two Doublet Higgs Model (II) with $`𝒞𝒫`$violation. If $`\mathrm{tan}\beta 1`$, contribution of $`b`$ quarks in loops is enhanced, which gives rise to the large imaginary part of the amplitudes. For the second mechanism complex $`g_i`$ could be signal of fundamental breaking of unitarity in theory.
We assume that future observations will reveal a picture close to $`𝒮`$ and therefore anomalies $`g_i`$ will be small. In the first mechanism of $`𝒞𝒫`$violation with $`\mathrm{\Lambda }<M_H`$ smallness of anomalies is related to small values of corresponding mixing angles $`\alpha _m`$, $`g_i\alpha _m`$. In the second mechanism it is related to large scale of New Physics $`\mathrm{\Lambda }`$, i.e. $`g_i=(v/\mathrm{\Lambda }_i)^2`$ with $`\mathrm{\Lambda }_i\mathrm{\Lambda }`$. The relation between parameters $`\mathrm{\Lambda }_i`$ and $`\mathrm{\Lambda }`$ depends on the nature of New Physics.
(A) The simplest extension of the $`𝒮`$consists in adding new charged heavy particles with mass $`M_n`$ that is not generated by a Higgs mechanism (like in MSSM). They will circulate in loops and give rise to anomalous effective $`H\gamma \gamma `$ and $`HZ\gamma `$ vertices, with $`\mathrm{\Lambda }^24\pi M_n^2/\alpha `$.
(B) If the heavy particle is a point-like Dirac monopole, then $`\mathrm{\Lambda }^2\alpha M_n^2`$.
(C) If New Physics is determined by higher dimension (Kaluza–Klein) mechanism, the quantity $`\mathrm{\Lambda }`$ is close to the energy scale at which the extra dimensions come into play.
For the second mechanism the anomalous amplitude is often described with the aid of Effective Lagrangian with operators of dimension 6, which has the same form as our effective vertices (1). Our particular parameterization can be readily linked to that used in other papers (e.g. ). For example, correspondence of our parameters $`g_i`$ to constants $`d_i`$ used in ref. reads $`d_{Z\gamma }=2g_{Z\gamma }/(c_Ws_W)`$, $`\overline{d}_{Z\gamma }=g_{PZ}/(c_Ws_W)`$.
Finally, we undertake a study where both $`|g_i|`$ and $`\xi _i`$ are treated as independent parameters. This is done in contrast to other similar investigations, where the complexity was not an explicitly free parameter, but fixed by the particular model considered. We argue that our approach accounts for the most wide range of possible anomalies. Determination of both sets of parameters should be considered primarily as an experimental task<sup>2</sup><sup>2</sup>2 Certainly, only phase differences are measurable for entire effective couplings. Expecting relatively small magnitude of anomaly, one can conclude that the phases of entire quantities $`G_\gamma `$, $`G_Z`$ are close to their $`𝒮`$values $`\xi _\gamma ^{SM}`$ and $`\xi _Z^{SM}`$ and the effect of anomaly itself is reduced by factor $`\mathrm{cos}(\xi _\gamma \xi _\gamma ^{SM})`$..
About figures and notation there. Currently, due to the large number of new model parameters, a thorough investigation of regions of the parameter space, achievable in future experiments, makes little sense. Instead of that we present in our figures examples for some values of parameters, which illustrate that the study of these effects at the Photon Colliders is indeed possible and profitable.
There are no doubts that relatively large anomalies will be discovered easily. Therefore, we concentrate our efforts on the case when the anomalous effects are relatively small as compared with basic $`𝒮`$effects. In this case the effects of anomalies will be seen mainly in the interference with the $`𝒮`$effects, and contributions of different anomalies in the observed cross sections are additive with good accuracy. This is why we treat each anomaly separately, assuming all other anomalies absent (the corresponding $`g_i=0`$).
## 3 Process $`𝜸𝜸\mathbf{}𝑯`$
Let us denote by $`\sigma ^{SM}_{np}`$ the $`𝒮`$Higgs boson production cross section in unpolarized photon collisions averaged over a certain small effective mass interval (see e.g. ). Then the cross section of the Higgs boson production can be written in the form:
$$\begin{array}{c}\sigma (\lambda _i,\mathrm{}_i,\psi )=\sigma ^{SM}_{np}T(\lambda _i,\mathrm{}_i,\psi );\\ T(\lambda _i,\mathrm{}_i,\psi )=\frac{\left|G_\gamma \right|^2}{\left|G_\gamma ^{SM}\right|^2}(1+\lambda _1\lambda _2+\mathrm{}_1\mathrm{}_2\mathrm{cos}2\psi )+\frac{\left|\stackrel{~}{G}_\gamma \right|^2}{\left|G_\gamma ^{SM}\right|^2}\left(1+\lambda _1\lambda _2\mathrm{}_1\mathrm{}_2\mathrm{cos}2\psi \right)\\ +2\frac{Re\left(G_\gamma ^{}\stackrel{~}{G}_\gamma \right)}{\left|G_\gamma ^{SM}\right|^2}(\lambda _1+\lambda _2)+2\frac{Im\left(G_\gamma ^{}\stackrel{~}{G}_\gamma \right)}{\left|G_\gamma ^{SM}\right|^2}\mathrm{}_1\mathrm{}_2\mathrm{sin}2\psi .\end{array}$$
(3)
Here $`\lambda _i`$ and $`\mathrm{}_i`$ ($`i=1,2`$) are degrees of circular and linear polarization respectively of the photon beams and $`\psi `$ is the polar angle between the linear polarization vectors of the two photon beams.
In the $`𝒮`$case we have only the first item in this sum. (Note that the $`\gamma \gamma b\overline{b}`$ background is practically independent on linear polarization of photons.)
An important feature here is interference terms. They give rise to the inequality of the two directions of rotation and to the modification of the $`\psi `$-dependence, which is entirely due to the $`𝒞𝒫`$odd admixture to $`𝒞𝒫`$even Lagrangian. Owing to these modifications, a number of experimentally measurable quantities appear that can help study $`𝒞𝒫`$even and odd anomalies separately.
It is useful to introduce five different asymmetries:
$$\begin{array}{c}T_\pm =\frac{\sigma \left(\lambda _i,\mathrm{}_i=0\right)\pm \sigma \left(\lambda _i,\mathrm{}_i=0\right)}{2\sigma ^{SM}_{np}}\{\begin{array}{c}(1+\lambda _1\lambda _2)(|\stackrel{~}{G}_\gamma |^2+|G_\gamma |^2),\hfill \\ 2(\lambda _1+\lambda _2)Re(\stackrel{~}{G}_\gamma ^{}G_\gamma );\hfill \end{array}\\ T_{}=\frac{\sigma \left(\lambda _i=0,\mathrm{}_i,\psi =0\right)}{\sigma ^{SM}_{np}}\left[|G_\gamma |^2(1+\mathrm{}_1\mathrm{}_2)+|\stackrel{~}{G}_\gamma |^2(1\mathrm{}_1\mathrm{}_2)\right],\\ T_{}=\frac{\sigma \left(\lambda _i=0,\mathrm{}_i,\psi =\pi /2\right)}{\sigma ^{SM}_{np}}\left[|G_\gamma |^2(1\mathrm{}_1\mathrm{}_2)+|\stackrel{~}{G}_\gamma |^2(1+\mathrm{}_1\mathrm{}_2)\right],\\ T_\psi =\frac{\sigma \left(\lambda _i=0,\mathrm{}_i,\psi =3\pi /4\right)\sigma \left(\lambda _i=0,\mathrm{}_i,\psi =\pi /4\right)}{\sigma ^{SM}_{np}}2\mathrm{}_1\mathrm{}_2Im(\stackrel{~}{G}_\gamma ^{}G_\gamma ),\end{array}$$
(4)
whose $`𝒮`$values are
$$T_+^{SM}=1+\lambda _1\lambda _2,T_{}^{SM}=0,T_{}^{SM}=1+\mathrm{}_1\mathrm{}_2,T_{}^{SM}=1\mathrm{}_1\mathrm{}_2,T_\psi ^{SM}=0.$$
The quantities $`T_+`$, $`T_{}`$ and $`T_{}`$ are combinations of $`|G_\gamma |^2`$ and $`|\stackrel{~}{G}_\gamma |^2`$ with different weights. These asymmetries are sensitive to the $`𝒞𝒫`$even anomaly and its phase $`\xi _\gamma `$ via its interference with the $`𝒮`$contribution. The best quantity for this study is of course $`T_+`$, which is illustrated by Fig. 2. Certainly, curves for $`𝒞𝒫`$even anomaly effects at $`\xi _\gamma =0`$ are the same as obtained in Ref. (modulo to reparameterization of anomalous terms). These three quantities include also the $`𝒞𝒫`$odd anomaly in the form $`|\stackrel{~}{G}_\gamma |^2`$, which is $`g_{P\gamma }^2`$, i.e. small and independent of $`\xi _{P\gamma }`$ (the corresponding $`g_{P\gamma }`$ dependence was studied in ref. ). Even in the case of $`T_{}`$, where the contribution of $`|\stackrel{~}{G}_\gamma |^2`$ is enhanced, it is difficult to see the effect of $`𝒞𝒫`$-odd anomalies at reasonably small $`g_{P\gamma }`$, Fig.4.
The remaining two quantities — $`T_{}`$ and $`T_\psi `$ — are much more useful for study of $`𝒞𝒫`$violating effects in $`\gamma \gamma H`$ interaction. Their study supplements each other. Both of them differ from zero only if the $`𝒞𝒫`$parity is broken. They derive from the interference of the $`𝒞𝒫`$odd and $`𝒞𝒫`$even items in (1). Fig. 2 shows the $`T_{}`$ dependence on $`|g_{P\gamma }|`$ and phase $`\xi _{P\gamma }`$ for different values of the Higgs boson mass. At $`M_H<160`$ GeV ($`WW`$ threshold) the basic quantity $`G_\gamma ^{SM}`$ is practically real. Therefore, the quantity $`T_{}`$ has maximum at $`\xi _{P\gamma }=0`$. Above this threshold the imaginary part of $`G_\gamma ^{SM}`$ becomes substantial, and the position of maximum is shifted to $`\xi _{P\gamma }0`$. Fig. 4 shows that the $`𝒞𝒫`$odd anomaly effect is strong in this asymmetry as well, and exhibits a remarkable dependence of $`T_\psi `$ on the value of phase $`\xi _{P\gamma }`$. With measurement of $`T_{}`$ and $`T_\psi `$ one can extract from the data both $`|g_{P\gamma }|`$ and $`\xi _{P\gamma }`$ since $`T_{}`$ and $`T_\psi `$ represent the real and imaginary part of the same quantity.
## 4 Process $`𝒆𝜸\mathbf{}𝒆𝑯`$
The process $`e\gamma eH`$ is considered here as a good tool for study of $`HZ\gamma `$ anomalous interactions provided $`H\gamma \gamma `$ anomalies are known from the experiments in the $`\gamma \gamma `$ mode. This process was studied within $`𝒮`$in detail in refs. <sup>3</sup><sup>3</sup>3 The production of the pseudoscalar Higgs boson in such a reaction was studied e.g. in ref. , see also ref. for the $`𝒮𝒮`$ case.. It is described by diagrams of three types — those with photon exchange in $`t`$-channel, with $`Z`$ exchange in $`t`$ channel and box diagrams. This subdivision is approximately gauge invariant with accuracy $`m_e/M_Z`$ . The difference in the cross sections $`\sigma ^L`$ and $`\sigma ^R`$ for the left-hand and right-hand polarized electrons is obliged to interference between photon and $`Z`$ exchange amplitudes.
The main contribution into the total cross section is given by diagrams with photon exchange in $`t`$-channel. Therefore, this total cross section is sensitive to the $`H\gamma \gamma `$ anomalies and weakly sensitive to the $`HZ\gamma `$ anomalies, which are our major concern here (the difference $`\sigma ^L\sigma ^R`$ is small as compared with the unpolarized cross section). This picture is improved with the growth of transverse momentum of the scattered electron $`p_{}`$. Indeed, with this growth photon exchange contribution is strongly reduced, while $`Z`$–boson exchange contribution changes only marginally at $`p_{}<M_Z`$. At transverse momenta of the scattered electrons $`p_{}>30`$ GeV and for longitudinally polarized initial electrons the effect of Z-exchange should be seen well . To feel the scale of observed effects, we present in Fig. 5 the $`𝒮`$cross sections $`\sigma ^L`$ and $`\sigma ^R`$ integrated over the region $`Q^2>1000`$ GeV<sup>2</sup> and averaged over initial photon polarizations. We use this limitation in $`Q^2`$ everywhere below.
We denote the particle momenta as $`p`$ for the incident electron, $`k`$ for the photon, $`p^{}=pq`$ for the scattered electron and $`Q^2=q^2`$. In our calculations far from the photon pole in $`t`$-channel we neglect the electron mass. We also denote: $`u=2kp^{}=M_H^2+Q^2s`$, $`x=2kq/s(M_H^2+Q^2)/s`$, $`E_{tot}=\sqrt{s}`$. The collision axis is labeled as $`z`$ axis and $`x`$ axis is chosen along the direction of the photon linear polarization vector $`\stackrel{}{\mathrm{}}`$. Finally, angle $`\varphi `$ is the azimuthal angle of the scattered electron relative to so-defined $`x`$ axis. The values $`\zeta =1`$ or $`\zeta =+1`$ correspond to left-hand or right-hand polarized initial electrons. We use superscripts $`L`$ and $`R`$ to label quantities referring to these polarizations.
The qualitative features of the observable effect could be understood taking into account that the quantities below could be treated as the average helicity $`\lambda _V`$ and degree of linear polarization $`\mathrm{}_V`$ of an exchanged virtual photon or $`Z`$ boson:
$$\lambda _V=\frac{s^2u^2}{s^2+u^2}\zeta =\frac{xx^2/2}{1x+x^2/2}\zeta ,\mathrm{}_V=\frac{2s|u|}{s^2+u^2}=\frac{1x}{1x+x^2/2},$$
(5)
with vector of linear polarization $`\stackrel{}{\mathrm{}}_V`$ lying in the electron scattering plane . Since usually $`x1`$, we have $`\lambda _V1`$ and $`\mathrm{}_V1`$. Therefore, joining the results of the previous section and those from ref., one can conclude that the effect of $`𝒞𝒫`$odd $`HZ\gamma `$ interaction can be seen in the dependence on angle $`\varphi `$ in the experiments with left– and right– polarized electrons and in the study of dependence on the sign of the incident photon helicity. These dependencies have not been studied earlier.
Helicity amplitudes of the process are calculated just as in Ref. . With notations for the box contributions from that paper we have (in these equations helicities $`\lambda ,\zeta =\pm 1`$)
$$\begin{array}{c}=\frac{4\pi \alpha }{M_Ws_W}\sqrt{\frac{Q^2}{2}}\left\{s\frac{1+\zeta \lambda }{2}+(sM_H^2Q^2)\left[\frac{1\zeta \lambda }{2}\mathrm{cos}2\varphi +\frac{\zeta \lambda }{2}i\mathrm{sin}2\varphi \right]\right\}\\ \times (\lambda K+\stackrel{~}{K}),(K=V\zeta A+B_+,\stackrel{~}{K}=\stackrel{~}{V}\zeta \stackrel{~}{A}+\zeta B_{}).\end{array}$$
(6)
Here $`V`$ and $`A`$ stand for vector and axial $`t`$–channel exchange contributions, $`B_\pm `$ are the box contributions which are composed from items related to the $`W`$ or $`Z`$ circulating in box<sup>4</sup><sup>4</sup>4 The box diagrams contribution (and their interference with other diagrams) is small in comparison with other contributions.:
$$\begin{array}{c}V=\frac{G_\gamma }{Q^2}+\frac{v_eG_Z}{4s_Wc_W\left(Q^2+M_Z^2\right)},A=\frac{G_Z}{4s_Wc_W\left(Q^2+M_Z^2\right)},\\ \stackrel{~}{V}=\frac{\stackrel{~}{G}_\gamma }{Q^2}+\frac{v_e\stackrel{~}{G}_Z}{4s_Wc_W\left(Q^2+M_Z^2\right)},\stackrel{~}{A}=\frac{\stackrel{~}{G}_Z}{4s_Wc_W\left(Q^2+M_Z^2\right)};\\ B_\pm =\frac{\alpha M_W^2}{4\pi s_W^2}\left[\frac{W(s,u)\pm W(u,s)}{2}+\frac{Z(s,u)\pm Z(u,s)}{2}\right].\end{array}$$
(7)
The amplitude squared for an arbitrarily polarized photon beam can be written in terms of helicity amplitudes and the photon density matrix $`\rho `$ written in helicity basis as
$$||^2=_a^{}\rho _{ab}_b,a,b=+,,\rho =\frac{1}{2}\left(\begin{array}{ccc}1+\lambda & & \mathrm{}\\ \mathrm{}& & 1\lambda \end{array}\right).$$
(8)
So that the cross section reads (here $`\zeta =\pm 1`$):
$$\begin{array}{cc}\hfill d\sigma =& \frac{\pi \alpha ^2}{2M_W^2s_W^2}\frac{d\varphi }{2\pi }Q^2dQ^2\frac{s^2+u^2}{2s^2}\left(U_0+\lambda U_\lambda +\mathrm{}\mathrm{cos}2\varphi U_{}\mathrm{}\mathrm{sin}2\varphi U_\psi \right);\hfill \\ & U_0=\left(|K|^2+|\stackrel{~}{K}|^2\right)+\lambda _V2Re\left(K\stackrel{~}{K}^{}\right),U_{}=\mathrm{}_V\left(|K|^2|\stackrel{~}{K}|^2\right)\hfill \\ & U_\lambda =2Re\left(K\stackrel{~}{K}^{}\right)+\lambda _V\left(|K|^2+|\stackrel{~}{K}|^2\right),U_\psi =2Im\left(K\stackrel{~}{K}^{}\right).\hfill \end{array}$$
(9)
With notations (5) it becomes evident that this equation reproduces term by term the polarization dependencies of $`\gamma \gamma H`$ process (3), in particular, $`T_+,T_{}U_0`$, $`T_{}U_{}`$, $`T_{}U_\lambda `$, $`T_\psi U_\psi `$. Therefore, the similar studies of $`HZ\gamma `$ interaction are possible here. However, there is a difference between effects of linear photon polarization in these two reactions. In the $`\gamma \gamma `$ collisions we can control linear polarizations and choose their relative orientation to study specific contribution. In the $`\gamma e`$ collision we cannot control relative orientation of linear polarizations, so that some Fourier-type analysis is necessary to see contributions under interest.
Different asymmetries. The quantities $`U_0`$ and $`U_{}`$ are weakly sensitive to the $`\stackrel{~}{G}_Z`$.The sensitivity of $`U_0`$ to the $`𝒞𝒫`$even anomalous interaction was studied, in fact, in refs. .
The quantities $`U_\lambda `$ and $`U_\psi `$ are most sensitive to the $`𝒞𝒫`$odd anomalies. Thus, we consider asymmetries
$$\begin{array}{c}V_\lambda ^{L,R}=\frac{{\displaystyle 𝑑\sigma ^{L,R}\left(\lambda \right)}{\displaystyle 𝑑\sigma ^{L,R}\left(\lambda \right)}}{\left|\lambda \right|{\displaystyle 𝑑\sigma _{np}^{SM}}}U_\lambda ^{L,R},\\ V_\psi ^{L,R}=\frac{{\displaystyle 𝑑\sigma ^{L,R}\mathrm{sin}2\varphi }}{\left|\mathrm{}\right|{\displaystyle 𝑑\sigma _{np}^{SM}}}U_\psi ^{L,R},\end{array}$$
(10)
with integrations spanning over the region $`Q^2>Q_0^2=1000`$ GeV<sup>2</sup> and the whole region of $`\varphi `$ for the left–hand and right–hand polarized initial electrons. (The integrals in denominators are calculated for the nonpolarized initial particles.) It happens that the cross sections for the left-hand polarized electrons are much higher than those for the right-handed electrons (see Fig. 5). Therefore, we present graphs for the left-handed electrons only. The anomalous effect for the right-handed electrons is also small in its absolute value. We have not encountered any case where $`\sigma ^R`$ would be a useful source of additional information, despite the relative value of anomaly contribution can be higher here.
The quantity $`V_\lambda ^L`$ describing the helicity asymmetry is analogous to $`T_{}`$ in the $`\gamma \gamma `$ case with accuracy to contribution $`(|K|^2+|\overline{K}|^2)`$ entering with small coefficient $`\lambda _V`$. This contribution results in non-zero $`V_\lambda `$ even in $`𝒮`$. Figs. 7 shows dependence of this quantity on $`|g_{PZ}|`$. For the purposes of comparison, the effect of a $`g_{P\gamma }=0.310^3`$ $`H\gamma \gamma `$ anomaly is also shown. We see that the values of this helicity asymmetry are large enough. Note that the signal/background ratio improves with the growth of energy since the $`𝒮`$contribution into the discussed quantity decreases approximately $`\lambda _Vs^1`$ while the anomaly effect increases weakly, $`\mathrm{ln}(s/M_Z^2)`$.
The same figure depicts also the quantity $`V_\psi ^L`$ at different values of $`|g_{PZ}|`$. Again we also draw a comparison with a $`H\gamma \gamma `$ $`𝒞𝒫`$-odd anomaly. This quantity is intrinsically smaller than $`V_\lambda ^L`$, so the $`𝒞𝒫`$-odd $`HZ\gamma `$ anomaly can be seen only at $`|g_{PZ}|>10^3`$.
The dependence of $`V_\lambda `$ and $`V_\psi `$ on the phase of $`HZ\gamma `$ anomaly $`\xi _{PZ}`$ is shown in Fig. 7. (The dependence of these quantities on the parameters of $`H\gamma \gamma `$ anomaly has the same form but the magnitude is somewhat larger.) These curves closely resemble dependencies of $`T_{}`$ and $`T_\psi `$ on $`\xi _{P\gamma }`$ in the $`\gamma \gamma H`$ case. We see the familiar phase dependence $`\mathrm{cos}(\xi _{PZ}\overline{\xi }_\gamma )`$ or $`\mathrm{sin}(\xi _{PZ}\overline{\xi }_\gamma )`$ (here $`\overline{\xi }_i`$ are phases of $`G_\gamma `$ and $`G_Z`$ which are close to their $`𝒮`$values). The effect of switching on of the imaginary part of the $`𝒮`$contribution at $`M_H160`$ GeV is clearly seen in these curves. In the phenomenological analysis, it is helpful that $`V_\lambda `$ and $`V_\psi `$ are intrinsically complementary: just as it was in $`\gamma \gamma H`$ case, $`V_\lambda `$ is the real part and $`V_\psi `$ is the imaginary part of the same quantity. Therefore, at any value of $`M_H`$ and $`\xi _{PZ}`$ either $`V_\lambda `$ or $`V_\psi `$ will deviate strongly from the $`𝒮`$value.
## 5 Scalar-pseudoscalar mixing within two doublet Higgs model
A specific case of $`𝒞𝒫`$violation takes place in the scalar–axial mixing within the two doublet Higgs model (2HDM). This model is described with the aid of mixing angle $`\beta `$ (defined via the ratio of v.e.v.’s for two basic scalar fields, $`\mathrm{tan}\beta =\varphi _1/\varphi _2`$) and three Euler mixing angles $`\alpha _1`$, $`\alpha _2`$, $`\alpha _3`$ (see, for example, ref. ). The observed neutral Higgs bosons are combined from the basic scalar fields as
$$\begin{array}{c}\left(\begin{array}{c}h_1\\ h_2\\ h_3\end{array}\right)=\sqrt{2}R\left(\begin{array}{c}Re\varphi _1^0\\ Re\varphi _2^0\\ Im\left(s_\beta \varphi _1^0c_\beta \varphi _2^0\right)\end{array}\right),\\ R=\left(\begin{array}{ccc}c_1& s_1c_2& s_1s_2\\ s_1c_3& c_1c_2c_3s_2s_3& c_1s_2c_3c_2s_3\\ s_1s_3& c_1c_2s_3+s_2c_3& c_1s_2s_3+c_2c_3\end{array}\right)\end{array}$$
(11)
Here $`c_i=\mathrm{cos}\alpha _i,s_i=\mathrm{sin}\alpha _i`$. Our definition differs from that used in Ref. by the sign minus in front of $`R`$ in (11). The $`𝒞𝒫`$conserving case is realized at $`\alpha _2=\alpha _3=0`$, the last angle $`\alpha _1`$ is related to the quantity $`\alpha `$ used for the case without $`𝒞𝒫`$violation as $`\alpha _1\pi /2\alpha `$, $`h_1h`$, $`h_2A`$, $`h_3H`$. Instead of $`\alpha _1`$, we use below the angle $`\delta =\beta (\pi /2\alpha _1)`$.
We consider only the lightest Higgs boson $`h_1`$ having in mind the decoupling regime where $`M_{H^\pm },M_{h_2},M_{h_3}M_{h_1}`$. Besides, we fix the only relevant free parameter of 2HDM in the Higgs self-interaction as $`\lambda _5=2M_{H^\pm }^2/v^2+g^2`$ (just as it is in MSSM, see for definition). This choice guarantees us negligibly small contribution of charged Higgs loops into the discussed couplings of Higgs boson with photons.
To describe couplings of the lightest Higgs boson $`h_1`$ with quarks and charged leptons we use the widespread ”Model II” in which the ratios of these couplings to those in the minimal $`𝒮`$ (one Higgs doublet) are
$$\begin{array}{c}\overline{u}h_1u(\mathrm{sin}\delta +\mathrm{cot}\beta \mathrm{cos}\delta )\mathrm{cos}\alpha _2i\gamma ^5\mathrm{cot}\beta \mathrm{cos}(\delta \beta )\mathrm{sin}\alpha _2,\\ \overline{d}h_1d,\overline{\mathrm{}}h_1\mathrm{}(\mathrm{sin}\delta \mathrm{tan}\beta \mathrm{cos}\delta )i\gamma ^5\mathrm{tan}\beta \mathrm{cos}(\delta \beta )\mathrm{sin}\alpha _2,\\ VVh_1\mathrm{sin}\delta \mathrm{sin}\beta \mathrm{cos}(\delta \beta )(1\mathrm{cos}\alpha _2).\end{array}$$
(12)
The effective couplings of Higgs boson with light $`G_i`$ (1) can be written via standard loop integrals and the above mixing angles (see for definitions).
$$\begin{array}{c}G^\gamma =G_{SM}^\gamma \mathrm{sin}\delta +\frac{\alpha }{12\pi }\mathrm{cos}\delta \left[\mathrm{\Phi }_{1/2}(b)\mathrm{tan}\beta +4\mathrm{\Phi }_{1/2}(t)\mathrm{cot}\beta \right]+\text{scalars}\\ \frac{\alpha }{12\pi }(1\mathrm{cos}\alpha _2)\left[3\mathrm{\Phi }_1^\gamma (W)\mathrm{sin}\beta \mathrm{cos}(\delta \beta )+4\mathrm{\Phi }_{1/2}(t)(\mathrm{sin}\delta +\mathrm{cot}\beta \mathrm{cos}\delta )\right],\\ \stackrel{~}{G}^\gamma =\frac{\alpha }{12\pi }\left[\mathrm{\Phi }_{1/2}^A(b)\mathrm{tan}\beta +4\mathrm{\Phi }_{1/2}^A(t)\mathrm{cot}\beta \right]\mathrm{cos}(\delta \beta )\mathrm{sin}\alpha _2,\\ G^Z=G_{SM}^Z\mathrm{sin}\delta +\frac{\alpha }{4\pi }\left[v_b\mathrm{\Phi }_{1/2}(b)\mathrm{tan}\beta +2v_t\mathrm{\Phi }_{1/2}\mathrm{cot}\beta \right]\mathrm{cos}\delta +\text{scalars}\\ \frac{\alpha }{4\pi }(1\mathrm{cos}\alpha _2)\left[\mathrm{\Phi }_1^Z(W)\mathrm{sin}\beta \mathrm{cos}(\delta \beta )+2v_t\mathrm{\Phi }_{1/2}(t)(\mathrm{sin}\gamma +\mathrm{cot}\beta \mathrm{cos}\delta )\right],\\ \stackrel{~}{G}^Z=\frac{\alpha }{4\pi }\left[2v_t\mathrm{\Phi }_{1/2}^A(t)\mathrm{cot}\beta v_b\mathrm{\Phi }_{1/2}^A(b)\mathrm{tan}\beta \right]\mathrm{cos}(\delta \beta )\mathrm{sin}\alpha _2;\\ v_b=\frac{34s_w^2}{12s_wc_w},v_t=\frac{38s_w^2}{12s_wc_w}\end{array}$$
(13)
The first lines in formulas for $`G_\gamma `$ and $`G_Z`$ give their form for the standard 2-doublet model without $`𝒞𝒫`$-mixing. At large $`\mathrm{tan}\beta `$ the imaginary part of all these couplings (arising from the $`b`$-quark contribution) becomes essential. It gives phases $`\xi _i`$ (2) which differ essentially from 0 or $`\pi `$. The corresponding values of $`g_i`$ and phases $`\xi _i`$ (2) could be calculated easily from these equations. The word scalars means charged Higgs loop contribution, it is negligibly small in the discussed case, so we will not write it below.
Finally, all box diagrams include $`VVh`$ vertex. Therefore the box contribution (7) to the amplitude changes as
$$B_\pm B_\pm ^{SM}\left[\mathrm{sin}\delta \mathrm{sin}\beta \mathrm{cos}(\delta \beta )(1\mathrm{cos}\alpha _2)\right].$$
(14)
To make new effects more manifest, we study the dependence on two parameters $`\alpha _2`$ and $`\beta `$ only, keeping main features of discussed Higgs boson $`h_1`$ as close as possible to the Higgs boson of $`𝒮`$. For this purpose we fix parameter $`\delta \pi /2`$ and consider small enough values of $`𝒞𝒫`$–violated mixing angle $`\alpha _2`$. According to eq. (12), in this case couplings of $`h`$ with quarks and gauge bosons are close to those in $`𝒮`$(see refs. for the detail discussion of this opportunity). In this case we have instead of previous equations
$$\begin{array}{c}\overline{u}h_1u\mathrm{cos}\alpha _2i\gamma ^5\mathrm{cos}\beta \mathrm{sin}\alpha _2,\overline{d}h_1d1i\gamma ^5\mathrm{tan}\beta \mathrm{sin}\beta \mathrm{sin}\alpha _2,\\ VVh_11\mathrm{sin}^2\beta (1\mathrm{cos}\alpha _2).\end{array}$$
(15)
$$\begin{array}{c}G^\gamma =G_{SM}^\gamma \frac{\alpha }{12\pi }(1\mathrm{cos}\alpha _2)\left[3\mathrm{\Phi }_1^\gamma (W)\mathrm{sin}^2\beta +4\mathrm{\Phi }_{1/2}(t)\right],\\ \stackrel{~}{G}^\gamma =\frac{\alpha }{12\pi }\left[\mathrm{\Phi }_{1/2}^A(b)\mathrm{tan}\beta +4\mathrm{\Phi }_{1/2}^A(t)\mathrm{cot}\beta \right]\mathrm{sin}\beta \mathrm{sin}\alpha _2,\\ G^Z=G_{SM}^Z\frac{\alpha }{4\pi }(1\mathrm{cos}\alpha _2)\left[\mathrm{\Phi }_1^Z(W)\mathrm{sin}^2\beta +2v_t\mathrm{\Phi }_{1/2}(t)\right],\\ \stackrel{~}{G}^Z=\frac{\alpha }{4\pi }\left[2v_t\mathrm{\Phi }_{1/2}^A(t)\mathrm{cot}\beta v_b\mathrm{\Phi }_{1/2}^A(b)\mathrm{tan}\beta \right]\mathrm{sin}\beta \mathrm{sin}\alpha _2,\\ B_\pm =B_\pm ^{SM}[1\mathrm{sin}\beta \mathrm{cos}\beta (1\mathrm{cos}\alpha _2)].\end{array}$$
(16)
Fig. 8 presents the overall dependence on $`\alpha _2`$. The strong oscillations might seem surprising. To explain them on the example of $`T_+`$, let us first note that at $`\alpha _2\pi /2`$ and $`\mathrm{tan}\beta 1`$ the boson $`h_1`$ becomes almost pseudoscalar. Next, it is well known that the two-photon decay width of the pseudoscalar is significantly smaller than the $`h\gamma \gamma `$ decay width. Therefore, quantity $`T_+`$ should be close to zero at $`\alpha _2\pi /2`$. For more details, one can consider this quantity $`T_+`$ for the case $`\mathrm{tan}\beta =3`$, for definiteness. In this case $`\mathrm{sin}^2\beta =0.9`$. By definition, $`T_+|G|^2+|\stackrel{~}{G}|^2`$, the $`W`$ contribution in the first term plays the dominant role everywhere except for narrow region near $`\alpha _2=\pi /2`$ and thus dictates the shape $`T_+[10.9(1\mathrm{cos}\alpha _2)]^2+`$ small remnants. At $`\alpha _2`$ slightly above $`\pi /2`$, when $`t`$-quark exactly cancels the remnant of $`W`$ boson contribution (and the real part of the $`b`$ contribution), $`T_+`$ is saturated by $`|\stackrel{~}{G}|^2`$, which is intrinsically smaller than $`|G_{SM}|^2`$ by two orders of magnitude. The shape of $`T_{}`$, etc. dependence on $`\alpha _2`$ can also be foreseen from Eq.(16) in the same way. Our calculations show that the quantities $`T_{}`$, $`T_\psi `$ as well as asymmetries $`V_i`$ of the $`e\gamma eH`$ reaction also exhibit a similar oscillatory dependence on $`\alpha _2`$. The principal features of the results remain the same for other values of Higgs boson masses, including region $`M_h>2M_W`$ above the WW threshold.
However, the case of strong $`𝒞𝒫`$mixing is obviously so prominent that it will be seen at other colliders. The opposite case — the ”weak mixing regime” (small values of $`\alpha _2`$) — looks especially interesting. The above equations show that in this region $`T_{},T_\psi ,V_\psi \alpha _2`$, $`V_\lambda c\lambda _V+\alpha _2`$, all other quantities differ from their values without $`𝒞𝒫`$mixing only a little, by a quantity $`\alpha _2^2`$. Therefore, the asymmetries $`T_{},T_\psi `$ for $`\gamma \gamma `$ collisions and $`V_\psi ,V_\lambda `$ for $`\gamma e`$collisions are most sensitive to the weak $`𝒞𝒫`$mixing, as it is seen in Figures.
The quantities $`T_{}`$ and $`T_\psi `$ are non-zero only due to $`𝒞𝒫`$violation. Their $`\mathrm{tan}\beta `$-dependence for different $`\alpha _2`$ is shown in Fig. 9. The measurements of both of these quantities supplement each other essentially: asymmetry $`T_\psi `$ is most sensitive to mixing effects at large $`\mathrm{tan}\beta `$, while in the small $`\beta `$ domain the best suited quantity is $`T_{}`$. This $`\mathrm{tan}\beta `$ dependence of the both quantities again can be traced form Eq.(16). Asymmetry $`T_{}`$, being proportional to $`Re(\stackrel{~}{G}_\gamma G_\gamma ^{})`$, borrows its $`\mathrm{tan}\beta `$ behavior from interplay of the $`b`$ and $`t`$ quark contributions to $`Re(\stackrel{~}{G}_\gamma )`$: the $`b`$ contribution, initially small, grows with $`\mathrm{tan}\beta `$. It compensates the $`t`$ loop at $`\mathrm{tan}\beta 10`$ and becomes dominant later on. At the same time, $`T_\psi `$ has $`\mathrm{tan}\beta `$–dependence similar to $`Im(\stackrel{~}{G}_\gamma )`$, where we have only $`b`$ quark loop contribution. Thus, the whole asymmetry $`T_\psi `$ scales as $`\mathrm{tan}\beta `$.
For the $`\gamma e`$ collision we present only the quantities arising from $`𝒞𝒫`$ non-conservation, they are $`\alpha _2`$ at small $`\alpha _2`$ (Fig. 9). Just as for $`\gamma \gamma `$ reaction the studies of both these quantities supplement each other. The effect of circular polarization $`V_\lambda ^L`$ (which is an analogue to $`T_{}`$) is relatively large at $`\mathrm{tan}\beta 1`$, the $`t`$/$`b`$ quark loop compensation point diminish this effect with growth of $`\mathrm{tan}\beta `$ (it becomes zero at large $`\mathrm{tan}\beta `$). Thus, in the whole $`\mathrm{tan}\beta `$ domain under investigation the $`t`$ quark loop in $`\stackrel{~}{G}_i`$ is dominant and therefore makes $`V_\lambda ^L`$ behave roughly as $`\mathrm{cot}\beta `$. On the contrary, the effect of linear photon polarization $`V_\psi ^L`$ (which is similar to $`T_\psi `$) is very small at $`\mathrm{tan}\beta 1`$ but it grows with $`\mathrm{tan}\beta `$. Nevertheless, it stays below $`0.05`$ and seems thus hardly measurable.
The obtained results describe also production of the lightest Higgs boson in the $`𝒮𝒮`$ in the decoupling regime (when all superparticles are heavy enough). It is necessary to note in this respect that the modern calculations in the $`𝒮𝒮`$ need to fix many subsidiary parameters. In the standard choice, the variation of Higgs mass and $`\mathrm{tan}\beta `$ shifts also quantity $`\delta `$, so that curves of Ref. , for example, present simultaneous dependence on parameters $`\alpha _2`$, $`\beta `$ and $`\delta `$. That is why numerical results of obtained for the specific problems discussed there differ from our Figs. 8,9. The numerical experiments show that simple variation of $`𝒮𝒮`$ parameters $`A`$ and $`\mu `$ allows one to have $`𝒮`$ like value $`\mathrm{sin}\delta 1`$ at $`M_h=105125`$ GeV . Our curves correspond to this very case of $`𝒮𝒮`$.
## 6 Discussion
In this work, together with , we gave detailed answers to the questions what is the whole experimentally available information about photon-Higgs boson anomalous interactions and how to extract it in a reasonable way from future experiments at Photon Colliders. In this problem, the comparative simultaneous analysis of both reactions $`\gamma \gamma H`$ and $`e\gamma eH`$is useful. Due to the absence of $`𝒮`$couplings of the Higgs boson with photons at tree level, the signal of non-standard phenomena can appear clean in Higgs boson production in photon collisions. The high sensitivity of reactions $`\gamma \gamma H`$ and $`e\gamma eH`$ to the admixture of various anomalous interactions makes these processes very useful in exploring the New Physics beyond TeV scale. With new degrees of freedom (2) in the parametric space, the unique opportunities of Photon Colliders in the variation of the initial photon polarization provide a new route to studying different anomalies in details and confident separation of different contributions.
In our investigation we treat anomalies in a universal manner, regardless of the particular mechanism of the $`𝒞𝒫`$violation phenomenon. This is possible because, as we showed, various sources of $`𝒞𝒫`$violation are indistinguishable in the two reactions discussed having relatively large cross sections. These mechanisms are, in principe, distinguishable via the study of such processes as $`\gamma \gamma HH`$ or $`\gamma \gamma H^{}ZZ`$ at $`sM_H^2`$. However they have very low cross sections and will hardly help.
Aiming at the most wide class of anomalous interactions, we parameterized the amplitudes in a very general way, treating the absolute values $`|g_i|`$ and phases $`\xi _i`$ of anomalies as independent parameters. The results presented shows the range of effects that could be resolved from the data, it is close to that for the $`𝒞𝒫`$-even case . They are $`g_\gamma ,g_{P\gamma }0.5÷110^4`$ for $`H\gamma \gamma `$ anomalies and $`g_Z,g_{PZ}510^4`$ for $`HZ\gamma `$ anomalies (in terms of $`\mathrm{\Lambda }_i`$ introduced in they read $`\mathrm{\Lambda }_\gamma ,\mathrm{\Lambda }_{P\gamma }40÷60`$ TeV and $`\mathrm{\Lambda }_Z,\mathrm{\Lambda }_{PZ}20`$ TeV). Effects depend strongly on the phase of anomaly. The comparative study of effects with circularly and linearly polarized photons is necessary to separate effects of amplitude and phase of anomaly ($`|g_i|`$ and $`\xi _i`$). Future simulations based on final versions of collider and detector will show the exact discovery limits before actual experiments.
Next, we analyzed some specific cases of anomalies: the presence of new particles within $`𝒮`$(for $`𝒞𝒫`$even anomalies, )<sup>5</sup><sup>5</sup>5 Note that ”existence of extra chiral generations with all fermions heavier than $`M_Z`$ is strongly disfavored by the precision electroweak data. However the data are fitted nicely even by a few extra generations, if one allows neutral leptons to have masses close to 50 GeV” and scalar-pseudoscalar mixing in the $`2𝒟`$. Their important feature is definite relation among the anomalous signals in $`\gamma \gamma `$and $`\gamma e`$collisions. In particular, the study of both $`\gamma \gamma `$ and $`\gamma e`$ reactions is essential to test if we deal with either $`𝒞𝒫`$violating mixing in 2HDM with definite relation among $`H\gamma \gamma `$ and $`HZ\gamma `$ anomalies or with some other mechanism of $`𝒞𝒫`$violation with now unpredicted relation between these two anomalies. The specific feature of result is that signals of small mixing ($`\mathrm{sin}\alpha _20.1`$) are seen well in effects with circular photon polarization at small and large $`\mathrm{tan}\beta `$ (but not at intermediate, $`\mathrm{tan}\beta 10`$), whereas effects with linear photon polarization can be seen well at intermediate and large values of $`\mathrm{tan}\beta `$.
Last, it is useful to note one more advantage of analysis of polarization asymmetry in the production of Higgs bosons. There is a possibility in the $`2𝒟`$ and $`𝒮𝒮`$ that the heavier scalar Higgs boson $`H`$ and its pseudoscalar counterpart $`A`$ are almost degenerate within the mass resolution without $`𝒞𝒫`$violation. In this case the study of polarization asymmetries in Higgs boson production like those discussed above can answer whether $`𝒞𝒫`$is violated or not. Contrary to this, the study of asymmetries of decay products cannot distinguish the true $`𝒞𝒫`$violation from accidental overlapping of $`H`$ and $`A`$ resonance curves.
We are thankful to V. Ilyin, M. Krawczyk, V. Serbo and P. Zerwas for discussions. IPI is thankful to Prof. J. Speth for hospitality at Forschungszentrum Jülich and IFG is thankful to Prof. A. Wagner for hospitality in DESY, where the paper was finished. This work was supported by grants RFBR 99-02-17211 and 00-15-96691, grant “Universities of Russia” 015.0201.16 and grant of Sankt-Petersburg Center of fundamental studies.
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# Isotropization of two-component fluids
## I Introduction
Since the Bianchi models were introduced into cosmology , they have been the most studied generalizations of the spatially homogeneous and isotropic Friedmann-Lemaître (FL) models. The Bianchi models are spatially homogeneous but anisotropic, and thus well suited for studying the effects of anisotropic expansion on the evolution of the universe. The complexity of a Bianchi model is determined by its three-dimensional symmetry group $`G_3`$, in conjunction with the chosen matter model. Since the present universe seems to be very well described by a FL model, our main interest is models that can be “close to a FL model” at late times. In this paper we will restrict our consideration to models that may approach the open FL model at late times. Within the class of Bianchi models, only models of Bianchi type V and VII<sub>h</sub> can possibly have this behavior (see Collins & Hawking ). This follows since the open FL metric admits a $`G_3`$ of these particular Bianchi types. Here we will focus on the type V models. The notion of a model approaching a FL model at late times is often referred to as late-time isotropization.
As regards the matter description, most studies of Bianchi cosmologies use non-tilted barotropic perfect fluids with a linear equation of state. The book by Wainwright & Ellis is basically devoted to such models. In the so-called non-tilted models, the fluid 4-velocity is orthogonal to the orbits of the isometry group $`G_3`$. Models in which the 4-velocity of the fluid is not orthogonal to the group orbits are referred to as tilted, and were introduced by King & Ellis . These models have been studied by, for example, Collins & Ellis , Hewitt & Wainwright , and Harnett . Since the non-tilted models are a subset of the tilted models, the latter can be viewed as simple generalizations of the former. Another generalization of non-tilted models is achieved by allowing the source of the gravitational field to be a combination of two non-interacting, non-tilted perfect fluids. Since these models can take into account both a radiation-dominated as well as a matter-dominated epoch of the universe, they may be considered as more physically relevant than single-fluid models. The qualitative behavior of two-fluid models were studied by Coley & Wainwright who showed that the models are dominated at early times by the stiffer fluid, and at late times by the softer fluid. This is the expected behavior of a universe filled with radiation and dust.
The next step is to allow one of the fluids in a two-fluid model to be tilted, and subsequently, to let both fluids be tilted. Since the dynamics of models containing tilted fluids is more difficult to analyze than that of non-tilted models, we will consider the combination of one non-tilted and one tilted fluid. The non-tilted fluid is by definition irrotational, but tilted fluids can rotate. We chose, however, to consider the subset of irrotational tilted fluids (For a discussion of general single-fluid Bianchi type V models, see Harnett ). We are thus focusing on the subclass of irrotational Bianchi type V models. Both fluids are assumed to satisfy linear barotropic equations of state,
$`p_\mathrm{o}`$ $`=`$ $`(\gamma _\mathrm{o}1)\mu _\mathrm{o},`$ (2)
$`p_\mathrm{t}`$ $`=`$ $`(\gamma _\mathrm{t}1)\mu _\mathrm{t},`$ (3)
where $`p_{\mathrm{o},\mathrm{t}}`$ are the pressures of the fluids, $`\mu _{\mathrm{o},\mathrm{t}}`$ the energy densities, and $`\gamma _{\mathrm{o},\mathrm{t}}`$ are constants that satisfy $`0\gamma _{\mathrm{o},\mathrm{t}}2`$ with $`\gamma _\mathrm{o}\gamma _\mathrm{t}`$. We exclude, however, the specific values $`2/3`$ and $`2`$ since these models behave qualitatively different. The indices $`\mathrm{o}`$ and $`\mathrm{t}`$ refer to the orthogonal and the tilted fluid respectively and will be used throughout the paper. We will also assume that the energy densities are non-negative, i.e.,
$`\mu _\mathrm{o}0,`$ (5)
$`\mu _\mathrm{t}0.`$ (6)
To discuss isotropization in detail, we need a well-defined notion of when a model is “close to a FL model”. For models with a single non-tilted fluid, a vanishing fluid shear defines the FL models (see, for example, section 2.4 in ). It is, however, not sufficient to demand that the fluid shear itself should approach zero at late times, since this will occur in any single-fluid Bianchi model, irrespectively of whether the model isotropizes or not. As realized already by Kristian & Sachs , the appropriate quantity to consider is the dimensionless ratio formed by normalizing the fluid shear with the fluid expansion, the so-called dimensionless shear. It measures the dynamical importance of the shear compared to the expansion of the fluid. Since there are now two fluids present, the notion of isotropization at late times needs to be generalized as follows. We say that a two-fluid model isotropizes at late times if the dimensionless shear of both fluids vanish in this limit. This issue was partially addressed by Goliath & Ellis , who considered models with a tilted fluid and a non-zero cosmological constant. Such models are contained within the models studied in this paper if we set $`\gamma _\mathrm{o}`$ equal to zero. We will comment on aspects of the behavior of these models that were not addressed in . In general, vanishing dimensionless shear is not sufficient for a model to isotropize, since models can be Weyl-dominated in the future, as is the case for Bianchi type VII<sub>0</sub> and Bianchi type VIII models . However, for Bianchi type V this is not the case .
The plan of the paper is as follows. In Sec. II, the field equations are rewritten as a first-order system of autonomous ordinary differential equations. Reduced dimensionless variables are then introduced leading to a compact state space. In Sec. III we perform a local analysis of the system of equations. In Sec. IV, Bianchi type I and type V models with two orthogonal fluids are studied, while the locally rotationally symmetric type V models are considered in Sec. V. The question of isotropization is discussed in detail. We end with a discussion in Sec. VI. In App. A, the kinematical properties of both fluids are given. The relationship to the variables used in is discussed in App. B. We use units such that $`c=8\pi G=1`$. Orthonormal frame indices are denoted by latin letters $`a,b,\mathrm{}`$ .
## II The gravitational field equations
The line element for the irrotational Bianchi type V models can be written
$$ds^2=dt^2+D_1(t)^2dx^2+\mathrm{e}^{2x}\left[D_2(t)^2dy^2+D_3(t)^2dz^2\right],$$
(7)
(see, for example, Ellis & MacCallum ). In the models we consider there are two preferred timelike congruences. First we have the congruence associated with the normal to the spatial symmetry surfaces, which is also, by definition, the congruence of the 4-velocity of the non-tilted fluid, $`u_\mathrm{o}^a`$. The time variable $`t`$ in (7) is chosen so that the normal to the symmetry surfaces is $`\frac{}{t}`$. The other preferred congruence is that of the tilted fluid. The specific form of the line element, Eq. (7), guarantees that this fluid is irrotational. Hence, the 4-velocity of this fluid, $`u_\mathrm{t}^a`$, is constrained by the field equations to be of a particular form, and can conveniently be parameterized in terms of the so-called tilt variable $`v`$ according to
$$u_\mathrm{t}^a=\mathrm{\Gamma }(1,v,0,0),\mathrm{\Gamma }=(1v^2)^{1/2},$$
(8)
where the value $`v=0`$ corresponds to a non-tilted fluid. The energy-momentum tensor of each fluid is
$`T_\mathrm{o}^{ab}`$ $`=`$ $`\left[\gamma _\mathrm{o}u_\mathrm{o}^au_\mathrm{o}^b+(\gamma _\mathrm{o}1)\eta ^{ab}\right]\mu _\mathrm{o},`$ (10)
$`T_\mathrm{t}^{ab}`$ $`=`$ $`\left[\gamma _\mathrm{t}u_\mathrm{t}^au_\mathrm{t}^b+(\gamma _\mathrm{t}1)\eta ^{ab}\right]\mu _\mathrm{t}.`$ (11)
The assumption that the two fluids are non-interacting leads to equations of motion of the form
$$_aT_\mathrm{o}^{ab}=0=_aT_\mathrm{t}^{ab}.$$
(12)
We choose to parameterize the gravitational field using the expansion, $`\theta `$, and the non-vanishing components of the shear-tensor $`\sigma _\pm `$, of the normal congruence (and thus of the non-tilted fluid). They are given by
$$\theta =\frac{d}{dt}\mathrm{ln}\left(D_1D_2D_3\right),\sigma _+=\frac{1}{2}\frac{d}{dt}\mathrm{ln}\left(\frac{D_1^2}{D_2D_3}\right),\sigma _{}=\frac{\sqrt{3}}{2}\frac{d}{dt}\mathrm{ln}\left(\frac{D_2}{D_3}\right).$$
(13)
The expansion and shear of the tilted fluid can be written as functions of these variables, see App. A. The gravitational field equations,
$$G_{ab}=T_\mathrm{o}^{ab}+T_\mathrm{t}^{ab},$$
(14)
and the equations of motion of the fluids, Eq. (12), become
Evolution equations
$`\dot{\theta }`$ $`=`$ $`\frac{1}{3}\theta ^2\frac{2}{3}\left(\sigma _{+}^{}{}_{}{}^{2}+\sigma _{}^{}{}_{}{}^{2}\right)\frac{1}{2}(3\gamma _\mathrm{o}2)\mu _\mathrm{o}{\displaystyle \frac{1}{2}}{\displaystyle \frac{(3\gamma _\mathrm{t}2)+(2\gamma _\mathrm{t})v^2}{1+(\gamma _\mathrm{t}1)v^2}}\mu _\mathrm{n},`$ (16)
$`\dot{\sigma }_+`$ $`=`$ $`\left(\theta 2vB_1\right)\sigma _+,`$ (17)
$`\dot{\sigma }_{}`$ $`=`$ $`\theta \sigma _{},`$ (18)
$`\dot{B}_1`$ $`=`$ $`\frac{1}{3}\left(\theta 2\sigma _+\right)B_1,`$ (19)
$`\dot{v}`$ $`=`$ $`{\displaystyle \frac{v(1v^2)}{3\left[1(\gamma _\mathrm{t}1)v^2\right]}}\left[2\sigma _++(3\gamma _\mathrm{t}4)\theta 6(\gamma _\mathrm{t}1)B_1v\right],`$ (20)
$`\dot{\mu }_\mathrm{o}`$ $`=`$ $`\gamma _\mathrm{o}\theta \mu _\mathrm{o}.`$ (21)
Constraint equation
$`0`$ $`=`$ $`\gamma _\mathrm{t}v\mu _\mathrm{n}+2(1+(\gamma _\mathrm{t}1)v^2)B_1\sigma _+.`$ (22)
Defining equation for $`\mu _\mathrm{n}`$
$`3\mu _\mathrm{n}`$ $`=`$ $`\theta ^2\sigma _{+}^{}{}_{}{}^{2}\sigma _{}^{}{}_{}{}^{2}9B_1^23\mu _\mathrm{o},`$ (23)
where we have introduced
$`B_1`$ $`=`$ $`D_1^1,`$ (25)
$`\mu _\mathrm{n}`$ $`=`$ $`{\displaystyle \frac{1+(\gamma _\mathrm{t}1)v^2}{1v^2}}\mu _\mathrm{t},`$ (26)
and a dot denotes differentiation with respect to $`t`$.
From the form of Eq. (23), it is now clear why our choice of variables is a good one. The assumption in Eq. (I), in conjunction with Eq. (23) and Eq. (26), implies that $`\theta `$ cannot change sign. Therefore we can, without loss of generality, assume that $`\theta 0`$, i.e., we restrict ourselves to expanding models. From Eq. (23) it also follows that $`\theta `$ is a dominant quantity. We then introduce bounded dimensionless “$`\theta `$-normalized” variables for which the system of equations (14), is reduced as far as possible, and for which the state space is compact. These variables are defined by
$$\mathrm{\Sigma }_\pm =\frac{\sigma _\pm }{\theta },A=\frac{3B_1}{\theta },\mathrm{\Omega }_\mathrm{o}=\frac{3\mu _\mathrm{o}}{\theta ^2},\mathrm{\Omega }_\mathrm{n}=\frac{3\mu _\mathrm{n}}{\theta ^2}.$$
(27)
The subsequent introduction of a dimensionless time variable $`\tau `$, which satisfies
$$\frac{d\tau }{dt}=\frac{\theta }{3},$$
(28)
leads to a decoupling of the $`\theta `$-equation, which can be written on the form
$$\theta ^{}=(1+q)\theta ,q=22A^2\frac{3(2\gamma _\mathrm{o})}{2}\mathrm{\Omega }_\mathrm{o}\frac{3(2\gamma _\mathrm{t})+(5\gamma _\mathrm{t}6)v^2}{2\left[1+(\gamma _\mathrm{t}1)v^2\right]}\mathrm{\Omega }_\mathrm{n},$$
(29)
where the prime denotes differentiation with respect to $`\tau `$. The parameter $`q`$ is the deceleration parameter associated with the normal congruence and the non-tilted fluid. The remaining equations can now be written on dimensionless form as follows.
Evolution equations
$`\mathrm{\Sigma }_+^{}`$ $`=`$ $`(2q2Av)\mathrm{\Sigma }_+,`$ (31)
$`\mathrm{\Sigma }_{}^{}`$ $`=`$ $`(2q)\mathrm{\Sigma }_{},`$ (32)
$`A^{}`$ $`=`$ $`(q+2\mathrm{\Sigma }_+)A,`$ (33)
$`v^{}`$ $`=`$ $`{\displaystyle \frac{v(1v^2)}{1(\gamma _\mathrm{t}1)v^2}}\left[2\mathrm{\Sigma }_++3\gamma _\mathrm{t}42(\gamma _\mathrm{t}1)Av\right],`$ (34)
$`\mathrm{\Omega }_{\mathrm{o}}^{}{}_{}{}^{}`$ $`=`$ $`\left[2q(3\gamma _\mathrm{o}2)\right]\mathrm{\Omega }_\mathrm{o}.`$ (35)
Constraint equation
$`0`$ $`=`$ $`\gamma _\mathrm{t}v\mathrm{\Omega }_\mathrm{n}+2[1+(\gamma _\mathrm{t}1)v^2]A\mathrm{\Sigma }_+.`$ (36)
Defining equation for $`\mathrm{\Omega }_\mathrm{n}`$
$`\mathrm{\Omega }_\mathrm{n}`$ $`=`$ $`1\mathrm{\Sigma }_+^2\mathrm{\Sigma }_{}^2A^2\mathrm{\Omega }_\mathrm{o}.`$ (37)
The set of equations (29) shows that the irrotational Bianchi type V models with one orthogonal and one tilted fluid is governed by a system of five autonomous ordinary differential equations subject to one constraint. The dimension of the state space is thus four. The set (29) is invariant under the discrete transformations
$$(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},A,v,\mathrm{\Omega }_\mathrm{o})(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},A,v,\mathrm{\Omega }_\mathrm{o}),(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},A,v,\mathrm{\Omega }_\mathrm{o})(\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},A,v,\mathrm{\Omega }_\mathrm{o}),$$
(38)
so we can, without loss of generality, restrict ourselves to the invariant set defined by $`A0`$ and $`\mathrm{\Sigma }_{}0`$. The variables $`\mathrm{\Sigma }_+,\mathrm{\Sigma }_{},A,v,`$ and $`\mathrm{\Omega }_\mathrm{o}`$ therefore satisfy $`0\mathrm{\Sigma }_+^2,\mathrm{\Sigma }_{},A,v^2,\mathrm{\Omega }_\mathrm{o}1`$. There is a number of important subsets of Eqs. (29):
1. The three-dimensional subset defined by $`v=0`$, which describes a Bianchi type V universe with two non-tilted fluids. This case is thus included in the general study of Coley & Wainwright . We will, however, consider these models in Sec. IV for the purpose of comparison with the models where one fluid is tilted. The subset $`v=0`$ also contains Bianchi type I models and the open FL model as special cases.
2. The three-dimensional subset $`\mathrm{\Sigma }_{}=0`$, which corresponds to locally rotationally symmetric (LRS) models (see, for example, ). This subset turns out to be very important for the evolution of the general irrotational type V models and will be considered in detail in Sec. V.
The next step in the analysis is to consider the equilibrium points of Eq. (29). This is done in the next section.
## III Qualitative analysis
In Table I, we present the equilibrium points of the system (29). The corresponding eigenvalues are given in Table II. Since the constraint, Eq. (36), cannot be solved globally in an analytic way, it will be treated locally (see, for example, ).
Some of the equilibrium points correspond to exact solutions of the field equations. For example, the equilibrium points $`F_\mathrm{t}^0`$ and $`F_\mathrm{o}^0`$ correspond to flat FL models in which the tilted fluid and the non-tilted fluid is dominant, respectively (The “tilted” fluid is in fact non-tilted at $`F_\mathrm{t}^0`$ since $`v=0`$, but we refer to it as the tilted fluid for simplicity). The point $`M^0`$ is the Milne model, while the equilibrium set $`K^0`$ corresponds to Kasner-like models. The point $`\stackrel{~}{M}`$ corresponds to flat space and coincides with $`M^0`$ for $`\gamma _\mathrm{t}=4/3`$. We also note the appearance of the equilibrium set $`F_\mathrm{o}^v`$ for the specific value $`\gamma _\mathrm{t}=\frac{4}{3}`$. It is associated with a line bifurcation that transfers stability between the equilibrium points $`F_\mathrm{o}^\pm `$ and $`F_\mathrm{o}^0`$. The points $`𝒞^\pm `$ correspond to particular Kasner solutions. There is also a number of equilibrium points for which the tilt is extreme ($`v^2=1`$). Whether these equilibrium points correspond to exact Bianchi solutions or not seems to be an open question (see comment on p. 4245 in ). We note that the constraint, Eq. (36), is degenerate ($`G=0`$) at the point $`F_\mathrm{o}^0`$, allowing all five eigenvector directions to be physical at this point.
## IV Two orthogonal fluids
Setting $`v=0`$ in the constraint Eq. (36), implies either $`A=0,\mathrm{\Sigma }_+0`$ (Bianchi type I models) or $`\mathrm{\Sigma }_+=0,A0`$ (Bianchi type V models). Without loss of generality we can assume that $`\gamma _\mathrm{t}>\gamma _\mathrm{o}`$. The boundary subsets of the state space for these two classes of models are given by the two invariant submanifolds $`\mathrm{\Omega }_\mathrm{o}=0`$ and $`\mathrm{\Omega }_\mathrm{n}=0`$, which describe the corresponding one-fluid models. The dynamics of the Bianchi I state space is shown in Fig. 2, while the dynamics of the type V models is shown in Fig. 2.
For type I models there are two sources. The equilibrium point $`F_\mathrm{t}^0`$ gives rise to a single orbit ending at $`F_\mathrm{o}^0`$. It corresponds to a flat Friedmann model with two orthogonal fluids. The other source is the equilibrium set $`K^0`$, of which each point is associated with a one-parameter set of orbits. Therefore this equilibrium set describes the generic behavior at early times. The future attractor of these orbits is the point $`F_\mathrm{o}^0`$. Thus, from Eqs. (A3) and (A9), all orthogonal two-fluid Bianchi type I models isotropize.
For type V models there are three sources. The equilibrium point $`F_\mathrm{t}^0`$ is associated with a one-parameter set of orbits (characterized by $`\mathrm{\Sigma }_{}=0`$), which corresponds to open Friedmann models with two orthogonal fluids. They all end at the equilibrium point $`M^0`$. The other two sources are the two equilibrium points belonging to the set $`K^0`$ with $`\mathrm{\Sigma }_{}=\pm 1`$. Both of these are associated with two-parameter sets of orbits, and thus describe the generic early-time behavior. They are future attracted to $`M^0`$. Consequently, from Eqs. (A3) and (A9), all orthogonal two-fluid Bianchi type V models isotropize.
From the analysis of the orthogonal two-fluid models, it is clear that the generic behavior is described by equilibrium points associated with two-parameter sets. Thus, in what follows, we will focus on such equilibrium points. Note that the variable $`\chi `$ of Coley & Wainwright (see App. B) is a monotone function when both the fluids are orthogonal. This is no longer the case when one of the fluids is tilted.
## V Bianchi type V LRS models with one tilted and one orthogonal fluid
As Bianchi type I models do not allow the combination of one non-tilted and one tilted fluid as a source due to the constraint Eq. (36), we focus on the Bianchi type V models. Since $`\mathrm{\Omega }_\mathrm{o},\mathrm{\Omega }_\mathrm{n}>0`$ imply $`q<2`$ by Eq. (29), the evolution equation for $`\mathrm{\Sigma }_{}`$ implies that $`\mathrm{\Sigma }_{}`$ is a monotone decreasing function along all orbits with $`\mathrm{\Sigma }_{}0`$ and $`\mathrm{\Omega }_\mathrm{o},\mathrm{\Omega }_\mathrm{n}>0`$. This fact significantly restricts the evolution at late times. It implies that $`lim_\tau \mathrm{}\mathrm{\Sigma }_{}=0`$, for all orbits with $`\mathrm{\Omega }_\mathrm{o},\mathrm{\Omega }_\mathrm{n}>0`$. This can be proven along the lines used to prove the similar statement for tilted single-fluid models, see . The asymptotic behavior as $`\tau \mathrm{}`$ is thus contained in the three-dimensional invariant set $`\mathrm{\Sigma }_{}=0`$ corresponding to the LRS models.
The three Kasner circles $`K^0`$ and $`K^\pm `$ each reduce to two equilibrium points in the LRS submanifold, namely points for which $`\mathrm{\Sigma }_+=\pm 1`$. We denote these equilibrium points $`K_\pm ^0`$, $`K_\pm ^+`$, and $`K_\pm ^{}`$, where the subscript distinguishes between the two signs of $`\mathrm{\Sigma }_+`$. The equilibrium points for which $`v`$ and $`\mathrm{\Sigma }_+`$ have the same sign (collectively denoted $`K_\pm ^\pm `$) are not located on the boundary of the interior of the LRS submanifold. Consequently, we do not need to consider them when studying the dynamics.
The state space for the LRS submanifold with $`2/3<\gamma _\mathrm{o}<\gamma _\mathrm{t}<2`$ is presented in Figs. 66. The effect of changing the equation-of-state parameters so that $`\gamma _\mathrm{o}>\gamma _\mathrm{t}`$ is that the flow along the orbit between $`F_\mathrm{o}^0`$ and $`F_\mathrm{t}^0`$ is reversed. Similarly, $`0\gamma _\mathrm{o}<2/3`$ results in a stability change along the orbits from the equilibrium points $`F_\mathrm{o}^{}`$ to $`M^{}`$, where $`\{0,+,\}`$. Note that the special case $`\gamma _\mathrm{o}=0`$ corresponds to a cosmological constant. Consequently, the LRS state space for $`0\gamma _\mathrm{o}<2/3`$ is contained in Figs. 9–11 of . Finally, if $`0\gamma _\mathrm{t}<2/3`$ the flow changes along the orbits $`F_t^0`$$`M^0`$ and $`K_+^0`$$`K_+^{}`$.
The sources and sinks for various equations of state are summarized in Table III. The isotropization properties, which can be found from Eqs. (A3) and (A6), are also listed. For $`\gamma _\mathrm{t}<6/5`$, all models isotropize. for $`\gamma _\mathrm{o}>2/3`$, $`6/5<\gamma _\mathrm{t}<4/3`$ there is a class of solutions of non-zero measure that does not isotropize, namely the class of solutions associated with $`M^{}`$. For $`\gamma _\mathrm{o}>2/3`$, $`\gamma _\mathrm{t}=4/3`$, all models isotropize. For $`\gamma _\mathrm{t}>4/3`$ no models isotropize. Note that the physically interesting combination of one dust fluid and one radiation fluid always isotropizes, regardless of which fluid is tilted.
When $`\gamma _\mathrm{o}<2/3,\gamma _\mathrm{t}>4/3`$, the future attractors are $`F_\mathrm{o}^\pm `$. Solutions corresponding to orbits approaching these equilibrium points do not isotropize, see Table III, even though the final states are inflationary in the sense that the deceleration parameter associated with the normal congruence $`q=\frac{1}{2}(3\gamma _\mathrm{o}2)`$ is negative. Thus, in some sense it seems to be misleading to refer to solutions corresponding to orbits approaching these equilibrium points as “asymptotically Friedmann models”, although the solutions corresponding to the points themselves are Friedmann models. In particular, for the case of a cosmological constant $`\gamma _\mathrm{o}=0`$, the solutions corresponding to $`F_\mathrm{o}^\pm `$ are de Sitter models . This is consistent with the cosmic “no-hair” theorem for Bianchi models . However, the theorem does not guarantee that the tilt tends to zero, as pointed out by Raychaudhuri & Modak . From our analysis, it is clear that this cautionary note is crucial for these models. As the tilted fluid does not become orthogonal at late times, the expansion-normalized shear of the tilted fluid does not vanish at late times. This stresses the point made in that one should be cautious about the isotropization of tilted models. As seen in Table III, this is in fact the generic behavior for models with $`0\gamma _\mathrm{o}<2/3`$ and $`\gamma _\mathrm{t}4/3`$.
For general irrotational Bianchi type V models, there is no a priori reason that $`\mathrm{\Sigma }_\mathrm{t}0`$ when $`\mathrm{\Sigma }_{\mathrm{t}+}0`$. However, this is indeed the case for the equilibrium points in question, as can be seen from Eq. (A8), noting that $`\mathrm{\Sigma }_{}0`$. Thus, the analysis of isotropization for the LRS submanifold holds for the general class of models as well.
## VI Discussion
In this paper we have continued the study of irrotational Bianchi type V cosmologies, using the dynamical systems approach initiated by Hewitt & Wainwright and Coley & Wainwright . The source of the gravitational field has been taken to be two non-interacting fluids, one orthogonal and one tilted. Such models can describe a universe where one of the fluids models the contribution of radiation to the energy density of the universe, and the other the matter content.
We have found that, although the orthogonal two-fluid models isotropize, this is not necessarily the case when one of the fluids is tilted. Thus, depending on the equation-of-state parameters $`\gamma _\mathrm{o}`$ and $`\gamma _\mathrm{t}`$, it is possible to find cases for which all or a subset of the solutions are anisotropic to the future. In particular, there are models which are inflationary, but do not isotropize. However, we emphasize that the cases of dust plus radiation always isotropize.
It should be stressed that the relevant quantities when determining whether the shear dies away are shear quantities normalized with respect to the expansion associated with each fluid.
## Acknowledgements
We thank Claes Uggla for useful comments. MG was supported by a grant from C F Liljevalch J:ors stipendiefond. USN was supported by Gålöstiftelsen, Svenska Institutet, Stiftelsen Blanceflor and the University of Waterloo.
## A Fluid properties
Here we present the kinematical properties of the two fluids in terms of the variables used to parameterize the gravitational field. For the non-tilted fluid we have
Fluid expansion
$$\theta _\mathrm{o}=\theta ,$$
(A1)
Fluid shear
$$\sigma _{\mathrm{o}}^{}{}_{}{}^{2}=\sigma _{+}^{}{}_{}{}^{2}+\sigma _{}^{}{}_{}{}^{2}$$
(A2)
The dimensionless shear $`\mathrm{\Sigma }_\mathrm{o}`$ of the non-tilted fluid is thus
$$\mathrm{\Sigma }_\mathrm{o}^2=\frac{\sigma _{\mathrm{o}}^{}{}_{}{}^{2}}{\theta _\mathrm{o}^2}=\mathrm{\Sigma }_+^2+\mathrm{\Sigma }_{}^2.$$
(A3)
The fluid properties of the tilted fluid are as follows.
Fluid expansion
$$\theta _\mathrm{t}=\frac{1}{\sqrt{1v^2}}\left[\frac{3\theta 6vB_1+v^2(2\sigma _+\theta )}{3\left[1(\gamma _\mathrm{t}1)v^2\right]}\right]=\frac{1}{\sqrt{1v^2}}\left[\frac{32vA+v^2(2\mathrm{\Sigma }_+1)}{3\left[1(\gamma _\mathrm{t}1)v^2\right]}\right]\theta ,$$
(A4)
Fluid shear
$$\sigma _\mathrm{t}^2=\frac{1}{1v^2}\left[\sigma _{+}^{}{}_{}{}^{2}+\sigma _{}^{}{}_{}{}^{2}B_1v(2\sigma _+B1v)\frac{2v\dot{v}}{1v^2}(\sigma _+B_1v)+\frac{v^2\dot{v}^2}{(1v^2)^2}\right].$$
(A5)
The dimensionless shear $`\mathrm{\Sigma }_\mathrm{t}`$ of the tilted fluid can be written
$$\mathrm{\Sigma }_\mathrm{t}^2=\frac{\sigma _\mathrm{t}^2}{\theta _\mathrm{t}^2}=\mathrm{\Sigma }_{\mathrm{t}+}^2+\mathrm{\Sigma }_\mathrm{t}^2,$$
(A6)
where
$`\mathrm{\Sigma }_{\mathrm{t}+}`$ $`=`$ $`1+{\displaystyle \frac{3(1+\mathrm{\Sigma }_+vA)\left[1(\gamma _\mathrm{t}1)v^2\right]}{32vA+(2\mathrm{\Sigma }_+1)v^2}},`$ (A7)
$`\mathrm{\Sigma }_\mathrm{t}`$ $`=`$ $`{\displaystyle \frac{3\left[1(\gamma _\mathrm{t}1)v^2\right]\mathrm{\Sigma }_{}}{32vA+(2\mathrm{\Sigma }_+1)v^2}},`$ (A8)
(see Eqs. (A3), (A4) and (A6) in ). Note that when both fluids are orthogonal ($`v=0`$), Eq. (A6) reduces to
$$\mathrm{\Sigma }_\mathrm{t}^2=\mathrm{\Sigma }_+^2+\mathrm{\Sigma }_{}^2.$$
(A9)
## B The variable $`\chi `$
In their study of two orthogonal fluids, Coley & Wainwright introduced the following variable instead of $`\mathrm{\Omega }_0`$:
$$\chi :=\frac{\mu _\mathrm{o}\mu _\mathrm{t}}{\mu _\mathrm{o}+\mu _\mathrm{t}}=\frac{\mathrm{\Omega }_\mathrm{o}\mathrm{\Omega }_\mathrm{t}}{\mathrm{\Omega }_\mathrm{o}+\mathrm{\Omega }_\mathrm{t}}.$$
(B1)
The evolution equation for $`\chi `$ becomes
$$\chi ^{}=\frac{\left(1\chi ^2\right)}{1+(\gamma _\mathrm{t}1)v^2}\left\{3\gamma _\mathrm{o}\left[1+(\gamma _\mathrm{t}1)v^2\right]\gamma _\mathrm{t}\left[3+v^22v(A+v\mathrm{\Sigma }_+)\right]\right\}.$$
(B2)
For the submanifold $`v=0`$ corresponding to two orthogonal fluids, the above equation simplifies to
$$\chi ^{}=\frac{1}{2}(1\chi ^2)(\gamma _\mathrm{o}\gamma _\mathrm{t}).$$
(B3)
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# Line bundles of type (1,…,1,2,…,2,4,…,4) on Abelian Varieties
## 1 Introduction
Let $`L`$ be an ample line bundle of type $`\delta =(\delta _1,\delta _2,\mathrm{},\delta _g)`$ on a $`g`$-dimensional abelian variety $`A`$. Consider the associated rational map $`\varphi _L:AIPH^0(A,L)`$.
When $`g=2`$, Birkenhake, Lange and van Straten ( see ) have studied line bundles of type $`(1,4)`$ on abelian surfaces. Suppose $`L`$ is an ample line bundle of type $`(1,4)`$ on an abelian surface $`A`$. Then there is a cyclic covering $`\pi :AB`$ of degree $`4`$ and a line bundle $`M`$ on $`B`$ such that $`\pi ^{}M=L`$. Let $`X`$ denote the unique divisor in $`|M|`$ and put $`Y=\pi ^1(X)`$. Their main theorem is
###### Theorem 1.1
1) $`\varphi _L:AA^{}IP^3`$ is birational onto a singular octic $`A^{}`$ in $`IP^3`$ if and only if $`X`$ and $`Y`$ do not admit elliptic involutions compatible with the action of the Galois group of $`\pi `$.
2)In the exceptional case $`\varphi _L:AA^{}IP^3`$ is a double covering of a singular quartic $`A^{}`$, which is birational to an elliptic scroll.
Here we generalise this situation to higher dimensions and show
###### Theorem 1.2
Suppose $`L`$ is an ample line bundle of type $`\delta =(1,\mathrm{},1,2,\mathrm{},2,4,\mathrm{},4)`$ on a $`g`$-dimensional abelian variety $`A`$, $`g3`$, such that $`1`$ and $`4`$ occur equally often and atleast once in $`\delta `$. Then, for a generic pair $`(A,L)`$, the following holds.
a) The associated morphism $`\varphi _L:AIPH^0(A,L)`$ is birational onto its image.
b) When $`g=3`$, the image $`\varphi _L(A)`$, can be described as follows,
there are $`4`$ curves $`C_i`$ on the image $`\varphi _L(A)`$ such that the restricted morphism $`\varphi _L:\varphi _L^1(C_i)C_i\varphi _L(A)`$ is of degree $`2`$.
Birkenhake et.al (see , Proposition 1.7, p.631) have shown the existence of the following commutative diagram
$`\begin{array}{ccccc}A& \stackrel{\varphi _L}{}& \varphi _L(A)& & IP^3=IPH^0(L)\\ \pi & & & & p\\ B& \stackrel{\varphi _{M^2}}{}& 𝒦(B)& & IP^3=IPH^0(M^2)\end{array}`$
where $`p(z_0:z_1:z_1:z_3)=(z_0^2:z_1^2:z_2^2:z_3^2)`$ and the pair $`(B,M)`$ is a principally polarized abelian surface. This diagram explains the geometry of the image $`\varphi _L(A)`$ from the geometry of the Kummer surface $`𝒦(B)`$ and it also gives the explicit equation of the surface $`\varphi _L(A)`$ in $`IP^3`$.
Similarly, when $`g3`$ and the pair $`(A,L)`$ as in 1.2, we show that there is a commutative diagram:
$`\begin{array}{ccccc}A& \stackrel{\varphi _L}{}& \varphi _L(A)& & IP^{2^g1}=IPH^0(L\\ \pi & & & & p\\ B& \stackrel{\varphi _{M^2}}{}& 𝒦(B)& & IP^{2^g1}=IPH^0(M^2)\end{array}`$
where $`p(z_0:\mathrm{}:z_{2^g1})=(z_0^2:\mathrm{}:z_{2^g1}^2)`$ and $`\pi `$ is an isogeny of degree $`2^g`$ and the pair $`(B,M)`$ is a principally polarized abelian variety. This will explain the birationality of the map $`\varphi _L`$ and the geometry of the image $`\varphi _L(A)`$, when $`g=3`$, as asserted in 1.2. Since $`deg(\varphi _{M^2}\pi )=2^{g+1}`$ and from the birationality of $`\varphi _L`$, it follows that $`deg(p|_{\varphi _L(A)})=2^{g+1}`$. But since $`degp=2^{2^g1}`$ the inverse image of the Kummer variety in $`IPH^0(L)`$ has components other than the image $`\varphi _L(A)`$. Hence the image $`\varphi _L(A)`$ will be defined by forms other than those coming from those forms which define the variety $`𝒦(B)`$.
We study the situation when $`g=3`$, in detail. Consider a pair $`(A,L)`$, with $`L`$ being an ample line bundle of type $`(1,2,4)`$ on an abelian threefold $`A`$. Consider an isogeny $`AB=A/G`$, where $`G`$ is a maximal isotropic subgroup of $`K(L)`$ of the type $`\frac{ZZ}{2ZZ}\times \frac{ZZ}{2ZZ}\times \frac{ZZ}{2ZZ}`$. Then $`B`$ is a principally polarized abelian threefold. If $`B`$ is isomorphic to the Jacobian variety of $`C`$, $`J(C)`$, where $`C`$ is a smooth non-hyperelliptic curve of genus $`3`$, then the situation becomes interesting because of the following results due to Narasimhan and Ramanan.
###### Theorem 1.3
(See , Main Theorem, p.416) If $`C`$ is a non-hyperelliptic curve of genus $`3`$, then the moduli space $`SU_C(2)`$ is isomorphic to a quartic hypersurface in $`IP^7`$.
( Here $`IP^7=|2\theta |`$, where $`\theta `$ is the canonical principal polarization on the Jacobian $`J(C)`$ and $`SU_C(2)`$ is the moduli space of rank $`2`$ semi-stable vector bundles with trivial determinant on the curve $`C`$).
###### Theorem 1.4
( See ) The Kummer variety $`𝒦`$ is precisely the singular locus of $`SU_C(2)`$, if $`g(C)3`$.
The quartic hypersurface, $`F=0`$, is classically called the $`\mathrm{𝐶𝑜𝑏𝑙𝑒}\mathrm{𝑞𝑢𝑎𝑟𝑡𝑖𝑐}`$ and is $`𝒢(2\theta )`$-invariant in the linear system $`|2\theta |`$. We identify the group of projective transformations, $`H`$, of order $`8`$, which acts on $`\pi ^1𝒦(C)`$, (see 3.7). The $`𝒢(L)`$-invariant octic hypersurface $`R`$, given as $`F(z_0^2:\mathrm{}:z_7^2)=0`$ in $`IPH^0(L)`$, then contains the components $`h(\varphi _L(A)),hH`$ in its singular locus.
Now we use the geometry of the moduli space $`SU_C(2)`$ in the linear system $`|2\theta |`$, which has been extensively studied ( see , for instance), to get analogous results in $`IPH^0(L)`$.
We show
###### Theorem 1.5
Consider a pair $`(A,L)`$, as above. Let $`aK(L)`$ be an element of order $`2`$ such that $`e^L(a,g)=1`$, for all $`gG`$, (here $`e^L`$ is the Weil form on the group $`K(L)`$). Let $`IPW_a`$ be an eigenspace in $`IPH^0(L)`$, for the action of $`a`$. Then there is a polarized abelian surface $`(Z,N)`$, $`N`$ is ample of type $`(1,4)`$ and a commutative diagram
$`\begin{array}{cccccc}Z& \stackrel{\varphi _N}{}& \varphi _N(Z)& & IPH^0(N)& IPW_a\\ f& & & & q& p\\ P_a& \stackrel{\varphi _{2\theta _a}}{}& 𝒦(P_a)& & IPH^0(2\theta _a)& IPV_a\end{array}`$
Here $`(P_a,\theta _a)`$ is the Prym variety associated to the $`2`$-sheeted unramified cover of the curve $`C`$, given by $`\pi (a)`$ and $`IPV_a`$ is the eigenspace in $`IPH^0(2\theta )`$, for the action of $`\pi (a)`$. The isomorphisms above are Heisenberg equivariant and the morphism $`q`$ is given as $`(r_0:r_1:r_2:r_3)(r_0^2:r_1^2:r_2^2:r_3^2)`$.
We thus obtain the situation described by Birkenhake et.al in the case $`g=2`$, nested in the case $`g=3`$.
Moreover, the $`𝒢(N)`$-invariant octic surface $`\varphi _N(Z)`$ is mapped isomorphically onto the $`a^{}/a(Heis(4))`$-octic $`RIPW_a`$ and we identify the set $`_{hH}h(\varphi _L(A))`$ with the set of all pinch points and the coordinate points in $`\varphi _N(Z)`$, occurring in each of the eigenspace $`IPW_a`$, (see 5.6). Finally, we make some remarks on the moduli space $`𝒜^{(1,2,4)}`$.
$`\mathrm{𝐴𝑐𝑘𝑛𝑜𝑤𝑙𝑒𝑑𝑔𝑒𝑚𝑒𝑛𝑡𝑠}`$: We thank W.M.Oxbury and B.van Geemen for making useful comments in an earlier version. We are grateful to Christian Pauly for suggestions during revision. We also thank the French Ministry of National Education, Research and Technology, for their support.
$`\mathrm{𝐍𝐨𝐭𝐚𝐭𝐢𝐨𝐧}:`$ Suppose $`L`$ is a symmetric line bundle i.e. $`Li^{}L`$ for the involution $`i:AA`$, $`aa`$.
The $`\mathrm{𝑓𝑖𝑥𝑒𝑑}\mathrm{𝑔𝑟𝑜𝑢𝑝}`$ of $`L`$ is $`K(L)=\{aA:Lt_a^{}L\},t_a:AA,xa+x`$.
The $`\mathrm{𝑡ℎ𝑒𝑡𝑎}\mathrm{𝑔𝑟𝑜𝑢𝑝}`$ of $`L`$ is $`𝒢(L)=\{(a,\varphi ):L\stackrel{\varphi }{}t_a^{}L\}`$.
$`K_1(\delta )=\frac{ZZ}{d_1ZZ}\times \mathrm{}\times \frac{ZZ}{d_gZZ}`$, and $`\widehat{K_1(\delta )}=Hom(K_1(\delta ),IC^{})`$.
The $`\mathrm{𝐻𝑒𝑖𝑠𝑒𝑛𝑏𝑒𝑟𝑔}\mathrm{𝑔𝑟𝑜𝑢𝑝}`$ of type $`\delta `$, $`Heis(\delta )=IC^{}\times K_1(\delta )\times \widehat{K_1(\delta )}`$ and $`V(\delta )=\{f:f:K_1(\delta )IC\}`$.
The $`\mathrm{𝑊𝑒𝑖𝑙}\mathrm{𝑓𝑜𝑟𝑚}e^L:K(L)\times K(L)IC^{}`$, is the commutator map $`(x,y)x^{}y^{}x^1y^1`$, for any lifts $`x^{},y^{}𝒢(L)`$ of $`x,yK(L)`$.
For any $`aK(L)`$, $`a^{}=\{xK(L):e^L(a,x)=1\}`$.
Consider the semi-direct product, $`𝒢(L)(\times (i)`$, of the theta group associated to $`L`$ and the group generated by the involution $`i`$. Let $`\gamma 𝒢(L)(\times (i)`$ be an element of order 2.
$`H^0(L)_\gamma ^\pm =(\pm 1)`$eigenspace of $`H^0(L)`$ for the action of $`\gamma `$.
$`h^0(L)_\gamma ^\pm =dimH^0(L)_\gamma ^\pm .`$
$`Q(V)=`$ function field of a variety $`V`$.
## 2 Birationality of the map $`\varphi _L`$.
Let $`L`$ be an ample line bundle of type $`\delta =(1,\mathrm{..2},..,4)`$ on a $`g`$-dimensional abelian variety $`A`$. Here number of 2’s= number of 4’s in $`\delta `$. Let $`K(L)=\{aA:t_a^{}LL\}`$, where $`t_a`$ denotes translation by $`a`$ on $`A`$. Choose a maximal isotropic subgroup $`G`$ of $`K(L)`$ w.r.t. the Weil form $`e^L`$, containing $`2K(L)`$ and having only elements of order 2. Then $`G\frac{ZZ}{2ZZ}\times \mathrm{}\times \frac{ZZ}{2ZZ}`$, $`g`$-times. Consider the exact sequence
$$1IC^{}𝒢(L)K(L)0.$$
Let $`G^{}`$ be a lift of $`G`$ in $`𝒢(L)`$. Consider the isogeny $`A\stackrel{\pi }{}B=A/G`$. Then $`L`$ descends to a principal polarization $`M`$ on $`B`$. By Projection formula and using the fact that $`\pi _{}𝒪_A=_{\chi \widehat{G}}L_\chi `$, where $`L_\chi `$ denotes the line bundle corresponding to the character $`\chi `$, we deduce that
$$H^0(L)=_{\chi \widehat{G}}H^0(ML_\chi ).$$
Hence $`\{s_\chi H^0(ML_\chi ):\chi \widehat{G}\}`$ is a basis for the vector space $`H^0(L)`$ and since $`M^2L_\chi ^2M^2`$, $`s_\chi ^2=s_\chi s_\chi H^0(M^2)\chi \widehat{G}`$.
Consider the homomorphism $`ϵ_2:𝒢(L)𝒢(L^2),(x,\varphi )(x,\varphi ^2)`$ and the inclusion $`K(L)K(L^2)`$.
Then the subgroup $`GK(L^2)`$ is isotropic for the Weil form $`e^{L^2}`$. Moreover, if $`xK(L)`$ and $`gG`$, then
$$e^{L^2}(x,g)=e^L(x,g).e^L(x,g)=1.$$
Hence $`ϵ_2(𝒢(L))𝒵(ϵ_2(G^{^{}}))`$ and $`\pi (K(L))K(M^2)`$. ( Here $`𝒵(ϵ_2(G^\mathrm{`}))=\{a𝒢(L^2):a.g^{^{}}=g^{^{}}.a,g^{^{}}ϵ_2(G^{^{}})\}`$).
Now $`𝒢(M^2)=𝒵(ϵ_2(G^{^{}}))/ϵ_2(G^{^{}})`$ and $`H^0(M^2)=H^0(L^2)^G^{^{}}`$, where $`H^0(L^2)^G^{^{}}`$ denotes the vector subspace of $`ϵ_2(G^{^{}})`$-fixed sections of $`H^0(L^2)`$. For $`g^{}G^{}`$ and $`\chi \widehat{G}`$, $`g^{}(s_\chi ^2)=\chi ^2(g).s_\chi ^2=s_\chi ^2`$. Hence $`s_\chi ^2H^0(L^2)^G^{}`$, for all $`\chi \widehat{G}`$.
We now show that $`\{s_\chi ^2:\chi \widehat{G}\}`$ is a basis for $`H^0(M^2)`$, for a generic pair $`(A,L)`$..
In fact, we show that the homomorphism
$$\underset{\chi \widehat{G}}{}H^0(ML_\chi ).H^0(ML_\chi )\stackrel{\rho }{}H^0(M^2)\mathrm{}()$$
is an isomorphism, for a generic pair $`(A,L)`$.
Consider the pair $`(A,L)=(E_1\times \mathrm{}\times E_r\times ,A_1\times \mathrm{}A_s,p_1^{}L_1\mathrm{}p_{r+s}^{}L_{r+s})`$, where $`r`$ is the number of 2’s occurring in $`\delta `$, $`E_1,\mathrm{},E_r`$ are elliptic curves with line bundles $`L_i`$ on $`E_i`$ of degree 2 and $`A_j`$ are simple abelian surfaces with line bundles $`L_j`$ on $`A_j`$ of type $`(1,4)`$ ( by 1.1, $`\varphi _{L_j}(A_j)|L_j|`$ is an octic surface).
In this case, one can easily see that the homomorphism
$$S=Sym^2H^0(L_1)\mathrm{}Sym^2H^0(L_{r+s})H^0(L_1^2)\mathrm{}H^0(L_{r+s}^2)=H^0(L_1^2\mathrm{}L_{r+s}^2)$$
is injective. Here, $`(B,M)=(F_1,M_1)\times \mathrm{}\times (F_r,M_r)\times (B_1,M_1^{})\times \mathrm{}(B_s,M_s^{})`$, where $`(F_j,M_j)`$ are polarised elliptic curves of degree 1 and $`(B_j,M_j)`$ are principally polarised abelian surfaces. Also, the group $`G`$ is generated by elements of the type $`(e_1,..,e_r,a_{r+1}^{},..,a_g^{})`$, where each of $`e_j`$ and $`a_j^{}`$ are non-trivial 2 torsion elements of $`E_j`$ and $`A_j`$, respectively. Now it is easy to see that $`_{\chi \widehat{G}}H^0(ML_\chi ).H^0(ML_\chi )S`$ and $`H^0(M^2)H^0(L^2)`$ and (\*) is an isomorphism.
Hence, for a generic pair $`(A,L)`$ as above, (\*) is an isomorphism.
As a consequence, we obtain the following
###### Proposition 2.1
Consider a generic principally polarized abelian variety $`(B^{},M^{})`$ of dimension $`g`$. Let $`H`$ be a subgroup of 2-torsion points of $`B^{}`$, of order g. Then the image of $`H`$ in $`𝒦(B^{})`$ generates the linear system $`|2M^{}|`$.
( This is well known if $`H`$ consists of all the 2-torsion points of $`B^{}`$, for any principally polarised pair $`(B^{},M^{})`$.)
Proof: Since the map $`B^{}\stackrel{\varphi _{2M^{}}}{}|2M^{}|`$ is given by $`at_a^{}\theta +t_a^{}\theta `$, where $`\theta `$ is the unique divisor in $`|M^{}|`$, the assertion is equivalent to showing the surjectivity of the multiplication map
$$\underset{\chi \widehat{H}}{}H^0(M^{}L_\chi )H^0(M^{}L_\chi )\stackrel{\rho }{}H^0(M^2)..(!).$$
Here $`\widehat{H}`$ is the dual image of $`H`$ in $`Pic^0(B^{})`$. But we showed above this isomorphism, if $`\widehat{H}`$ gives rise to a $`g`$-sheeted cover $`(A^{},L^{})`$ of $`(B^{},M^{})`$, where $`L^{}`$ is of type $`(1,..,2,..,4)`$. Otherwise, $`\widehat{H}`$ gives a cover $`(A^{},L^{})`$ where $`L^{}`$ is of type $`(2,2,\mathrm{},2)`$. By similar argument used in proving (\*), (!) is still true when $`A^{}=E_1\times \mathrm{}\times E_g`$ and $`L^{}=L_1\times L_2\mathrm{}\times L_g`$, where $`L_j`$ are line bundles of degree 2 on the elliptic curves $`E_j`$. Hence our assertion is true for a generic pair $`(B^{},M^{})`$. $`\mathrm{}`$
So, for a generic pair $`(A,L)`$, the map $`IPH^0(L)IPH^0(M^2)`$, given as $`(\mathrm{},s_\chi ,\mathrm{})(\mathrm{},s_\chi ^2,\mathrm{})`$ is a morphism and we obtain a commutative diagram (I),
$`\begin{array}{ccccc}A& \stackrel{\varphi _L}{}& \varphi _L(A)& & IP^{2^g1}=IPH^0(L)\\ \pi & & & & p\\ B=A/G& \stackrel{\varphi _{M^2}}{}& 𝒦(B)& & IP^{2^g1}=IPH^0(M^2)\end{array}`$
where $`p(\mathrm{},s_\chi ,\mathrm{})=(\mathrm{},s_\chi ^2,\mathrm{})`$.
###### Remark 2.2
Since $`\varphi _{M^2}\pi `$ is a morphism, $`\varphi _L`$ is a morphism i.e. $`L`$ is base point free.
###### Lemma 2.3
Consider a pair $`(A,L)`$ as in 1.2. Let $`\gamma 𝒢(L)(\times (i)`$ be an element of order 2. Then $`H^0(L)H^0(L)_\gamma ^\pm `$.
Proof: Case 1: Suppose $`\gamma =g𝒢(L)`$. Then the action of $`\gamma `$ is fixed point free on $`A`$. Hence by Atiyah- Bott fixed point theorem,
$$h^0(L)_\gamma ^+=h^0(L)_\gamma ^{}=h^0(L)/2.$$
Case 2: Suppose $`\gamma =i`$. Then
$$h^0(L)_i^\pm =h^0(L)/2\pm 2^{gs1}$$
( see , 4.6.6), where $`s`$ is the number of odd integers occurring in the type of $`L`$.
Case 3: Suppose $`\gamma =i.g`$ and $`H^0(L)=H^0(L)_\gamma ^+`$, where $`g𝒢(L)`$ is an element of order 2. Let $`sH^0(L)_g^{}`$. Then $`\gamma (s)=s`$ gives $`i(s)=s`$, i.e. $`sH^0(L)_i^{}`$. Hence $`H^0(L)_g^{}H^0(L)_i^{}`$ . But this contradicts the fact that $`h^0(L)_g^{}=2^{g1}`$ and $`h^0(L)_i^{}=2^{g1}2^{gs1}`$ (here $`s>1`$). Similarly $`H^0(L)H^0(L)_\gamma ^{}`$. $`\mathrm{}`$
Suppose $`\varphi _L`$ is not birational and is a finite morphism of degree $`d`$, $`d>1`$. Notice that $`A\stackrel{\varphi _{M^2}\pi }{}𝒦(B)`$ is a Galois covering with Galois group $`(G,i)(\frac{ZZ}{2ZZ})^{g+1}`$ and we have the extension of fields, $`Q(𝒦(B))Q(\varphi _L(A))Q(A)`$. Hence the Galois group of $`Q(A)`$ over $`Q(\varphi _L(A))`$ is a subgroup of $`(G,i)`$, say $`H`$, of order $`d`$. Let $`\gamma H`$. Then $`\gamma `$ is an involution on $`A`$, given as $`aϵa+g`$ where $`ϵ=\pm 1`$, $`gG`$ and it induces an involution $`\gamma ^{^{}}`$ on $`H^0(L)`$.
Hence $`\varphi _L`$ factorizes as $`A\stackrel{\psi _1}{}A/(\gamma )\stackrel{\psi _2}{}\varphi _L(A)IP^{2^g1}`$. This means that the morphism $`\psi _2`$ is given by the pair $`(N,H^0(L)_\gamma ^{^{}}^+)`$ or $`(N^{^{}},H^0(L)_\gamma ^{^{}}^{})`$, where $`N`$ and $`N^{^{}}`$ are line bundles on $`A/(\gamma )`$ whose pullback to $`A`$ is $`L`$. By 2.3, $`H^0(L)H^0(L)_\gamma ^{^{}}^\pm `$ and hence $`\varphi _L(A)`$ is a degenerate variety in $`IP^{2^g1}`$. This contradicts the fact that the morphism $`\varphi _L`$ is given by a complete linear system. Hence $`\varphi _L`$ is a birational morphism.
## 3 Configuration when $`g=3`$
Assume $`g=3`$. Choose a $`\mathrm{𝑡ℎ𝑒𝑡𝑎}\mathrm{𝑠𝑡𝑟𝑢𝑐𝑡𝑢𝑟𝑒}`$ $`f:𝒢(L)Heis(2,4)`$, ( i.e. $`f`$ is an isomorphism which restricts to identity on $`IC^{}`$.) This induces an isomorphism $`H^0(L)V(2,4)`$ and a level structure $`K(L)\frac{ZZ}{2ZZ}\frac{ZZ}{4ZZ}\frac{ZZ}{2ZZ}\frac{ZZ}{4ZZ}`$. Let $`\sigma _1,\tau _1,\sigma _2,\tau _2`$ be the generators of the summands such that $`o(\sigma _i)=2`$ and $`o(\tau _i)=4`$. The Weil form $`e^L`$ is given as
$$e^L(\sigma _1,\sigma _2)=1$$
$$e^L(\tau _1,\tau _2)=i$$
$$e^L(\sigma _i,\tau _j)=1.$$
Then we see that the subgroup $`G=<\sigma _1,\tau _1^2,\tau _2^2>`$ of $`K(L)`$ is maximal isotropic for the form $`e^L`$.
We may assume $`L`$ is strongly symmetric (see , Remark 2.4., p.160), i.e., $`e_{}^L(g)=1`$ for all $`gK(L)_2`$, after choosing a normalized isomorphism $`\psi :Li^{}(L)`$, i.e. $`\psi (0)=+1`$. Here $`e_{}^L:A_2\{\pm 1\}`$ is a quadratic form whose value at an element $`a`$, of order 2 is the action of $`\psi `$ at the fibre of $`L`$ at $`a`$.
Consider the exact sequence
$$1IC^{}𝒢(L)K(L)0$$
and the homomorphism $`\delta _1:𝒢(L)𝒢(L)`$, $`zizi`$. Then $`\delta _1(z)=\alpha z^1`$ for some $`\alpha IC^{}`$.
By , Proposition 2.3, p.141, we further assume that $`f`$ is a $`\mathrm{𝑠𝑦𝑚𝑚𝑒𝑡𝑟𝑖𝑐}\mathrm{𝑡ℎ𝑒𝑡𝑎}\mathrm{𝑠𝑡𝑟𝑢𝑐𝑡𝑢𝑟𝑒}`$, i.e. $`f\delta _1=D_1f`$, where $`D_1:Heis(\delta )Heis(\delta )`$ is the homomorphism $`(\alpha ,x,l)(\alpha ,x,l)`$.
###### Lemma 3.1
If $`z𝒢(L)`$ is an element of order 2 and $`z\pm 1`$ then $`\delta _1(z)=e_{}^L(z)z`$.
Proof: : See , Proposition 3, p.309. $`\mathrm{}`$
###### Remark 3.2
Let $`\sigma _1^{^{}},\sigma _2^{^{}},\tau _1^{^{}},\tau _2^{^{}}𝒢(L)`$ be lifts of $`\sigma _1,\sigma _2,\tau _1,\tau _2`$ such that $`o(\sigma _i^{^{}})=2,o(\tau _i^{^{}})=4.`$ Since $`\tau _i^2G`$, $`e_{}^L(\tau _i^2)=1`$, hence by 3.1, $`\delta _1((\tau _i^{^{}})^2)=(\tau _i^{^{}})^2`$. Hence $`\delta _1(\tau _i^{^{}})=c.\tau _{}^{^{}}{}_{i}{}^{1},c=\pm 1`$. We may assume $`c=+1`$, by suitably altering the lift $`\tau _i^{^{}}`$.
Let $`G^{^{}}=<\sigma _1^{^{}},(\tau _1^{^{}})^2,(\tau _2^{^{}})^2>𝒢(L).`$
Then $`L`$ descends to a principal polarization $`M`$ on $`B=A/G`$.
As remarked in Section 2,
$$H^0(L)=_{\chi \widehat{G}}H^0(ML_\chi )$$
and $`\{s_\chi H^0(ML_\chi ),\chi \widehat{G}\}`$ form a basis of $`H^0(L)`$.
Consider the commutative diagram,
$`\begin{array}{ccc}A& \stackrel{\psi _L}{}& Pic^0(A)\\ \pi & & \widehat{\pi }\\ B& \stackrel{\psi _M}{}& Pic^0(B)\end{array}`$
where $`\psi _L(a)=t_a^{}LL^1`$ and $`\psi _M(b)=t_b^{}MM^1`$. Then $`\psi _M`$ is an isomorphism and since $`\widehat{\pi }(L_\chi )=0`$, we have $`\pi ^1\psi _M^1(L_\chi )K(L)\chi \widehat{G}`$. Hence $`ML_\chi t_b^{}M`$ where $`b\pi (K(L))`$. The basis elements $`\{s_\chi \}_{\chi \widehat{G}}`$ can be written as $`s_0,s_1=\sigma _2^{^{}}(s_0),s_2=\tau _1^{^{}}(s_0),s_3=\tau _2^{^{}}(s_0),s_4=\sigma _2^{^{}}\tau _1^{^{}}(s_0),s_5=\sigma _2^{^{}}\tau _2^{^{}}(s_0),s_6=\tau _1^{^{}}\tau _2^{^{}}(s_0),s_7=\sigma _2^{^{}}\tau _1^{^{}}\tau _2^{^{}}(s_0)`$.
###### Lemma 3.3
If $`aK(L)_2`$, then $`a.i=i.a`$.
Proof: By 3.1, $`\delta _1(a)=e_{}^L(a)a`$. Since $`e_{}^L(a)=1,a.i=i.a`$. $`\mathrm{}`$
In particular, $`g^{^{}}i(s_0)=ig^{^{}}(s_0)`$, for all $`g^{^{}}G^{^{}}`$. Since $`g^{^{}}s_0=s_0`$, $`i(s_0)H^0(M)`$. This implies that $`i(s_0)=\pm s_0`$. We may assume $`i(s_0)=s_0`$.
###### Lemma 3.4
a) $`i\sigma _2^{^{}}(s_0)=\sigma _2^{^{}}(s_0)`$.
b) $`i\tau {}_{}{}^{}{}_{j}{}^{}(s_0)=\tau _j^{^{}}(s_0)`$.
c) $`i\sigma _2^{^{}}\tau _j^{^{}}(s_0)=\sigma _2^{^{}}\tau _j^{^{}}(s_0)`$.
d) $`i\tau _1^{^{}}\tau _2^{^{}}(s_0)=\tau _1^{^{}}\tau _2^{^{}}(s_0)`$.
e)$`i\sigma _2^{^{}}\tau _1^{^{}}\tau _2^{^{}}(s_0)=\sigma _2^{^{}}\tau _1^{^{}}\tau _2^{^{}}(s_0)`$
Proof: We will use 3.3 and the fact that $`g^{^{}}(s_0)=s_0,`$ for all $`g^{^{}}G^{^{}}`$.
a) $`i\sigma _2^{}(s_0)=\sigma _2^{}i(s_0)=\sigma _2^{^{}}(s_0)`$.
b) $`i\tau {}_{}{}^{}{}_{j}{}^{}(s_0)=\tau _j^1i(s_0)=\tau _j^3(s_0)=\tau _j^{^{}}(s_0)`$, ( since $`\tau _j^2G^{}`$).
c) $`i\sigma _2^{^{}}\tau _j^{^{}}(s_0)=\sigma _2^{}i\tau _j^{}(s_0)=\sigma _2^{^{}}\tau _j^{^{}}(s_0)`$.
d) $`i\tau _1^{^{}}\tau _2^{^{}}(s_0)=\tau _1^1i\tau _2^{}(s_0)=\tau _1^{}\tau _1^2\tau _2^{}(s_0)=\tau _1^{}\tau _2^{}\tau _1^2(s_0)=\tau _1^{^{}}\tau _2^{^{}}(s_0)`$ ( since $`e^L(\tau _1^2,\tau _2^{})=1,\tau _1^2G^{}`$).
e) $`i\sigma _2^{^{}}\tau _1^{^{}}\tau _2^{^{}}(s_0)=\sigma _2^{}i\tau _1^{^{}}\tau _2^{^{}}(s_0)\sigma _2^{^{}}\tau _1^{^{}}\tau _2^{^{}}(s_0)`$ $`\mathrm{}`$
Hence we have shown the following.
###### Proposition 3.5
The vector subspace $`H^0(L)_i^+`$ of $`H^0(L)`$ is generated by the sections $`s_0,s_1,s_2,s_3,`$ $`s_4,s_5`$ and the subspace $`H^0(L)_i^{}`$ of $`H^0(L)`$ is generated by the sections $`s_6`$ and $`s_7`$.
We then have the commutative diagram,
$`\begin{array}{ccccc}A& \stackrel{\varphi _L}{}& \varphi _L(A)& & IP(H^0(L))\\ \pi & & & & p\\ B=A/G& \stackrel{\varphi _{M^2}}{}& 𝒦(B)& & IP(H^0(M^2))\end{array}`$ …(I).
Here $`degree(p)=2^7`$ and $`degree(\pi )=8`$. Since we have shown that $`\varphi _L`$ is a birational morphism, $`degree(\varphi _L)=1`$ and hence $`degree(p|_{\varphi _L(A)})=2^4`$. The ramification locus of $`p|_{\varphi _L(A)}`$ is $`_{i=0}^7(H_i\varphi _L(A))`$, where $`H_i`$ is the hyperplane $`\{s_i=0\}`$ in $`IP(H^0(L))`$, $`0i7`$.
Consider the group $`J`$ generated by the projective transformations $`\alpha _i`$,
$$(s_0,\mathrm{},s_i,\mathrm{},s_7)(s_0,\mathrm{},s_i,\mathrm{},s_7)$$
for $`i=1,\mathrm{},7`$.
Then $`order(J)=2^7`$ and the group $`J`$ is the Galois group of the finite morphism $`p`$.
###### Proposition 3.6
The group $`G^{}\times <i>`$ can be identified as a subgroup of $`J`$.
Proof: : Since the action of $`gG`$ on the abelian threefold is fixed point free, the $`\pm 1`$-eigenspaces of $`H^0(L)`$ under the transformation $`gG^{^{}}`$ are equidimensional. Also, $`g(s_\chi )=\chi (g).s_\chi `$, for all $`\chi \widehat{G}`$, implies that $`g=\alpha _i\alpha _j\alpha _k\alpha _lJ`$, for some $`0i<j<k<l7`$. Here $`\alpha _0=\alpha _1\alpha _2\mathrm{}\alpha _7`$. By 3.5, $`i(s_0:\mathrm{}:s_7)=(s_0:\mathrm{}:s_5:s_6:s_7)`$. Hence the involution $`i=\alpha _6.\alpha _7`$. Hence we can identify $`G^{^{}}\times <i>`$ as a subgroup of $`J`$. $`\mathrm{}`$
Moreover, since the Galois group of the morphism $`p`$, $`Gal(p)=J`$ and the subgroup $`G^{}\times <i>J`$, leaves the image $`\varphi _L(A)`$ invariant in $`IPH^0(L)`$, we have the following
###### Proposition 3.7
Consider the commutative diagram (I). The inverse image of the variety, $`𝒦(B)`$, has eight distinct components $`h(\varphi _L(A))`$, where $`hJ/(G^{}\times <i>)`$.
In Section 2, we have seen that $`\{t_0=s_0^2,t_1=\sigma _2^{}(s_0^2),t_2=\tau _1^{}(s_0^2),t_3=\tau _2^{}(s_0^2),t_4=\sigma _2^{}\tau _1^{}(s_0^2),t_5=\sigma _2^{}\tau _2^{}(s_0^2),t_6=\tau _1^{}\tau _2^{}(s_0^2),t_7=\sigma _2^{}\tau _1^{}\tau _2^{}(s_0^2)\}`$
form a basis of $`H^0(M^2)`$.
###### Remark 3.8
( We use the same notations for the elements in $`K(L)`$ and their images in $`K(M^2)`$.) The elements $`\sigma _2^{},\tau _1^{},\tau _2^{}`$ of $`𝒢(M^2)`$ act on these sections as follows.
| | $`\sigma _2^{^{}}`$ | $`\tau _1^{^{}}`$ | $`\tau _2^{^{}}`$ |
| --- | --- | --- | --- |
| $`t_0`$ | $`t_1`$ | $`t_2`$ | $`t_3`$ |
| $`t_1`$ | $`t_0`$ | $`t_4`$ | $`t_5`$ |
| $`t_2`$ | $`t_4`$ | $`t_0`$ | $`t_6`$ |
| $`t_3`$ | $`t_5`$ | $`t_6`$ | $`t_0`$ |
| $`t_4`$ | $`t_2`$ | $`t_1`$ | $`t_7`$ |
| $`t_5`$ | $`t_3`$ | $`t_7`$ | $`t_1`$ |
| $`t_6`$ | $`t_7`$ | $`t_3`$ | $`t_2`$ |
| $`t_7`$ | $`t_6`$ | $`t_5`$ | $`t_4`$ |
Now let $`H_i=\{s_i=0\}`$ denote the coordinate hyperplanes in $`IPH^0(L)`$, for $`i=0,1,\mathrm{},7`$. Consider the curve $`C=H_6H_7\varphi _L(A)`$. Then the involution $`i`$ acts trivially on the curve $`C`$ and hence the degree of the restricted morphism $`\varphi _L^1(C)C`$ is at least $`2`$.
###### Proposition 3.9
The restricted morphism $`\varphi _L^{}:\varphi _L^1(C)C`$ is of degree $`2`$.
Proof: : Consider the commutative diagram
$`\begin{array}{ccc}\varphi _L^1(C)& \stackrel{\varphi _L^{}}{}& C\\ \pi ^{}& & p^{}\\ \varphi _{M^2}^1(p(C))& \stackrel{\varphi _{M^2}^{}}{}& p(C)\end{array}`$
Suppose the degree of the restricted morphism $`\varphi _L^{}`$ is greater than $`2`$. Since the Galois group of the morphism $`\varphi _{M^2}^{}\pi ^{}`$ is the group $`G\times <i>`$, the Galois group of $`\varphi _L^{}`$ contains an element $`gG`$. Hence the element $`g`$ acts trivially on the curve $`C`$. This means that $`C`$ is contained in one of the eigenspaces $`IPW^\pm `$ of $`IPH^0(L)`$, for the action of $`g`$. We claim that the intersection $`\varphi _L(A)IPW^\pm `$ is at most a finite set of points. This will give a contradiction.
If $`g^{}=\{aK(L):e^L(a,g)=1\}`$, then $`\frac{g^{}}{<g>}Heis(1,1,4)`$ or $`Heis(1,2,2)`$ and the group $`\frac{g^{}}{<g>}`$ acts on the linear space $`IPW^\pm `$. Hence projecting from $`IPW^\pm `$ gives a map $`\varphi _g:\frac{A}{<g>}IPW^{}`$, which is base point free in the first case ( by ) and has a finite base locus in the second case ( by ). This proves our claim. $`\mathrm{}`$
Now, the group $`G`$ leaves the curve $`C`$ invariant and moreover since $`\sigma _2(H_6)=H_7`$, we get $`\sigma _2(C)=C`$. Hence the curves
$$\tau _1(C)=H_3H_5\varphi _L(A)$$
$$\tau _2(C)=H_2H_4\varphi _L(A)$$
$$\tau _1.\tau _2(C)=H_0H_1\varphi _L(A)$$
are also invariant for the action of $`\sigma _2`$ and since for $`xC`$, $`i(x)=x`$, $`i.\tau _j^2(\tau _j(x))=\tau _j^2.\tau _j^1i(x)=\tau _j(x)`$. By $`K(L)`$-invariance of the image $`\varphi _L(A)`$, we get
###### Corollary 3.10
The morphism $`\varphi _L`$ restricts to a morphism of degree $`2`$ on the curves $`\varphi _L^1(C),\varphi _L^1(\tau _1(C)),\varphi _L^1(\tau _2(C))`$ and $`\varphi _L^1(\tau _1.\tau _2(C))`$, onto their respective images. Moreover, the Galois groups of these restricted morphisms are $`<i>,<i.\tau _1^2>,<i.\tau _2^2>`$ and $`<i.\tau _1^2.\tau _2^2>`$, respectively.
Let $`A_2^+`$ denote the set of points of order 2 on $`A`$ where the involution $`i`$ acts on the fibre of $`L`$ at those points as $`+1`$ and $`A_2^{}`$ denote the set of points where $`i`$ acts as $`1`$. By , Remark 4.7.7, $`cardinality(A_2^+)=48`$ and $`cardinality(A_2^{})=16`$. Hence if $`aA_2^{}`$ and $`sH^0(L)_i^+`$, then $`s(a)=0`$. This implies that for $`aA_2^{}`$,$`\varphi _L(a)=(0:0:\mathrm{}:0:c_1:c_2)IPH^0(L)`$, for some $`c_1,c_2IC`$.
###### Proposition 3.11
Let $`aA_2^+`$( respectively $`A_2^{}`$) and $`gK(L)_2`$. Then $`a+gA_2^+`$( respectively $`A_2^{}`$).
Proof: : Let $`gK(L)_2`$ and $`(g,\varphi )𝒢(L)`$ be a lift of order $`2`$ and $`\psi :Li^{}(L)`$ be the normalized isomorphism. By , Proposition 3, p.309,
$$\delta _1(g,\varphi )=(g,(t_g^{}\psi )^1i^{}\varphi \psi )$$
$$=e_{}^L(g).(g,\varphi )$$
$$=(g,\varphi )(sinceLisstronglysymmetric).$$
Hence the following diagram commutes
$`\begin{array}{ccc}L& \stackrel{\psi }{}& i^{}(L)\\ \varphi & & i^{}(\varphi )\\ t_g^{}L& \stackrel{t_g^{}(\psi )}{}& i^{}t_g^{}L=t_g^{}(i^{}L)\end{array}`$
Evaluating at $`aA_2^+`$( respectively $`A_2^{}`$), gives $`\psi (a)=t_g^{}(\psi )(a)=\psi (a+g)`$, i.e. $`a+gA_2^+`$( respectively $`A_2^{}`$). $`\mathrm{}`$
Now let $`aA_2^{}`$ then $`\varphi _L(a)=(0:\mathrm{}:c_1,c_2)`$ for some $`c_1,c_2IC`$. Then $`\sigma _2\varphi _L(a)=(0:\mathrm{}:c_2:c_1)`$. We may assume $`c_20`$. Let $`P_0=\varphi _L(a)=(0:\mathrm{}:c:1)`$ and $`Q_0=p(P_0)=(0:\mathrm{}:c^2:1)`$, for some $`cIC`$.
###### Proposition 3.12
The points $`h(P_0),hK(L)/<\tau _1^2,\tau _2^2>`$ are of degree $`4`$ on the image $`\varphi _L(A)`$.
Proof: : By 3.11, the action of $`G`$ on the set $`A_2^{}`$ has two distinct orbits, namely $`O_1=\{a+g:gG\}`$ and $`O_2=\{a+\sigma _2+g:gG\}`$. Then $`\varphi _{M^2}\pi (O_1)=Q_0`$ and $`\varphi _{M^2}\pi (O_2)=\sigma _2(Q_0)`$. Notice that $`P_0\tau _1(C)\tau _2(C)\tau _1.\tau _2(C)`$. Hence, by 3.10, $`\varphi _L^1(P_0)=\{a,a+2\tau _1,a+2\tau _2,a+2\tau _1+2\tau _2\}`$. The assertion now follows from the $`K(L)`$-invariance of the image $`\varphi _L(A)`$. $`\mathrm{}`$
###### Corollary 3.13
The points $`b(Q_0)`$, where $`b<\pi (\sigma _2),\pi (\tau _1),\pi (\tau _2)>`$, lie on the Kummer $`𝒦(B)`$.
## 4 Prym Varieties
We recall few facts on Prym varieties ( see , , , for details).
Let $`C`$ be a smooth projective curve of genus $`g`$. We will assume $`C`$ has no vanishing theta nulls. In particular, when $`g=3`$, this means $`C`$ is a non-hyperelliptic curve. A point of order $`2`$, in $`X=Jac(C)`$, say $`x`$, defines an unramified $`2`$\- sheeted cover $`C_x`$ of $`C`$, $`q_x:C_xC`$. Let $`P_x=Ker(Nm(q_x):Jac(C_x)X)^o`$, where $`\mathrm{`}o\mathrm{`}`$ denotes the connected component containing $`0Jac(C_x)`$. Here $`Nm(q_x)(𝒪(r_iP_i))=𝒪(r_iq_x(P_i))`$ is the norm map. This defines a principally polarized abelian variety $`(P_x,\theta _{P_x})`$, of dimension $`g1`$. Since the kernel of the dual map $`q_x^{^{}}:XJac(C_x)`$ is generated by the element $`x`$, $`q_x^{}`$ induces an isomorphism $`x^{}/xP_x[2]`$. Since $`q_x𝒪_{C_x}𝒪_Cx`$, we have $`detq_x𝒪_{C_x}x`$. Hence $`det(q_x(p))`$ is also $`x`$, for any $`pker(Nm(q_x))`$.
Fix a $`zX`$ with $`z^2x`$. This gives a map
$$\psi _x:Ker(Nm(q_x))P_xP_xSU_C(2).$$
where $`\psi _x(p)=(q_xp)z`$.
The image of $`\psi _x`$ is independent of the choice of $`z`$. Recall the map
$$SU_C(2)\stackrel{\varphi }{}|2\theta _C|IP(H^0(SU_C(2),))$$
where $``$ generates $`Pic(SU_C(2))ZZ`$.
Let $`IPV_x^+`$ and $`IPV_x^{}`$ be the two eigenspaces for the action of $`x`$ on $`|2\theta _C|`$. Then there is one component of $`Ker(Nm(q_x))`$ in each eigenspace. So we get a map $`\varphi _x:P_xIPV_x`$.
###### Proposition 4.1
The map $`\varphi _x:P_xIPV_x`$ is the natural map
$$P_x𝒦(P_x)IP(H^0(P_x,2\theta _{P_x})IPV_x.$$
Proof: : See , Proposition 1, p.745.
###### Proposition 4.2
For any curve $`C`$ and any $`x`$ in $`X[2]\{0\}`$, we have $`𝒦(C)IPV_x=𝒦(P_x[2])`$, ( the Schottky Jung relations).
Proof: : See , Proposition 2 (1), p.746.
## 5 Situation in $`IP(H^0(L))`$, when $`g=3`$.
We now assume $`B=J(C)`$, where $`J(C)`$ is the Jacobian of a non-hyperelliptic curve $`C`$ of genus 3. (This is the generic situation, since the dimension of the moduli space of principally polarized abelian threefolds is $`6`$ which equals the dimension of the moduli space of curves of genus 3.) Recall the results of Narasimhan and Ramanan ( $`Theorem1.3,Theorem1.4`$), to obtain a morphism
$$J(C)\stackrel{\varphi _{2\theta }}{}𝒦(C)F|2\theta |$$
where
1) $`F`$ is a quartic hypersurface and is the isomorphic image of the moduli space $`SU_C(2)`$ and
2) the Kummer variety $`𝒦(C)`$ is precisely the singular locus of $`F`$.
We will use the following
###### Proposition 5.1
Let $`L`$ be an ample line bundle of type $`\delta =(d_1,d_2,\mathrm{},d_g)`$ on an abelian variety $`A`$. Then the set of irreducible representations of the theta group $`𝒢(L)`$, where $`\alpha IC^{}`$ acts as multiplication by $`\alpha ^n`$( called as of ’weight $`n`$’), is in bijection with the set of characters on the subgroup of $`n`$torsion elements, $`K(L)_n`$, of $`K(L)`$. Moreover, the dimension of any such representation is $`\frac{d_1.d_2\mathrm{}d_g}{(n,d_1)\mathrm{}(n,d_g)}`$. ( $`(n,d_i)`$ denotes the greatest common divisor of $`n`$ and $`d_i`$.)
Proof: : When $`n=2`$, the statement is proved in , Proposition 3.2, p.142. The same proof holds when $`n>2`$, by choosing a section over the subgroup of $`n`$-torsion elements, $`K(L)_n`$, of $`K(L)`$ in the exact sequence
$$1IC^{}𝒢(L)K(L)0$$
in the proof of , Proposition 3.2. $`\mathrm{}`$
###### Corollary 5.2
The quartic $`F`$ in $`|2\theta |`$ is $`𝒢(2\theta )`$-invariant and the linear span of the eight cubics $`\{\frac{dF}{dt_i}\}`$ for $`i=0,1,\mathrm{},7`$ form an irreducible $`𝒢(2\theta )`$-module where $`\alpha IC^{}`$ acts as multiplication by $`\alpha ^3`$.
Proof: : Consider the multiplication maps $`Sym^nH^0(2\theta )\stackrel{\rho _n}{}H^0(2n\theta )`$. Then $`I_n=Ker(\rho _n)=`$ vector space of degree $`n`$ forms containing the image $`𝒦(B)`$ in $`IPH^0(2\theta )`$. Since the vector spaces $`Sym^nH^0(2\theta )`$ and $`H^0(2n\theta )`$ ( via the homomorphism $`𝒢(2\theta )\stackrel{ϵ_n}{}𝒢(2n\theta )`$) are $`𝒢(2\theta )`$-modules, of weight $`n`$ and $`\rho _n`$ is equivariant for the $`𝒢(2\theta )`$-action, $`I_n`$ is also a $`𝒢(2\theta )`$-module of weight $`n`$. Now the homogenous polynomial $`FI_4`$ and the partial derivatives $`\frac{dF}{dt_i}I_3`$. By 5.1, it follows that $`F`$ is $`𝒢(2\theta )`$-invariant , upto scalars. If $`z𝒢(2\theta )`$, then $`z\frac{dF}{dt_i}=\frac{d(zF)}{d(zt_i)}=\alpha \frac{dF}{d(zt_i)}W=IC\{\frac{dF}{dt_i}\}_{i=0}^7`$, for some scalar $`\alpha `$. Hence $`W`$ is a $`𝒢(2\theta )`$-module of weight $`3`$. By 5.1, dimension of such an irreducible representation is $`8`$. This proves our assertion. $`\mathrm{}`$
Similarly, we see that $`R=F(s_0^2,\mathrm{},s_7^2)`$ is a $`𝒢(L)`$-invariant octic hypersurface in $`IPH^0(L)`$, by applying 5.1.
Recall the Weil form $`e^L`$ on $`K(L)`$ and the isotropic subgroup $`G=<\sigma _1,\tau _1^2,\tau _2^2>K(L)`$. Then $`e^L(\sigma _2+g,\sigma _1)=1`$, for all $`gG`$. Let $`a=\sigma _2+g`$, for $`gG`$ and $`a^{}=\sigma _2^{}+g^{}𝒢(L)`$.
Recall the basis $`\{s_0,s_1,\mathrm{},s_7\}`$ of $`H^0(L)`$ and $`\{s_0^2,\mathrm{},s_7^2\}`$ of $`H^0(M^2)`$, (see Section 3). Let $`W_a^+`$ and $`W_a^{}`$ denote the eigen spaces in $`H^0(L)`$, for the action of $`a^{}`$. Now $`IPW_a^\pm =\{s=0:sW_a^{}\}`$ and $`IPV_a^+=\{t=0:tH^0(M^2)_a^{}\}`$. Now $`W_{\sigma _2}^\pm =IC\{s_0\pm s_1,s_2\pm s_4,s_3\pm s_5,s_6\pm s_7\}`$ and $`H^0(M^2)_{\sigma _2}^{}=IC\{s_0^2s_1^2,s_2^3s_4^2,s_3^2s_5^2,s_6^2s_7^2\}`$.
Then $`p`$ restricts on $`IPW_{\sigma _2}^\pm IPV_{\sigma _2}^+`$ as $`(s_0;s_2:s_3,s_6)(s_0^2:s_2^2:s_3^2:s_6^2)`$, of degree $`2^3`$. Similarly, one checks that if $`a=\sigma _2+g,gG`$ then $`p`$ restricts to $`IPW_a^\pm IPV_{\sigma _2}^+`$ as $`(z_0:\mathrm{}:z_3)(z_0^2:\mathrm{}:z_3^2)`$ of degree $`2^3`$.
###### Proposition 5.3
Consider a principally polarized abelian surface $`(Y,P)`$, which is not a product of elliptic curves. Let $`y_1,y_2Y`$ be elements of order $`2`$, such that $`e^{P^2}(y_1,y_2)=1`$. Then we have the following.
1) There is a polarized abelian surface $`(Z,N)`$, such that $`N`$ is strongly symmetric of type $`(1,4)`$ and there is a covering map $`f:ZY`$ with the Galois group of the map $`f`$ being isomorphic to $`ZZ/2ZZ\times ZZ/2ZZ`$.
2) The vector space $`H^0(N)`$ can be written as
$$H^0(N)=H^0(P)H^0(t_{y_1}^{}P)H^0(t_{y_2}^{}P)H^0(t_{y_1+y_2}^{}P).$$
and there is a commutative diagram
$`\begin{array}{ccccc}Z& \stackrel{\varphi _N}{}& \varphi _N(Z)& & IP^3=IPH^0(N)\\ f& & & & q\\ Y& \stackrel{\varphi _{P^2}}{}& 𝒦(Y)& & IP^3=IPH^0(M^2)\end{array}`$
where $`q(r_0:r_1:r_2:r_3)=(r_0^2:r_1^2:r_2^2:r_3^2)`$. Here $`\{r_0,r_1,r_2,r_3\}`$ is a basis obtained from above decomposition of $`H^0(N)`$, such that $`r_0,r_1,r_3H^0(N)_i^+`$ and $`r_3H^0(N)_i^{}`$.
Proof: : 1) Consider the isomorphism $`\varphi _P:YPic^0(Y)`$, $`bt_b^{}PP^1`$. Let $`L_{y_1}`$ and $`L_{y_2}`$ denote the images of $`y_1`$ and $`y_2`$ under this map. These two line bundles define an unramified cover, $`f:ZY`$, whose Galois group is isomorphic to $`ZZ/2ZZ\times ZZ/2ZZ`$, as asserted.
Then $`N=f^{}P`$ is an ample line bundle and $`dimH^0(N)=4`$. So to see that $`N`$ is of type $`(1,4)`$, it is enough to show that $`K(N)`$ has an element of order $`4`$. Consider the commutative diagram
$`\begin{array}{ccc}Z& \stackrel{\psi _N}{}& Pic^0(Z)\\ f& & \widehat{f}\\ Y& \stackrel{\psi _M}{}& Pic^0(Y)\end{array}`$
Then $`\widehat{f}\psi _M(y_i)=0`$. This implies that if $`z_1`$ and $`z_2`$ are in $`Z`$ such that $`f(z_i)=y_i`$, then $`z_1,z_2K(N)`$. Moreover, since $`e^{P^2}(y_1,y_2)=1`$ and $`N^2f^{}(P^2)`$, we have $`e^{N^2}(z_1,z_2)=1`$. This gives $`e^N(z_1,z_2)=\pm i`$. Hence the elements $`z_1,z_2K(N)`$ are of order $`4`$.
2) Clearly, $`f_{}N=P(PL_{y_1})(PL_{y_2})(PL_{y_1+y_2})`$. Now, in the algebraic equivalence class of $`N`$, there are strongly symmetric line bundles. Hence, by tensoring $`P`$ with a suitable line bundle of order $`2`$, we may assume that $`N=f^{}P`$ is strongly symmetric and $`r_0H^0(P)`$ is such that $`i(r_0)=r_0`$.
Since $`N`$ is strongly symmetric, by 3.1, $`\delta _1(z_j^{^{}})^2=(z_j^{^{}})^2`$, for some lifts $`z_j^{^{}}𝒢(N)`$ of $`z_jK(N)`$. We may further choose the lifts such that $`\delta _1(z_j^{^{}})=(z_j^{^{}})^1`$, ( as in 3.2). In particular, the descent data of $`N`$ to $`P`$ is $`K^{^{}}=<(z_1^{^{}})^2,(z_2^{^{}})^2>𝒢(N)`$, which is a splitting over $`K=<z_i^2,z_2^2>K(N)`$ in the exact sequence
$$1IC^{}𝒢(N)K(N)0.$$
This means $`(z_j^{^{}})^2r_0=r_0`$. Also this gives
As in 3.5, we see that
$$i.z_j^{^{}}(r_0)=z_j^{^{}}(r_0)$$
and
$$i.z_1^{^{}}.z_2^{^{}}(r_0)=z_1^{^{}}.z_2^{^{}}(r_0).$$
Thus $`r_0,r_1=z_1^{^{}}(r_0),r_2=z_2^{^{}}(r_0)H^0(N)_i^+`$ and $`r_3=z_1^{^{}}.z_2^{^{}}(r_0)H^0(N)_i^{}`$.
Hence one sees as earlier that $`Gal(q)=<z_1^2,z_2^2,i>`$, with a commuatative diagram as in 5.3. $`\mathrm{}`$
###### Proposition 5.4
Let $`a=\sigma _2+g,gG`$ and $`IPW_a`$ denote an eigenspace of $`a`$ in $`IPH^0(L)`$. Then there is an abelian surface $`Z`$ and a symmetric line bundle $`N`$ on $`Z`$ of type $`(1,4)`$ such that $`Z\stackrel{\varphi _N}{}IP(H^0(N))\stackrel{Heis(4)}{}IPW_aIPH^0(L)`$. Moreover, under this isomorphism, the image $`\varphi _N(Z)`$ is mapped onto the $`Heis(4)`$-invariant surface $`S=RIPW_a`$, where $`R`$ is the $`Heis(2,4)`$\- invariant hypersurface of degree 8 in $`IPH^0(L)`$, defined by $`F(s_0^2:s_1^2:,,,:,s_7^2)=0`$. ( $`F`$ being the Coble quartic).
Proof: : Consider the restricted morphism $`p:IPW_aIPV_a`$, given as $`(z_0:\mathrm{}:z_3)(z_0^2:\mathrm{}:z_3^2)`$. Then $`a`$ acts trivially on $`IPW_a`$ and $`a^{}/a(Heis(4)`$) acts on $`IPW_a`$, (here $`a^{}=\{yK(L):e^L(a,y)=1\}`$). Hence there is a $`Heis(4)`$\- action on $`IPW_a`$ and similarly a $`Heis(2,2)`$\- action on $`IPV_a`$. By 4.1, there is a principally polarized abelian surface $`(P_a,\theta _{C_a})`$, ($`P_a`$ being the Prym variety associated to the element $`\pi (a)K(M^2))`$, such that
$$P_a𝒦(P_a)|2\theta _{C_a}|IPV_a.$$
Consider the images of $`\tau _1,\tau _2`$, which are elements of order 2 in $`J(C)`$. Since $`e^{L^2}(\tau _i,a)=1`$, for the Weil form $`e^{2\theta }`$ on $`J(C)[2]`$, $`\pi (\tau _1),\pi (\tau _2)\pi (a)^{}/\pi (a)`$. Moreover, $`e^{2\theta }(\pi (\tau _1),\pi (\tau _2))=1`$. By 4.2, the points $`\varphi _{M^2}\pi (\tau _i)`$, are nodes in the Kummer of the Prym variety $`P_a`$. These nodes correspond to elements of order $`2`$ in $`P_a`$, say $`\beta _1`$ and $`\beta _2`$. Since the Weil form $`e^{2\theta _{C_a}}`$ on $`P_a[2]`$ is induced from the Weil form $`e^{2\theta }`$, we have $`e^{2\theta _{C_a}}(\beta _1,\beta _2)=1`$. By 5.3, there is a polarized abelian surface $`(Z,N)`$ of type $`(1,4)`$, such that the following diagram commutes
$`\begin{array}{ccccc}Z& \stackrel{\varphi _N}{}& \varphi _N(Z)& & IPH^0(N)\\ f& & & & q\\ P_a& \stackrel{\varphi _{2\theta _{C_a}}}{}& 𝒦(P_a)& & |2\theta _{C_a}|\end{array}`$
and for the choice of basis $`\{r_0,r_1,r_2,r_3\}`$, in 5.3 2), the morphism $`q`$ is defined as $`(r_0:r_1:r_2:r_3)(r_0^2:r_1^2:r_2^2:r_3^2)`$, with $`Gal(q)=<z_1^2,z_2^2,i>`$, ($`z_j`$ as in 5.3).
Now, $`R`$ is the $`Heis(2,4)`$-invariant octic $`F(s_0^2:\mathrm{}:s_7^2)=0`$, where $`F`$ is the Coble quartic. Note that $`S=RIPW_a`$ is $`a^{}/a`$-invariant and is mapped onto the Kummer, $`K(P_a)`$, under the restriction morphism. Moreover, the Galois group of $`p_{|S}`$ is $`<\tau _1^2,\tau _2^2,i>`$ which is isomorphic to the Galois group of $`q`$. Hence there is a $`Heis(4)`$\- isomorphism $`IPH^0(N)IPW_a`$, such that the Heisenberg invariant octic surface $`\varphi _N(Z)`$ is mapped onto the Heis(4)-invariant octic surface $`S=RIPW_a`$. This proves the assertion. $`\mathrm{}`$
It is known that the Kummer $`𝒦(P_a)`$, has $`6`$ of its nodes in each of the coordinate hyperplane, namely the coordinate points and $`3`$ other distinct points. The preimages of the coordinate points are the coordinate points in $`IPH^0(N)`$ and $`q`$ is etale over the other $`3`$ points which are the pinch points of $`\varphi _N(Z)`$ in the respective coordinate hyperplane.
###### Proposition 5.5
$`\varphi _N(Z)`$ has exactly $`48`$ pinch points, $`12`$ in each coordinate hyperplane.
Proof: : See , Proposition 2.2, p.633.
Let $`T_a`$ denote the set of pinch points and the coordinate points in $`\varphi _N(Z)`$.
###### Proposition 5.6
The components $`h(\varphi _L(A)),hH`$ (here $`H=J/(G^{^{}}\times i)`$) and $`IPW_a`$ intersect at the subset $`T_a`$ of $`\varphi _N(Z)`$. In particular $`_{hH}h(\varphi _L(A))=_{a=\sigma _2+g,gG}T_a`$.
Proof: : Since $`\pi ^1𝒦(C)=_{hH}h(\varphi _L(A))`$, by 4.2 and 5.5, we conclude that $`h(\varphi _L(A))IPW_a=T_a`$, for all $`hH`$. This gives the assertion. $`\mathrm{}`$
## 6 Some remarks
a) Consider the moduli space $`𝒜_{(1,2,4)}^l`$ of triples $`(A,c_1(L),f)`$, where $`f:K(L)ZZ/DZZ\times ZZ/DZZ`$ is a level structure, ( here $`D=(1,2,4)`$). Consider the subset of $`𝒜_{(1,2,4)}^l`$, $`𝒜_{(1,2,4)}^{lo}`$, parametrizing triples which admit a $`(ZZ/2ZZ)^3`$isogeny to the Jacobian of a non-hyperelliptic curve.
Since $`dim𝒜_{(1,2,4)}^{lo}=dim𝒜_{(1,2,4)}^l=6`$ and $`c_1(L)`$ gives a birational morphism , $`𝒜_{(1,2,4)}^{lo}`$ is an open subset of $`𝒜_{(1,2,4)}^l`$.
Consider a triple $`(A,c_1(L),f)𝒜_{(1,2,4)}^{lo}`$. We have seen that there is a $`Heis(2,4)`$-invariant octic hypersurface $`R`$, defined by $`F(s_0^2:s_1^2:\mathrm{}:s_7^2)=0`$, ( $`F`$ being the Coble quartic), such that $`\varphi _L(A)RIPV(2,4)`$. In fact $`h(\varphi _L(A))Sing(R),`$ for all $`hH`$, ( $`H`$ as in 5.6).
Now $`F`$ is a $`Heis(2,2,2)`$-invariant quartic polynomial in $`IPV(2,2,2)`$. Since the space of $`Heis(2,2,2)`$-invariant quartics is $`14`$-dimensional, ( see , p.186\]), the space of $`Heis(2,4)`$-invariant octics in $`IP^7`$ which are of the form $`R=F(s_0^2:\mathrm{}:s_7^2)`$ where $`F`$ is a $`Heis(2,2,2)`$-invariant quartic, is also $`14`$-dimensional. Call this space as
$`P(Sym^8V(2,4)^{Heis(2,4)^{}})=IP^{14}`$.
So there is a morphism
$$𝒜_{(1,2,4)}^{lo}\stackrel{T}{}IP^{14}$$
where $`T`$ is defined as $`(A,c_1(L),f)R`$.
One may try to study this morphism, from a moduli point of view.
b) Consider the special basis $`\{s_0^2,\mathrm{},s_7^2\}`$ ( which is different from the usual $`\mathrm{𝐻𝑒𝑖𝑠𝑒𝑛𝑏𝑒𝑟𝑔}`$ basis) of $`H^0(2\theta )`$ and the action of the elements of the subgroup $`<\sigma _2,\tau _1^2,\tau _2^2>K(2\theta )`$ on this basis ( see 3.8).
Also, by 3.12, the points $`b(P_0)\varphi _L(A)`$, where $`b<\sigma _2,\tau _1,\tau _2>K(L)`$, $`P_0=(0:\mathrm{}:0:c:1)`$ and the point $`Q_0=(0:\mathrm{}:0:c^2:1)𝒦(C)`$, for some non-zero $`cIC`$. With these data, in addition to knowing the geometry of $`SU_C(2)`$ in $`|2\theta |`$\- linear system one may try to know the equation of the $`\mathrm{𝐶𝑜𝑏𝑙𝑒}\mathrm{𝑞𝑢𝑎𝑟𝑡𝑖𝑐}`$, in terms of this basis $`\{s_0^2,\mathrm{},s_7^2\}`$.
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# General Relativistic Constraints on Emission Models of Anomalous X-ray Pulsars
## 1. INTRODUCTION
Among pulsating compact X-ray sources are a small subset with pulsation periods between 6 and 12 seconds, soft spectra, and no identifiable companions. The first suggestion that these objects might form a separate class of neutron stars – later called the anomalous X-ray pulsars (AXPs) – was made by Mereghetti & Stella (1995), who proposed that they might be powered by accretion from a very low mass companion (this was also earlier hinted at by Hellier 1994). The lack of optical counterparts, however, and the absence of observable Doppler shifts in the frequency of the X-ray pulses led van Paradijs, Taam, & van den Heuvel (1995) to favor a different accretion model for AXPs. According to their suggestion, material from a fossil accretion disk, possibly the debris of a disrupted binary companion after a period of common-envelope evolution, is being accreted by a solitary neutron star. Recently, a similar model has been proposed by Chatterjee et al. (2000), in which the accreting material is supplied by the post-supernova fallback material from the neutron star progenitor itself.
In a different class of models, AXPs are considered to be isolated neutron stars, spinning down by magnetic dipole radiation. Because of their unusually high period derivatives, a simple application of the relationship between the dipole magnetic field strength and the spin-down rate implies a field strength for these objects of $`10^{14}10^{15}`$ G (see, e.g., Thompson & Duncan 1996). Two main types of models that rely upon the presumed high magnetic field of the stars have been proposed. Thompson & Duncan (1996; see also Duncan & Thompson 1992) suggested that the released energy may be drawn from the decay of the magnetic field itself and from differential movements in the stellar crust. This model also serves to explain the bursts observed from soft gamma-ray repeaters as being produced by larger-scale magnetospheric phenomena. Because of such models, AXPs are often called “magnetars” and are grouped into the same class of sources as the soft gamma-ray repeaters (see, e.g, Thompson & Duncan 1995, 1996; also Hurley 2000 for a review of SGR properties and models). In the alternative high magnetic-field model, Heyl & Hernquist (1997a, 1997b) suggested that AXPs draw their energy from the residual thermal energy of the star itself.
All these models face a number of difficulties. For example, if AXPs are powered by accretion from a stellar companion, the absence of detectable Doppler shifts in the arrival of X-ray pulses cannot be easily explained(Mereghetti et al. 1998). If accretion is from either a companion or a fossil disk, the optical fluxes one would expect directly from the disk or due to X-ray reprocessing are too high compared to the observed upper limits (see, e.g., Perna, Hernquist, & Narayan 2000; Hulleman et al. 2000). On the other hand, in the magnetar model the absence of bright faster AXPs and the observed variations in spin-down rates are hard to account for (see, e.g., Baykal & Swank 1996; Marsden et al. 2000; see, however, Heyl & Hernquist 1999; Melatos 1999).
In all the above models of AXP, nearly all of the X-ray emission is produced at the surface of the neutron star. It is well known that strong gravitational fields tend to smooth out the variability produced by a spinning compact star, even if the emission is highly localized in bright spots (Pechenick, Ftaclas, & Cohen 1983). Indeed, the X-ray pulse amplitudes of the three radio pulsars that show thermal emission from their surfaces are only $`30`$% (see, e.g., the discussion in Page 1995; Harding & Muslimov 1998). This is in contrast to the non-thermal emission from radio pulsars (which is magnetospheric in origin) and from accretion-powered X-ray pulsars (which is from collimated accretion columns) which often show pulse amplitudes as high as $`90`$% (see, e.g., Nagase 1989). In this respect, AXPs are similar to accretion-powered pulsars, showing X-ray pulse amplitudes anywhere between $`10`$% and $`70`$%.
We study in this paper a set of variability diagnostics that may be used in constraining emission models of AXPs. We examine three parameters: the pulse fraction observed at infinity, which is a measure of the overall amplitude of the variations, and the Fourier amplitudes at the first and second harmonics of the neutron-star spin frequency. We find tight constraints on the properties of magnetar models. Most importantly, we are unable to reproduce the observed variability properties of AXPs with thermal cooling models.
## 2. FORMALISM
In order to determine how relativistic effects suppress or enhance variability amplitudes, we need to consider curved photon paths from the surface of the star to an observer at infinity. Since the objects under consideration are rotating slowly, we use the Schwarzschild spacetime, which is appropriate for a non-spinning mass, and ignore effects such as relativistic frame dragging, which are important only for rapidly spinning objects. In this section, we outline the basic ingredients of our method, drawing on the work of Pechenick et al. (1983).
For each model we specify both the brightness distribution over the surface of the star and the effective beaming of radiation, i.e., we specify the specific intensity integrated over photon energy, $`I(\theta ,\varphi ,\theta ^{})`$, that emerges from each point on the stellar surface with polar coordinates $`(\theta ,\varphi `$) at an angle $`\theta ^{}`$ with respect to the normal.
In the current analysis we consider a number of different, physically motivated, mathematical expressions for the dependence of the emerging specific intensity on $`\theta ^{}`$ (hereafter called the beaming function). This allows us to explore a large parameter space and draw conclusions that do not depend strongly on any particular emission model. We consider isotropic emission, i.e., no dependence on $`\theta ^{}`$, as well as the beaming described by the Hopf function (Chandrasekhar 1950, eq. \[III.50\])
$$I(\theta ,\varphi ,\theta ^{})=I_0(\theta ,\varphi )\left(\underset{a=1}{\overset{3}{}}\frac{L_a}{1+k_a\mathrm{cos}\theta ^{}}\mathrm{cos}\theta ^{}+Q\right),$$
(1)
where the parameters $`L_a`$, $`k_a`$, and $`Q`$ are given in Chandrasekhar (1950; Table III.VII). The Hopf function describes the beaming of radiation emerging from a scattering atmosphere heated from below and is, therefore, suitable for a weakly-magnetic H–He atmosphere at energies $`1`$ keV (Zavlin et al. 1998).
For photon energies near the cyclotron energy ($`E_{\mathrm{cyc}}11.6[B/10^{12}`$ G$`]`$ keV), the beaming function for a magnetic atmosphere does not always decrease monotonically away from the radial direction. Indeed the beaming function may have a local minimum at small angles from the normal (Zavlin et al. 1995). However, for a dipole magnetic field this effect is not very significant and results in a rather small pulse fraction ($`30`$%) even when general relativistic effects are not taken into account (see Zavlin et al. 1995). Moreover, the cooling models of AXPs discussed here
require that the magnetic field of the neutron star is $`10^{14}10^{15}`$ G and hence that the cyclotron energy is $`110`$ MeV, i.e, much larger than the $`1`$ keV photon energies that are of interest here. For such energies, it is reasonable to consider beaming functions that decrease monotonically away from the normal to the surface (Zavlin et al. 1995).
In models of accretion columns, the interaction of radiation with a plasma in the strong magnetic field of the neutron star, as well as the possibility of radiation being obscured by the accretion column (see, e.g., Riffert et al. 1993) leads to a much sharper beaming. For this reason, we also consider beaming functions of the form
$$I(\theta ,\varphi ,\theta ^{})=I_0(\theta ,\varphi )\mathrm{cos}^n\theta ^{}.$$
(2)
According to Nagel (1981; see also Meszaros & Nagel 1985), the radiation pattern emerging from an accretion column at low accretion rates is described quite well with $`n23`$.
Given a model for the X-ray emission from the stellar surface, the flux measured by an observer at distance $`d`$, whose polar coordinates with respect to the stellar rotation axis are $`(\beta ,\mathrm{\Phi })`$, is given by (cf. Pechenick et al. 1983)
$`F_{\mathrm{}}(\beta ,\mathrm{\Phi })`$ $`=`$ $`\left({\displaystyle \frac{R}{d}}\right)^2\left({\displaystyle \frac{M}{R}}\right)^2\left(1{\displaystyle \frac{2M}{R}}\right)^2`$ (3)
$`\times {\displaystyle _{x=0}^{x_{\mathrm{max}}}}{\displaystyle _{y=0}^{2\pi }}I(\theta ,\varphi ,\theta ^{})xdxdy.`$
In the above equation $`\theta `$, $`\varphi `$, and $`\theta ^{}`$ depend implicitly on the angles $`x`$ and $`y`$ as described in Pechenick et al. (1983), $`x_{\mathrm{max}}(R/M)(12M/R)^{1/2}`$, $`M`$ and $`R`$ are the neutron star mass and radius, and we have set $`c=G=1`$. The double integral has an integrable pole at $`x_{\mathrm{max}}`$. In the calculations presented here we have evaluated this integral to an accuracy of $`10^3`$ using Romberg integration of the fifth order.
| AXP PROPERTIES | | | |
| --- | --- | --- | --- |
| Label | Source Name | Pulse Fraction | $`I_2/I_1`$ |
| A | 1E 1048.1–5937 | 0.76 | 0.15 |
| B | 1E 1841–045 | 0.15 | 0.54 |
| C | AX J1845.0–0258 | 0.63 | 0.14 |
| D | 1RXS J170849.0$``$400910 | 0.50 | 0.40 |
| E | 4U 0142$`+`$61 | 0.17 | 0.70 |
| F | 1E 2259$`+`$586 | 0.35 | 1.35 |
References.— A: Oosterbroek et al. 1998 (SAX); B: Gotthelf et al. 1999 (ASCA+SAX); C: Torri et al. 1998 (ASCA); D: Sugizaki et al. 1997 (ASCA); E: Israel et al. 1999 (SAX); F: Iwasawa et al. 1992 (GINGA).
The degree of suppression of the pulsation amplitude depends sensitively on the compactness of the neutron star. Figure 1 shows the ratio $`pRc^2/2GM`$ for different neutron-star masses and equations of state. Based on this figure, we limit our parameter study to $`p=2,3`$, and 4; larger values of $`p`$ would correspond to unrealistically light neutron stars, even for the stiffest proposed equations of state (i.e., $`1.4M_{}`$ even for equation of state M).
## 3. PULSATION AMPLITUDES
The pulsation amplitudes observed in anomalous X-ray pulsars allow us to place constraints on models of X-ray emission from their surface. In this paper, we do not attempt to fit particular observed pulse profiles but rather try to set general constraints on large classes of models. For this reason, we only consider the pulse fraction, defined as (cf. eq.)
$$PF\frac{F_{\mathrm{}}^{\mathrm{max}}F_{\mathrm{}}^{\mathrm{min}}}{F_{\mathrm{}}^{\mathrm{max}}+F_{\mathrm{}}^{\mathrm{min}}},$$
(4)
and the Fourier amplitudes of the harmonics of the pulse frequency, defined by
$$F_{\mathrm{}}(\beta ,\mathrm{\Phi })=I_0(\beta )+I_1(\beta )\mathrm{cos}(\mathrm{\Phi })+I_2(\beta )\mathrm{cos}(2\mathrm{\Phi })+\mathrm{}$$
(5)
Table I shows the observed pulse fractions and harmonic content for six AXPs, as inferred approximately from published lightcurves. We see that half of the known systems have pulse fractions of 0.5 or larger, with the largest being $`0.7`$. As we discuss below, this fact provides strong constraints on some models of AXPs.
### 3.1. Cooling of Magnetic Neutron Stars
The brightness distribution on the surface of a strongly-magnetic cooling neutron star depends on the local magnetic field strength $`B`$ and its angle $`\psi `$ with respect to the local radial direction (see, e.g., Heyl & Hernquist 1998). The flux emerging from a spot on the stellar surface is $`B^m\mathrm{cos}^2\psi `$, with $`m0.4`$ approximating well the numerical results (Heyl & Hernquist 1998). For a dipole stellar field we, therefore, use
$$I_0(\theta ,\varphi )I_0\frac{\mathrm{cos}^2\theta }{(3\mathrm{cos}^2\theta +1)^{1m/2}},m=0.4,$$
(6)
where ($`\theta `$,$`\varphi `$) are polar coordinates on the stellar surface with respect to the magnetic axis.
Figure 2 shows the predicted pulse fractions for a cooling neutron-star model, for different choices of the angle $`\alpha `$ between the magnetic dipole axis and the rotation axis and the angle $`\beta `$ between the light of sight to the observer and the rotation axis, as well as for different beaming functions. Some general features of this figure are worth noting. For example, the predicted pulse fraction remains unchanged when the magnetic inclination and the inclination of the observer are interchanged. This is true because the flux at infinity measured at any give pulse phase $`\mathrm{\Phi }`$ depends only on the angular distance
$$\theta _0=\mathrm{cos}^1(\mathrm{sin}\alpha \mathrm{sin}\beta \mathrm{cos}\mathrm{\Phi }+\mathrm{cos}\alpha \mathrm{cos}\beta )$$
(7)
between the magnetic axis and the direction to the observer, which is symmetric in $`\alpha `$ and $`\beta `$. Furthermore, the predicted pulse fraction is maximum when $`\alpha =\beta `$. The overall maximum occurs when $`\alpha =\beta =90^{}`$ and the pulse fraction is zero when either $`\alpha =0^{}`$ or $`\beta =0^{}`$.
The most important result from Figure 2 is that even for the most favorable geometry (i.e., $`\alpha =\beta =90^{}`$) and for strong beaming (i.e., $`n=3`$), the pulse fraction does not exceed 37%. The reason is that the surface brightness distribution (6) is too smooth. Assuming a less relativistic neutron star ($`R/2M=4`$), or even a significantly stronger dependence of the emerging flux on the local magnetic field ($`m=1`$ in eq. ) results in only a modest increase of the maximum pulse fraction (51% and 50%, respectively, for $`n=3`$). We thus conclude that neutron star cooling models with the surface emission described by equation (6) cannot reproduce the large pulse fractions observed from AXPs.
### 3.2. Localized Thermal Emission
Thermal emission from the neutron-star surface can, in principle, be more localized than indicated in equation (6), if it is confined mostly around the stellar magnetic poles. For example, metallicity gradients on the surface of a cooling neutron star, possibly produced by magnetically-channelled fallback material during the supernova explosion, can lead to larger effective temperatures near the magnetic poles than in the magnetic equator (see, e.g., Pavlov et al. 2000). Alternatively, non-uniform heating of the neutron-star atmosphere, e.g., caused by magnetic field decay or crustal heating (see, e.g., Thompson & Duncan 1996), can lead to more localized thermal emission from its surface.
We estimate the magnitude of the first effect using the analytic expressions for neutron-star atmosphere models given by Chabrier, Potekhin, & Yakovlev (1997). We assume that a fraction $`(\mathrm{\Delta }M/M)(\mathrm{\Omega }/4\pi )`$ of light element material has accumulated only over a fraction ($`\mathrm{\Omega }/4\pi `$) of the stellar surface. We neglect the fact that the surface layers of the neutron star are in the liquid phase and hence lateral diffusion may smooth out the metallicity gradient (L. Hernquist, private communication). For a given core temperature we then calculate the ratio $`f`$ of the bolometric flux emerging from the light-element region of the surface to the bolometric flux from the region consisting purely of iron. This flux contrast is shown in Figure 3, for different values of the mass fraction of the accreted atmosphere. The flux contrast does not further increase, for $`(\mathrm{\Delta }M/M)(\mathrm{\Omega }/4\pi )10^7`$, because the base of the accreted layer reaches densities that are high enough ($`10^{10}`$ g cm<sup>-3</sup>) for it to be part of the isothermal core. Although the above models are strictly valid only for weakly magnetic neutron stars, the flux ratios are significantly more sensitive to the composition of the neutron-star envelope than the magnetic field strength (cf. Heyl & Hernquist 1997b).
According to Figure 3, the expected flux contrast between the polar caps and the rest of the stellar surface can be at most $`f9`$ and the contrast depends very weakly on the amount of accreted material. For the flux contrast to attain its maximum value, the polar caps must accrete $`\mathrm{\Delta }M10^7M_{}`$, while the rest of the surface must accumulate more than three orders of magnitude less material. Even if these conditions were met and even if the photons emerged from the stellar surface strictly radially, the pulse fraction would be at most $`(f1)/(f+1)0.8`$.
Any realistic beaming function reduces this upper bound significantly below the highest observed pulse fraction for AXPs ($`0.7`$; see Table I), even for mildly relativistic neutron stars. This is demonstrated in Figure 4, which shows the maximum pulse fraction of radiation emerging from a neutron star with $`R/M=3`$, as a function of the angular radius of each polar cap, for various realistic beaming functions, and an emerging flux from the polar caps that is nine times larger than the rest of the stellar surface. For small polar-cap sizes, the radiation flux emerging from the caps is only a small perturbation to the total brightness of the star and hence the pulse fraction is small. With increasing cap size the relative contribution of the caps to the brightness of the star increases, leading to an increase of the pulse fraction, until the caps cover a large enough fraction of the stellar surface and the pulse fraction drops again. For the flux contrast and neutron star radius used here, the maximum pulse fraction occurs when the polar-cap size is $`40^{}`$ and this maximum value is $`55`$%. Clearly, no realistic model of this kind can fit the pulse fraction of $`0.7`$ seen in 1E 1048.1$``$5937.
The pulse fraction can be further enhanced if the fluxes emerging from two antipodal polar caps are unequal, e.g., because of uneven fallback or crustal heating. As the contrast between the two polar caps increases, the pulse fraction becomes larger, but only at the expense of the Fourier amplitudes at the harmonics of even order. This is shown in Figure 5, where the predicted ratio $`I_2/I_1`$ is plotted against the pulse fraction for $`R/M=3`$, two
antipodal caps with a half-opening angle of 40 degrees, and beaming described by the Hopf function. For these calculations, the two caps are assumed to have different emerging radiation fluxes, with the flux of the brightest set to 9 times the flux emerging from the rest of the neutron-star surface. For any given pulse fraction, there exists a maximum value that the ratio $`I_2/I_1`$ can attain, shown by the dashed line in Figure 5. As a result, detection of a source with a large pulse fraction and significant amplitude at the even order harmonics can exclude such a configuration.
The observed properties of some AXPs are not consistent even with a model with such unequal polar caps (Figure 5). For example, the source 1E 2259$`+`$586 is characterized by a large pulse fraction and high harmonic content that cannot be achieved by any of the configurations considered here. It does not, therefore, appear plausible for thermal emission from the stellar surface to account for the variability properties of AXPs. Achieving the kind of variability observed is possible only when the emission is both localized and strongly beamed.
### 3.3. Accretion Onto Magnetic Neutron Stars
In accretion models of AXPs (see, e.g., Mereghetti & Stella 1995; van Paradijs et al. 1996; Chatterjee et al. 1999), a large fraction of the X-ray emission may be produced mainly at localized “hot spots” where the accretion columns meet the stellar surface. In order to take such configurations into account, we describe the surface brightness distribution with two circular antipodal spots, with a brightness that is constant over their surface area. We denote by $`\alpha `$ the angular distance of the center of each spot from its closest rotation pole.
If AXPs are powered by accretion from a geometrically thin disk, the half-opening angle of each polar cap can be very small. We estimate the size of each polar cap as the angular distance on the stellar surface from the magnetic axis of the footpoint of the last magnetic field line that
penetrates the accretion disk. Given that all AXPs are observed to be spinning down, their accretion disks must be truncated near the outer corotation radius (see, e.g., Ghosh & Lamb 1979)
$`R_{\mathrm{co}}`$ $``$ $`\left({\displaystyle \frac{GMP^2}{4\pi ^2}}\right)^{1/3}`$ (8)
$``$ $`550R\left({\displaystyle \frac{R}{10^6\text{cm}}}\right)^1\left({\displaystyle \frac{M}{1.4M_{}}}\right)^{1/3}\left({\displaystyle \frac{P}{6\text{s}}}\right)^{2/3},`$
where $`P`$ is the spin period of the pulsar. For a dipole magnetic field, the quantity $`\mathrm{sin}^2\theta /r`$, with $`\theta `$ measured from the magnetic pole, remains constant along a field line, and therefore the half-opening angle of the polar cap is $`(R/R_{\mathrm{co}})^{1/2}2.5^{}`$.
Figure 6 shows the pulse fraction calculated for a model with two identical hot spots, a half-opening angle of 5 degrees, and various beaming functions (see eq. ). Because of the general relativistic deflection of light, we see that even such a small spot cannot produce a $`70`$% pulse fraction, unless there is significant beaming ($`n2`$). On the other hand, if the emerging radiation is strongly beamed towards the radial direction, a large pulse fraction can be achieved for a wide range of polar cap sizes ($`40^{}`$; Fig. 7) and for all realistic neutron-star radii (Fig. 8).
## 4. CONCLUSIONS
In this paper we have used the high pulse amplitudes observed from a number of AXPs to constrain the properties of their emission mechanism. We find that the observations can be accounted for only if the surface emission is localized (half-opening angle $`<40^{}`$) and strongly beamed
($`n2`$ in eq. ), as summarized quantitatively in Figure 9. These constraints are a consequence of the compactness of the neutron stars and the resulting strong general relativistic deflection of photon trajectories. Our conclusions are valid for all realistic neutron star masses and radii.
The properties of individual sources offer a number of additional clues. For example, the double-peaked pulse profile of 1E 2259$`+`$586 requires that the emission is localized around two antipodal spots on the neutron-star surface, probably associated with the magnetic poles. Furthermore, the change in the relative strength of the two peaks observed with GINGA (Iwasawa et al. 1992) implies that the pulse shape cannot be solely due to geometric effects but should also reflect a flux contrast between the two antipodal spots. Furthermore, this flux contrast should be variable and, therefore, cannot be caused by a non-dipolar magnetic field configuration or a non-uniform fallback of low metallicity material.
Such arguments, together with the constraints presented in Figures 2, 4, 5, and 9 appear to rule out thermal cooling models for AXPs. They are also inconsistent with those magnetar models in which most of the X-ray flux originates from heating in the deep surface layers of the neutron star. On the other hand, the localized emission and beaming predicted by accretion models seem to be consistent with the observations. A magnetospheric model in which the neutron-star surface is heated over localized spots by particle bombardment may also be viable, though the beaming properties of such a model are unknown.
We thank Deepto Chakrabarty, Lars Hernquist, Vicky Kaspi, Jessica Lackey, and Feryal Özel for many useful discussions. This work was supported in part by NSF grant AST 9820686. D. P. acknowledges the support of a postdoctoral fellowship of the Smithsonian Institution.
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# Non-Abelian Einstein-Born-Infeld Black Holes
## 1 Introduction
Non-linear electrodynamics was proposed in the thirties to remove singularities associated with charged pointlike particles . More recently such non-linear theories were considered in order to remove singularites associated with charged black holes .
Among the non-linear theories of electrodynamics Born-Infeld (BI) theory is distinguished, since BI type actions arise in many different contexts in superstring theory . The non-abelian generalization of BI theory yields an ambiguity in defining the Lagragian from the Lie-algebra valued fields. While the superstring context favors a symmetrized trace , the resulting Lagranian is so far only known in a perturbative expansion . However, the ordinary trace structure has also been suggested .
Motivated by this strong renewed interest in BI and non-abelian BI theory, recently non-perturbative non-abelian BI solutions were studied, both in flat and in curved space. In particular, non-abelian BI monopoles and dyons exist in flat space as well as in curved space together with dyonic BI black holes.
Likewise one expects regular and black hole solutions in pure SU(2) Einstein-Born-Infeld (EBI) theory, representing the BI generalizations of the regular Bartnik-McKinnon solutions and their hairy black hole counterparts . This expectation is further nurtured by the observation, that SU(2) BI theory possesses even in flat space a sequence of regular solutions .
Here we construct these SU(2) EBI regular and black hole solutions, employing the ordinary trace as in . The set of EBI equations depends essentially on one parameter, $`\gamma `$, composed of the coupling constants of the theory. The solutions are labelled by the node number $`n`$ of the gauge field function. We construct solutions up to node number $`n=5`$ and discuss the limiting solutions, obtained for $`n\mathrm{}`$, for various values of the horizon radius $`x_\mathrm{H}`$ and the parameter $`\gamma `$. We also address thermodynamic properties of the black hole solutions.
## 2 SU(2) EBI Equations of Motion
We consider the SU(2) EBI action
$$S=S_G+S_M=L_G\sqrt{g}d^4x+L_M\sqrt{g}d^4x$$
(1)
with
$$L_G=\frac{1}{16\pi G}R,L_M=\beta ^2(1\sqrt{1+\frac{1}{2\beta ^2}F_{\mu \nu }^aF^{a\mu \nu }\frac{1}{16\beta ^4}(F_{\mu \nu }^a\stackrel{~}{F}^{a\mu \nu }})^2),$$
(2)
field strength tensor
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+eϵ^{abc}A_\mu ^bA_\nu ^c,$$
(3)
gauge coupling constant $`e`$, gravitational constant $`G`$ and BI parameter $`\beta `$.
To construct static spherically symmetric regular and black hole solutions we employ Schwarzschild-like coordinates and adopt the spherically symmetric metric
$$ds^2=g_{\mu \nu }dx^\mu dx^\nu =A^2Ndt^2+N^1dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(4)
with the metric functions $`A(r)`$ and
$$N(r)=1\frac{2m(r)}{r}.$$
(5)
The static spherically symmetric and purely magnetic Ansatz for the gauge field $`A_\mu `$ is
$$A_0^a=0,A_i^a=ϵ_{aik}r^k\frac{1w(r)}{er^2}.$$
(6)
For purely magnetic configurations $`F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }=0`$.
With the ansatz (4)-(6) we find the static action
$`S=e{\displaystyle 𝑑r\frac{A}{2}\left[1+N\left(1+2r\left(\frac{A^{^{}}}{A}+\frac{N^{}}{2N}\right)\right)\right]}\alpha ^2\beta ^2r^2A\left[1\sqrt{1+{\displaystyle \frac{2Nw_{}^{2}{}_{}{}^{}}{\beta ^2e^2r^2}}+{\displaystyle \frac{(1w^2)^2}{\beta ^2e^2r^4}}}\right],`$ (7)
where $`\alpha ^2=4\pi G`$.
Introducing dimensionless coordinates and a dimensionless mass function
$$x=\sqrt{e\beta }r,\mu =\sqrt{e\beta }m$$
(8)
as well as the dimensionless parameter
$$\gamma =\frac{\alpha ^2\beta }{e},$$
(9)
we obtain the set of equations of motion
$`\mu ^{^{}}`$ $`=`$ $`\gamma x^2\left[1\sqrt{1+{\displaystyle \frac{2Nw_{}^{2}{}_{}{}^{}}{x^2}}+{\displaystyle \frac{(1w^2)^2}{x^4}}}\right],`$ (10)
$`{\displaystyle \frac{A^{^{}}}{A}}`$ $`=`$ $`\gamma {\displaystyle \frac{2w_{}^{2}{}_{}{}^{}}{x\sqrt{1+\frac{2Nw_{}^{2}{}_{}{}^{}}{x^2}+\frac{(1w^2)^2}{x^4}}}},`$ (11)
$`\left(AN{\displaystyle \frac{w^{^{}}}{\sqrt{1+\frac{2Nw_{}^{2}{}_{}{}^{}}{x^2}+\frac{(1w^2)^2}{x^4}}}}\right)^{^{}}`$ $`=`$ $`{\displaystyle \frac{Aw(w^21)}{x^2\sqrt{1+\frac{2Nw_{}^{2}{}_{}{}^{}}{x^2}+\frac{(1w^2)^2}{x^4}}}}.`$ (12)
In eq. (12) the metric function $`A`$ can be eliminated by means of eq. (11).
We consider only asymptotically flat solutions, where the metric functions $`A`$ and $`\mu `$ both approach a constant at infinity. Here we choose
$$A(\mathrm{})=1.$$
(13)
For magnetically neutral solutions the gauge field configuration approach a vacuum configuration at infinity
$$w(\mathrm{})=\pm 1.$$
(14)
Globally regular solutions satisfy at the origin the boundary conditions
$$\mu (0)=0,w(0)=1,$$
(15)
whereas black hole solutions with a regular event horizon at $`x_\mathrm{H}`$ satisfy there
$$\mu (x_\mathrm{H})=\frac{x_\mathrm{H}}{2},N^{}w^{}|_{x_\mathrm{H}}=\frac{w\left(w^21\right)}{x^2}|_{x_\mathrm{H}}.$$
(16)
## 3 Embedded Abelian BI Solutions
Before discussing the non-abelian BI solutions, let us briefly recall the embedded abelian BI solutions . With constant functions $`A`$ and $`w`$,
$$A=1,w=0,$$
(17)
they carry one unit of magnetic charge, and their mass function satisfies
$$\mu ^{}=\gamma \left(\sqrt{x^4+1}x^2\right),$$
(18)
or upon integration
$$\mu (x)=\mu (0)+\gamma _0^x\left(\sqrt{x_{}^{}{}_{}{}^{4}+1}x_{}^{}{}_{}{}^{2}\right)𝑑x^{}.$$
(19)
Their mass $`\mu _{\mathrm{}}`$ is thus given by
$$\mu _{\mathrm{}}=\mu (0)+\gamma \frac{\pi ^{3/2}}{3\mathrm{\Gamma }^2(3/4)}.$$
(20)
The solutions are classified according to the integration constant $`\mu (0)`$ . For $`\mu (0)>0`$ black hole solutions with one non-degenerate horizon $`x_\mathrm{H}`$ are obtained. For $`\mu (0)=0`$ black hole solutions with one non-degenerate horizon $`x_\mathrm{H}`$ are obtained for $`\gamma >1/2`$, otherwise the solutions possess no horizon. For $`\mu (0)<0`$ black hole solutions with two non-degenerate horizons, extremal black hole solutions with one degenerate horizon or solutions with no horizons are obtained, similarly to the Reissner-Nordstrøm (RN) case.
For black hole solutions with event horizon at $`x_\mathrm{H}=2\mu (x_\mathrm{H})`$ the integration constant $`\mu (0)`$ is given by
$$\mu (0)=\frac{x_\mathrm{H}}{2}\gamma _0^{x_\mathrm{H}}\left(\sqrt{x^4+1}x^2\right)𝑑x.$$
(21)
Since extremal black hole solutions satisfy
$$\mu ^{}(x_\mathrm{H})=\frac{\mu (x_\mathrm{H})}{x_\mathrm{H}}=\frac{1}{2},$$
(22)
eq. (18) yields for the degenerate horizon of extremal black holes
$$x_\mathrm{H}^{\mathrm{ex}}=\sqrt{\gamma \frac{1}{4\gamma }}.$$
(23)
Thus we find a critical value $`\gamma _{\mathrm{cr}}=1/2`$, where $`x_{\mathrm{H},\mathrm{cr}}^{\mathrm{ex}}=0`$ and $`\mu _{\mathrm{cr}}(0)=0`$. Finally, the Hawking temperature $`T`$ of the black holes is given by
$$T=\frac{1}{4\pi x_\mathrm{H}}\left[12\gamma \left(\sqrt{x_\mathrm{H}^4+1}x_\mathrm{H}^2\right)\right].$$
(24)
## 4 Regular Solutions
Recently, a sequence of non-abelian regular BI solutions was found in flat space by Gal’tsov and Kerner , labelled by the node number $`n`$ of the gauge field function $`w`$. We now consider this sequence of BI solutions in the presence of gravity, i.e. for finite $`\gamma `$ (and thus finite $`\alpha `$), and compare the non-abelian regular BI solutions to the regular Einstein-Yang-Mills (EYM) solutions .
With increasing $`\gamma `$ the regular BI solutions evolve smoothly from the corresponding flat space BI solutions. Let us first consider the sequence of regular BI solutions for $`\gamma _{\mathrm{cr}}`$. In Fig. 1 we show the gauge field function $`w`$ of the regular BI solutions with node numbers $`n=1`$, 3 and 5.
The metric function $`N`$ of these BI solutions is shown in Fig. 2 together with the metric function $`N`$ of the extremal abelian BI solution with $`x_\mathrm{H}^{\mathrm{ex}}=0`$ and one unit of magnetic charge. Clearly, with increasing $`n`$ the metric function of the non-abelian regular BI solutions tends to the metric function of the extremal abelian solution.
In Fig. 3 we show the charge function $`P(x)`$
$$P^2(x)=\frac{2x}{\gamma }\left(\mu _{\mathrm{}}\mu (x)\right),$$
(25)
obtained from an expansion of the abelian mass function for general charge $`P`$, for these regular BI solutions together with the constant charge $`P=1`$ of the extremal abelian BI solution with $`x_\mathrm{H}^{\mathrm{ex}}=0`$.
Furthermore, for $`\gamma =\gamma _{\mathrm{cr}}`$ the masses of the non-abelian regular BI solutions converge exponentially to the mass of the extremal abelian BI solution with $`x_\mathrm{H}^{\mathrm{ex}}=0`$ and unit magnetic charge, which represents the limiting solution of this sequence.
Let us now consider $`\gamma \gamma _{\mathrm{cr}}`$. For $`\gamma <\gamma _{\mathrm{cr}}`$ the sequence of non-abelian regular BI solutions tends to an abelian BI solution without horizon and unit magnetic charge, and with integration constant $`m(0)=0`$.
For $`\gamma >\gamma _{\mathrm{cr}}`$ we must distinguish two spatial regions, $`x<x_\mathrm{H}^{\mathrm{ex}}`$ (given in eq. (23)), and $`x>x_\mathrm{H}^{\mathrm{ex}}`$. Only in the region $`x>x_\mathrm{H}^{\mathrm{ex}}`$ the metric function $`N`$ of the sequence of non-abelian regular BI solutions tends to the metric function of the extremal abelian BI black hole solution with horizon $`x_\mathrm{H}^{\mathrm{ex}}`$ and unit magnetic charge. For $`x<x_\mathrm{H}^{\mathrm{ex}}`$ it tends to a non-singular limiting function. This is demonstrated in Fig. 4, where the metric function $`N`$ is shown for $`n=5`$ and $`\gamma =0.01`$, $`\gamma =1`$, and $`\gamma =100`$, together with the metric function $`N`$ of the corresponding limiting abelian BI solutions.
Thus the non-abelian regular BI solutions are very similar to their EYM counterparts. In particular, there with increasing node number $`n`$ the sequence of neutral $`SU(2)`$ EYM solutions also tends to a limiting charged solution, which for radius $`x1`$ is the extremal embedded RN solution with magnetic charge $`P=1`$ .
## 5 BI Black Hole Solutions
Imposing the boundary conditions eq. (16), leads to non-abelian BI black hole solutions with regular horizon. Like their non-abelian EYM counterparts, non-abelian BI black hole solutions exist for arbitrary value of the horizon radius $`x_\mathrm{H}`$. In both cases the sequences of neutral non-abelian BI black hole solutions tend to limiting solutions with unit magnetic charge.
For $`\gamma <\gamma _{\mathrm{cr}}`$ the limiting solution of the non-abelian BI black hole solutions is always the abelian BI black hole solution with the same horizon. The same holds true for $`\gamma =\gamma _{\mathrm{cr}}`$.
For $`\gamma >\gamma _{\mathrm{cr}}`$ and $`x_\mathrm{H}<x_\mathrm{H}^{\mathrm{ex}}`$ the limiting solution is the extremal abelian BI black hole solution with horizon $`x_\mathrm{H}^{\mathrm{ex}}`$
only in the region $`x>x_\mathrm{H}^{\mathrm{ex}}`$, but differs from it in the region $`x_\mathrm{H}<x<x_\mathrm{H}^{\mathrm{ex}}`$. For $`x_\mathrm{H}>x_\mathrm{H}^{\mathrm{ex}}`$ the limiting solution is always the abelian BI black hole solution with the same horizon. This is demonstrated in Fig. 5 for $`\gamma =1`$ and horizon radii $`x_\mathrm{H}=0.1`$, 0.2, 0.5, and 10.
The temperature of the non-abelian BI black holes is obtained from
$$T=\frac{1}{4\pi }AN^{}.$$
(26)
In Fig. 6 we show the inverse temperature of the non-abelian BI black hole solutions with node numbers $`n=1`$, 3 and 5 as a function of their mass for $`\gamma =0.01`$ and $`\gamma =1`$. Also shown is the inverse temperature of the corresponding limiting abelian BI solutions. Again, with increasing $`n`$ rapid convergence towards the limiting values is observed, analogously to the EYM and EYM-dilaton case .
Further details will be given elsewhere .
## 6 Conclusions
We have constructed sequences of regular and black hole solutions in SU(2) EBI theory. The solutions are labelled by the node number $`n`$. With increasing node number these sequences of non-abelian neutral solutions tend to limiting solutions, corresponding (at least in the outer spatial region) to abelian BI solutions with unit magnetic charge. These features are similar to those observed for non-abelian EYM and EYM-dilaton regular and black hole solutions . This similarity is also observed for the Hawking temperature.
By generalizing the framework of isolated horizons to non-abelian gauge theories , recently new results were obtained for EYM black hole solutions. In particular nontrivial relations between the masses of EYM black holes and regular EYM solutions were found, and a general testing bed for the instability of non-abelian black holes was suggested . Application of such considerations to non-abelian BI black holes appears to be interesting.
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# 1 Introduction
## 1 Introduction
Recent experimental results concerning the structure of the deuteron have led to the speculations that manifestation of the quark-gluon degrees of freedom are present even at relatively large distances between nucleons. Measurements of the cross section of the inclusive deuteron breakup $`A(d,p)X`$ reaction on carbon with the proton emitted at a zero degree have shown the relatively broad shoulder at internal nucleon momenta $`k0.35`$ GeV/c in the deuteron defined in the light-cone dynamics -. This enhancement has been observed later at different initial energies and for different $`A`$ values of the target -. This shoulder could not be reproduced by the calculations within relativistic impulse approximation (IA) using the standard deuteron wave functions , as well as by the inclusion of the rescattering corrections . Theoretical work of Kobushkin and Vizireva led to the possibility of existing of a $`6q`$ admixture in the deuteron wave function . This $`6q`$ amplitude, arising from the $`S`$ configurations of six quarks, must be added to the $`S`$ component of the standard deuteron wave function (DWF) with a relative phase, $`\chi `$. The fit of the experimental data gave the probability of the $`6q`$ configuration about $`4\%`$ and relative angle $`\chi `$ $`82^0`$ and $`61^0`$ for Paris and Reid Soft Core (RSC) $`NN`$ potentials, respectively. Admixture of a $`6q`$ state of about $`3.4\%`$ was imposed also in to describe the tail of momentum spectrum of the $`{}_{}{}^{12}C(d,p)X`$ reaction .
One of the important features of this hybrid wave function is that an additional $`6q`$ admixture masks the node of the $`NN`$ $`S`$ wave, what drastically reflects on the behaviour of polarization observables. For instance, the data on tensor analyzing power $`T_{20}`$ and cross section in inclusive deuteron breakup at zero degree and at $`2.1`$ GeV of the initial energy obtained at Saclay were explained by the hybrid wave function with $`4\%`$ of $`|6q>`$ configuration probability with $`55^0`$ of relative phase between $`|6q>`$ and $`S`$ component from RSC potential .
Recent measurements of tensor analyzing power $`T_{20}`$ for deuteron inclusive breakup at $`0^0`$ performed in Saclay and in Dubna at different energies and for different targets have shown the strong deviation from the IA predictions at $`k0.2`$ GeV/c. The behaviour of the polarization transfer coefficient from vector polarized deuteron to proton $`\kappa _0`$ also disagrees with the calculations using conventional DWFs at $`k0.2`$ GeV/c. On the other hand, both $`T_{20}`$ and $`\kappa _0`$ data demonstrate a weak dependence on $`A`$ value of the target, as well as an approximate energy independence, i.e. features of IA. Considering of the mechanisms additional to IA can not explain the experimental data.
Most intriguing feature of the experimental data is that tensor analyzing power $`T_{20}`$ in deuteron inclusive breakup and deuteron–proton backward elastic scattering show at high internal momenta of proton the same negative value $`0.3÷0.4`$ , incompatible with the predictions using any reasonable nucleon–nucleon potential. Various attempts were undertook to explain the $`T_{20}`$ data taking into account the nonnucleon degrees of freedom in the deuteron. An asymptotic negative limit of $`T_{20}`$ was obtained in framework of the QCD motivated approach based on reduced nuclear amplitude method . The results of calculations with the hybrid DWF allowed to describe satisfactory the $`T_{20}`$ data up to $`k1`$ GeV/c . Recently the data on $`T_{20}`$ and $`\kappa _0`$ in the $`{}_{}{}^{12}C(d,p)X`$ reaction at $`0^0`$ were reasonably reproduced within a model which incorporates multiple scattering and Pauli principle at the quark level . The additional account of the negative parity nucleon resonance exchanges improves the accordance of calculations with the experimental data on $`T_{20}`$ in backward elastic $`dp`$ scattering . The tensor analyzing power $`A_{yy}`$ in deuteron inclusive breakup obtained up to 600 MeV/c of a proton transverse momenta also disagrees with the calculations within hard scattering model using conventional DWFs. However, the sign of $`A_{yy}`$ at large proton momenta as at a zero angle <sup>1</sup><sup>1</sup>1At a zero angle emission angle $`A_{yy}=T_{20}/\sqrt{2}`$, as well as at a $`90^0`$ in the rest frame of the deuteron is the same as predicted by QCD motivated approach .
These peculiarities of the experimental data and relative successful attempts to describe them by the considering of the nonnucleon degrees of freedom stimulate to measure additional polarization observables crucial to the quark degrees of freedom in the deuteron.
In our previous paper we have considered the using of the polarized proton target and proton polarimeter to study the deuteron structure at short distances. Here we propose to study $`T`$-odd polarization observables in deuteron exclusive breakup in the collinear geometry and $`dp`$ backward elastic scattering in order to identify the exotic $`6q`$ configurations inside the deuteron.
## 2 Matrix elements of the $`dpppn`$ and $`dppd`$ reactions
In this section we analyze the polarization effects in two processes : deuteron breakup in the strictly collinear geometry, $`d+pp(0^0)+p(180^0)+n`$, and deuteron–proton backward elastic scattering, $`d+pp+d`$, using the hybrid DWF with the complex $`6q`$ admixture.
This function can be presented in the momentum space in the following form:
$`\mathrm{\Phi }_d\left(𝐩\right)={\displaystyle \frac{i}{\sqrt{2}}}{\displaystyle \frac{1}{\sqrt{4\pi }}}\psi _p^{\alpha +}\left[\left(U\left(p\right)\left(\stackrel{}{\sigma }\stackrel{}{\xi }\right){\displaystyle \frac{W\left(p\right)}{\sqrt{2}}}\left(3\left(\widehat{𝐩}\stackrel{}{\xi }\right)\left(\stackrel{}{\sigma }\widehat{𝐩}\right)\left(\stackrel{}{\sigma }\stackrel{}{\xi }\right)\right)\right)\sigma _y\right]_{\alpha \beta }\psi _n^{\beta +},`$ (1)
where $`\psi _p`$ and $`\psi _n`$ are the proton and neutron spinors, respectively, $`\stackrel{}{\xi }`$ is the deuteron polarization vector, defined in a standard manner:
$`\stackrel{}{\xi }_1={\displaystyle \frac{1}{\sqrt{2}}}(1,i,0)\stackrel{}{\xi }_1={\displaystyle \frac{1}{\sqrt{2}}}(1,i,0)\stackrel{}{\xi }_0=(0,0,1),`$ (2)
$`𝐩`$ is the relative proton–neutron momentum inside the deuteron, $`\widehat{𝐩}=𝐩/|𝐩|`$ is the unit vector in the $`𝐩`$ direction. Here $`S`$ and $`D`$ components are defined as
$`U(p)`$ $`=`$ $`u(p)+v_o(p)e^{i\chi },`$
$`W(p)`$ $`=`$ $`w(p)+v_2(p)e^{i\chi },`$ (3)
where $`u(p)`$ and $`w(p)`$ are $`S`$ and $`D`$ components of the standard deuteron wave function based on the $`NN`$ potentials; $`v_o(p)e^{i\chi }`$ and $`v_2(p)e^{i\chi }`$ are the complex $`6q`$ admixtures to the $`S`$ and $`D`$ components of the standard DWF, respectively.
Using parity conservation, time reversal invariance and the Pauli principle we can write the matrix of $`NN`$ elastic scattering in terms of 5 independent complex amplitudes (when isospin invariance is assumed):
$`M(𝐤^{},𝐤)={\displaystyle \frac{1}{2}}`$ $`((a+b)+(ab)(\stackrel{}{\sigma }_1𝐧)(\stackrel{}{\sigma }_2𝐧)+`$
$`+`$ $`(c+d)(\stackrel{}{\sigma }_1𝐦)(\stackrel{}{\sigma }_2𝐦)+`$
$`+`$ $`(cd)(\stackrel{}{\sigma }_1𝐥)(\stackrel{}{\sigma }_2𝐥)+e((\stackrel{}{\sigma }_1+\stackrel{}{\sigma }_2)𝐧)),`$ (4)
where $`a`$, $`b`$, $`c`$, $`d`$ and $`e`$ are the scattering amplitudes, $`\stackrel{}{\sigma }_1`$ and $`\stackrel{}{\sigma }_2`$ are the Pauli $`2\times 2`$ matrices, k and k are the unit vector in the direction of the incident and scattered particles, respectively, and center-of-mass basis vectors n, m, l are defined as:
$`𝐧={\displaystyle \frac{𝐤^{}\times 𝐤}{|𝐤^{}\times 𝐤|}},𝐥={\displaystyle \frac{\stackrel{}{𝐤}^{}+𝐤}{|𝐤^{}+𝐤|}},𝐦={\displaystyle \frac{𝐤^{}𝐤}{|𝐤^{}𝐤|}}.`$ (5)
However, at a zero angle there are only 3 independent amplitudes and the matrix element (2) can be written as
$`(0)={\displaystyle \frac{1}{2}}\left(A+B(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2)+C(\stackrel{}{\sigma }_1𝐤)(\stackrel{}{\sigma }_2𝐤)\right),`$ (6)
where amplitudes $`A`$, $`B`$ and $`C`$ are related to the amplitudes defined in ref. as follows
$`A=a(0)+b(0),B=c(0)+d(0),C=2d(0).`$ (7)
We consider deuteron breakup reaction in the special kinematics, i.e. with the emission of the spectator proton at zero angle, while the neutron interacts with the proton target and the products of this interaction go along the axis of the reaction.
Using of both (1) and (6) expressions the matrix element of deuteron breakup process in collinear geometry can be written as
$`=`$ $`{\displaystyle \frac{i}{2\sqrt{2}}}{\displaystyle \frac{1}{\sqrt{4\pi }}}\psi _1^+\psi _2^+\left[(U(p)(\stackrel{}{\sigma }\stackrel{}{\xi }){\displaystyle \frac{W(p)}{\sqrt{2}}}(3(\widehat{𝐩}\stackrel{}{\xi })(\stackrel{}{\sigma }\widehat{𝐩})(\stackrel{}{\sigma }\stackrel{}{\xi })))\sigma _y\right]\times `$
$`\times `$ $`\left(A+B(\stackrel{}{\sigma }_1,\stackrel{}{\sigma }_2)+C(\stackrel{}{\sigma }_1,𝐤)(\stackrel{}{\sigma }_2,𝐤)\right)\psi _1^{}\psi _2.`$ (8)
The matrix element of $`dp`$ backward elastic scattering within framework of one nucleon exchange has the following form
$`={\displaystyle \frac{1}{8\pi }}\psi _f^+`$ $`(U^{}(p)(\stackrel{}{\sigma _f}\stackrel{}{\xi }_f^{}){\displaystyle \frac{W^{}(p)}{\sqrt{2}}}(3(\widehat{𝐩}\stackrel{}{\xi }_f^{})(\stackrel{}{\sigma _f}\widehat{𝐩})(\stackrel{}{\sigma _f}\stackrel{}{\xi }_f^{})))\times `$
$`\times `$ $`\left(U(p)(\stackrel{}{\sigma _i}\stackrel{}{\xi }_i){\displaystyle \frac{W(p)}{\sqrt{2}}}(3(\widehat{𝐩}\stackrel{}{\xi }_i)(\stackrel{}{\sigma _i}\widehat{𝐩})(\stackrel{}{\sigma _i}\stackrel{}{\xi }_i))\right)\psi _i`$ (9)
## 3 Six-quark configurations
In this section we consider different models considering the quark (or baryon–baryon) degrees of freedom inside the deuteron.
In the hybrid model of the DWF the $`6q`$ amplitude, arising from the $`s^6`$ configurations of six quarks, must be added to the $`S`$ component of the standard DWF according the following expression
$`U(k)=\sqrt{1\beta ^2}u(k)+\beta v_0(k)e^{i\chi },`$ (10)
where parameter $`\beta `$ and phase $`\chi `$ represent value of $`6q`$ admixture in deuteron and the degree of non–orthogonality between $`np`$ and $`6q`$ components of the DWF, respectively.
The $`6q`$ admixture has the following form
$`v_0(k)=I\sqrt{10}2^{2/3}\left({\displaystyle \frac{2}{1+\sqrt{2}}}\right)^6\left({\displaystyle \frac{2}{3\pi \omega }}\right)^{3/4}e^{k^2/3\omega }`$ (11)
Factor $`I0.332`$ is the overlap factor of color spin–isospin wave fuctions and $`\omega `$ defines the root-mean-square radius of the $`6q`$ configuration $`r^2=5/4\omega `$, $`k`$ is internal momentum of a nucleon in the deuteron defined in the light-cone dynamics -.
The parameters of the $`6q`$ admixture $`r`$, $`\beta `$ and $`\chi `$ were obtained in from the fit of the experimental data on the momentum density of the nucleon in deuteron $`\varphi ^2(k)`$ , tensor analyzing power $`T_{20}`$ and polarization transfer coefficient $`\kappa _0`$ for deuteron inclusive breakup reaction with the emission of the proton at a zero angle using standard DWFs . The results of the fit are given in Table 1 and shown in Fig.1 by the solid, dashed and dotted lines for RSC , Paris and Bonn (version C) DWFs, respectively. One can see the satisfactory description of the experimental data. The probability of the $`6q`$ admixture is found to be $`34\%`$. The relative phase is $`40^0`$ for RSC and Paris DWFs and $`55^0`$ for Bonn C DWF . The radius is $`r0.6`$ fm for all used DWFs. The parameters are comparable with the results obtained in using RSC DWF .
In the nonnucleon degrees of freedom ($`NN^{}`$, $`NN\pi `$ and higher components of the Fock space) were taken into account in the following way
$`\mathrm{\Phi }^2(\alpha ,k_t)=(1\beta ^2)\varphi _{NN}^2(\alpha ,k_t)/(2\alpha (1\alpha ))+\beta ^2G_d(\alpha ,k_t),`$ (12)
where $`\mathrm{\Phi }^2(\alpha ,k_t)`$ is the distribution of constituents in the deuteron;
$`\varphi _{NN}^2(\alpha ,k_t)`$ is the relativized standard DWF; $`G_d(\alpha ,k_t)`$ is the distribution of $`NN^{}`$, $`NN\pi `$ …, or $`6q`$ component in the deuteron. Parameter $`\beta ^2`$ gives the probability of this nonnucleon component. The relativistic form of the DWF $`\varphi _{NN}(\alpha ,k_t)`$ can be written according - as
$`\varphi _{NN}(\alpha ,k_t)=\left({\displaystyle \frac{m_p^2+k_t^2}{4\alpha (1\alpha )}}\right)^{1/4}\varphi (k),`$ (13)
where $`\varphi (k)`$ is the standard DWF (for instance, ) and internal momentum $`k`$ and longitudinal momentum fraction $`\alpha `$ are defined as -
$`k^2`$ $`=`$ $`{\displaystyle \frac{m_p^2+k_t^2}{4\alpha (1\alpha )}}m_p^2,`$
$`\alpha `$ $`=`$ $`{\displaystyle \frac{k_{||}+\sqrt{k_{||}^2+m_p^2}}{2\sqrt{k_{||}^2+m_p^2}}}.`$ (14)
Here $`k_{||}`$ is the longitudinal momentum in infinite momentum frame and $`m_p`$ is the nucleon mass.
The expression for nonnucleon component $`G_d(\alpha ,k_t)`$ is written as
$`G_d(\alpha ,k_t)=b^2/(2\pi )G_1(\alpha )e^{bk_t}`$ (15)
with
$`G_1(\alpha )=\mathrm{\Gamma }(A_2+B_2+2)/(\mathrm{\Gamma }(A_2+1)\mathrm{\Gamma }(B_2+1))\alpha ^{A_2}(1\alpha )^{B_2},`$ (16)
where $`\mathrm{\Gamma }(\mathrm{})`$ denotes the $`\mathrm{\Gamma }`$-function. The parameter $`b`$ is chosen to be 5 GeV/c. We assume, that nonnucleon component (15)-(16) has the relative phase $`\chi `$ with the $`S`$ wave of the standard DWF .
The results of the fit of the experimental data are given in Table 2 and shown in Fig.2 by the solid, dashed and dotted lines for RSC , Paris and Bonn (version C) DWFs, respectively. The probability of the nonnucleon component is found to be also $`3\%`$. The relative phase $`\chi `$ between $`NN`$ and nonnucleon components is $`4060^o`$. The parameters $`A_2`$ and $`B_2`$ are found to be approximately the same for Paris , RSC and Bonn C DWFs. Note, all the used $`NN`$ DWFs provide satisfactory agreement with the existing data, however, the using of RSC DWF gives better description of the polarization transfer coefficient $`\kappa _0`$.
## 4 T-odd polarization effects
Let us define the general spin observable of the third order in terms of Pauli $`2\times 2`$ spin matrices $`\sigma `$ for protons and a set of spin operators $`S_\lambda `$ for the spin 1 particle for both reactions as
$`C_{\alpha ,\lambda ,\beta ,0}={\displaystyle \frac{Tr(\sigma _\alpha ^pS_\lambda ^d^+\sigma _\beta ^p)}{Tr(^+)}},`$ (17)
where indices $`\alpha `$ and $`\lambda `$ refer to the initial proton and deuteron polarization, index $`\beta `$ refers to the final proton, respectively.
We use a righthand coordinate system, defined in accordance with Madison convention . This system is specified by a set of three orthogonal vectors $`\stackrel{}{L}`$, $`\stackrel{}{N}`$ and $`\stackrel{}{S}`$, where $`\stackrel{}{L}`$ is the unit vector along the momenta of the incident particle, $`\stackrel{}{N}`$ is taken to be orthogonal to $`\stackrel{}{L}`$, $`\stackrel{}{S}=\stackrel{}{N}\times \stackrel{}{L}`$.
In this paper we consider $`T`$-odd polarization observables, namely: tensor- vector spin correlations $`C_{N,SL,0,0}`$ due to tensor polarization of the beam and polarization of the initial proton and polarization transfer coefficient $`C_{0,SL,N,0}`$ from tensor polarized deuteron to proton in the $`dppd`$ and $`dpp(0^0)+p(180^0)+n`$ reactions. Note, that such observables must be zero in framework of one nucleon exchange using standard deuteron wave functions, however, they do not vanish with the existing of $`6q`$ admixture in the DWF.
Using the formulas for the matrix elements of the $`p(d,p)pn`$ and $`p(d,p)d`$ reactions (2) and (2), respectively, one can obtain the expression for the polarization transfer coefficient $`C_{0,SL,N,0}`$
$`C_{0,SL,N,0}={\displaystyle \frac{3}{\sqrt{2}}}{\displaystyle \frac{wv_0sin\chi }{u^2+w^2+v_0^2+2uv_0cos\chi }}.`$ (18)
One can see that $`C_{0,SL,N,0}`$ does not depend on the initial energy and is defined only by the interference between $`D`$ wave of the standard DWF and $`6q`$ admixture. The results of the calculations with the use of Paris , RSC and Bonn C DWFs are presented in Fig. 3 a, b and c for two different models of the $`6q`$ admixture: and given by the solid and dashed lines, respectively. These two types of the hybrid DWFs give quite similar behaviour of the $`C_{0,SL,N,0}`$ up to $`k800`$ MeV/c, however, they differ at higher momenta. Both models predict the smooth variation of the $`C_{0,SL,N,0}`$ of about $`1`$ at $`k600`$ MeV/c. The dependence on the used $`NN`$ deuteron wave function occurs only at high $`k`$ of about $`900`$ MeV/c, therefore, the observation of large negative value of $`C_{0,SL,N,0}`$ could indicate that quark degrees of freedom play quite important role in the deuteron at large $`k`$.
Spin correlation parameter $`C_{N,SL,0,0}`$ due to tensor polarization of the beam and polarization of the initial proton for the $`dpp(0^0)+p(180^0)+n`$ process can be written as
$`C_{N,SL,0,0}={\displaystyle \frac{3}{\sqrt{2}}}{\displaystyle \frac{wv_0sin\chi }{u^2+w^2+v_0^2+2uv_0cos\chi }}A_{oonn}(0^o),`$ (19)
where $`A_{oonn}(0^o)`$ is spin correlation of neutron–proton elastic scattering at a zero angle for vertically polarized particles (see notations used in ). Therefore, the behaviour of $`C_{N,SL,0,0}`$ in the $`dpp(0^0)+p(180^0)+n`$ reaction is defined both the DWF and $`np`$ elementary amplitude which is energy dependent. The calculation of $`C_{N,SL,0,0}`$ for the deuteron initial energy of 2.1 GeV and 1.25 GeV using the results of phase-shift analysis performed in are shown in Figs 4 and 5, respectively. One can see that $`C_{N,SL,0,0}`$ is positive at 2.1 GeV up to $`k550`$ MeV/c and negative at 1.25 GeV at $`k300÷400`$ MeV/c. The difference between two models of $`6q`$ admixture shown by the solid and dashed lines in a, b and c figures for Paris , RSC and Bonn C DWFs is not dramatic at both energies.
Spin correlation parameter $`C_{N,SL,0,0}`$ in deuteron–proton backward elastic scattering is given in the following form
$`C_{N,SL,0,0}={\displaystyle \frac{1}{\sqrt{2}}}{\displaystyle \frac{wv_0sin\chi \left((u+v_0cos\chi \sqrt{2}w)^2+v_0^2sin^2\chi \right)}{(u^2+w^2+v_0^2+2uv_0cos\chi )^2}}.`$ (20)
The behaviour of this observable for different types of $`6q`$ admixture in the DWFs is shown in Fig. 6 a, b and c for Paris , RSC and Bonn C DWFs by the solid and dashed lines, respectively. One can see that the spin correlation $`C_{N,SL,0,0}`$ has a small negative value at low $`k`$, then it approaches a minima of $`0.7÷0.8`$ at $`k400`$ MeV/c and afterwards it goes smoothly to a zero for both models of $`6q`$ component. However, the use of DWFs with the $`6q`$ admixture adopted in gives systematically more negative value of spin correlation at internal momenta of $`300`$ MeV/c. The use of different $`NN`$ potentials (see Fig.6 a, b and c, respectively) gives slightly different behaviour of $`C_{N,SL,0,0}`$ for both models of $`6q`$ admixture. Nevertheless, one can conclude that the measurements of spin correlation $`C_{N,SL,0,0}`$ in $`dp`$ backward elastic scattering can help to distinguish between these two models.
Note that non–orthogonality in the deuteron wave function results in the $`T`$-invariance violation, which contradicts the experiment. However, $`NN`$ and $`6q`$ components can be orthogonalized following the procedure described in ref.. Such a procedure only slightly changes the probability of $`6q`$ admixture , but does not affect on the behaviour of the considered observables. For instance, the probability of $`6q`$ component changes from $`2.96\%`$ to $`3.31\%`$ and from $`3.42\%`$ to $`4.17\%`$ for the models and , respectively, in the case of the use of Paris DWF .
The six-quark wave function of the deuteron has been calculated recently not only from $`s^6`$ , but also from $`s^4p^2`$ configurations . Such configurations are orthogonal to the usual $`S`$ and $`D`$ waves in the deuteron. Tensor analyzing power $`T_{20}`$ and polarization transfer coefficient $`\kappa _0`$ in deuteron inclusive breakup at a zero proton emission angle have been qualitatively reproduced at large internal momenta using the results of these calculations . The probability of the $`D`$ wave originated from $`s^4p^2`$ configurations was found to be about $`0.5\%`$ (a small part of $`D`$ wave probability $`6\%`$). The results on the polarization transfer $`C_{0,SL,N,0}`$ and spin correlation $`C_{N,SL,0,0}`$ in $`dp`$ backward elastic scattering using Paris DWF and the $`6q`$ projection on $`NN`$ component from are given in Fig.7 a and b, respectively. The behaviour of these observable differ significantly from the results shown in Figs. 3 and 6. This deviation is due to presence of $`D`$\- wave in six-quark wave function. The results on tensor- vector spin correlation $`C_{SL,N,0,0}`$ in the reaction $`dpp(0^0)+p(180^0)+n`$ at the initial deuteron energy of 2.1 and 1.25 GeV are shown in Fig.8 a and b, respectively. The behaviour is qualitatively the same as shown in Figs. 4 and 5, however, the value of $`C_{SL,N,0,0}`$ at 2.1 GeV and $`k300`$ MeV/c is as twice as much than that in the case of absence of $`D`$ wave.
Note, the one of interesting features of QCD is the possible existing of resonances in the dibaryon system corresponding to six-quark states which are dominantly hidden color, i.e., orthogonal to the usual $`np`$ states. The rich structure in the behaviour of the tensor analyzing power $`T_{20}`$ in $`dp`$ backward elastic scattering can be an indication of such dibaryon resonances .
Of course, the mechanisms additional to ONE can contribute to $`C_{N,SL,0,0}`$ and $`C_{0,SL,N,0}`$. However, the calculations taking into account such mechanisms show that their contribution is small at large internal momenta. Thus, the observation of a large values of $`C_{N,SL,0,0}`$ and $`C_{0,SL,N,0}`$ at momenta higher 600 MeV/c could be rather clear indication of the exotic $`6q`$ configurations.
## 5 Conclusion
We have considered $`T`$-odd observables in deuteron exclusive breakup and $`dp`$ backward elastic scattering, namely, tensor- vector polarization transfer coefficient $`C_{0,SL,N,0}`$ and tensor- vector spin correlation $`C_{N,SL,0,0}`$. These observables, which are associated with the tensor polarization of the deuteron and polarization of the proton, show their sensitivity to the quark degrees of freedom in the deuteron and their spin structure. The calculations give a sizeable effects at large internal momenta, which could be measured with the existing experimental techniques.
Measurements of these observables could be performed at COSY at Zero Degree Facility (ANKE) using internal polarized target with the detection of two charged particles in case of deuteron breakup and with the detection of the fast proton in case of $`dp`$ backward elastic scattering.
Such experiments could be also performed at the Laboratory for High Energies of Joint Institute for Nuclear Researches. The rotation of the primary deuteron spin could be provided by the magnetic field of the beam line upstream of the target or by the special spin-rotating magnet.
Acknowledgments
I’m indebted to Dr. N.B.Ladygina for critically revising of this manuscript and for permanent help and support. Author is grateful to Prof. A.P.Kobushkin for helpful comments and discussions. I also thank I.M.Sitnik for estimation of the bending angle of the magnetic system upstream of the polarized target installed now at the LHE of JINR.
Table 1. Parameters of the $`6q`$ admixture in the hybrid model for different standard DWFs .
| DWF | $`\beta ^2,\%`$ | $`\chi `$ | $`r,fm`$ |
| --- | --- | --- | --- |
| | $`3.42\pm 0.09`$ | $`47.2^0\pm 0.6^0`$ | $`0.578\pm 0.009`$ |
| | $`4.07\pm 0.10`$ | $`40.1^0\pm 0.6^0`$ | $`0.590\pm 0.009`$ |
| | $`2.79\pm 0.09`$ | $`55.1^0\pm 0.7^0`$ | $`0.595\pm 0.010`$ |
Table 2. Parameters of the $`6q`$ admixture for different standard DWFs
.
| DWF | $`\beta ^2,\%`$ | $`A_2`$ | $`B_2`$ | $`\chi `$ |
| --- | --- | --- | --- | --- |
| | $`2.96\pm 0.18`$ | $`10.0^{}`$ | $`20.23\pm 0.42`$ | $`47.0^0\pm 0.6^0`$ |
| | $`3.70\pm 0.43`$ | $`10.0\pm 0.9`$ | $`19.87\pm 2.72`$ | $`40.2^0\pm 0.6^0`$ |
| | $`2.67\pm 0.19`$ | $`10.0^{}`$ | $`19.46\pm 0.48`$ | $`54.6^0\pm 0.7^0`$ |
$``$ \- parameter is fixed.
Figure captions
Fig.1. Momentum density $`\mathrm{\Phi }^2(k)`$ , tensor analyzing power $`T_{20}`$ and polarization transfer coefficient $`\kappa _0`$ (open squares), (full triangles) and (full circles and squares) versus internal momentum $`k`$ in deuteron inclusive breakup with the emission of proton at $`0^0`$. Full, dashed and dotted lines correspond to calculations with hybrid wave function using RSC , Paris and Bonn C DWFs, respectively.
Fig.2. Momentum density $`\mathrm{\Phi }^2(k)`$ , tensor analyzing power $`T_{20}`$ and polarization transfer coefficient $`\kappa _0`$ versus internal momentum $`k`$ in deuteron inclusive breakup with the emission of proton at $`0^0`$. Full, dashed and dotted lines correspond to calculations with the wave function adopted in using RSC , Paris and Bonn C DWFs, respectively. The symbols are the same as in Fig.1.
Fig.3. Tensor-vector polarization transfer coefficient $`C_{0,SL,N,0}`$ in deuteron exclusive breakup in the collinear geometry and $`dp`$ backward elastic scattering using $`6q`$ admixture adopted in and and given by the solid and dashed lines, respectively. The curves in Figs. a, b and c are obtained with the use of Paris , RSC and Bonn C DWFs, respectively.
Fig.4. Tensor-vector spin correlation parameter $`C_{N,SL,0,0}`$ in deuteron exclusive breakup in the collinear geometry at 2.1 GeV of the deuteron initial energy using $`6q`$ admixture adopted in (solid lines) and (dashed lines). The curves in Figs. a, b and c are obtained with the use of Paris , RSC and Bonn C DWFs, respectively.
Fig.5. Tensor-vector spin correlation parameter $`C_{N,SL,0,0}`$ in deuteron exclusive breakup in the collinear geometry at 1.25 GeV of the deuteron initial energy using $`6q`$ admixture adopted in and and given by the solid and dashed lines, respectively. The curves in Figs. a, b and c are obtained with the use of Paris , RSC and Bonn C DWFs, respectively.
Fig.6. Tensor-vector spin correlation parameter $`C_{N,SL,0,0}`$ in deuteron- proton backward elastic scattering using $`6q`$ admixture adopted in and and given by the solid and dashed lines, respectively. The curves in Figs. a, b and c are obtained with the use of Paris , RSC and Bonn C DWFs, respectively.
Fig.7. a) Tensor-vector polarization transfer coefficient $`C_{0,SL,N,0}`$ in deuteron exclusive breakup in the collinear geometry and $`dp`$ backward elastic scattering and b) tensor-vector spin correlation parameter $`C_{N,SL,0,0}`$ in deuteron–proton backward elastic scattering using results of and Paris DWF .
Fig.8. Tensor- vector spin correlation parameter $`C_{N,SL,0,0}`$ in deuteron exclusive breakup in the collinear geometry at a) 2.1 GeV and b) at 1.25 GeV of the deuteron initial energy results of and Paris DWF .
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# Many-Polaron System Confined to a Quantum Dot: Ground-State Energy and Optical Absorption
## I Introduction
The path integral method for indistinguishable particles developed in Refs. seems an adequate tool for the investigation of both the equilibrium and the non-equilibrium properties of interacting quantum many-body systems. The advantages of the path integral representation are particularly manifest when investigating systems with a fixed (few or many) number of particles. As demonstrated in Refs. , the thermodynamical properties of systems with a limited number of particles might deviate substantially from those obtained in the thermodynamical limit.
Electrons confined to a quantum dot can provide a typical example of a system with a fixed number of indistinguishable particles. In recent years, the quantum states and the optical properties of these systems have received considerable attention. Multi-electron states in 3D and 2D quantum dots (without the electron-phonon interaction) have, e. g., been treated in Refs. . It has been shown in Ref. that the interaction between the charge carriers strongly influences the optical spectra of few-particle quantum dots. The bipolaron ground state in quantum dots has been studied in Ref. . Cooperative effects for a gas of polarons and bipolarons were treated in Ref. on the basis of a model introduced by Friedberg and Lee . To the best of our knowledge, the effects of the electron-phonon interaction on the optical spectra of multi-electron quantum dots with a fixed number of electrons have hitherto not been investigated theoretically.
Experimentally observed optical absorption spectra of high-$`T_c`$ cuprates are very promising to reveal new manifestations of the interaction between electrons (holes) and the longitudinal optical (LO) phonons. A possible role of polarons in high-T<sub>c</sub> superconductivity has been analyzed by several authors (see, e. g., Refs. ). Some aspects of the recently observed optical absorption spectra of high-$`T_c`$ superconductors can be interpreted using the polaron theory . It has been shown in Ref. that certain characteristics of the polaron optical absorption at intermediate Fröhlich coupling constant $`\alpha `$ appear in the aforementioned spectra . Bipolaron optical absorption in bulk materials has, e. g., been treated in Refs. . In view of the relatively high concentration of electrons (holes) in high-$`T_c`$ superconductors, the development of an all-coupling and all-concentration theory of the optical response of many-polaron systems is an urgent problem.
## II Ground-state energy
In the present communication, we treat the ground state energy and the optical absorption spectra of a multi-electron (multi-polaron) parabolic quantum dot for arbitrary electron-phonon coupling strength, using a recently developed path integral formalism for the quantum statistical treatment of identical particles and considering the expected quite large polaron coupling constant in those materials. We consider a system consisting of $`N`$ electrons in a parabolic confinement potential characterized by the frequency parameter $`\mathrm{\Omega }_0.`$ These electrons are assumed to interact with each other and with the LO phonons. The electron subsystem is subdivided into two groups ($`N=_\sigma N_\sigma ,`$ where $`N_\sigma `$ is the number of electrons with definite spin projection $`\sigma =\pm 1/2`$).
In the present work the ground state energy of $`N`$ polarons confined to a quantum dot has been determined for the first time. This calculation has been performed within the *extended Jensen-Feynman variational principle* taking into account the symmetry properties of the electrons with respect to permutations. The validity of this extension of the Jensen-Feynman inequality for systems of indistinguishable particles, which is not obvious, and which is of crucial importance for the present work, has been demonstrated in Ref. . We have used an auxiliary model system of particles in a harmonic confinement potential with elastic interparticle interactions as studied in Refs. . The parameters of the auxiliary model system have been treated as variational parameters. We have worked out the variational procedure to obtain those variational parameters at arbitrary temperature for the physical system of $`N`$ polarons confined to a parabolic potential. The calculation has been performed for the case of closed shells (i. e. of a non-degenerate ground state) and for the case of open shells.
We have calculated the ground state energy of a parabolically confined many-polaron system as a function of the total spin of the system $`S=\frac{1}{2}\left|N_{+1/2}N_{1/2}\right|`$. For relatively small polaron coupling constant and all values of $`\eta \epsilon _{\mathrm{}}/\epsilon _0`$ ($`\epsilon _{\mathrm{}}`$ and $`\epsilon _0`$ are the high-frequency and the static dielectric constants, respectively), except $`\eta 1`$, the electrons tend to maximally fill up open shells with the *same* spin projection (cf. the first Hund’s rule). In this domain of the parameters $`\alpha `$ and $`\eta `$, the ground state is obtained when at least one of the numbers $`N_{\pm 1/2}`$ equals the number of states $`M_n`$ in a set of shells labelled by the index $`k`$ from 0 to $`n`$ (see Ref. ):
$$M_n=\underset{k=0}{\overset{n}{}}g_k=\frac{\left(n+1\right)\left(n+2\right)\left(n+3\right)}{6}(n=0,1,2,\mathrm{}),$$
(1)
where $`g_k\frac{1}{2}\left(k+1\right)\left(k+2\right)`$ is the degeneracy of the $`k`$th energy level of a three-dimensional oscillator. For $`n=0,1,2,3\mathrm{}`$ the values of the number $`M_n`$ are $`1,4,10,20\mathrm{}`$.
This filling scheme is *broken* for strongly polar substances ($`\eta 1`$), where we find that the lowest variational energy corresponds to the minimal possible total spin ($`S=0`$ for $`N`$ even, and $`S=\frac{1}{2}`$ for $`N`$ odd), or for $`\alpha 1`$ (except the region $`\eta 1`$), where the lowest variational energy corresponds to the maximal possible total spin ($`S=N/2`$).
In Table 1, the variational ground state energy per particle $`E_0/N`$ is shown for $`N=20`$ polarons in a quantum dot with confinement frequency parameter $`\mathrm{\Omega }_0=0.5`$ at $`\alpha =5`$, for different values of $`\eta .`$
Table 1. The variational ground state energy per particle $`E_0/N`$ (in units of the LO phonon energy $`\mathrm{}\omega _{\mathrm{LO}}`$) for different $`\eta `$ for $`N=20`$ polarons in a quantum dot with the confinement frequency parameter $`\mathrm{\Omega }_0=0.5`$ (in units of $`\omega _{\mathrm{LO}})`$ at $`\alpha =5.`$
| $`\eta `$ | $`E_0/N\left(\begin{array}{c}N_{+\frac{1}{2}}=20,\\ N_{\frac{1}{2}}=0\end{array}\right)`$ | $`E_0/N\left(\begin{array}{c}N_{+\frac{1}{2}}=10,\\ N_{\frac{1}{2}}=10\end{array}\right)`$ |
| --- | --- | --- |
| 0.01 | $`\mathrm{3.\hspace{0.17em}9404}`$ | $`4.1078`$ |
| 0.081 | $`2.0867`$ | $`2.0873`$ |
| 0.082 | $`2.0626`$ | $`2.0620`$ |
| 0.4 | $``$$`5.8487`$ | $``$$`6.0482`$ |
As follows from this table, a ferromagnetic-to-nonmagnetic transition takes place at a value of $`\eta `$ in the interval between 0.081 and 0.082. It will be analyzed below how this transition influences the optical absorption spectra.
## III Optical absorption spectra
In order to investigate the optical properties of the confined many-polaron system, we use the formalism developed in Refs. . An alternative derivation can be found in Ref. . Within this technique, the absorption spectrum for a many-polaron system in a parabolic confinement potential is given by the expression
$$\mathrm{\Gamma }\left(\omega \right)\frac{\text{Im}\chi \left(\omega \right)}{\left[\omega \mathrm{\Omega }_0^2/\omega \text{Re}\chi \left(\omega \right)\right]^2+\left[\text{Im}\chi \left(\omega \right)\right]^2},$$
(2)
where the memory function is
$`\chi \left(\omega \right)={\displaystyle \underset{𝐪}{}}{\displaystyle \frac{2\left|V_𝐪\right|^2q^2}{3N\mathrm{}\omega }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}\left(e^{i\omega t}1\right)\text{Im}\left[T_{\omega _{\mathrm{LO}}}^{}\left(t\right)\rho _𝐪\left(t\right)\rho _𝐪\left(0\right)_{S_M}\right]𝑑t.`$ (3)
The function $`T_\omega \left(t\right)=\mathrm{cos}\left[\omega \left(ti\beta /2\right)\right]/\mathrm{sinh}\left(\beta \omega /2\right)`$ with $`\beta =\mathrm{}/k_BT`$ describes the phonon dynamics. The time-dependent correlation function $`\rho _𝐪\left(t\right)\rho _𝐪\left(0\right)_{S_M}`$ is the path integral average for an auxiliary model action functional $`S_M`$ of $`N`$ identical electrons and $`N_B`$ identical fictitious particles (simulating the influence of the LO phonon bath). As emphasized above, the parameters of the auxiliary model system are determined using the variational procedure for the physical system of $`N`$ confined polarons under consideration. Note that, formally, the term $`\mathrm{\Omega }_0^2/\omega `$ plays a role similar to the cyclotron frequency $`\omega _c`$ in the theory of the cyclotron resonance of polarons . In the zero temperature limit, we have obtained the following analytic expression for the memory function:
$`\chi \left(\omega \right)={\displaystyle \frac{2\alpha m^{}}{3\pi N\omega }}\left({\displaystyle \frac{\omega _{\mathrm{LO}}}{A}}\right)^{3/2}{\displaystyle \underset{p_1=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p_2=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p_3=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\right)^{p_3}}{p_1!p_2!p_3!}}\left({\displaystyle \frac{a_1^2}{N\mathrm{\Omega }_1A}}\right)^{p_1}\left({\displaystyle \frac{a_2^2}{N\mathrm{\Omega }_2A}}\right)^{p_2}\left({\displaystyle \frac{1}{NwA}}\right)^{p_3}`$ (4)
$`\times `$ $`\{{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\sigma }{}}N_{m,\sigma }(1N_{n,\sigma })[{\displaystyle \frac{1}{\omega \omega _{\mathrm{LO}}\left[p_1\mathrm{\Omega }_1+p_2\mathrm{\Omega }_2+\left(p_3m+n\right)w\right]+\mathrm{i}\epsilon }}`$ (5)
$``$ $`{\displaystyle \frac{1}{\omega +\omega _{\mathrm{LO}}+p_1\mathrm{\Omega }_1+p_2\mathrm{\Omega }_2+\left(p_3m+n\right)w+\mathrm{i}\epsilon }}+2{\displaystyle \frac{𝒫}{\omega _{\mathrm{LO}}+p_1\mathrm{\Omega }_1+p_2\mathrm{\Omega }_2+\left(p_3m+n\right)w}}]`$ (6)
$`\times `$ $`{\displaystyle \underset{l=0}{\overset{m}{}}}{\displaystyle \underset{k=nm+l}{\overset{n}{}}}{\displaystyle \frac{\left(1\right)^{nm+l+k}\mathrm{\Gamma }\left(p_1+p_2+p_3+k+l+\frac{3}{2}\right)}{\left(k\right)!\left(l\right)!}}`$ (7)
$`\times `$ $`\left({\displaystyle \frac{1}{wA}}\right)^{l+k}\left(\begin{array}{c}n+2\\ nk\end{array}\right)\left(\begin{array}{c}2k\\ kln+m\end{array}\right)+[{\displaystyle \frac{1}{\omega \omega _{\mathrm{LO}}\left(p_1\mathrm{\Omega }_1+p_2\mathrm{\Omega }_2+p_3w\right)+\mathrm{i}\epsilon }}`$ (12)
$``$ $`{\displaystyle \frac{1}{\omega +\omega _{\mathrm{LO}}+p_1\mathrm{\Omega }_1+p_2\mathrm{\Omega }_2+p_3w+\mathrm{i}\epsilon }}+2{\displaystyle \frac{𝒫}{\omega _{\mathrm{LO}}+p_1\mathrm{\Omega }_1+p_2\mathrm{\Omega }_2+p_3w}}]{\displaystyle }_{m=0}^{\mathrm{}}{\displaystyle }_{n=0}^{\mathrm{}}{\displaystyle }_{\sigma ,\sigma ^{}}N_{m,\sigma }N_{n,\sigma ^{}}`$ (13)
$`\times `$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \underset{l=0}{\overset{m}{}}}{\displaystyle \frac{\left(1\right)^{k+l}\mathrm{\Gamma }\left(p_1+p_2+p_3+k+l+\frac{3}{2}\right)}{k!l!}}\left({\displaystyle \frac{1}{wA}}\right)^{k+l}\left(\begin{array}{c}n+2\\ nk\end{array}\right)\left(\begin{array}{c}m+2\\ ml\end{array}\right)\},`$ (18)
where $`\epsilon +0`$, $`𝒫`$ denotes the principal value, $`A\left[_{i=1}^2a_i^2/\mathrm{\Omega }_i+\left(N1\right)/w\right]/N.`$ $`\mathrm{\Omega }_1,`$ $`\mathrm{\Omega }_2`$ and $`w`$ are the eigenfrequencies of the model system ($`\mathrm{\Omega }_1`$ is the frequency of the relative motion of the center of mass of the electrons with respect to the center of mass of the fictitious particles; $`\mathrm{\Omega }_2`$ is the frequency related to the center of mass of the entire model system; $`w`$ is the frequency of the internal degrees of freedom), $`a_1`$ and $`a_2`$ are the coefficients of a canonical transformation which diagonalizes the model Lagrangian, and $`N_{n,\sigma }`$ is the number of electrons with spin projection $`\sigma `$ in the $`n`$th single-particle level (shell).
## IV Discussion of results
In Fig. 1, the optical absorption spectra of a many-polaron system in a quantum dot with parabolic confinement potential are plotted for $`\mathrm{\Omega }_0=0.5`$ (all frequencies are measured in units of $`\omega _{\mathrm{LO}}`$), $`\alpha =1`$ and for $`N=5,10,14`$ (panels $`a,b,c,`$ respectively). Due to the confinement in all three dimensions, the electron motion is fully quantized. Hence, when a photon is absorbed, the electron recoil can be transferred only by discrete quanta. As a result, the absorption spectrum consists of a series of $`\delta `$-like peaks as distinct from the absorption spectrum of a bulk polaron. In this and subsequent figures, the height of each peak represents its intensity.
Fig. 2 shows, for reference, the optical absorption spectrum for a *single* polaron confined to a quantum dot with $`\mathrm{\Omega }_0=0.2`$ for $`\alpha =5,`$ and reveals several essential elements. Following the nomenclature of Refs. we distinguish a central (zero-phonon) peak, peaks due to transitions to the relaxed excited state (RES), and peaks associated with transitions to the Franck-Condon state (FC). Also it is worth mentioning that there is the “one-phonon shoulder” in the optical absorption spectrum with threshold at $`\omega 1`$. The envelope of the absorption spectrum for a polaron in a quantum dot is very similar to the bulk polaron absorption spectra for the same $`\alpha `$, as first obtained in Ref. . It follows that the obtained optical absorption spectra for a polaron confined to a quantum dot are consistent with those for a bulk polaron .
As follows from Figs. 1 and 2, the $`\delta `$-like “central peak” for a system of polarons, confined to a quantum dot with parabolic confinement potential, is positioned at $`\omega \mathrm{\Omega }_0.`$ Without the electron-phonon interaction, the “central peak” for a system of electrons, confined to a quantum dot with parabolic confinement, is exactly at the confinement parameter: $`\omega =\mathrm{\Omega }_0.`$
Fig. 3 shows the evolution of the absorption spectrum as a function of the confinement parameter $`\mathrm{\Omega }_0,`$ for $`\mathrm{\Omega }_0=0.5,`$ $`\mathrm{\Omega }_0=0.8,`$ $`\mathrm{\Omega }_0=1`$ and $`\mathrm{\Omega }_0=1.2`$. This calculation is performed for $`N=4,`$ $`\alpha =3,`$ $`\eta =0.3.`$ Near $`\mathrm{\Omega }_0=1,`$ the zero-phonon and the one-phonon peaks are of comparable oscillator strength. A formal analogy (particularly manifest for sufficiently small $`\alpha )`$ exists between the resonance condition $`\mathrm{\Omega }_0=1`$ discussed here and the wee-established magnetophonon resonance around $`\omega _c=1`$ (Refs. ). We suggest the designation “confinement-phonon resonance”.
Fig. 4 illustrates the ferromagnetic-to-nonmagnetic transition induced by the increase of the electron-phonon interaction for a many-polaron system confined to a quantum dot. In this figure, optical absorption spectra are plotted for $`N=20`$ polarons in a quantum dot with confinement parameter $`\mathrm{\Omega }_0=0.5,`$ and with electron-phonon coupling constant $`\alpha =5.`$ The parameter $`\eta `$ varies from $`\eta =0.01`$ to $`\eta =0.4,`$ and we find that, as a consequence, the total spin of the ground state jumps from $`S=0`$ to $`S=10`$. As seen from Table 1, for $`\eta 0.081,`$ the ground state with $`S=0`$ is energetically favorable. The optical absorption spectra for this case are shown on panels *a* ($`\eta =0.01`$) and *b* ($`\eta =0.081`$) of Fig. 4. A series of pronounced peaks clearly related to the internal polaron excitations (cf. RES of a single polaron in Fig. 2) is seen in the spectrum for the case of a strongly polar substance ($`\eta =0.01`$). With increasing $`\eta ,`$ the relative intensity of these “RES” (many-polaron) peaks with respect to that of the zero-phonon peak decreases as is seen from Fig. 4. This effect is due to weakening of the electron-phonon interaction with increasing parameter $`\eta .`$
In the specific case under consideration ($`N=20`$), when $`\eta `$ varies from $`\eta =0.081`$ to $`\eta =0.082`$, the total spin abruptly changes from $`S=0`$ \[$`N_{1/2}=10`$, $`N_{1/2}=10`$\] to its maximal value $`S=10`$ \[$`N_{1/2}=20`$ (0), $`N_{1/2}=0`$ (20)\], i. e., the electrons become fully spin-polarized. Such a jump in the magnitude of the total spin can also be realized by varying the confinement parameter at fixed $`\alpha `$ and $`\eta .`$ Comparing the optical absorption spectra in panels *b* and *c* of Fig. 4, one observes that the optical absorption spectrum as a whole shows an abrupt change at the ferromagnetic-to-nonmagnetic transition, though the shift of the most intense phonon sidebands towards lower frequencies (with increasing $`\eta `$) continuously proceeds when $`\eta `$ passes through the value 0.081. Note also the shift to lower frequencies of the most intense phonon sidebands (again related to the “many-polaron RES”) in Fig. 4$`b`$ compared to Fig. 4$`a`$.
## V Moments of the optical absorption spectra
We have calculated (both for closed-shell and open-shell systems) the first frequency moment $`M_1`$ of the optical absorption spectrum for a $`N`$-polaron system confined to a quantum dot with parabolic confinement potential
$`\omega {\displaystyle \frac{\mu _1}{\mu _0}}={\displaystyle \frac{_0^{\mathrm{}}\omega \mathrm{\Gamma }\left(\omega \right)𝑑\omega }{_0^{\mathrm{}}\mathrm{\Gamma }\left(\omega \right)𝑑\omega }}`$ (19)
with $`\mathrm{\Gamma }\left(\omega \right)`$ the optical absorption coefficient, and the parameter $`\sigma \sqrt{\left(\omega \omega \right)^2}`$, where
$`\left(\omega \omega \right)^2=\omega ^2\omega ^2`$ (20)
is the normalized second frequency moment of the optical absorption spectrum.
The results for the first frequency moment of the optical ansorption spectrua for $`N`$ polarons in a quantum dot with parabolic confinement are shown in Fig. 5$`a`$. The first frequency moment is plotted as a function of the effective density $`N/𝒱`$, where $`𝒱=\left(4\pi /3\right)\left(\mathrm{}/2m\mathrm{\Omega }_0\right)^{3/2}.`$ The points corresponding to definite $`N`$ are shown by filled circles. The number of particles ranges from $`N=1`$ to $`N=20.`$ The points indicated by arrows are related to the closed-shell systems.
The function $`\omega \left(N/𝒱\right)`$ turns out to be non-monotonous. There is a maximum of the first frequency moment $`\omega `$ at $`N=2`$. The first frequency moment has a minimum when the number of particles takes the value $`N=14`$ corresponding to the closed-shell system ($`N_{+1/2}=10,`$ $`N_{1/2}=4`$).
By its general trend, the first frequency moment of the optical absorption spectrum as a function of the effective density strikingly resembles the first moment of optical conductivity spectrum in Nd$`_{2\text{x}}`$Ce$`_\text{x}`$CuO<sub>4-y</sub> recently observed experimentally (Fig. 5$`b`$ and Ref. ). Note that our theory describes a 3D confined system, while the experiment relates to a quasi-2D translationally invariant system. Therefore, the theory can be expected to reveal only a qualitative trend of the first frequency moment as a function of concentration in comparison with experiment.
## VI Conclusions
We have presented the ground state energy and the optical absorption spectra calculated for a system of $`N`$ polarons in a parabolic confinement potential for any strength of the polaron coupling. Path integral formalism for the quantum statistical physics of indistinguishable particles has allowed to develop the variational procedure for the ground state energy for a finite number of polarons. For the first time, the ground state energy and the optical absorption spectra have been analyzed for $`N`$ electrons (holes) interacting with each other and with the longitudinal optical (LO) phonons at an arbitrary electron-phonon coupling strength $`\alpha `$ to a parabolic confinement potential. A new type of transition for $`N`$ polarons confined in a parabolic potential (ferromagnetic-to-nonmagnetic transition) is found between states with different total spin, which is related to the competition between the *Coulomb repulsion* and the *phonon-mediated attraction* between the electrons. For relatively weak polaron coupling constant, the electrons are shown to maximally fill up shells with the same spin projection (cf. the first Hund’s rule). This filling scheme is demonstrated to be broken for strongly polar substances ($`\eta 1`$), where we find that the lowest variational energy corresponds to the minimal possible total spin. The present analysis has been executed for closed-shell and open-shell systems.
The optical absorption spectra have been calculated here using the memory function approach as applied to path integrals for a many-polaron system confined to a quantum dot with different number of polarons. The dependence of the optical absorption spectra on the confinement parameter $`\mathrm{\Omega }_0`$ reveals a resonant behavior for $`\mathrm{\Omega }_01`$, especially if $`\alpha `$ is small. The polaron RES are seen to also influence the optical absorption spectra of $`N`$ confined polarons.
We have analysed also the first frequency moment of the optical absorption spectrum for a $`N`$-polaron system in a parabolic quantum dot for both closed-shell and open-shell systems.
###### Acknowledgements.
This work has been supported by the BOF NOI (UA-UIA), GOA BOF UA 2000, IUAP, FWO-V. projects G.0287.95, 9.0193.97, and the W.O.G. WO.025.99N (Belgium). We are indebted to P. Calvani for fruitful interactions and communication of experimental data. Discussions with L. F. Lemmens on the many body aspects and a discussion with K. H. Michel on phase transitions are gratefully acknowledged.
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# Mirror Dark Matter
## I Introduction
Dark matter constitutes the bulk of the matter in the universe and a proper understanding of the nature of the new particle that plays this role has profound implications not only for cosmology but also for particle physics beyond the standard model. It is therefore not surprising that one of the major areas of research in both particle physics and cosmology continues to be the physics of dark matter.
Apart from the simple requirement that the right particle physics candidate must have properties that allow it to have the requisite relic density and mass to dominate the mass content of the universe, it should be required to provide a satisfactory resolution of three puzzles of dark matter physics: (i) why is it that the contribution of baryons to the mass density ($`\mathrm{\Omega }`$) of the universe is almost of the same order as the contribution of the dark matter to it ? (ii) how does one understand the dark objects with mass $`0.5M_{}`$ observed in the MACHO experiment, which are supposed to constitute up to 20% of the mass of the halo of the Milky way galaxy and presumably be connected to the dark constituent that contributes to $`\mathrm{\Omega }`$ ? and, finally (iii) what explains the density profile of dark matter in galactic halos; in particular, the evidence that the core densities of galactic halos remain constant as the radius goes to zero.
There are many particle physics candidates for dark constituent of the universe. Generally speaking, the prime consideration that leads to such candidates is that they yield the right order of magnitude for the relic density and mass necessary to get the desired $`\mathrm{\Omega }_{DM}0.21`$. This is, of course, the minimal criterion. It requires that the annihilation cross section of the particles must be in a very specific range correlated with their mass. The most widely discussed candidates are the lightest supersymmetric particle (LSP) and the Peccei-Quinn particle, the axion. The first one is expected to have a mass in the range of 100 GeV whereas the mass of the second would be in the range of $`10^6`$ eV. The present consensus seems to be that value of $`\mathrm{\Omega }_{CDM}`$ is around $`0.20.3`$, with $`\mathrm{\Omega }_\mathrm{\Lambda }`$ making up the rest of the energy density of the universe at the moment. Compare these values with $`\mathrm{\Omega }_B0.05`$. The CDM contribution and the baryon contribution to $`\mathrm{\Omega }`$ are roughly of same order. On the other hand, the nucleon mass is very different from the masses of either the axion or the SUSY LSP. So to understand within the SUSY or the axion models why $`\mathrm{\Omega }_B\mathrm{\Omega }_{DM}`$, one needs to work in a special range for the parameters of the theory. In either of these pictures, the MACHO observations must have a separate explanation.
Furthermore, in recent years it has been emphasized that the LSP and the axion may also have difficulty in explaining the observed core density behaviour of dwarf speroidal galaxies which are known to be dark matter dominated. The point here is that both the axions and neutralinos, being collisionless and nonrelativistic, accumulate at the core of galactic halos, leading to a core density $`\rho (R)`$ which goes like $`R^2`$ rather than a constant which seems to fit data better. We will refer to this as the core density puzzle.
Spergel and Steinhardt have recently revived an old idea, that dark matter may be strongly self interacting, to resolve this puzzle. They argue that, for the right range of the parameters of the particle, it may lead to a halo core which is much less dense and hence in better agreement with observations. Specifically they note that if the dark matter particle has mean free collision path of about a kpc to a Mpc, then the core on this scale cannot “keep on accumulating” dark matter particles, since these will now scatter and “diffuse out”. For typical dark matter particle densities of order of one particle per cm<sup>3</sup>, this requires a cross section for scattering of $`\sigma 10^{21}10^{24}`$ cm<sup>2</sup>. Furthermore the properties of the dark matter particle must be such that it must not allow for dissipation of the thermal energy via emission of light particles; otherwise,the galactic core would cool and lead to an increase in core density. If these considerations stand the test of time, a theoretical challenge would be to look for alternative dark matter candidates (different from the popular ones described above) and the associated scenarios for physics beyond the standard model.
A class of models known as mirror universe models have recently been discussed. These are motivated theoretically by string theory and experimentally by neutrino physics. They predict the existence of a mirror sector of the universe with matter and force content identical to the familiar sector (prior to symmetry breaking). Symmetry breaking might either keep the mirror symmetry exact or break it. This leads to two classes of mirror models: the symmetric mirror model, where all masses and forces in the two sectors remain the same after symmetry breaking and the asymmetric mirror model where the masses in the mirror sector are larger than those in the familiar sector. The mirror particles interact with the mirror photon and not the familiar photon so that they remain dark to our observations. Since the the lightest particles of the mirror sector (other than the neutrinos), the mirror proton and the mirror electron (like those in the familiar sector) are stable and will have abundances similar to the familiar protons and electrons, the proton being heavier could certainly qualify as a dark matter candidate. We will show in the next section they can indeed play this role.
It has been shown that, consistent with the cosmological constraints of the mirror universe theory, the mirror baryons have the desired relic density to play the role of dark matter of the universe. The additional neutrinos of the mirror sector are the sterile neutrinos that appear to be needed in order to have a simultaneous understanding of the three different observed neutrino oscillations i.e. solar, atmospheric and the LSND observations. In fact, one view of neutrino oscillation explanations of these phenomena fixes the ratio of familiar particle mass to mirror particle mass thereby narrowing down the freedom in mirror sector parameters. If indeed sterile neutrinos turn out to be required, mirror universe theory is one of the few models where they appear naturally with masses in the desired range. If we denote the mass ratio $`m_p^{}/m_p=\zeta `$, then a value of $`\zeta 1030`$ is required to explain the neutrino puzzles. What is more interesting is that for the same parameter range required to solve the neutrino puzzles, mirror matter can also provide an explanation of the microlensing observations\- in particular why the observed MACHOs have a mass very near the solar mass and are still dark.
## II Asymmetric mirror model in brief
Let us start with a brief overview of the mirror matter models . The basic idea of the model is extremely simple: duplicate the standard model or any extension of it in the gauge symmetric Lagrangian and allow for the possibility that symmetry breaking may be different in the two sectors. (See table below) There is an exact mirror symmetry connecting the Lagrangians (prior to symmetry breaking) describing physics in each sector. Clearly the $`W^{}s,\gamma ^{}s`$ etc in each sector differ from those in the other as do the quarks and leptons. When the symmetry breaking scale is different in the two sectors, we will call this the asymmetric mirror model. The QCD scale being an independent scale in the theory could be arbitrary. We will allow both the weak scale and the QCD scale of the mirror sector to differ from those of the familiar sector and assume the same common ratio $`\zeta `$ for both scales i.e. $`<H^{}>/<H>=\mathrm{\Lambda }^{}/\mathrm{\Lambda }\zeta `$.
| $`u,d,e,\nu _e`$ | $``$ | $`u^{},d^{},e^{},\nu _e^{}`$ |
| --- | --- | --- |
| $`W,Z,\gamma ,G`$ | $``$ | $`W^{},Z^{},\gamma ^{},G^{}`$ |
| $`\varphi ,\nu _R,\mathrm{}`$ | $``$ | $`\varphi ^{},\nu _R^{},\mathrm{}`$ |
| | $``$Gravity$``$ | |
It is assumed that the two sectors in the universe are connected by only gravitational interactions. It was shown in that gravity induces nonrenormalizable operators that generate mixings between the familiar and the mirror neutrinos. This is one of the ingredients in the resolution of neutrino puzzles. To get an idea of how this works, note that the lepton operators induced due to nonperturbative gravitational effects have the form $`\frac{LHLH}{M_{Pl}}`$, $`\frac{LHL^{}H^{}}{M_{Pl}}`$ and $`\frac{L^{}H^{}L^{}H^{}}{M_{Pl}}`$. After spontaneous breakdown (i.e. $`<H>=v`$ and $`<H^{}>=v^{}`$, we get for the mass matrix mixing the first generation neutrinos from each sector to have the form (in the basis $`\nu _e,\nu _e^{}`$):
$`M={\displaystyle \frac{v^2}{M_{Pl}}}\left(\begin{array}{cc}1& \zeta \\ \zeta & \lambda \zeta ^2\end{array}\right)`$ (3)
where $`\zeta v^{}/v`$ defined above. To solve the solar neutrino puzzle via small angle MSW solution, we will choose $`\lambda 1`$ and $`\zeta 1030`$. This gives the sterile neutrino mass of order $`10^3`$ eV or so, choosing $`M_{Pl}10^{18}`$ GeV, which is also in the right range to explain the solar neutrino puzzle. Note that these are meant to indicate that the model leads to numbers in the right ball park. Emboldened by this result, we will consider the asymmetric version and look at its cosmological implications.
In discussing cosmology, we first note that both sectors of the universe will evolve according to the rules of the usual big bang model except that the cosmic soups in the two sectors may have different temperature. In fact the constraints of big bang nucleosynthesis require that the post inflation reheat temperature in the mirror sector $`T_R^{}`$ be slightly lower than that in the familiar sector $`T_R`$ (define $`\beta T_R^{}/T_R`$) so that the contribution of the light mirror particles such as $`\nu ^{},\gamma ^{}`$ etc. to nucleosynthesis is not too important. This is called asymmetric inflation and can be implemented in different ways. We will see that, if we want the mirror nucleons to play the role of the dark matter, we will need a definite value of $`\beta `$ depening on the choice of $`\zeta `$. This in turn will help us to predict a value for the equivalent extra neutrinos at the BBN epoch (i.e. $`\delta N_\nu `$).
Before detailed discussion, let us first note the impact of the asymmetry on physical parameters and processes. First it implies that $`m_i\zeta m_i`$ with $`i=n,p,e,W,Z`$. This has important implications for the nuclear and atomic physics of the mirror sector . For instance, the binding energy of mirror hydrogen is $`\zeta `$ times larger so that the recombination in the mirror sector takes place much earlier than in the visible sector. With $`\beta T_R^{}/T_R`$ as above, the mirror recombination occurs when the temperature of the familiar sector is $`\zeta /\beta T_r`$ where $`T_r`$ is the recombination temperature in the familiar sector. The mirror sector recombination takes place before familiar sector recombination; this means that density inhomogeneities in the mirror sector begin to grow earlier and familiar matter can fall into them later as in typical cold dark matter scenarios.
## III Mirror nucleon as dark matter
One can also compute the contribution of mirror baryons to the mass density of the universe as follows:
$`{\displaystyle \frac{\mathrm{\Omega }_B^{}}{\mathrm{\Omega }_B}}\beta ^3\zeta `$ (4)
Here we have assumed that baryon to photon ratio in the familiar and the mirror sectors are the same as would be expected since the dynamics are same in both sectors due to mirror symmetry. Eq. (2) implies that both the baryonic and the mirror baryon contribution to $`\mathrm{\Omega }`$ are roughly of the same order, as observed. This provides a resolution of the first conceptual puzzle. Furthermore if we take $`\mathrm{\Omega }_B0.05`$, then $`\mathrm{\Omega }_B^{}0.2`$ would require that $`\beta =(4/\zeta )^{1/3}`$. From this one can calculate the effective $`\delta N_\nu `$ using the following formula:
$`\delta N_\nu =3\beta ^4+{\displaystyle \frac{4}{7}}\beta ^4({\displaystyle \frac{11}{4}})^{4/3}`$ (5)
where the last factor $`(11/4)^{4/3}`$ is due to the reheating of the mirror photon gas subsequent to mirror $`e^+^{}e^{^{}}`$ annihilation. For $`\zeta =20`$, this implies $`\delta N_\nu 0.6`$ and it scales with $`\zeta `$ as $`\zeta ^{4/3}`$. Thus in principle the idea that mirror baryons are dark matter could be tested by more accurate measurements of primordial He<sup>4</sup>, Deuterium and Li<sup>7</sup> abundances which determine $`\delta N_\nu `$.
Clearly to satisfy the inflationary constraint of $`\mathrm{\Omega }_{TOT}=1`$, we need $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$. These kinds of numbers for cold dark matter density apparently emerge from current type I supernovae observations. It is interesting to note that if one were to require that $`\mathrm{\Omega }_{CDM}=1`$, the mirror model would require that $`\zeta `$ be much larger (more than 100) which would then create difficulties in understanding both the neutrino data and the microlensing anomalies. Thus mirror baryons seem to have just the right properties to be the cold dark matter of the universe.
## IV Explaining the microlensing anomaly
Let us now turn to show how the mirror model accounts for the microlensing observations. We start with the four equations of stellar structure:
$`dP/dr=G\rho (r)M(r)/r^2`$ (6)
$`dM(r)/dr=4\pi r^2\rho (r)`$ (7)
$`L(r)/4\pi r^2=(16/3)\sigma _{SB}(T^3/\rho \kappa )dT/dr`$ (8)
$`dL/dr=4\pi r^2ϵ(r)\rho (r)`$ (9)
where $`\kappa (r)`$ is the opacity (cross section per unit mass) at radius $`r`$, $`\sigma _{SB}`$ the Stefan-Boltzmann constant, L(r) the luminosity at radius $`r`$, and $`ϵ(r)`$, the rate of energy generation per unit mass at radius $`r`$. We will need three terms in the equation of state (below) taken one or two at a time:
$`P=(\rho /m)kT+(4\sigma _{SB}/3c)T^4+(h^2/2m_e)(3/8\pi )^2(\rho /m)^{5/3}`$ (10)
where the three terms represent gas pressure, radiation pressure, and (non-relativistic) degenerate electron pressure. $`m`$ is the nucleon mass, $`m_e`$ that of the electron. We have neglected such niceties as keeping track of how many objects there are for each $`m`$ of gas (2 for H, 3/4 for He, etc)
We will make standard, illuminating if crude, approximations in order to understand the $`\zeta `$ behavior of the solutions to the above equations. First we write
$`P=\rho GM/R,\rho =3M/4\pi R^3`$ (11)
where $`P`$ and $`\rho `$ are roughly core averages. Here $`M`$ and $`R`$ are mass and radius of the star.This gives the useful relation
$`P=(4\pi /3)^{1/3}GM^{2/3}\rho ^{4/3}`$ (12)
To find the maximum mass of a (main sequence) mirror star, which is of interest to us, we note that as the mass of the star gets bigger, the core temperature rises. Therefore, of the three terms in the expression for the pressure in the Equation 10, we expect $`P_g`$ and $`P_r`$ to dominate. Following Phillips, we parameterize them as fractions of the total pressure $`P`$ as below:
$`P_g=\beta P,P_r=(1\beta )P`$ (13)
We eliminate $`T`$ and solve for P from this parameterization to obtain
$`\beta P=[(\rho k/m)^4(\beta ^11)/(4\sigma _{SB}/3c)]^{1/3}`$ (14)
Using Equation (12) again then gives
$`M_{max}[(1\beta )c/\sigma _{SB}]^{1/2}G^{3/2}(k/m)^2/\beta ^2`$ (15)
As $`\beta `$ approaches $`1`$, the energy density is increasingly dominated by photons (relativistic particles) and stars become unstable. Taking a cutoff around $`\beta 1/2`$ gives a maximum stellar mass around $`70M_{}`$. Thus the range for stars is roughly $`0.07M_{}`$ to $`70M_{}`$. From Equation (15) one sees, in the approximation that instability sets in at the same $`\beta `$ independent of $`\zeta `$, that $`M_{max}`$ varies as $`\zeta ^2`$. For $`\zeta =15`$, we get for the maximum mass of the mirror star $`0.5M_{}`$, which is of the same order of mass as the MACHO microlensing events. Our model therefore provides a resolution of the microlensing anomaly that avoids the strong constraints of Freese et al for familiar sector white dwarfs.
We want to point out here that we do not expect all of mirror dark matter to condense to form mirror stars. Instead, we would expect it to be in the form of a mixture of mirror dust and mirror stars. In this connection, it has been noted recently that current upper limits to scattering optical depths for Thomson scattering in early universe suggests that compact objects of any kind cannot be the main dark matter constituent. This would also suggest a mixed picture of the kind mentioned.
## V Self interaction of mirror matter and halo core density problem
As noted in the introduction, there appear to be indications from the core density profile of dwarf and low surface brightness galaxies that the dark matter may need to be endowed with a significant self interaction. According to the analysis of Spergel and Stenhardt, the self interaction must be such that the collision cross section of dark matter particles should be of order $`10^{23}(m/GeV)`$ cm<sup>2</sup> for a one GeV particle corresponding to a DM mean free path of order of $`\lambda _{ss}=300`$ kpc in a gas with number density given by $`\rho _{DM}/m`$. Also we note that the cross section would scale linearly with the mass of the dark matter particle. The mean free path requirement cannot be met by the neutralino or the axion but is met quite naturally by the mirror dark matter forces.
Two obvious kinds of self interactions are the self interactions due to inter-atomic forces and nuclear forces. The latter are not very effective as can be seen by the following crude argument. From pure $`\zeta `$ scaling, we can infer that $`\sigma _{N^{}N^{}}\sigma _{NN}\zeta ^2`$. The value of $`\zeta =20`$ would put the mirror nucleon cross section to be of order $`10^{26}`$ cm<sup>2</sup>, which is much too small. Note that for $`\zeta =20`$, one would need a cross section of order $`10^{22}`$ cm<sup>2</sup> to get $`\lambda =\lambda _{ss}`$. On the other hand, one would expect that there can be scattering of the dark matter particles (mirror atoms) due to inter-atomic forces. A crude estimate of such cross sections is given by $`\sigma _{H^{}H^{}}\pi (a_0^{})^2`$, where $`a_0^{}`$ is the Bohr radius of the mirror hydrogen atom. One can then estimate the mean free path of the dark matter particles in the mirror model to be $`\lambda _{DM}3\pi \zeta ^3\times 10^{16}`$ cm. For $`\zeta =20`$, this gives $`\lambda _{DM}0.3`$ kpc, which is not far below the required values for explaining the halo core density profile. Note that Spergel and Steinhardt require $`\lambda _{DM}`$ kpc to mpc. On the other hand, one could take the value of $`\zeta =30`$ and thereby get $`\lambda `$ one kpc.
We thus feel that mirror dark matter presents the best scenario for understanding the halo core density profile using a self interacting dark matter model. Of course more detailed numerical work is needed to confirm these qualitative conclusions.
## VI Structure Formation
Finally, we use the remainder of this article to make plausible that, in spite of the $`\zeta ^2`$ decrease in cross sections for most processes, the facts that (a) structure formation begins earlier for the mirror sector (because recombination occurs before matter-radiation equality) and (b) the higher mirror temperatures for the same processes, than familiar temperatures, permit formation of galactic and smaller structures. In doing this, we will make use of our previous work in and , as well as that of Tegmark et al .
Much of the work of can be carried over to the present work, after suitable modification to take into account the fact that, in the current model, the proton mass scales as $`\zeta `$. Here, we will assume that primordial perturbations are ”curvature” or ”adiabatic” perturbations. This means that the scale of the largest structures are set by mirror Silk damping . $`\gamma ^{}`$ diffusion wipes out inhomogeneities until the $`\gamma ^{}`$ mean free path,
$$\lambda ^{}=[\sigma _T\zeta ^2n_e^{^{}}]^1$$
(16)
where $`\zeta ^2\sigma _T`$ is the mirror matter Thomson cross section and $`n_e^{}`$ is its electron number density, becomes greater than one third the horizon distance ($`ct`$). Silk damping turns off because the $`\lambda ^{}`$ increases as $`z^3`$ while $`ct`$ only increases as $`z^2`$.
First, we compute, from Silk damping, the masses of the largest structures in this picture. Structure formation starts with mirror sector particles, and familiar sector particles later fall into these. For numerical values below, we will take, $`h=0.7`$ and $`\mathrm{\Omega }_B^{}=0.2`$. We pick $`t(z_1/z)^2s`$ with $`z_1=4\times 10^9`$ and $`n_B^{}=\mathrm{\Omega }_B^{}\rho _cz^3/(\zeta m_p)`$ with $`\rho _c=1.9h^2\times 10^{29}g/cm^3`$. Silk damping stops at around $`z_{sd}8\zeta ^3`$ which gives
$`\lambda _{sd}2.5\times 10^{27}/\zeta ^6cm`$ (17)
$`M_{sd}10^{54}/\zeta ^9gm`$ (18)
Note that, for $`\zeta 10`$, this is about the mass (and size) of a large galaxy. This coincidence could be an important factor in understanding galaxy sizes should this model correspond to reality.
As in , we parametrize the separation of $`M_{sd}`$ from the expanding universe as taking place at
$$z_{stop}=z_Mz_{sd}$$
(19)
with
$$R_G=\lambda _{sd}/z_M$$
(20)
After violent relaxation we have for the proton temperature
$$T_p=GM_G\zeta m_p/R_G10^4z_M/\zeta ^2ergs$$
(21)
with $`\rho `$, outside the central plateau, given by
$$\rho (R)=A/R^210^{26}z_M/(\zeta ^3R^2)gm/cm^3$$
(22)
We now turn to the question of whether this isothermal sphere is likely to fragment and form mirror stars. For this we compute the amount of mirror molecular hydrogen since it is its collisional excitation (and subsequent radiation) that is believed to be the chief mechanism that provides cooling for formation of stars. If the rate for this mechanism is faster than the rate for free fall into the mass of the structure at issue, we can expect local regions to cool fast enough to result in fragmentation of that structure. We do here a very rough estimate of mirror galaxy fragmentation into mirror globular clusters, using the results of , but leave to a more detailed work further fragmentation into the $`0.5M_{}`$ structures predicted in .
Reference give a useful approximation to their numerical results for the fraction of molecular hydrogen, $`f_2`$ ($`f_0`$ denotes its primordial value):
$$f_2(t)=f_0+(k_m/k_1)ln[1+x_0nk_1t)$$
(23)
where, as a first try, $`k_m`$ can be taken as just the rate for $`H+e^{}H^{}+\gamma `$ (which they conveniently give as about $`2\times 10^{18}T^{.88}cm^3s^1`$), while $`k_1`$ is the rate for $`H^++e^{}H+\gamma `$ ($`2\times 10^{10}T^{0.64}cm^3s^1`$). Equation (8) is the result of $`H_2`$ production from the catalytic reactions $`H+e^{}H^{}+\gamma `$ followed by $`H+H^{}H_2+e^{}`$ competing against the recombination reaction that destroys the catalyst, free electrons, (approximately) as $`1/t`$ (assuming constant density). Our goal here is to show from Equation (20) that it is plausible that $`f_2`$, the fraction of molecular hydrogen, rises from its primordial value of $`10^6`$ to the region above $`10^4`$ where cooling tends to be competitive with free fall.
First, we note that, if $`k=<v\sigma >AT^\gamma cm^3/s`$, for familiar e and p, we expect that, for mirror e and p, scaling with $`\zeta `$ to go as $`\zeta ^{(2+\gamma )}AT^\gamma `$, since $`\sigma `$ must go as $`\zeta ^2`$, all factors of $`T`$ must be divided by some combination of $`m_e`$ and $`M_p`$, both of which go as $`\zeta `$ (making this model much easier to compute from than that of ).
We now estimate fragmentation. From Equation (18) we see that the galactic temperature should begin at about 10 eV at a time when the cosmic temperature is about 1 eV and the cosmic gamma number density is about $`10^9/cm^3`$. The rate for “compton cooling” is very fast at this high density (unlike at later times for the familiar case) and there should be rapid cooling to about 1 eV. We can now compute the Jeans length for fragments as a function of distance R from galaxy central. We use
$$\rho _J=(T/Gm)^3/M^2$$
(24)
If we set the Jeans mass, $`M`$, to $`4\pi r^3\rho _J/3`$, we can solve for $`r`$ obtaining (if we are careful to convert $`T`$ in Equation (6) from ergs)
$$r=R[10^7\zeta ^2/z_M]^{1/2}10^2R$$
(25)
Now inserting into Equation (20) gives the coefficient of the log term on the order of $`10^{3.5}`$ and the coefficient of $`t`$ in the argument varying from $`10^{13}`$ to $`10^{17}`$ as $`R`$ varies from 1 to 100 kiloparsecs while the free fall time ($`(G\rho )^{1/2}`$) varies from $`10^{14}`$ to $`10^{16}`$. This would appear to indicate the likelihood of fragmentation of the original Silk damping structure into smaller units, and the eventual formation of the $`0.5M_{odot}`$ black holes that explain the microlensing events of .
## VII Conclusion
We have argued that the asymmetric mirror model, originally proposed to solve neutrino puzzles and subsequently advocated as providing an alternative dark matter candidate has the advantage of resolving the microlensing anomaly and possibly the core density problem of dark halos. Possible tests of these models are to narrow the allowed values of $`\delta N_\nu `$ from more accurate observation of deuterium, He<sup>4</sup> and Li<sup>7</sup> and observing whether further accumulation of MACHO candidates lie in the mass between 0.1- 1 solar mass. Needless to say that if the underground searches for the cold dark matter now under way lead to a positive signal, mirror matter cannot be the dominant component of the dark matter of the universe.
The work of R. N. M. is supported by the National Science Foundation grant under no. PHY-9802551 and the work of V. L. T. is supported by the DOE under grant no. DE-FG03-95ER40908.
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# Non-Abelian Aharonov-Bohm Scattering of Spin 1/2 Particles*footnote **footnote *Copyright by The American Physical Society
## I INTRODUCTION
Starting from different perspectives, a scalar non-Abelian Aharonov-Bohm (AB) effect has been discussed by several authors. This subject has interesting implications to the physics of peculiar objects as cosmic string and black holes; it has also applications to some aspects of gravitation in 2+1dimensions. Cosmic strings, for example, may have trapped non-Abelian magnetic flux tubes so that the scattering of charged particles by these strings is just a manifestation of the non-Abelian AB effect.
The study of the AB effect was started through the exact calculation of the scattering amplitude of scalar particles by a thin magnetic flux tube at the origin. As it is nowadays well-known, in that situation the perturbative Born approximation fails to reproduce the expansion of the exact result and, moreover, the second term of the Born series is divergent. This discrepancy is due to the fact that the perturbative wave function does not satisfy the same boundary condition as the exact one. Actually, in a perturbative treatment for a nonrelativistic field theory describing spinless Abelian particles scattered through a Chern-Simons field, it was shown that to eliminate the divergences, to recuperate the scale invariance and to reproduce the result of the expansion of the exact solution, it is necessary to add a contact term $`(\varphi ^{}\varphi )^2`$.
Recently, the perturbative treatment was applied to relativistic non-Abelian scalar particles minimally coupled to non-Abelian CS field. By considering the low momentum limit, it was shown that, up to leading order, the same results, got through the calculation of a non-Abelian nonrelativistic field theory, is obtained. In the next to leading approximation, new corrections appear which are absent in the direct nonrelativistic approach. These corrections also differ from the ones got in the Abelian theory.
By analyzing the Abelian AB effect, it has been verified that new features appear if spin is introduced . For example, the Pauli’s magnetic term plays the role of a contact interaction and no quartic self-interaction is needed. Besides that, as shown in new effective low momentum interactions are induced if one starts from a fully relativistic theory.
Completing our study of the non-Abelian AB effect began in, in this work we analyze the AB scattering for non-Abelian spin 1/2 particles. We start by calculating the AB scattering in a nonrelativistic setting. We then consider the AB scattering from a more basic standpoint, starting from a relativistic quantum field theory, and then taking the appropriate nonrelativistic limit of the scattering amplitudes. One of the advantages of such procedure is in the fact that it automatically incorporates quantum radiative corrections as the vacuum polarization and induced magnetic moment. To take the nonrelativistic limit most easily, we use an intermediate auxiliary cutoff separating the low and high loop momenta in the Feynman integrals. As it happened in our previous studies, it is also convenient to work in the Coulomb gauge, since in this gauge the Chern-Simons propagator depends only on the spatial part of the loop momentum variable.
## II NONRELATIVISTIC THEORY
We consider the non-Abelian Pauli–Schrödinger model for fermions minimally coupled to a non-Abelian Chern-Simons field specified by the Lagrangian
$``$ $`=`$ $`\mathrm{\Theta }\epsilon ^{\alpha \beta \lambda }\text{tr}\left(A_\alpha _\beta A_\lambda +{\displaystyle \frac{2g}{3}}A_\alpha A_\beta A_\lambda \right)+\psi ^{}\left[i_t+{\displaystyle \frac{(\mathrm{𝐠𝐀})^\mathrm{𝟐}}{2m}}+igA_0i{\displaystyle \frac{g}{2m}}B\right]\psi `$ (1)
$``$ $`{\displaystyle \frac{1}{\xi }}\text{tr}(.𝐀)^\mathrm{𝟐}𝐜^𝐚(\delta _{\mathrm{𝐚𝐛}}^\mathrm{𝟐}+𝐠\epsilon _{\mathrm{𝐚𝐛𝐜}}𝐀^𝐜.)𝐜^𝐛,`$ (2)
where $`\psi `$ is a one-component anticommuting field, belonging to the fundamental representation of the $`SU(2)`$ group, and $`A_\mu =A_\mu ^aT^a`$, with $`T^a`$ being the generator of the Lie Algebra of $`SU(2)`$ satisfying
$$[T_a,T_b]=\epsilon _{abc}T^c,$$
(3)
and normalized such that
$$T^aT^b=\frac{\delta _{ab}}{4}I+\frac{1}{2}\epsilon ^{abc}T_c$$
(4)
The term containing the “magnetic” field $`B`$, is the Pauli term and $`c`$ is the ghost field needed to guarantee unitarity. For convenience, we will work in a strict Coulomb gauge obtained by letting $`\xi 0`$.
We will use a graphical notation where the CS field, the matter field and the ghost field propagators are represented by wavy, continuous and dashed lines respectively. The analytic expression for the $`A_\mu `$ free propagator is
$$D^{\mu \nu }(k)_{ba}=D^{\mu \nu }(k)\delta _{ba}=\frac{1}{\mathrm{\Theta }}\epsilon ^{\mu \nu \lambda }\frac{\overline{k}_\lambda }{𝐤^2}\delta _{ba},$$
(5)
where $`\overline{k}_\lambda (0,𝐤)`$. The matter field propagator is
$$S(p)_{nm}=S(p)\delta _{nm}=\frac{i}{p_0\frac{𝐩^2}{2m}+iϵ}\delta _{nm},$$
(6)
whereas the ghost field propagator is
$$G(p)_{ba}=G(p)\delta _{ba}=\frac{i}{𝐩^2}\delta _{ba}.$$
(7)
Since the $`B`$ field occurs in Eq. (2) it is convenient to have at hand
$$\mathrm{\Delta }_B^{ba}(x)=<TB^b(x)A_0^a(0)>=\frac{i}{\mathrm{\Theta }}\delta ^{(3)}(x)\delta ^{ba},$$
(8)
which is the only nonvanishing propagator involving the $`B`$ field; graphically it will be represented by a dotted line. Expression (8 shows that the Pauli term, i . e., the interaction $`\psi ^{}B\psi `$, plays the same role as the quartic term $`(\varphi ^{}\varphi )^2`$ in the scalar case.
The graphical representation for the vertices is given in Fig. 1 and the corresponding analytical expression are
$`\mathrm{\Gamma }_{nm}^{a,0}(p,p^{})=g(T^a)_{nm},`$ (9)
$`\mathrm{\Gamma }_{nm}^{a,i}(p,p^{})={\displaystyle \frac{g}{2m}}(T^a)_{nm}(p^i+p^i),`$ (10)
$`\mathrm{\Gamma }_{nm}^{ab,ij}(p,p^{})=i{\displaystyle \frac{g^2}{2m}}(T^aT^b+T^bT^a)_{nm}g^{ij},`$ (11)
$`\mathrm{\Gamma }_{nm}^{a,B}(p,p^{})={\displaystyle \frac{g}{2m}}(T^a)_{nm},`$ (12)
$`\mathrm{\Gamma }^{abc,\mu \nu \lambda }(p,p^{})=ig\mathrm{\Theta }\epsilon ^{abc}\epsilon ^{\mu \nu \lambda },`$ (13)
$`\mathrm{\Gamma }_{nm}^{abc,i}(p,p^{})=g\epsilon ^{abc}p^i\delta _{nm}.`$ (14)
In the tree approximation and in the center-of-mass frame the two body scattering amplitude is given by
$$(\theta )=\frac{ig^2}{m\mathrm{\Theta }}[T^aT_a]\left[1+i\frac{\mathrm{sin}\theta }{(1\mathrm{cos}\theta )}\right],$$
(15)
where $`\theta `$ is the scattering angle. Here and in what follows we employ a simplified notation where the isospin indices are omitted. If the incoming and outgoing particles have isospin ($`n,m`$) and ($`n^{},m^{}`$) the total amplitude for the process is given by
$$_{n^{}m^{};nm}=n^{},m^{}(\theta )n,mm^{},n^{}(\theta +\pi )n,m.$$
(16)
The one-loop contribution to AB scattering is depicted in the Fig. 2. The incoming and outgoing fermions are assumed to have momenta $`p_1=(𝐩_1^2/2m,𝐩_1)`$, $`p_2=(𝐩_2^2/2m,𝐩_2)`$ and $`p_3=(𝐩_3^2/2m,𝐩_3)`$, $`p_4=(𝐩_4^2/2m,𝐩_4)`$, respectively. We work in the center-of-mass frame where $`𝐩_1=𝐩_2=𝐩`$, $`𝐩_3=𝐩_4=𝐩^{}`$ and $`|𝐩|=|𝐩^{}|`$. For the first graph, Fig. 2(a), we get
$`_\text{a}(\theta )={\displaystyle \frac{d^3k}{(2\pi )^3}}`$ $`[\mathrm{\Gamma }^{d,\alpha }(p_1+p_2k,p_4)S(p_1+p_2k)\mathrm{\Gamma }^{c,\nu }(p_2,p_1+p_2k)`$ (18)
$`D_{\mu \nu }^{ac}(kp1)D_{\alpha \beta }^{db}(kp_3)\mathrm{\Gamma }^{b,\beta }(k,p_3)S(k)\mathrm{\Gamma }^{a,\mu }(p_1,k)].`$
After performing the $`k_0`$ integration, this gives
$$_\text{a}(\theta )=\frac{4ig^4}{m\mathrm{\Theta }^2}[T^bT^aT_bT_a]\frac{d^2𝐤}{(2\pi )^2}\frac{1}{𝐩^2𝐤^2+iϵ}\left[\frac{(𝐩_1𝐤)(𝐩_3𝐤)}{(𝐤𝐩_1)^2(𝐤𝐩_3)^2}\right].$$
(19)
As a general rule, whenever dealing with divergent spatial integrals we will introduce a nonrelativistic cutoff $`\mathrm{\Lambda }_{NR}`$. However, in Eq. (19) such regulator is not necessary as the integral is ultraviolet finite. The final result is
$$_\text{a}(\theta )=\frac{ig^4}{4\pi m\mathrm{\Theta }^2}[T^bT^aT_bT_a]\left\{\mathrm{log}\left[\frac{𝐪^2}{𝐩^2}\right]+i\pi \right\}.$$
(20)
where $`𝐪=𝐩_\mathrm{𝟑}𝐩_\mathrm{𝟏}`$ is the momentum transferred.
The same procedure can be used to calculate the other graphs in Fig. 2. Here the spatial integrals are logarithmically divergent and are done after the introduction of the aforementioned cutoff. Graph 2(b) gives
$$_\text{b}(\theta )=\frac{g^4}{m^2\mathrm{\Theta }^2}[T^bT^aT_bT_a]\frac{d^3k}{(2\pi )^3}S(p_1+p_2k)S(k),$$
(21)
from which we obtain
$$_\text{b}(\theta )=\frac{ig^4}{4\pi m\mathrm{\Theta }^2}[T^bT^aT_bT_a]\left\{\mathrm{log}\left[\frac{\mathrm{\Lambda }_{NR}^2}{𝐩^2}\right]+i\pi \right\}.$$
(22)
Similarly, graph 2(c) corresponds to
$$_\text{c}(\theta )=2\frac{d^3k}{(2\pi )^3}\mathrm{\Gamma }^{cd,ij}D_{0i}^{ac}(k)D_{j0}^{db}(k+q)\mathrm{\Gamma }^{b,0}\mathrm{\Gamma }^{a,0}.$$
(23)
The $`k^0`$ integration is straightforward and gives
$$_\text{c}(\theta )=\frac{ig^4}{2m\mathrm{\Theta }^2}[(T^aT^b+T^bT^a)T_bT_a]\frac{d^2k}{(2\pi )^2}\frac{𝐤(𝐤+𝐪)}{𝐤^2(𝐤+𝐪)^2}.$$
(24)
Effectuating the remaining integral produces
$$_\text{c}(\theta )=\frac{ig^4}{4\pi m\mathrm{\Theta }^2}[T^bT^aT_bT_a+\frac{1}{2}\epsilon ^{cab}T_cT_bT_a]\left\{\mathrm{log}\left[\frac{𝐪^2}{\mathrm{\Lambda }_{NR}^2}\right]\right\}.$$
(25)
The last diagram, graph 2(d) gives
$`_\text{d}(\theta )`$ $`=`$ $`2{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}[\mathrm{\Gamma }^{b,\nu }(p_2,p_4)D_{\nu \mu }^{ab}(q)\mathrm{\Gamma }^{ac^{}d^{},\mu \rho \sigma }D_{\sigma \alpha }^{d^{}d}(kp_3)`$ (26)
$`\mathrm{\Gamma }^{d,\alpha }(k,p_3)S(k)\mathrm{\Gamma }^{c,\beta }(p_1,k)D_{\beta \rho }^{cc^{}}(kp_1)]`$ (27)
so that, after the $`k^0`$ integration,
$$_\text{d}(\theta )=\frac{ig^4}{m\mathrm{\Theta }^2}[\epsilon ^{cab}T_cT_bT_a]\frac{d^2k}{(2\pi )^2}\frac{[𝐪𝐤𝐩_1𝐩_3](𝐪𝐤)}{𝐪^2(𝐤𝐩_1)^2(𝐤𝐩_3)^2}$$
(28)
leading to
$$_\text{d}(\theta )=\frac{ig^4}{8\pi m\mathrm{\Theta }^2}[\epsilon ^{cab}T_cT_bT_a]\left\{1\mathrm{log}\left[\frac{𝐪^2}{\mathrm{\Lambda }_{NR}^2}\right]\right\}.$$
(29)
Thus, the sum of the one-loop contribution is
$$_{1loop}(\theta )=\frac{ig^4}{8\pi m\mathrm{\Theta }^2}[\epsilon ^{cab}T_cT_bT_a]=\frac{ig^4}{8\pi m\mathrm{\Theta }^2}[T^aT_a].$$
(30)
It happens that the nonvanishing result in the last equation is only due to the regularization used. Really, as the original expression was logarithmically divergent, different regularization schemes will produce results that for the finite part will differ at most by a constant. This remark holds even for the sum of the Feynman integrals which is only conditionally convergent and leads to different results depending on the way it is treated. In particular, had we used the dimensional regularization, as it was done in the reference for the scalar case, Eq. (30) would be zero. Our constant term in that result may be eliminated through a redefinition of the cutoff $`\mathrm{\Lambda }_{NR}`$ in the Eq. (29) or by adding a counterterm of the form $`(\psi ^{}T^a\psi )^2`$ to the original Lagrangian. In the relativistic theory the divergences are milder, the graphs are individually finite and no such counterterms are needed.
## III RELATIVISTIC THEORY
We will now consider the non-Abelian scattering within the full relativistic context. The Lagrangian describing the model is
$``$ $`=`$ $`\mathrm{\Theta }\epsilon ^{\alpha \beta \lambda }\text{tr}\left(A_\alpha _\beta A_\lambda +{\displaystyle \frac{2g}{3}}A_\alpha A_\beta A_\lambda \right)+i\overline{\mathrm{\Psi }}(\overline{)}Dm)\mathrm{\Psi }`$ (31)
$``$ $`{\displaystyle \frac{1}{\xi }}\text{tr}(.𝐀)^\mathrm{𝟐}𝐜^𝐚(^\mathrm{𝟐}+𝐠\epsilon _{\mathrm{𝐚𝐛𝐜}}𝐀^𝐜.)𝐜^𝐛.`$ (32)
where $`D_\mu =_\mu +gA_\mu `$ and $`\mathrm{\Psi }`$ is a two-component Dirac field belonging to the fundamental representation of the SU(2) gauge group. $`\psi `$ represents particles and anti-particles with the same spin and we take $`m`$ to be positive so that. Our graphical notation is specified in Fig 3. The corresponding analytical expressions for the gauge and ghost field propagators are the same as in the previous section. The matter field propagator and the vertices, however, are now given by
$`S(p)_{nm}=S(p)\delta _{nm}={\displaystyle \frac{i(\overline{)}p+m)}{p^2m^2+iϵ}}\delta _{nm},`$ (33)
$`\mathrm{\Gamma }_{nm}^{a,\mu }(p,p^{})=g(T^a)_{nm}(\gamma ^\mu ),`$ (34)
$`\mathrm{\Gamma }^{abc,\mu \nu \lambda }(p,p^{})=ig\mathrm{\Theta }f^{abc}\epsilon ^{\mu \nu \lambda },`$ (35)
$`\mathrm{\Gamma }_{nm}^{abc,i}(p,p^{})=g\epsilon ^{abc}p^i\delta _{nm}.`$ (36)
The model is renormalizable. Actually, without the matter field it was found that there are no radiative corrections to the Green functions . We can therefore restrict our study of one-loop renormalization to superficially divergent graphs arising from the coupling to the matter field, i. e., the 1-loop correction to the self-energy, vacuum polarization and vertex corrections. The nonvanishing self-energy graph depicted in Fig. 4 is given by
$`\mathrm{\Sigma }(p)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\left[\mathrm{\Gamma }^{a,\mu }(p+k,p)S(p+k)\mathrm{\Gamma }^{b,\nu }(p,p+k)D_{\nu \mu }^{ab}(k)\right]}`$ (37)
$`=`$ $`{\displaystyle \frac{ig^2}{\mathrm{\Theta }}}[T^aT_a]{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{[\gamma ^\mu (\overline{)}p+\overline{)}k+m)\gamma ^\nu ]\epsilon _{\mu \nu \lambda }\overline{k}^\lambda }{[(p+k)^2m^2+iϵ]𝐤^2}},`$ (38)
so that the inverse of the complete fermion propagator is written as $`𝒮^1(p)=\overline{)}pm+i\mathrm{\Sigma }`$. Notice that the self-energy is diagonal in isospin space. After doing the $`k_0`$ integration we obtain
$$\mathrm{\Sigma }(p)=\frac{ig^2}{8\pi \mathrm{\Theta }}[T^aT_a]_0^{\mathrm{\Lambda }_0^2}𝑑𝐤^2\frac{1}{w_k}\left\{\frac{m}{𝐩^2}\gamma 𝐩[1ϵ(𝐤^2𝐩^2)]+[1+ϵ(𝐤^2𝐩^2)]\right\},$$
(39)
where $`ϵ(x)`$ is the signal function, $`w_k=\sqrt{𝐤^2+m^2}`$ and a cutoff $`\mathrm{\Lambda }_0`$ was introduced to take care of the ultraviolet divergence of the integral. The integral is easily done and gives
$$\mathrm{\Sigma }(p)=\frac{ig^2}{2\pi \mathrm{\Theta }}[T^aT_a]\left[\gamma 𝐩\frac{m}{𝐩^2}(w_pm)+\sqrt{\mathrm{\Lambda }_0^2+m^2}w_p\right]$$
(40)
and so, for $`\mathrm{\Lambda }_0\mathrm{}`$,
$$\mathrm{\Sigma }(p)=\frac{ig^2}{2\pi \mathrm{\Theta }}[T^aT_a]\left\{\frac{m\gamma 𝐩𝐩^2}{w_p+m}m+\mathrm{\Lambda }_0\right\}.$$
(41)
The linear ultraviolet divergence may be eliminated through the imposition of an adequate renormalization condition. Due to our use of the Coulomb gauge, a convenient condition is the one adopted in the work ; denoting the renormalized propagator by $`𝒮_R`$, this condition reads $`𝒮_R(p_0,𝐩=0)=S(p_0,𝐩=0)`$. Proceeding in this way, we get for the renormalized propagator
$$𝒮_R(p)=i\frac{(\overline{)}p+m)+\alpha (mw_p)[1+\frac{m}{𝐩^2}\gamma .𝐩]}{(p^2m^2)}.$$
(42)
where $`\alpha =g^2[T^aT_a]/(2\pi \mathrm{\Theta })`$.
Let us now turn our attention to the vacuum polarization correction. The only graph that contributes is the one drawn in Fig. 5. As this graph is gauge independent, the would be linear divergence may be eliminated if one employs a gauge invariant regularization scheme. Use of dimensional regularization gives
$$\mathrm{\Pi }^{\mu \nu }(q)=\frac{ig^2}{4\pi }\text{tr}[T^aT^b]\left[\left(g^{\mu \nu }\frac{q^\mu q^\nu }{q^2}\right)\mathrm{\Pi }_1(q^2)+im\epsilon ^{\mu \nu \lambda }q_\lambda \mathrm{\Pi }_2(q^2)\right],$$
(43)
with
$`\mathrm{\Pi }_1(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑x{\displaystyle \frac{2q^2x(1x)}{[m^2q^2x(1x)]^{1/2}}}{\displaystyle \frac{q^2}{3|m|}}`$ (44)
$`\mathrm{\Pi }_2(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑x{\displaystyle \frac{1}{[m^2q^2x(1x)]^{1/2}}}{\displaystyle \frac{1}{|m|}},`$ (45)
where the expressions on the right of these equations are the leading approximations for low momenta $`q`$. From these results, we see that for low momentum a Yang-Mills term may be induced, as one would expect on general grounds.
The 1-loop corrections to the CS-matter field vertex are given by the graphs in Figs. 6(a)and 6(b). The on-shell analytic expression associated to the graph 6(a) is
$$\overline{u}(p^{})\mathrm{\Gamma }_\text{a}^{a,\mu }u(p)=\frac{g^3}{\mathrm{\Theta }}[T^bT^aT_b]\frac{d^3k}{(2\pi )^3}\frac{\epsilon _{\rho \sigma \lambda }\overline{k}^\lambda \overline{u}(p^{})[\gamma ^\sigma (\overline{)}p^{}\overline{)}k+m)\gamma ^\mu (\overline{)}p\overline{)}k+m)\gamma ^\rho ]u(p)}{[(pk)^2m^2+iϵ][(p^{}k)^2m^2+iϵ][𝐤^2]}.$$
(46)
Here and in what follows the isospin indices $`(n,m)`$ will omitted. Up to the group factor $`T^bT^aT_b`$, this agrees with the vertex for the Abelian theory discussed in . Using dimensional regularization, the result can be read from that reference but for general momenta it is not particularly illuminating. Nevertheless, for small momenta (i. e., for $`|𝐩||𝐩^{}|m`$) a great simplification occurs and one finds ($`\eta =\frac{|𝐩|}{m}`$)
$$\overline{u}(p^{})\mathrm{\Gamma }_\text{a}^{a,0}u(p)=𝒪(\eta ^2),$$
(47)
$$\overline{u}(p^{})\mathrm{\Gamma }_\text{a}^{a,i}u(p)=\frac{g^3}{4\pi \mathrm{\Theta }}[T^bT^aT_b]\frac{1}{2m}[P^ii\epsilon ^{ij}q_j]+𝒪(\eta ^2),$$
(48)
where $`P^i=p^i+p^i`$ and $`q=p^ip^i`$
Similarly, the graph 6(b) which corresponds to
$$\overline{u}(p^{})\mathrm{\Gamma }_\text{b}^{a,\mu }u(p)=\frac{g^3}{\mathrm{\Theta }}[\epsilon ^{abc}T_bT_c]\frac{d^3k}{(2\pi )^3}\frac{\overline{u}(p^{})[\gamma ^\sigma (\overline{)}k+m)\gamma ^\beta ]u(p)\epsilon ^{\mu \sigma \rho }\epsilon _{\sigma \beta \lambda }\overline{(pk)^\lambda }\epsilon _{\alpha \rho \xi }\overline{(p^{}k)^\xi }}{[k^2m^2+iϵ](𝐩𝐤)^2(𝐩^{}𝐤)^2}$$
(49)
gives for small momenta the result
$$\overline{u}(p^{})\mathrm{\Gamma }_\text{b}^{a,0}u(p)=𝒪(\eta ^2)$$
(50)
and
$`\overline{u}(p^{})\mathrm{\Gamma }_\text{b}^{a,i}u(p)={\displaystyle \frac{g^3}{8\pi m\mathrm{\Theta }}}[\epsilon ^{abc}T_cT_b]\left\{P^i+i\epsilon ^{ij}q_j\left[1+\mathrm{log}\left({\displaystyle \frac{4m^2}{𝐪^2}}\right)\right]\right\}.`$ (51)
The renormalized vertex part is defined by $`\mathrm{\Gamma }_R^{a\mu }=Z_1\mathrm{\Gamma }^{a,\mu }`$ where $`\mathrm{\Gamma }^{a\mu }=gT^a\gamma ^\mu +\mathrm{\Lambda }^{a\mu }`$ is the unrenormalized one. Fixing the vertex renormalization constant $`Z_1`$ by the condition that for $`𝐩=𝐩^{}=0`$ and $`p^0=p^0=m`$
$$\overline{u}\mathrm{\Gamma }_R^{a\mu }u=gT^ag^{0\mu }$$
(52)
we get $`Z_1=1`$, so that up to 1-loop there is no coupling constant renormalization. This result is also in accord with the computation of the correction to the trilinear CS vertex shown in Fig 6(c); simple symmetry considerations shown that the result is finite and no counterterm is necessary. Actually, the graph 6(c) plus the graphs with four external gauge lines and the polarization tensor give an induced Yang-Mills term (and also a finite correction for the Chern-Simons term), as commented before. However, up to one-loop the graph 6(c) does not contribute to the scattering and for that reason it will not be considered any longer.
Summarizing, up to one-loop one needs just a mass renormalization counterterm to fix the fermion mass. There are neither vertex nor wave function renormalizations.
Although not 1PI, we have drawn in Fig.6 the graphs 6(d) and 6(e) which are needed to compute the anomalous magnetic moment of the fermions. At low momenta these graphs give the contributions
$`\overline{u}\mathrm{\Gamma }_\text{d}^{a,0}u`$ $`=`$ $`\overline{u}\mathrm{\Gamma }_\text{e}^{a,0}u=𝒪(\eta ^2)`$ (53)
$`\overline{u}\mathrm{\Gamma }_{\text{de}}^{a,i}u`$ $`=`$ $`\overline{u}\mathrm{\Gamma }_\text{d}^{a,i}u+\overline{u}\mathrm{\Gamma }_\text{e}^{a,i}u={\displaystyle \frac{g^3}{4\pi \mathrm{\Theta }}}[(T^bT_b)T^a]\left\{{\displaystyle \frac{1}{2m}}[P^i+i\epsilon ^{ij}q_j]\right\}.`$ (54)
In the Abelian situation the contribution in Eq. (51) is absent and, in the expressions corresponding to Eqs.(48) and (54) the $`P^i`$ dependent part is exactly canceled. Here, due to the group factors, to get cancellation it is necessary to take into account the new contribution arising from (51). This can be easily verified using the identity $`T^bT^aT_b=T^a(T^bT_b)+\epsilon ^{abc}T_cT_b`$. The remaining local parts occurring in $`\mathrm{\Gamma }_{\text{a-e}}^{a\mu }`$ will contribute to the (matrix) magnetic moment and we get
$`\mu _{1loop}^a={\displaystyle \frac{ig^3}{4\pi m\mathrm{\Theta }}}[T^a(T^bT_b)].`$ (55)
This expression only differs from the corresponding result in the Abelian case by the group factor. In the “Abelian limit” ($`g=e\sqrt{2}`$ e $`T=i/\sqrt{2}`$) the result of is recovered.
To complete our discussion of the one-loop properties of the model one still has to calculate the fermion-fermion scattering. Fig. 7 shows the contributing graphs. The only tree level graph, depicted in Fig. 7(a), furnishes
$`_\text{a}(\theta )=[\overline{u}(𝐩_\mathrm{𝟒})\mathrm{\Gamma }^{b,\nu }(p_2,p_4)u(𝐩_\mathrm{𝟐})]D_{\nu \mu }^{ba}(q)[\overline{u}(𝐩_\mathrm{𝟑})\mathrm{\Gamma }^{a,\mu }(p_1,p_3)u(𝐩_\mathrm{𝟏})].`$ (56)
To the leading order of $`𝐩/m`$, this gives
$`_\text{a}(\theta )={\displaystyle \frac{ig^2}{m\mathrm{\Theta }}}[T^aT_a]\left\{1+i{\displaystyle \frac{\mathrm{sin}\theta }{(1\mathrm{cos}\theta )}}\right\},`$ (57)
which exactly agrees with that obtained previously for the nonrelativistic theory.
The one-loop graphs are represented in Figs. 7(b)-7(h). To facilitate our computation we will use an intermediate cutoff $`\mathrm{\Lambda }`$, satisfying $`|𝐩|<<\mathrm{\Lambda }<<m`$, which separates the loop integrals in two regions. In the low $`(L)`$ region $`(0|𝐤|^2\mathrm{\Lambda }^2)`$ the integrand is expanded in power of $`1/m`$, and in high $`(H)`$ region $`(|𝐤|^2\mathrm{\Lambda }^2)`$ we make a Taylor series of the integrand around $`|𝐩|0`$. We will retain terms up to order $`\eta =\frac{|𝐩|}{m}\left(\frac{\mathrm{\Lambda }}{m}\right)^2\left(\frac{|𝐩|}{\mathrm{\Lambda }}\right)^2`$.
Using that
$`S(p)=i\left[{\displaystyle \frac{u(𝐩)\overline{u}(𝐩)}{p^0w_p+iϵ}}+{\displaystyle \frac{v(𝐩)\overline{v}(𝐩)}{p^0+w_piϵ}}\right],`$ (58)
we may decompose the amplitude for the graph in Fig. 7(b)
$`_\text{b}(\theta )`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}[\overline{u}(𝐩_4)\mathrm{\Gamma }^{c,\alpha }(t,p_4)S(t)\mathrm{\Gamma }^{d,\beta }(p_2,t)u(𝐩_2)]D_{\alpha \mu }^{ca}(l^{})}`$ (59)
$`[\overline{u}(𝐩_3)\mathrm{\Gamma }^{a,\mu }(r,p_3)S(r)\mathrm{\Gamma }^{b,\nu }(p_1,r)u(𝐩_2)]D_{\nu \beta }^{bd}(l)`$ (60)
where $`l=(k^0,𝐤𝐩_1)`$, $`l^{}=(k^0,𝐤𝐩_3)`$, $`r=(w_p+k^0,𝐤)`$ and $`t=(w_p+k^0,𝐤)`$, into a sum of terms
$$_b^{uu}+_b^{vv}$$
(61)
where $`_b^{uu}`$ and $`_b^{vv}`$ designate the contributions of the $`u`$ and $`v`$ fermion wave functions to the two internal lines of the graph. The mixed contributions in which one has $`u`$ in one line and $`v`$ in the other vanish. After integrating in $`k^0`$ we obtain
$$_b^{uu}=\frac{ig^4}{2}[T^aT^bT_aT_b]\frac{d^2k}{(2\pi )^2}\left[\frac{T(k,p_1)T^{}(k,p_3)}{w_kw_p}\right]$$
(62)
and
$$_b^{vv}=\frac{ig^4}{2}[T^aT^bT_aT_b]\frac{d^2k}{(2\pi )^2}\left[\frac{H(p_3,k)H^{}(p,k)}{w_k+w_p}\right],$$
(63)
where
$`T(k,p)`$ $`=`$ $`[\overline{u}(𝐤)\gamma ^\nu u(𝐩)]D_{\nu \beta }(kp)[\overline{u}(𝐤)\gamma ^\nu u(𝐩)],`$ (64)
$`H(p,k)`$ $`=`$ $`[\overline{u}(𝐩)\gamma ^\nu v(𝐤)]D_{\nu \beta }(kp)[\overline{u}(𝐩)\gamma ^\nu v(𝐤)].`$ (65)
Introducing the intermediate cutoff to separate the low and high parts we get
$`_{\text{b}Low}^{uu}(\theta )`$ $`=`$ $`{\displaystyle \frac{ig^4}{4\pi m\mathrm{\Theta }^2}}[T^aT^bT_aT_b]\left\{\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{𝐪^2}}\right)+𝒪(\eta )\right\},`$ (66)
$`_{\text{b}High}^{uu}(\theta )`$ $`=`$ $`{\displaystyle \frac{ig^4}{4\pi m\mathrm{\Theta }^2}}[T^aT^bT_aT_b]\left\{\mathrm{log}\left({\displaystyle \frac{2m^2}{\mathrm{\Lambda }^2}}\right)+𝒪(\eta )\right\},`$ (67)
$`_{\text{b}Low}^{vv}(\theta )`$ $`=`$ $`{\displaystyle \frac{ig^4}{4\pi m\mathrm{\Theta }^2}}[T^aT^bT_aT_b]\left\{𝒪(\eta )\right\},`$ (68)
$`_{\text{b}High}^{vv}(\theta )`$ $`=`$ $`{\displaystyle \frac{ig^4}{4\pi m\mathrm{\Theta }^2}}[T^aT^bT_aT_b]\left\{\mathrm{log}\left(2\right)+𝒪(\eta )\right\}.`$ (69)
Putting these results together we arrive at
$`_\text{b}(\theta )={\displaystyle \frac{ig^4}{4\pi m\mathrm{\Theta }^2}}[T^aT^bT_aT_b]\left\{\mathrm{log}\left({\displaystyle \frac{4m^2}{𝐪^2}}\right)\right\}.`$ (70)
as the leading contribution.
For the crisscross graph, Fig. 7(c), we proceed analogously and obtain (in this case what survives are the mixed $`uv`$ and $`vu`$ contributions)
$`_{\text{c}Low}(\theta )`$ $`=`$ $`{\displaystyle \frac{ig^4}{2\pi m\mathrm{\Theta }^2}}[T^bT^aT_aT_b]\left\{{\displaystyle \frac{1}{2}}\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{𝐪^2}}\right)+𝒪(\eta )\right\},`$ (71)
$`_{\text{c}High}(\theta )`$ $`=`$ $`{\displaystyle \frac{ig^4}{2\pi m\mathrm{\Theta }^2}}[T^bT^aT_aT_b]\left\{1+{\displaystyle \frac{1}{2}}\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{4m^2}}\right)+𝒪(\eta )\right\}.`$ (72)
i.e.,
$`_\text{c}(\theta )={\displaystyle \frac{ig^4}{2\pi m\mathrm{\Theta }^2}}[T^bT^aT_aT_b]\left\{1+{\displaystyle \frac{1}{2}}\mathrm{log}\left({\displaystyle \frac{𝐪^2}{4m^2}}\right)+𝒪(\eta )\right\}.`$ (73)
The graph 7(d) does not exist in the Abelian theory but here it is essential to cancel the extra contribution coming through group factors in other graphs. It corresponds to
$`M_d`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}\{[\overline{u}(p_4)\mathrm{\Gamma }^{b,\nu }u(p_2)]D_{\nu \mu }^{ba}(q)\mathrm{\Gamma }^{ac^{}d^{},\mu \sigma \rho }D_{\sigma \alpha }^{dd^{}}(kp_3)`$ (74)
$`D_{\beta \rho }^{cc^{}}(kp_1)[\overline{u}(p_3)\mathrm{\Gamma }^{d,\alpha }S(k)\mathrm{\Gamma }^{c,\beta }u(p_1)]\}`$ (75)
has a low and high momentum parts given by
$`M_{\text{d}Low}`$ $`=`$ $`{\displaystyle \frac{ig^4}{8\pi m\mathrm{\Theta }^2}}[\epsilon ^{acd}T_aT_dT_c]\{{\displaystyle \frac{}{}}1+i{\displaystyle \frac{\mathrm{sin}\theta }{1\mathrm{cos}\theta }}`$ (76)
$`+`$ $`[\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{𝐪^2}}\right){\displaystyle \frac{\mathrm{\Lambda }^2}{2m^2}}(1+2\mathrm{cos}\theta ){\displaystyle \frac{𝐩^2}{\mathrm{\Lambda }^2}}]\},`$ (77)
and
$`M_{\text{d}High}={\displaystyle \frac{ig^4}{8\pi m\mathrm{\Theta }^2}}[\epsilon ^{acd}T_aT_dT_c]\left\{\mathrm{log}\left({\displaystyle \frac{\mathrm{\Lambda }^2}{4m^2}}\right)+{\displaystyle \frac{\mathrm{\Lambda }^2}{2m^2}}+(1+2\mathrm{cos}\theta ){\displaystyle \frac{𝐩^2}{\mathrm{\Lambda }^2}}\right\},`$ (79)
respectively. Summing Eqs. (LABEL:4.81) and (79) we get
$`M_\text{d}={\displaystyle \frac{ig^4}{8\pi m\mathrm{\Theta }^2}}[\epsilon ^{acd}T_aT_dT_c]\left\{1+i{\displaystyle \frac{\mathrm{sin}\theta }{1\mathrm{cos}\theta }}+\mathrm{log}\left({\displaystyle \frac{4m^2}{𝐪^2}}\right)\right\}.`$ (80)
Finally, incorporating the radiative correction, Figs. 7(e)-7(h), we obtain
$`_{\text{e-g}}(\theta )=`$ $``$ $`{\displaystyle \frac{ig^4}{4\pi m\mathrm{\Theta }^2}}[T^a(T^bT_b)T_a]`$ (81)
$``$ $`{\displaystyle \frac{ig^4}{8\pi m\mathrm{\Theta }^2}}[\epsilon ^{abc}T_aT_cT_b]\left\{1+i{\displaystyle \frac{\mathrm{sin}\theta }{(1\mathrm{cos}\theta )}}\right\}.`$ (82)
$`_\text{h}(\theta )={\displaystyle \frac{ig^4}{24\pi m\mathrm{\Theta }^2}}[T^aT_a].`$ (83)
Summing all these contributions and using the relation (4) to simplify the result, we get the total one-loop amplitude
$`_{1loop}(\theta )={\displaystyle \frac{ig^4}{4\pi m\mathrm{\Theta }^2}}\left\{{\displaystyle \frac{3}{8}}[II]+{\displaystyle \frac{2}{3}}[T^aT_a]\right\}.`$ (84)
## IV Conclusions
In this work we studied the scattering of isospin 1/2 fermionic particles interacting through a non-Abelian Chern-Simons field. In the nonrelativistic formulation we found that, up to a finite constant term, there is no one-loop correction to the tree approximation to the scattering amplitude. This is similar to what happens in the scalar theory where the constant one-loop contribution may be eliminated by a finite quartic counterterm .
We have also considered the same problem starting from the fully relativistic theory. After discussing the one-loop renormalizability of the model and determining anomalous contributions to the matrix magnetic moment of the fermions, we considered the low momenta limit of the two body scattering amplitude obtaining a nonvanishing one-loop contribution. This result, shown in Eq. (84), is a correction to the scattering which does not appear in the nonrelativistic theory. It is a leading order contribution and implies that the effective low momentum Lagrangian contains a four-fermion self interaction with a coupling which can be read from Eq. (84). These terms can not be eliminated by adding counterterms to the original Lagrangian (32) without destroying the renormalizability of the relativistic model. Furthermore, as also happens in the Abelian case, these new terms come from the high part of the original theory and could not be suspected in a direct nonrelativistic approach.
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# X-ray Spectral Diagnostics of Gamma-Ray Burst Environments
## 1 Introduction
The detection of a discrete spectral feature in the X-ray afterglow spectrum of GRB970508 was recently reported by Piro et al. (1999a,b). A similar feature in the afterglow of GRB970828 was reported by Yoshida et al. (1999). The redshift of the host galaxy of GRB970508 was determined to be $`z=0.835`$ (Metzger et al. 1997; Bloom et al. 1998). The apparent energy of the discrete feature in the X-ray spectrum is consistent with the redshifted energies of the $`n=21`$ transitions in all possible charge states of Fe, i.e., from Fe K$`\alpha `$ at 6.4 keV, to H-like Fe Ly$`\alpha `$ at 6.95 keV. Therefore, none of the possible line excitation mechanisms (fluorescence, collisional excitation in hot gas, or recombination in photoionized gas or a transient collisional plasma) is currently ruled out on the basis of the measured line energy alone.
The conditions required for each of these emission mechanisms (i.e. the presence of a given charge state of Fe) are all still compatible with the constraints on ionization parameter and gas temperature that can be derived from the fact that the line source must be located no more than about 1 light day from the ionizing source, and from the shape of the ionizing spectrum. All three mechanisms require high gas density ($`n_e>10^{11}`$ cm<sup>-3</sup>) for bound Fe to exist at all this close to the burst site. The presence of a large mass of dense gas close to the burst has led to the suggestion that a merging neutron star-neutron star binary is an unlikely site for this burst, and that a scenario involving some sort of stellar collapse is more likely (Piro et al. 1999b). In that case, the line source may be associated with debris from the stellar collapse.
Resolved X-ray spectroscopy can of course distinguish between the various line emission mechanisms, and lead to the correct characterization of the conditions in the source and the implied mass of radiating Fe. In the following, we will briefly recapitulate the constraints imposed by proximity and the energy budget, and then draw attention to an example of direct spectroscopic analysis of the X-ray spectrum: the fact that the measured centroid position, and to some extent even the shape, of the X-ray emission feature in the afterglow of GRB970508 already imply that recombination in afterglow-photoionized gas is unlikely to be the emission mechanism for this burst.
Our analysis applies to an optically thin, homogeneous medium. Complications arising from radiative transfer effects or external heat input by a relativistic shock may conspire to alter the shape of the spectrum, perhaps to the point that no definitive conclusions can be drawn from an intantaneous, low sensitivity spectrum. Such effects have been calculated by Weth et al. (1999) for the case of photoionization in the context of two specific geometries for the burst site, and by Vietri (1999) for a scenario in which the relativistic shock heats the medium shortly after the onset of the afterglow. With better data, it may hopefully be possible to reverse the argument, and infer the geometry, thermal history, and abundances from spectroscopy.
## 2 Ionization and Thermal Conditions in the GRB970508 X-ray Line Source
The discrete feature appears in the BeppoSAX Medium Energy Concentrator Spectrometer (MECS) spectrum immediately before a sudden rise in the afterglow flux, beginning at $`6\times 10^4`$ sec after the burst, and disappears thereafter (Piro et al. 1999b). If one attributes the feature to emission by Fe, one requires bound Fe to be present close to the source, which emits a very large flux of ionizing photons in the afterglow ($`L_{1\mathrm{Ryd}10\mathrm{keV}}10^{51}(t/1\mathrm{sec})^{1.1}`$ erg s<sup>-1</sup>; here, $`t`$ is the time since the burst \[Piro et al. 1999b\]).
The prompt burst and afterglow are intense enough that the photoionization timescale, $`t_{\mathrm{ion}}`$, is extremely short compared to the burst duration and the time during which the emission feature is visible. Specifically, the (inverse of the) ionization timescale for Fe in the afterglow radiation field is
$`t_{\mathrm{ion}}^1`$ $`=`$ $`{\displaystyle _\chi ^{\mathrm{}}}𝑑EF(E)E^1\sigma (E)=`$
$`=`$ $`1.4\times 10^4(t/1\mathrm{sec})^{1.1}r_{16}^2\mathrm{sec}^1,`$
with $`F(E)`$ the ionizing flux at the line source, $`\sigma (E)6\times 10^{18}Z^2(\chi /E)^3`$ cm<sup>2</sup> the K-shell photoionization cross section for neutral atoms of nuclear charge $`Z`$, $`\chi `$ the ionization potential, and $`r_{16}`$ the distance to the burst site in units $`10^{16}`$ cm. We have assumed a $`E^2`$ photon number spectrum to calculate $`F(E)`$ from the $`210`$ keV afterglow flux quoted by Piro et al. (1998), converted to luminosity assuming $`H_0=75\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and $`q_0=1/2`$, and we extrapolate the luminosity to 1 MeV. Ionization is essentially instantaneous.
The recombination timescale is given by $`t_{\mathrm{rec}}(n_e\alpha (T_e))^1`$, where $`\alpha (T_e)`$ is the recombination coefficient, $`n_e`$ the electron density, and $`T_e`$ the electron temperature. Since $`\alpha `$ depends on the temperature, we first need to consider the thermal evolution of the source subject to irradiation by the prompt burst and afterglow radiation fields.
Given sufficient interaction time, the gas will relax to the Compton temperature, $`T_\mathrm{C}`$, at which the Compton heating and cooling rates equal each other, and which is therefore determined only by the shape of the ionizing spectrum. In the non-relativistic limit (both the photon energies $`E`$ and the electron energies $`kT_em_ec^2`$), the Compton temperature is given by
$$kT_\mathrm{C}=\frac{1}{4}\frac{𝑑EEF(E)}{𝑑EF(E)}$$
(2)
(Ross 1979). As either the photon energies or the gas temperature approach the electron rest energy, the correct relativistic scattering cross section should be used, and the energy exchange between photons and electrons should be calculated to higher order than linear in $`E/m_ec^2`$ and $`kT_e/m_ec^2`$. In order to obtain a very rough estimate, we will assume here that we can mimic the effects of these modifications by simply cutting off the integrals at photon energies of order 1 MeV. If we assume a prompt burst spectrum of the form $`F(E)E^{1/2}`$ at low energies, and $`F(E)E^{3/2}`$ above a break energy of order 1 MeV (Conners et al. 1998), the Compton temperature becomes approximately $`kT_\mathrm{C}(\mathrm{MeV})(3/16)(E_{\mathrm{max}}/1\mathrm{MeV})^{1/2}`$, with $`E_{\mathrm{max}}`$ the cutoff energy. For $`E_{\mathrm{max}}`$ in the few MeV range, the Compton temperature is therefore in the several hundred keV range. The timescale for the Compton interactions to equilibrate should be of order (Rybicki & Lightman 1979)
$$t_{\mathrm{Compton}}=\frac{1}{n_e\sigma _\mathrm{T}c}\frac{m_ec^2}{4kT_e}7\times 10^4n_{10}^1T_8^1\mathrm{sec}$$
(3)
with $`\sigma _\mathrm{T}`$ the Thomson cross section. Unless the density is much higher than $`10^{10}`$ cm<sup>-3</sup>, this timescale is much longer than the burst timescale, which implies that the Compton temperature will not be reached. Moreover, the average photon energy in the burst spectrum is of order 1 MeV, which produces mildly relativistic electrons. The stopping timescale of these fast particles on stationary electrons is (e.g. Longair 1992)
$$t_{\mathrm{e}\mathrm{e}}\frac{\gamma m_e^2c^3}{2\pi e^4n_e\mathrm{ln}\mathrm{\Lambda }}=1.3\times 10^3n_{10}^1(\gamma /2)((\mathrm{ln}\mathrm{\Lambda })/10)^1\mathrm{sec},$$
(4)
(with $`\gamma `$ the electron Lorentz factor, and $`\mathrm{ln}\mathrm{\Lambda }`$ the Coulomb logarithm) which leads to the interesting conclusion that the plasma may not be in equilibrium at the end of the burst.
In any case, as the (much softer) afterglow begins, the plasma will rapidly cool by inverse Compton scattering, bremsstrahlung, and, at low temperatures, atomic emission. For an initial gas temperature of order $`10^8`$ K or hotter, the Comptonization timescale is of the order of, or shorter than the duration of the afterglow, and equilibrium is established at the afterglow Compton temperature, which for an $`E^2`$ photon spectrum is equal to $`kT_\mathrm{C}=(1/4)E_{\mathrm{max}}/\mathrm{ln}(E_{\mathrm{max}}/E_{\mathrm{min}})`$, with $`E_{\mathrm{min}}`$ the lowest photon energy in the afterglow spectrum; $`E_{\mathrm{max}}`$ is the upper end of the integration range, $`E_{\mathrm{max}}`$ 1 MeV. Assuming the afterglow spectrum flattens below $`E_{\mathrm{min}}<1`$ Ryd (Wijers & Galama 1999), we find $`T_\mathrm{C}<2\times 10^8`$ K. Kallman and McCray (1979) in their numerical models for X-ray photoionized nebulae, actually find that the electron temperature saturates at a lower temperature of approximately $`T_\mathrm{C}10^7`$ K, due to bremsstrahlung cooling, so most likely the gas is cooler than $`10^8`$ K in the afterglow.
The recombination coefficient as a function of temperature for recombination onto a bare nucleus is given by Seaton (1959); a powerlaw fit to the coefficient is $`\alpha 1.26\times 10^6T_e^{0.75}`$ cm<sup>3</sup> s<sup>-1</sup> (for Z = 26), accurate to $`<20\%`$ in the range $`T_e=10^52\times 10^8`$ K. Thus, the recombination timescale is
$$t_{\mathrm{rec}}=1.4n_{11}^1T_7^{0.75}\mathrm{sec}.$$
(5)
We conclude that at sufficiently high densities ($`n_{11}>1`$), the line source will be relatively cool, and close to ionization equilibrium at all times during the afterglow.
Assuming ionization equilibrium, we can derive a constraint on the density from the requirement that bound Fe be present in the line source. For Fe to recombine to the H-like charge state requires an ionization parameter in equilibrium of $`\xi L/nr^2<10^4`$, with $`L`$ the ionizing luminosity, $`n`$ the particle density, and $`r`$ the distance to the ionizing source. This value of the ionization parameter applies to an $`E^2`$ ionizing photon spectrum (Kallman & McCray 1982, their model 7). Inserting numbers and extrapolating the ionizing luminosity to 1 MeV, one infers a lower limit to the density at a distance of $`6\times 10^4`$ lt sec from the source, of
$$n>1.7\times 10^{11}(t/6\times 10^4\mathrm{sec})^{1.1}(r/6\times 10^4\mathrm{lt}\mathrm{sec})^2\mathrm{cm}^3$$
(6)
Similarly, for Fe to be recombined to Li-like or less ionized, necessary for fluorescence to operate, $`\xi <1000`$ in equilibrium, yielding a ten times higher density than the above estimate.
## 3 X-ray Spectroscopy
From the previous we conclude that if the interaction with the afterglow radiation is the only source of heating in the line emitting gas, then, later in the afterglow, the source is likely to be dense ($`>10^{11}`$ cm<sup>-3</sup>) and relatively cool. Line emission can be excited both by fluorescence and cascading following recombination. Only if there is an additional source of heat, sufficient to heat the gas to $`T>10^8`$ K, will collisional excitation play a role in driving Fe $`n=21`$ lines.
Obviously, X-ray spectroscopy of the discrete emission strongly distinguishes between these various possible emission mechanisms. If one is confident of the redshift of the emission complex, fluorescence can be distinguished from emission by highly ionized Fe, simply from the apparent energy of the emission. Unfortunately, the resolution of the MECS, and the statistical quality of the spectra under discussion are not sufficient to allow this distinction to be made.
Recombination and collisional excitation, however, can readily be distinguished in the following way. For a given degree of ionization, a photoionized medium in equilibrium is much cooler than the corresponding collisional plasma, in which the electron temperature has to be of order the ionization potential in order to support a given charge state. The low electron temperature in photoionization equilibrium has a dramatic spectroscopic signature. Since most of the free electrons have small kinetic energies compared to the ionization potential, photons resulting from radiative recombination will have a narrow energy distribution, bunched up just above the ionization potential, with a typical width of order $`\mathrm{\Delta }EkT_e\chi `$. In a collisional plasma on the other hand ($`kT_e\chi `$), the recombination rate is much reduced, and in addition, the recombination photons are spread out over a wide energy range above the ionization potential. The narrow radiative recombination continuum (RRC) in a photoionized source is easy to detect, because it contains an integrated photon flux comparable to that in the $`n=21`$ discrete transition; in fact, for H-like Fe, the ratio between the fluxes in the RRC to that in the Ly$`\alpha `$ line is approximately equal to $`0.93(kT_e/1\mathrm{keV})^{0.21}`$. So far, narrow RRC’s have only been seen in the spectrum of the massive binary Cyg X-3 (Liedahl & Paerels 1996; Kawashima & Kitamoto 1996).
## 4 Application to the Spectrum of GRB980508
We use the dataset shown by Piro et al. (1999b; their dataset 1a). At the resolution of the MECS ($`\mathrm{\Delta }E340`$ eV at 3 keV), the RRC is barely resolvable if the electron temperature is of order or less than 1 keV. For simplicity we use a model for the emission spectrum of cool, photoionized gas consisting of a narrow line plus exponentially decaying RRC, to represent the Ly$`\alpha `$ emission line and RRC at 6.95 and 9.28 keV, respectively, with equal photon fluxes. The He-like recombination spectrum would look much the same at the MECS resolution. To this, we added a simple power law with absorption by neutral gas, to represent the continuum. We let the redshift float, as well as the continuum parameters.
By fixing the electron temperature (i.e., the width of the RRC) and fitting for the remaining free parameters we obtain the minimum $`\chi ^2`$ values displayed in Figure 1. At either very low $`T_e`$ ($`kT_e<0.2`$ keV) or very high $`T_e`$ ($`kT_e>4`$ keV) we obtain $`\chi ^2`$ values of order $`\chi _{\mathrm{min}}^27`$ for 5 degrees of freedom. In themselves these values are acceptable (given the small number of degrees of freedom). However, at low temperatures the model has a best fitting redshift of $`z>1.6`$ (for Hydrogenic Fe; 68% confidence for one parameter of interest only). At high temperatures, the RRC becomes essentially undetectable and the model consists effectively of a single narrow line. This model will of course fit the data, but at these temperatures the plasma would be closer to collisional equilibrium, which is not consistent with the assumption we set out to test. Finally, at $`kT_e1`$ keV (which is about the temperature for a source in photoionization equilibrium with the afterglow), we find $`\chi _{\mathrm{min}}^212`$, which indicates a poor fit ($`\chi _{\mathrm{min}}^2`$ about two standard deviations away from the expected $`\chi _{\mathrm{min}}^2=5`$). By itself, this is probably not enough to confidently reject the model, and the fit can be improved by changing the temperature somewhat. But the implied redshift is still $`z1.20`$, higher than for a single line because we have an additional emission component at $`4/3`$ times the line energy. Unless we are willing to entertain the possibility that the line source is actually behind the galaxy at $`z=0.835`$, we conclude that the emission spectrum from a cool photoionized source is incompatible with the measured spectrum. In Figure 1, we have indicated the (68% confidence) limits on the redshift as a function of the assumed temperature of the gas, assuming H-like emission. At high temperature, the contrast in the RRC becomes very small, and the spectrum is dominated by a single line. In this limiting case, the implied redshift approaches $`1.05\pm 0.05`$ (68% confidence), appropriate for H-like Fe Ly$`\alpha `$. Widening the confidence interval to 90% confidence will produce agreement with the optical redshift of $`z=0.835`$ (Piro et al. 1999b).
Figure 2 displays the data and the lowest$`\chi ^2`$ recombination model with $`kT_e=1`$ keV; the values for the other spectral parameters are: power law photon index $`3.1`$, column density $`N_H=7.1\times 10^{21}`$ cm<sup>-2</sup>. This continuum spectral shape appears to be quite a bit steeper than the best fit given by Piro & al. 1999b, but in fact, the uncertainty on the index is large ($`1.1`$ for one parameter of interest, at 1$`\sigma `$), and so the continuum shapes are in fact consistent.
## 5 Conclusion
We find that Fe emission from a cool photoionized source in equilibrium with the afterglow cannot account for the X-ray spectrum of GRB970508. Fitting a spectral model for such a source either implies a large redshift, or indicates such high electron temperatures that the source would be closer to collisional ionization equilibrium.
This would imply that the discrete feature in the spectrum of GRB970508 arises from either Fe fluorescence or from collisional excitation, each of which implies a very different estimate for the total mass of radiating Fe ($`M_{\mathrm{Fe}}2.6\times 10^4M_{}`$ for fluorescence, $`M_{\mathrm{Fe}}7\times 10^2M_{}`$ for collisional excitation, which corresponds to $`M_{\mathrm{total}}8\times 10^2M_{}(A_{\mathrm{Fe}}/A_{\mathrm{Fe},})^1`$ and $`M_{\mathrm{total}}28M_{}(A_{\mathrm{Fe}}/A_{\mathrm{Fe},})^1`$, respectively (see, for instance, Meszaros & Rees , Piro et al. \[1999b\], Lazzati et al. ); here, $`A_{\mathrm{Fe}}`$ is the abundance of Fe).
With higher spectral resolution and sensitivity, we may be able to distinguish spectroscopically between a fluorescent and a collisional spectrum (a recombination spectrum is readily identified from the presence of the RRC). Working from the K-shell spectrum of a single ionization stage of Fe alone, one might use the $`n=21`$ and $`n=31`$ transitions, which are resolved at $`E/\mathrm{\Delta }E10`$. Intensity ratios do not effectively discriminate between collisional excitation and fluorescence: $`I(\mathrm{Ly}\beta )/I(\mathrm{Ly}\alpha )=0.100.15`$ (Hydrogenic Fe) for electron temperatures between $`3\times 10^7`$ and $`10^9`$ K (Mewe, Gronenschild, & van den Oord 1985), while for neutral Fe $`I(\mathrm{K}\beta )/I(\mathrm{K}\alpha )=0.13`$ (Kaastra & Mewe 1993). Instead, with somewhat higher resolution, one might use the fact that the ratio of the transition energies depends on the charge state: $`E(\mathrm{Ly}\beta )/E(\mathrm{Ly}\alpha )=1.18`$, while $`E(\mathrm{K}\beta )/E(\mathrm{K}\alpha )=1.10`$.
At higher resolution, one would resolve the $`1s2p_{1/2}`$ and $`1s2p_{3/2}`$ transitions ($`E/\mathrm{\Delta }E>500)`$. The ratios of the energies are $`E(\mathrm{Ly}\alpha _2)/E(\mathrm{Ly}\alpha _1)=1.0030`$ and $`E(\mathrm{K}\alpha _2)/E(\mathrm{K}\alpha _1)=1.0020`$, and a resolving power of 1000 is required to distinguish between these two.
A collisional source at $`kT<15`$ keV will also show emission from He-like Fe, and the simultaneous appearance of both the H- and the He-like charge states readily indicates a high degree of ionization of the source. For temperatures in the range $`3<kT<15`$ keV, the He-like emission dominates, which could still be identified as such if the characteristic ”triplet” structure of the $`n=21`$ lines is resolved ($`E/\mathrm{\Delta }E>300`$). The highest resolution available at energies above 2 keV is provided by the High Energy Transmission Grating Spectrometer on Chandra (Canizares et al. 2000). The Fe K spectrum will be resolved at $`E/\mathrm{\Delta }E150(1+z)`$, which would at least enable charge-state spectroscopy at moderately large redshift.
Finally, given sufficient sensitivity at low energies (and sufficiently low absorption in the source), we may of course detect emission from lower-$`Z`$ elements, either collisional or fluorescent, which would greatly facilitate the spectroscopic diagnosis of the excitation mechanism. With the grating spectrometers on Chandra and XMM, standard detailed plasma diagnostics can be invoked, at resolving powers of $`E/\mathrm{\Delta }E1001500`$ at $`E<2`$ keV, depending on photon energy (Canizares et al. 2000, Brinkman et al. 2000, Brinkman et al. 1998). Clearly, it would be very interesting to pursue more sensitive X-ray spectroscopy of gamma ray burst afterglows.
We gratefully acknowledge discussions with Marten van Kerkwijk. D. A. L. was supported in part by a NASA Long Term Space Astrophysics Program grant (LTSA S-92654-F). Work at LLNL was performed under the auspices of the U. S. Department of Energy, Contract No. W-7405-Eng-48.
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# References
IASSNS-HEP-00/35
NSF-ITP-00-36
PUPT-1927
Matrix Theory Interpretation of DLCQ String Worldsheets
G. Grignani <sup>a</sup><sup>1</sup><sup>1</sup>1Supported by INFN and MURST of Italy. Permanent Address: Dipartimento di Fisica and Sezione I.N.F.N., Università di Perugia, Via A. Pascoli I-06123, Perugia, Italia., P. Orland <sup>b</sup><sup>2</sup><sup>2</sup>2Supported by PSC-CUNY Research Award Program Grants nos. 668460, 61508-00, 69466-00, a CUNY Collaborative Incentive grant 91915-00-06 and the National Science Foundation under Grant No. PHY94-07194. Permanent address: Center for Theoretical Physics, The Graduate School and University Center and Baruch College, The City University of New York, New York, NY., L. D. Paniak <sup>c</sup><sup>3</sup><sup>3</sup>3Supported by NSF grant PHY98-02484 and NSERC of Canada. and G. W. Semenoff <sup>a</sup><sup>4</sup><sup>4</sup>4 Supported by the Ambrose Monell Foundation and NSERC of Canada. On leave from: Department of Physics and Astronomy, University of British Columbia, Vancouver, V6T 1Z1 Canada.
* Institute for Advanced Study, Einstein Drive, Princeton, NJ 08540.
* Institute for Theoretical Physics, University of California, Santa Barbara, CA 93106-4030.
* Joseph Henry Laboratories, Department of Physics, Princeton University,
Princeton, NJ 08544.
Abstract
We study the null compactification of type-IIA-string perturbation theory at finite temperature. We prove a theorem about Riemann surfaces establishing that the moduli spaces of infinite-momentum-frame superstring worldsheets are identical to those of branched-cover instantons in the matrix-string model conjectured to describe M-theory. This means that the identification of string degrees of freedom in the matrix model proposed by Dijkgraaf, Verlinde and Verlinde is correct and that its natural generalization produces the moduli space of Riemann surfaces at all orders in the genus expansion.
It is widely believed that each of the five known consistent string theories are limits of a single eleven-dimensional theory called M-theory. While this theory has not yet been fully mathematically formulated, there is an interesting proposal, namely the matrix model . This model is conjectured to describe M-theory in a particular kinematical region, the infinite-momentum frame. In the matrix model, the superstring is a composite object resembling a necklace of D0-branes.
A nontrivial check of the matrix model proposal would be to use it to obtain perturbative string theory. There are already convincing arguments that the non-interacting type-IIA strings emerge in the appropriate limit of the matrix model . There are several approaches to deriving superstring interactions along these lines -. In these, matrix-theory instantons interpolate between initial states and final states of strings through a Riemann surface which is a branched cover of the cylinder. Whether these branched-cover instantons account correctly for string-scattering amplitudes is an important question.
In this Letter, we show for the first time that in the appropriate context, perturbative string theory can be formulated using only branched-cover Riemann surfaces.
The matrix model is a formulation of M-theory in the infinite-momentum frame. Compactification of a spatial direction of M-theory produces type-IIA string theory. The matrix model becomes 1+1-dimensional, maximally-supersymmetric Yang-Mills theory.
To compare the matrix model directly with infinite-momentum-frame string theory, we study the string path integral with a compactified null direction. The natural quantization of the string in this frame is discrete light-cone quantization (DLCQ). We find it necessary to introduce a finite temperature by further compactifying Euclidean time. We prove a new theorem on Riemann surfaces, stating that to any order in the genus expansion, these compactifications restrict the worldsheets to branched covers of a torus. This differs from the moduli space of strings without these compactifications, which includes all Riemann surfaces up to conformal diffeomorphisms.
The same branched covers appear in the matrix model. According to Dijkgraaf, Verlinde and Verlinde , the string degrees of freedom are simultaneous eigenvalues of the matrices. At finite temperature, the matrices are defined on a torus and their eigenvalues, since they solve polynomial equations, are functions on branched covers of the torus. If the matrix model is to agree with perturbative string theory, these branched covers must be the full set of Riemann surfaces that contribute to the string path integral.
We compare only the degrees of freedom of the two theories and not the energy spectra. However, to one-loop order, the spectra are known to coincide . The coincidence of energy spectra once higher order corrections are included is still an open question.
The string path integral for the vacuum energy is
$$F=\underset{g,\sigma }{}g_s^{2g2}[dh_gdXd\mathrm{\Psi }]\mathrm{exp}\left(\frac{1}{4\pi \alpha ^{}}\sqrt{h}\left(h^{ab}_aX^\mu _bX^\mu 2\pi i\alpha ^{}\mathrm{\Psi }^\mu \gamma \mathrm{\Psi }^\mu \right)\right)$$
(1)
Here, we use, for example, the Neveu-Schwarz-Ramond superstring<sup>5</sup><sup>5</sup>5Most of our considerations apply to the bosonic sector of any string theory. We use the superstring as an example. For the relationship with matrix theory, however, supersymmetry is important and that case is more closely related to the Green-Schwarz superstring.. The string coupling constant is $`g_s`$ and its powers weight the genus, $`g=0,1,\mathrm{}`$, of the string’s worldsheet. There is also a sum over spin structures, $`\sigma `$ which, with the appropriate weights, imposes the GSO projection. For each value of the genus, $`g`$, $`[dh_g]`$ is an integration measure over all metrics of that genus and is normalized by dividing out the volume of the worldsheet re-parameterization and Weyl groups. We will assume that the metrics of both the worldsheet and the target spacetime have Euclidean signatures.
We wish to study the situation where the target space has particular compact dimensions. Two compactifications will be needed. The first compactifies the light-cone in Minkowski space by making the identification $`\frac{1}{\sqrt{2}}\left(tx^9\right)\frac{1}{\sqrt{2}}\left(tx^9\right)+2\pi R`$. In our Euclidean coordinates it is the identification
$$(X^0,\stackrel{}{X},X^9)(X^0+\sqrt{2}\pi iR,\stackrel{}{X},X^9\sqrt{2}\pi R)$$
(2)
With this compactification the GSO projection is unmodified. The factor of $`i`$ in the identification of $`X^0`$ might seem unnatural since it identifies a real integration variable periodically in a complex direction. However, we shall see that, for the path integral at genus 1, where we can check the result independently by using operator methods to compute the same partition function, this identification is indeed the correct thing to do. We shall postulate that it also gives the correct partition function at genus greater than one.
The second compactification that we shall need is that of Euclidean time,
$$(X^0,\stackrel{}{X},X^9)(X^0+\beta ,\stackrel{}{X},X^9)$$
(3)
This compactification, with the appropriate modification of the GSO projection to make space-time fermions anti-periodic, introduces a temperature, $`T=1/k_B\beta `$ where $`k_B`$ is Boltzmann’s constant, so that (1) computes the thermodynamic free energy.
In order to implement this compactification in the path integral, we assume that the worldsheet is a Riemann surface $`\mathrm{\Sigma }_g`$ of genus $`g`$ whose homology group $`H_1(\mathrm{\Sigma }_g)`$ is generated by the closed curves,
$`a_1,a_2,\mathrm{},a_g,b_1,b_2,\mathrm{},b_g`$
$`a_ia_j=\mathrm{},b_ib_j=\mathrm{},a_ib_j=\delta _{ij}`$ (4)
Furthermore, one may pick a basis of holomorphic differentials $`\omega _iH^1(\mathrm{\Sigma }_g)`$ with the properties
$$_{a_i}\omega _j=\delta _{ij},_{b_i}\omega _j=\mathrm{\Omega }_{ij}$$
(5)
where $`\mathrm{\Omega }`$ is the period matrix. It is complex, symmetric, $`\mathrm{\Omega }_{ij}=\mathrm{\Omega }_{ji}`$, and has positive definite imaginary part.
Compactification is implemented by including the possible windings of the string worldsheet on the compact dimensions. These form distinct topological sectors in the path integration in (1). In the winding sectors, the bosonic coordinates of the string should have a multi-valued part which changes by $`\beta `$integer or $`(i)\sqrt{2}R`$integer as it is moved along a homology cycle. The derivatives of these coordinates should be single-valued functions. It is convenient to consider their exterior derivatives which can be expressed as linear combinations of the holomorphic and anti-holomorphic 1-forms and exact parts,
$$dX^0=\underset{i=1}{\overset{g}{}}\left(\lambda _i\omega _i+\overline{\lambda }_i\overline{\omega }_i\right)+\mathrm{exact},dX^9=\underset{i=1}{\overset{g}{}}\left(\gamma _i\omega _i+\overline{\gamma }_i\overline{\omega }_i\right)+\mathrm{exact}$$
(6)
Then, we require
$$_{a_i}𝑑X^0=\beta n_i+\sqrt{2}\pi Rip_i,_{b_i}𝑑X^0=\beta m_i+\sqrt{2}\pi Riq_i$$
(7)
$$_{a_i}𝑑X^9=\sqrt{2}\pi Rp_i,_{b_i}𝑑X^9=\sqrt{2}\pi Rq_i$$
(8)
With (5), we use these equations to solve for the constants in (6). With the formula
$$\omega _i\overline{\omega }_j=\underset{k=1}{\overset{g}{}}\left(_{a_k}\omega _i_{b_k}\overline{\omega }_j_{b_k}\omega _i_{a_k}\overline{\omega }_j\right)=2i\left(\mathrm{\Omega }_2\right)_{ij}$$
(9)
we compute the part of the string action which contains the winding integers,
$`S={\displaystyle \frac{\beta ^2}{4\pi \alpha ^{}}}(n\mathrm{\Omega }^{}m)\mathrm{\Omega }_2^1(\mathrm{\Omega }nm)+2\pi i{\displaystyle \frac{\sqrt{2}\beta R}{4\pi \alpha ^{}}}{\displaystyle \frac{1}{2}}[(p\mathrm{\Omega }^{}q)\mathrm{\Omega }_2^1(\mathrm{\Omega }nm)`$
$`+(n\mathrm{\Omega }^{}m)\mathrm{\Omega }_2^1(\mathrm{\Omega }pq)]+\mathrm{}`$ (10)
Note that the integers $`p_i`$ and $`q_i`$ appear linearly in a purely imaginary term in the action. Furthermore, since they come from the compactification of the light cone, this is the only place that they will appear in the string path integral (unlike $`m_i`$ and $`n_i`$ which should appear in the weights of the sum over spin structures). When the action is exponentiated and summed over $`p_i`$ and $`q_i`$, the result will be periodic Dirac delta functions. It can be shown that these delta functions impose a linear constraint on the period matrix of the worldsheet. Thus, with the appropriate Jacobian factor, the net effect is to insert into the path integral measure the following expression,
$$\underset{mnrs}{}e^{\frac{\beta ^2}{4\pi \alpha ^{}}\left(n\mathrm{\Omega }^{}m\right)\mathrm{\Omega }_2^1\left(\mathrm{\Omega }nm\right)}\nu ^{2g}\left|det\mathrm{\Omega }_2\right|\underset{j=1}{\overset{g}{}}\delta \left(\underset{i=1}{\overset{g}{}}\left(n_i+i\nu r_i\right)\mathrm{\Omega }_{ij}\left(m_j+i\nu s_j\right)\right)$$
(11)
where $`\nu =4\pi \alpha ^{}/\sqrt{2}\beta R`$ is a fixed constant. Consequently, the integration over metrics in the string path integral is restricted to those for which the period matrix obeys the constraint
$$\underset{i=1}{\overset{g}{}}\left(n_i+i\nu r_i\right)\mathrm{\Omega }_{ij}\left(m_j+i\nu s_j\right)=0$$
(12)
for all combinations of the $`4g`$ integers $`m_i,n_i,r_i,s_i`$ such that $`\mathrm{\Omega }`$ is in a fundamental domain.
Since the columns of the period matrix are linearly independent vectors, these are $`g`$ independent complex constraints on the moduli space of $`\mathrm{\Sigma }_g`$. Thus its complex dimension $`3g3`$ is reduced to $`2g3`$ and there is further discrete data contained in the integers. One would expect that, when the compactifications are removed, either $`\beta \mathrm{}`$ or $`R\mathrm{}`$, the discrete data assembles itself to a “continuum limit” which restores the complex dimension of moduli space.
It is interesting to ask whether the Riemann surfaces with the constraint (12) can be classified in a sensible way. The answer to this question is yes, a Riemann surface obeys the constraint (12) if and only if it is a branched cover of the torus, $`T^2`$, with Teichmüller parameter $`i\nu `$. This is established through the
Theorem: $`\mathrm{\Sigma }_g`$ is a branched cover of $`T^2`$ if and only if the period matrix obeys (12), for some choice of integers $`m_i,n_i,r_i`$ and $`s_i`$.
Proof: The generators of the first homology group of $`T^2`$ are two closed loops $`(\alpha ,\beta )`$ which span the vector space $`H_1(T^2,\text{})`$. The dual vector space, the first cohomology group, $`H^1(T^2,\text{})`$ is spanned by the basis of holomorphic and anti-holomorphic differentials $`\gamma `$ and $`\overline{\gamma }`$. They can be normalized as,
$$_\alpha \gamma =1,_\beta \gamma =i\nu \mathrm{and}_\alpha \overline{\gamma }=1,_\beta \overline{\gamma }=i\nu $$
(13)
The Riemann surface $`\mathrm{\Sigma }_g`$ is a branched cover of $`T^2`$ if there exists a continuous, onto, holomorphic map $`f`$, such that
$$\mathrm{\Sigma }_g\stackrel{f}{}T^2$$
(14)
The map $`f`$ takes closed loops on $`\mathrm{\Sigma }_g`$ to closed loops on $`T^2`$. In particular, the generators (4) must map as
$$(a_i,b_j)\stackrel{f}{}(n_i\alpha +r_i\beta ,m_j\alpha +s_j\beta )$$
(15)
for some integers $`m_i,n_i,r_i,s_i`$. This gives a mapping between the vector spaces $`H_1(\mathrm{\Sigma }_g,\text{})`$ and $`H_1(T^2,\text{})`$. A mapping of vector spaces induces a pull-back on the dual vector spaces
$$H^1(T^2,\text{})\stackrel{f^{}}{}H^1(\mathrm{\Sigma }_g,\text{})$$
(16)
defined by its action on the basis,
$`a_if^{}(\gamma )={\displaystyle _{a_i}}f^{}(\gamma )f(a_i)\gamma =n_i{\displaystyle _\alpha }\gamma +r_i{\displaystyle _\beta }\gamma =n_i+i\nu r_i`$ (17)
$`b_jf^{}(\gamma )={\displaystyle _{b_j}}f^{}(\gamma )f(b_j)\gamma =m_j{\displaystyle _\alpha }\gamma +s_j{\displaystyle _\beta }\gamma =m_j+i\nu s_j`$ (18)
Consider the particular elements of $`H_1(\mathrm{\Sigma }_g,\text{})`$,
$$c_j=\underset{i=1}{\overset{g}{}}a_i\mathrm{\Omega }_{ij}b_j$$
(19)
It can be checked from the definition of the period matrix (4) (and the fact that it is symmetric) that any holomorphic differential, $`\eta `$, on $`\mathrm{\Sigma }_g`$ has the property
$$c_j\eta =\underset{i=1}{\overset{g}{}}\mathrm{\Omega }_{ij}_{a_i}\eta _{b_j}\eta =0$$
(20)
The holomorphic nature of the mapping, $`f`$, guarantees that $`f^{}(\gamma )`$ is a holomorphic differential on $`\mathrm{\Sigma }_g`$. Then it follows that
$$0=c_jf^{}(\gamma )=f(c_j)\gamma =\underset{i=1}{\overset{g}{}}(n_i+i\nu r_i)\mathrm{\Omega }_{ij}(m_j+i\nu s_j)$$
(21)
which is the constraint on the period matrix in eq.(12).
To prove the converse, we must show that if the constraint (12) is satisfied, then a covering map, $`f`$, exists. We will demonstrate this by explicit construction. Consider the line integral of a linear combination of the holomorphic differentials on $`\mathrm{\Sigma }_g`$,
$$z(P)=_{P_0}^P\underset{k=1}{\overset{g}{}}\lambda _k\omega _k$$
(22)
with $`P_0`$ a fixed base-point. We wish to choose the coefficients $`\lambda _k`$ such that this integral defines a map from points $`P\mathrm{\Sigma }_g`$ to the torus $`z(P)T^2`$ whose holomorphic coordinates are the complex numbers $`z(P)`$ with the identification $`zz+p+i\nu q`$ where $`p`$ and $`q`$ are integers. The integral depends on the path of integration. If the path is changed by a combination of homology cycles of $`\mathrm{\Sigma }_g`$, $`k_ia_i+l_ib_i`$, the integral on the right-hand side of (22) changes by
$$\delta z=\underset{j}{}\lambda _j\left(k_j+\underset{i}{}\mathrm{\Omega }_{ji}l_i\right)$$
$`\lambda _k`$ must be chosen so that this change is commensurate with the periods of $`T^2`$. This can easily be done if $`\mathrm{\Omega }`$ obeys (12). Then, with the choice $`\lambda _i=n_i+i\nu r_i`$,
$$\delta z=\underset{j}{}(n_j+i\nu r_j)k_j+\underset{i}{}(m_i+i\nu s_i)l_i=\mathrm{integer}+i\nu \mathrm{integer}$$
and we have constructed an explicit covering map, $`z(P)`$. Q.E.D.
As a concrete example, the constraint (12) can be solved explicitly for genus one. The leading contribution to the free energy of the finite temperature type II superstring (without DLCQ) is from the torus amplitude and is given by
$`{\displaystyle \frac{F}{V}}={\displaystyle _{}}{\displaystyle \frac{d^2\tau }{\tau _2}}{\displaystyle \underset{mn}{}}e^{\frac{\beta ^2|n\tau m|^2}{4\pi \alpha ^{}\tau _2}}\left({\displaystyle \frac{1}{4\pi ^2\alpha ^{}\tau _2}}\right)^5{\displaystyle \frac{1}{4\left|\eta (\tau )\right|^{24}}}[(\theta _2^4\overline{\theta }_2^4+\theta _3^4\overline{\theta }_3^4+\theta _4^4\overline{\theta }_4^4)(0,\tau )+`$
$`+e^{i\pi (m+n)}(\theta _2^4\overline{\theta }_4^4+\theta _4^4\overline{\theta }_2^4)(0,\tau )e^{i\pi n}(\theta _2^4\overline{\theta }_3^4+\theta _3^4\overline{\theta }_2^4)(0,\tau )e^{i\pi m}(\theta _3^4\overline{\theta }_4^4+\theta _4^4\overline{\theta }_3^4)(0,\tau )]`$ (23)
where $`\theta _k(0,\tau )`$ are Jacobi theta functions and $`\eta (\tau )`$ is the Dedekind eta-function. $``$ is the fundamental domain of the torus,
$$\left\{\tau =\tau _1+i\tau _2\right|\frac{1}{2}<\tau _1\frac{1}{2};|\tau |1;\tau _2>0\}$$
(24)
The modification of this formula by the null compactification can be found using (11),
$`{\displaystyle \frac{F}{V}}={\displaystyle \underset{\tau }{}}{\displaystyle \frac{\nu ^2}{m^2+\nu ^2n^2}}e^{\frac{\beta ^2|n\tau m|^2}{4\pi \alpha ^{}\tau _2}}\left({\displaystyle \frac{1}{4\pi ^2\alpha ^{}\tau _2}}\right)^5{\displaystyle \frac{1}{4\left|\eta (\tau )\right|^{24}}}[(\theta _2^4\overline{\theta }_2^4+\theta _3^4\overline{\theta }_3^4+\theta _4^4\overline{\theta }_4^4)(0,\tau )+`$
$`+e^{i\pi (m+n)}(\theta _2^4\overline{\theta }_4^4+\theta _4^4\overline{\theta }_2^4)(0,\tau )e^{i\pi n}(\theta _2^4\overline{\theta }_3^4+\theta _3^4\overline{\theta }_2^4)(0,\tau )e^{i\pi m}(\theta _3^4\overline{\theta }_4^4+\theta _4^4\overline{\theta }_3^4)(0,\tau )]`$ (25)
where the the solution of (12) yields the discrete Teichmüller parameter,
$$\tau =\frac{m+i\nu s}{n+i\nu r}$$
and one should sum over the integers so that $`\tau `$ is in the fundamental domain, $``$.
In ref. it was shown that this formula can also be obtained by operator methods by finding the spectrum of the non-interacting type II superstring in DLCQ and explicitly computing the thermodynamic free energy by summing over energy states. This provides a strong check of the path integral technique for DLCQ we have used in the present paper.
Modular transformations and identities for theta functions can be used to rewrite (25) as the Hecke operator acting on the partition function of a superconformal field theory, with torus worldsheet and target space $`R^8`$:
$$\frac{F}{V}=\frac{1}{\sqrt{2}\pi R\beta }[e^{\beta /\sqrt{2}R}]\left[\left(\frac{1}{4\pi ^2\alpha ^{}\tau _2}\right)^4\frac{1}{\left|\eta (\tau )\right|^{24}}\left|\theta _4(0,\tau )\right|^8\right]_{\tau =1/i\nu }$$
(26)
A modular transform can be used to write this as
$$\frac{F}{V}=\frac{1}{\sqrt{2}\pi R\beta }[e^{\beta /\sqrt{2}R}]\left[\left(\frac{1}{4\pi ^2\alpha ^{}\tau _2}\right)^4\frac{1}{\left|\eta (\tau )\right|^{24}}\left|\theta _2(0,\tau )\right|^8\right]_{\tau =i\nu }$$
(27)
The factor in front is the ratio of volumes of $`R^8`$ and $`R^9\times S^1`$ with compactified light cone. The action of $`[p]`$ on a function $`\varphi (\tau ,\overline{\tau })`$ is defined by
$$[p]\varphi (\tau ,\overline{\tau })=\underset{N=0}{\overset{\mathrm{}}{}}p^N\underset{\stackrel{kr=N,r\mathrm{odd}}{s\mathrm{mod}k}}{}\frac{1}{N}\varphi (\frac{r\tau +s}{k},\frac{r\overline{\tau }+s}{k})$$
(28)
In , this formula was shown to arise in the $`g_s0`$ limit of matrix string theory at finite temperature. In this limit, matrix string theory reduces to a theory of eigenvalues of matrices which naturally live on covers of the torus. Generally, these should be branched covers. In the leading order in $`g_s`$ only the covers without branch points contribute. It was shown that the thermodynamic free energy of the matrix model arising from summing over unbranched covers is identical to (27). The combinatorics of enumerating them is elegantly accounted for by the Hecke operator.
In matrix string theory the limit $`g_s0`$ corresponds to large gauge coupling. This limit projects the theory onto the zeros of the superpotential. These zeros occur when all of the matrices are simultaneously diagonalizable. This is why the matrix model reduces to a theory of eigenvalues . Beyond the leading order in strong coupling, the quadratic fluctuations of off-diagonal parts of the matrices can also be analyzed . It is found that the fluctuation determinants are ultra-local operators and almost cancel due to supersymmetry. Arguments are given that, when carefully treated, the gauge field sector produces the power of the string coupling constant $`g_s^{2g2}`$ accompanying the branch covers of genus $`g`$ (for a discussion of gauge fields, see refs.). The same argument would apply to an analysis of the matrix theory at finite temperature. Even though the supersymmetry is broken by temperature boundary conditions, since the fluctuations of non-diagonal fields are governed by ultra-local operators, the same supersymmetric cancellations should occur. The remaining theory of diagonal matrices is identical to the finite temperature Green-Schwarz superstring with worldsheets which are branched covers of the torus. In this paper we have shown that this is exactly what is realized in the DLCQ of string theory at finite temperature.
What must be done to demonstrate that this limit of matrix string theory describes the perturbative type-II superstring? It must be shown that the integration measures over the worldsheets in the string theory and over the collective coordinates of the branched covers, which can be regarded as instantons in matrix theory, are identical.
The authors thank M. Goresky, J. Harvey, I. Klebanov, J. Kollár, Y. Matsuo, D. Morrison, E. Verlinde and K. Zarembo for helpful conversations.
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# Multifractality in uniform hyperbolic lattices and in quasi–classical Liouville field theory
## I Introduction
The concept of multifractality consists in a scale dependence of critical exponents . It has been widely discussed in the literature in the context of various problems such as, for example, statistics of strange sets , diffusion limited aggregation , wavelet transforms , conformal invariance . This concept also proves to be useful in the context of disordered systems . It was recently found that the ground state wave function of two dimensional Dirac fermions in a random magnetic field has a multifractal behavior. The field theoretic investigation of the multifractality has been undertaken in the papers , while different interpretations of these field theoretic results from a geometrical and physical points of view were presented in and correspondingly. This problem was recently reanalyzed in the more general setting of systems caracterized by logarithmic correlations .
Our work is mainly inspired by the approach developed in where the authors obtain the multifractal exponents of the critical wave function by a mapping on the problem of directed polymers on a Cayley tree. However our starting point is different and we treat a deterministic model defined on a Cayley tree. We take advantage of the fact that the Cayley tree can be isometrically embedded in a space of constant negative curvature. We assume that each vertex of the tree carries a Boltzmann weight that depends on the hyperbolic distance from a given root point. The corresponding partition function is a sum over a finite number of tree vertices and has the form of a truncated Poincaré series. Its scaling dependence on the size of the system is controlled by the probability distribution of traces of $`2\times 2`$ matrices which belong to a discrete subgroup of $`PSL(2,\mathrm{𝖨𝖱})`$. This distribution, obtained by using the central limit theorem for Markov multiplicative processes , allows us to compute the multifractal exponents and discuss the termination of multifractality. The study of the convergence of the measure on the boundary reveals some interesting links with work of Gutzwiller and Mandelbrot on multifractal measures. Another interesting, although more speculative, aspect is connected with a geometric approach to Liouville field theory arising in the study of low dimensional disordered systems . We suggest that in two dimensions our model exhibits a new type of multifractal behavior which has a purely geometric origin.
This paper is organized as follows. In section II we introduce the geometrical model possessing the multifractal behavior, develop methods for its investigation and explicitly show the number theoretic origin of multifractality; section III is devoted to applications of these results to quasi–classical 2D Liouville field theory (LFT); the conclusion presents some speculations regarding the applicability of our geometric considerations to some other disordered physical systems.
## II The model
We begin with the investigation of geometrical properties of lattices uniformly embedded in the hyperbolic 2-space. Lattices under consideration are defined as follows: we construct the set of all possible orbits of a given root point under the action of a discrete subgroup of $`PSL(2,\mathrm{𝖨𝖱})`$ (group of motion of the hyperbolic 2-space). We restrict ourselves with the simplest example of 3–branching Bethe lattice (Cayley tree) which is generated by reflections of zero–angled curvilinear triangle—see fig.1.
The graph connecting the centers of the neighboring triangles forms a Cayley tree isometrically embedded in the Poincaré unit disc (a Riemann surface of constant negative curvature).
Consider the $`n^{\mathrm{th}}`$ generation of the vertices of the 3–branching Cayley tree. Denote by $`r_j(n)`$ the Euclidean distance of the vertex $`j`$ (which belongs to the $`n^{\mathrm{th}}`$ generation of the Cayley tree) from the center of the unit disc ($`1j3\times 2^{n1}`$). The corresponding hyperbolic (geodesic) distance $`d_j(n)`$ is given by:
$$d_j(n)=\mathrm{ln}\frac{1+r_j(n)}{1r_j(n)}\left(r_j(n)=\mathrm{tanh}\frac{d_j(n)}{2}\right)$$
(1)
Define the generating function $`𝒵(q,N)`$
$$𝒵(q,N)=\underset{n=1}{\overset{N}{}}\left(\underset{j=1}{\overset{3\times 2^{n1}}{}}e^{qd_j(n)}\right)$$
(2)
In a physical context $`𝒵(q,N)`$ may be interpreted as a partition function on the hyperbolic lattice with an action linear in the length of trajectory.
The Bethe lattice involved can be constructed by the action of the discrete group $`\mathrm{\Gamma }_\theta `$ which operates on the unit disc by a set of fractional–linear transformations. Despite the simple structure of the group it is believed that the techniques involved are quite general and could be easily generalized in order to cover more sophisticated lattices.
We are interested in the scaling dependence of the partition function $`𝒵(q,N)`$ as a function of the size $`N`$ of the system. Scaling considerations suggest the following behaviour
$$𝒵(q,N)=_{\mathrm{max}}^{\tau ^{}(q)}$$
(3)
where in our case $`_{\mathrm{max}}=3(2^N1)`$ is the total number of Cayley tree vertices in the bulk restricted by the generation $`n=N`$.
We show below that the critical exponent $`\tau (q)`$ defined as follows
$$\underset{N\mathrm{}}{lim}\frac{\mathrm{ln}𝒵(q,N)}{N}=\tau (q)$$
(4)
depends nonlinearly on $`q`$ i.e. exhibits the multicritical behavior (note that $`\tau ^{}=\tau /\mathrm{ln}2`$). Note that the free energy normalized per volume of the system $`f(q,N)=\frac{\mathrm{ln}𝒵(q,N)}{N}`$ coincides with the multifractal exponent $`\tau (q)`$:
$$\underset{N\mathrm{}}{lim}f(q,N)=\tau (q)$$
(5)
### A Numerical results
We first compute numerically the histogram, which counts the number of vertices belonging to generation $`n`$ (properly normalized), $`W_n(d)`$, lying in the shell $`[d,d+\delta d]`$.
In our particular computations we restrict ourselves with two cases depending on the length of the trajectories:
1. Short trajectories. We enumerate all trajectories and the computations have been carried out for all $`n[1,N]`$ up to $`N=25`$ generations. The figure fig.2a shows the histogram for the distribution of hyperbolic distances for $`n=25`$. The absolute value of number of events in the fig.2a depends on the particular choice of the width of the shell $`\delta `$. It can be seen from fig.2a that the corresponding plot is highly nonsymmetric with respect to the mean value $`d`$.
2. Long trajectories. For $`n=200`$ the enumeration of all different paths is very time consuming, therefore we compute numerically the histogram $`W_n(d)`$ developing partial ensemble of $`\mathrm{200\hspace{0.17em}000}`$ directed random walks of $`n=200`$ step each. As $`n\mathrm{}`$ the distribution function $`W_n(d)`$ becomes more and more symmetric in accordance with the statement that there exists a central limit theorem for such random walks on noncommutative groups (see the discussion below). The results of corresponding numerical computations are presented in Fig.2b. The distribution $`W_n(d)`$ is well fitted by a Gaussian function:
$$W_n(d)=A_0e^{\frac{(dd)^2}{2\mathrm{\Delta }^2}}$$
where for $`n=200`$ one has: $`A_01929.96`$ and depends on normalization; $`d159.18`$; $`\mathrm{\Delta }^217.01`$.
In spite of the fact that convergence to the Gaussian distribution is slow, the linear dependence in $`n`$ of the mean value $`d=\gamma n`$ and the variance $`(dd)^2\mathrm{\Delta }^2=\sigma ^2n`$ is numerically evident, which permits one to get an accurate estimate of $`\gamma `$ and $`\sigma ^2`$.
The numerical computation of the probability distribution $`W_n(d)`$ allows one to compute the multifractal exponent $`\tau (q)`$ following the definitions (3)–(5). The corresponding results are shown in fig.3, for $`N=40`$. Due to the slow convergence of the distribution, the discrepancy between numerical data (technically limited to $`N40`$) and the theoretical prediction can not be quantitatively taken into account. We here insist on the multifractal behaviour, shown by the non-linear depence on $`q`$.
### B Analytic results
Let us return to the definition of the model and recall that the group $`\mathrm{\Gamma }_\theta `$ acts in the hyperbolic Poincaré upper half–plane $`^2=\{z|𝖢,\mathrm{Im}(z)>0\}`$ by fractional–linear transforms<sup>*</sup><sup>*</sup>*It is convenient first to define the representation of the group $`\mathrm{\Gamma }_\theta `$ in the Poincaré upper half–plane and then use the conformal transform to the unit disc.. The matrix representation of the generators of the group $`\mathrm{\Gamma }_\theta `$ is well known (see, for example ), however for our purposes it is more convenient to take a framework that consists of the composition of standard fractional–linear transform and complex conjugacy. Namely, denoting by $`\overline{z}`$ the complex conjugate of $`z`$, we consider the following action
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right):z\frac{a\overline{z}+b}{c\overline{z}+d}$$
(6)
A possible set of generators is then:
$$h_0=\left(\begin{array}{cc}1& 2/\sqrt{3}\\ 0& 1\end{array}\right),h_1=\left(\begin{array}{cc}1& 2/\sqrt{3}\\ 0& 1\end{array}\right),h_2=\left(\begin{array}{cc}0& 1/\sqrt{3}\\ \sqrt{3}& 0\end{array}\right)$$
(7)
Choosing the point $`(x_0,iy_0)=(0,i)`$ as the tree root—see fig.4, any vertex on the tree is associated with an element $`M_n={\displaystyle \underset{k=1}{\overset{n}{}}}h_{\alpha _k}`$ where $`\alpha _k\{0,1,2\}`$ and is parametrized by its complex coordinates $`z_n=M_n\left((1)^ni\right)`$ in the hyperbolic plane.
Strictly speaking $`^2`$ should be identified with $`SL(2,\mathrm{𝖨𝖱})/SO(2)`$; we here identify an element with its class of equivalence of $`SO(2)`$. If one denotes by $`d(M_n)d(i,z_n)`$ the hyperbolic distance between $`i`$ and $`z_n`$, the following identity holds
$$2\mathrm{cosh}\left(d(M_n)\right)=\mathrm{Tr}(M_nM_n^{})$$
(8)
where dagger denotes transposition.
#### 1 Distribution function, invariant measure on the boundary and Lyapunov exponents
We are interested in the distribution function $`W_n(d)`$ which is the probability to find the tree vertices in generation $`n`$ at the distance $`d`$ from the root point. It means that we are looking for the distribution of the traces for matrices $`M_n`$ which are the irreducible products of $`n`$ generators. If we denote by $`l(M_n)`$ the irreducible length of the word represented by the matrix $`M_n`$, then $`M_n`$ is irreducible if and only if $`l(M_n)=n`$. Such word enumeration problem is simple in case of the group $`\mathrm{\Gamma }_\theta `$, because of its free product structure: $`\mathrm{\Gamma }_\theta \mathrm{𝖹𝖹}_2\mathrm{𝖹𝖹}_2\mathrm{𝖹𝖹}_2`$. Indeed, if $`M_n={\displaystyle \underset{k=1}{\overset{n}{}}}h_{\alpha _k}`$ one has $`l(M_n)=n`$ if and only if $`h_{\alpha _k}h_{\alpha _{k1}}k`$. Hence we have to study the behavior of the random matrix $`M_n`$, generated by the following Markovian process
$$M_{n+1}=M_nh_{\alpha _{n+1}}\mathrm{with}\alpha _{n+1}=\{\begin{array}{cc}(\alpha _n+1)\mathrm{mod}\mathrm{\hspace{0.17em}3}\hfill & \text{with probability }\frac{1}{2}\hfill \\ (\alpha _n+2)\mathrm{mod}\mathrm{\hspace{0.17em}3}\hfill & \text{with probability }\frac{1}{2}\hfill \end{array}$$
(9)
We use the standard methods of random matrices and consider the entries of the $`2\times 2`$–matrix $`M_n`$ as a 4–vector $`𝒱_n`$. The transformation $`M_{n+1}=M_nh_\alpha `$ reads
$$𝒱_{n+1}=\left(\begin{array}{cc}h_\alpha ^{}\hfill & 0\hfill \\ 0\hfill & h_\alpha ^{}\hfill \end{array}\right)𝒱_n$$
(10)
This block–diagonal form allows one to restrict ourselves to the study of one of two 2–vectors, composing $`𝒱_n`$, say $`v_n`$. Parametrizing $`v_n=(\varrho _n\mathrm{cos}\theta _n,\varrho _n\mathrm{sin}\theta _n)`$ and using the relation $`d(M_n)d_n2\mathrm{ln}\varrho _n`$ valid for $`n1`$, one gets a recursion relation $`v_{n+1}=h_\alpha ^{}v_n`$ in terms of hyperbolic distance $`d_n`$:
$$d_{n+1}=d_n+\mathrm{ln}\left[\frac{5}{3}+\frac{4}{3}\mathrm{cos}(2\theta _n+\phi _\alpha )\right]$$
(11)
where $`\phi _\alpha `$ depends on the transform $`h_\alpha `$ through $`\phi _\alpha =(2\alpha 1)\pi /3(\alpha =0,1,2)`$, while for the angles one gets straightforwardly
$$\mathrm{tan}\theta _{n+1}=h_{\stackrel{~}{\alpha }}\left(\mathrm{tan}(\theta _n)\right)$$
(12)
and the change $`\alpha \stackrel{~}{\alpha }`$ means the substitution $`(0,1,2)(1,0,2)`$. Action of $`h_\alpha `$ is still fractional–linear.
Define now three invariant measures $`\mu _\alpha (\theta )`$ corresponding to transformations of $`M_n`$ ($`n1`$) whose last step is given by a matrix $`h_\alpha `$. The form of (12) suggests to consider the corresponding $`\mu _\alpha (x)`$ with $`x=\mathrm{tan}\theta `$. We are then led to study the action of $`\mathrm{\Gamma }_\theta `$ restricted on the real line parametrized by $`x`$. Interesting properties of the average $`\overline{\mu }=(\mu _0+\mu _1+\mu _2)/3`$ have been discussed by Gutzwiller and Mandelbrot . In particular they pointed out the connexion with the arithmetic function $`\beta (\xi )`$ which maps some number $`\xi [0,1]`$ written as a continued fraction expansion
$$\frac{1}{n_1+{\displaystyle \frac{1}{n_2+\mathrm{}}}}$$
to the real number $`\beta `$ whose binary expansion is made by the sequence of $`n_11`$ times 0, followed by $`n_2`$ times 1, then $`n_3`$ times 0, and so on. To account for this fact, one has first to notice that the construction (9) of any word $`M`$ in $`\mathrm{\Gamma }_\theta `$ is exactly encoded by the binary representation of a real $`\xi `$, the $`n`$’s letter of this expansion being $`\alpha _{n+1}\alpha _n+1\mathrm{mod}\mathrm{\hspace{0.17em}3}`$. The second argument, due to Series , is that the real part of the vertex $`M(i)`$ is precisely the continued fraction $`\mu (\xi )`$. Therefore $`\beta (\xi )`$ has to be proportional to the “number” of vertices lying in the interval $`[0,\xi ]`$, that is to $`\overline{\mu }([0,\xi ])`$. Taking the limit $`n\mathrm{}`$ is not well defined. An alternative, which was used in this work, is to define $`\mu _\alpha (x)`$ as the limit of the following recurrency:
$$\mu _\alpha ^{(n+1)}(x)=\frac{1}{2}\left|\frac{dh_\alpha (x)}{dx}\right|\underset{\alpha ^{}\alpha }{}\mu _\alpha ^{}^{(n)}\left(h_\alpha ^{}(x)\right)$$
(13)
The symetry of such expression leads, after summing over $`\alpha `$, to the following relation admitting as fixed point $`\overline{\mu }(x)`$ at $`n\mathrm{}`$:
$$\mu ^{(n+1)}(x)=\frac{1}{3}\underset{\alpha =0}{\overset{2}{}}\mu ^{(n)}\left(h_\alpha (x)\right)\left|\frac{dh_\alpha (x)}{dx}\right|$$
(14)
The convergence $`\mu ^{(n)}(x)\overline{\mu }(x)`$ for $`n\mathrm{}`$ is assured by ergodic properties of such functional transform in case of equation (14), and has been successfully checked numerically by comparing to direct sampling of different orbits. Obtaining $`\mu _\alpha `$ for $`\alpha =\{0,1,2\}`$ from $`\overline{\mu }`$ is not difficult, taking into account the symmetric role that they play with respect to the three intervals $`I_0=]\mathrm{},1/\sqrt{3}],I_2=[1/\sqrt{3},1/\sqrt{3}],I_1=[1/\sqrt{3},+\mathrm{}[`$ (see fig.5). Contracting properties of $`h_\alpha (x)`$ allow convergence of (14) only if
$$\mu _k(x)=3\chi _{I_k}(x)\overline{\mu }(x)$$
(15)
where $`\chi _{I_k}`$ is the characteristic function of the interval $`I_k`$.
We would like to point out an interesting fact, even if far from being rigorous, which is very similar to the argument put forward in for justifying the connexion between the invariant measure and the arithmetic function $`\beta (\xi )`$. It has been shown in that the lattice under consideration can be isometrically embedded in a 2–manifold $`=\{c|\eta (z)|^4,z^2\}`$, where
$$\eta (z)=e^{\pi iz/12}\underset{n=1}{\overset{\mathrm{}}{}}(1e^{2\pi inz})$$
is the Dedekind $`\eta `$–function. The mountain range (relief) $``$ displays a very steep valley structure, and our tree lattice was defined as the ridges of this relief. The natural “counting” of vertices whose real part lies in $`[0,\xi ]`$ in is in our case equivalent to counting the number of maxima of $`|\eta (x+i0^+)|^4`$, that can be directly reexpressed as a density if one admits that all maxima are equivalent and well separated:
$$\overline{\mu }(x)\frac{|\eta (x+i0^+)|^4}{_0^1|\eta (t+i0^+)|^4𝑑t}$$
(16)
The intriguing fact is that $`\eta ^4`$ is an automorphic form of weight 2, what makes $`|\eta |^4`$ precisely a possible fixed point of Eq.(14). We recall the fundamental property of automorphic forms $`f`$ of weight 2 under the action of $`SL(2,\mathrm{𝖨𝖱})`$:
$$f(z)=\frac{e^{i\varphi (a,b,c,d)}}{(cz+d)^2}f\left(\frac{az+b}{cz+d}\right)$$
(17)
The main problem is that the boundary behavior of automorphic forms is far from trivial (see ), and (16) has no rigorous mathematical sense. In particular compatibility of (14) and (16) is not obvious even numerically. Nevertheless we insist on the fact that $`\mu _k`$ is defined with no ambiguity by (14), what enables us to compute the desired $`W_n(d)`$. The crucial point here, already required for convergence of $`\mu ^{(n)}`$, is ergodicity property of $`\theta _n`$. It means that for $`n1`$, the distribution of $`\theta _n`$ is exactly given by $`\overline{\mu }(\theta )`$, independently of $`n`$ and initial condition. Then, denoting by $`d_n^\alpha `$ the value $`d_n`$ obtained for a word ending with $`h_\alpha `$, one can transform (11) in the following way:
$$d_{n+1}=\frac{1}{3}\underset{\alpha =0}{\overset{2}{}}\left(d_n^\alpha +\mathrm{ln}\left[\frac{5}{3}+\frac{4}{3}\mathrm{cos}(2\theta _n+\phi _l)\right]\right)$$
(18)
with the condition $`l\alpha `$. Thus we obtain
$$e^{ikd_n}=\frac{1}{6}\underset{k=0}{\overset{2}{}}\underset{jk}{}\left[_{\pi /2}^{\pi /2}𝑑\theta \mu _k(\theta )\left(\frac{5}{3}+\frac{4}{3}\mathrm{cos}(2\theta +\phi _j)\right)^{ik}\right]^n$$
(19)
which finally leads to
$$W_n(d)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑ke^{ikd}\left[_0^{\frac{\pi }{3}}𝑑\theta \mu _1(\theta \frac{\pi }{6})\left(\frac{5}{3}+\frac{4}{3}\mathrm{cos}2\theta \right)^{ik}\right]^n$$
(20)
This form suggests that for $`n`$ large $`W_n(d)`$ satisfies a central limit theorem. Indeed such a theorem exists (see ) for Markovian processes, provided that the phase space is ergodic. We are then led to compute only the first two moments (Lyapunov exponents) which gives
$$\gamma =\frac{d}{n}=_0^{\frac{\pi }{3}}𝑑\theta \mu _1(\theta \frac{\pi }{6})\mathrm{ln}\left(\frac{5}{3}+\frac{4}{3}\mathrm{cos}2\theta \right)0.792$$
(21)
and
$$\sigma ^2=\frac{(dd)^2}{n}=\gamma _2\gamma ^2$$
(22)
with
$$\gamma _2=_0^{\frac{\pi }{3}}𝑑\theta \mu _1(\theta \frac{\pi }{6})\mathrm{ln}^2\left(\frac{5}{3}+\frac{4}{3}\mathrm{cos}2\theta \right)0.68$$
(23)
Numerical simulations presented in previous section yield $`\gamma 0.793`$ and $`\gamma _20.66`$, which finally allow us to conclude that for $`n1`$ $`W_n(d)`$ has a Gaussian behavior
$$W_n(d)=Ae^{\frac{(dn\gamma )^2}{2\sigma ^2n}}$$
(24)
centered at $`\gamma n`$ and of variance $`\sigma ^2n`$ ($`A`$ is the normalization).
The numerical values of the Lyapunov exponents $`\gamma `$ and $`\gamma _2`$ (see Eqs.(21) and (23)) are obtained by means of semi–numerical procedure which involves the numerical information about the invariant measure $`\mu _1(\theta )`$. However one can get the estimates for the Lyapunov exponents $`\gamma `$ and $`\gamma _2`$ by approximating the measure $`\mu _1(\theta )`$ on the interval $`0\theta \frac{\pi }{3}`$ in two different ways:
$$\begin{array}{c}\mu _1(\theta \frac{\pi }{6})\mu _1^A(\theta )=\frac{3}{\pi }\hfill \\ \mu _1(\theta \frac{\pi }{6})\mu _1^B(\theta )=\frac{3}{2}\mathrm{sin}(3\theta )\hfill \end{array}$$
(25)
Both measures $`\mu _1^A`$ and $`\mu _1^B`$ are properly normalized on the interval $`[0,\frac{\pi }{3}]`$. Substituting (25) in (21) and (23) and computing (analytically for $`\gamma `$) the corresponding integrals, one finally gets:
$$\mu _1^A:\{\begin{array}{c}\gamma 0.749\hfill \\ \gamma _20.665\hfill \end{array}\mu _1^B:\{\begin{array}{c}\gamma 0.792\hfill \\ \gamma _20.684\hfill \end{array}$$
(26)
As one can see, the agreement between numerical values of Lyapunov exponents obtained for the measures $`\mu _1`$ and its approximants $`\mu _1^{A,B}`$ is reasonable for $`\mu _1^A`$ and very good for $`\mu _1^B`$.
#### 2 Multifractal exponents
The partition function $`𝒵(q,N)`$ introduced in (2) can be defined for any discrete subgroup of $`PSL(2,\mathrm{𝖨𝖱})`$ of generic element $`\tau `$ by
$$𝒵(q,N)=\underset{\tau ,l(\tau )N}{}e^{qd(\tau )}$$
(27)
and the associated critical exponent is then
$$\tau (q)=\underset{N\mathrm{}}{lim}\frac{\mathrm{ln}𝒵(q,N)}{N}$$
(28)
The probability distribution (24) enables us to rewrite (27) for the group $`\mathrm{\Gamma }_\theta `$ in the limit $`n1`$ as follows
$$𝒵(q,N)=\underset{n=1}{\overset{N}{}}3\times 2^{n1}a_n$$
(29)
where
$$a_n=_0^{\mathrm{}}e^{qt}W_n(t)𝑑t=_0^{\mathrm{}}Ae^{\frac{(tn\gamma )^2}{2n\sigma ^2}}e^{qt}𝑑t$$
(30)
The following two cases should be distinguished:
* For $`q<\gamma /\sigma ^2`$, the minimum of the exponent in Eq.(30) is within the range of integration and
$$a_ne^{n\gamma q+n\sigma ^2q^2/2}$$
(31)
hence
$$𝒵(q,N)\frac{3}{2}\underset{k=1}{\overset{N}{}}e^{k(\mathrm{ln}2\gamma q+\sigma ^2q^2/2)}$$
(32)
The convergence of the sum (32) for $`N\mathrm{}`$ depends on the sign of the function in the exponent. For
$$q<q_0=\frac{\gamma \sqrt{\gamma ^22\sigma ^2\mathrm{ln}2}}{\sigma ^2}$$
(33)
one has $`\mathrm{ln}2+\gamma q\sigma ^2q^2/2>0`$ and the multifractal exponent is
$$\tau (q)=\mathrm{ln}2+\gamma q\sigma ^2q^2/2$$
(34)
while for $`q>q_0`$, the series $`𝒵(q,N)`$ converges and $`\tau (q)=0`$ what signals the termination of the multifractality. Note that $`q_0`$ is real at least in the case of the group $`\mathrm{\Gamma }_\theta `$.
* For $`q>\gamma /\sigma ^2`$ the minimum of the exponent in Eq.(30) is out of the range of integration and
$$\mathrm{ln}a_n\frac{n\gamma ^2}{2\sigma ^2}$$
(35)
hence $`𝒵(q,N)`$ is no longer extensive in $`N`$, which leads to $`\tau (q)=\tau (q_0)=0`$.
The nonlinear dependence on $`q`$ obtained above shows the multifractal behaviour of this model below the termination point $`q_0`$. It seems more transparent to summarize all these results in a table
#### 3 Conformal mapping approach to computation of partition function and multifractal exponent
We propose in this section a completely different approach allowing to get a closed analytic expression for the partition function similar to $`𝒵(q,N_{\mathrm{max}})`$ (see Eq.(2)). The construction presented below is a by-product of our former investigations of analytic structure of the covering Riemann space of the multi–punctured complex plane (see, for review ). We explore the properties of the Jacobian of conformal mapping of the infinite complex plane with a triangular lattice of punctures into the unit disc parametrized by $`w=re^{i\alpha }`$, which in this particular case represents the multi–sheeted universal covering space . Namely, we define two functions $`f(r,\alpha )`$ and $`g(r)`$:
$$f(r,\alpha )=c\frac{\left|\theta _1^{}(0,e^{i\pi \zeta (w)})\right|^{8/3}}{|1+iw|^4}c\frac{\left|\eta (\zeta (w))\right|^8}{|1+iw|^4},g(r)=\frac{1}{(1r^2)^2}$$
(36)
where
$$\{\begin{array}{c}\theta _1^{}(0,e^{i\pi \zeta })=2e^{i\frac{\pi }{4}\zeta }\underset{n=0}{\overset{\mathrm{}}{}}(1)^n(2n+1)e^{i\pi n(n+1)\zeta }\hfill \\ \zeta (w)=e^{i\pi /3}\frac{w+e^{i\pi /6}}{wi}\hfill \\ c=\left|\theta _1^{}(0,e^{i\pi \zeta (0)})\right|^{8/3}0.933293\hfill \end{array}$$
(37)
One can show that the functional equation
$$\frac{f(r,\alpha )}{g(r)}1=0$$
(38)
has a family of solutions $`(r_\mathrm{c},\alpha _\mathrm{c})`$ exactly at positions of 3–branching Cayley tree isometrically embedded in the hyperbolic unit disc (in the Klein’s model of the surface of constant negative curvature). In fig.6 we have plotted the 3D section of the function
$$u(r,\alpha )=\frac{f(r,\alpha )}{g(r)}$$
(39)
in polar coordinates $`(r,\alpha )`$ for $`0.9<u(r,\alpha )1`$. The function $`u(r,\alpha )`$ has local maxima with one and the same value $`u=1`$ only at the coordinates of isometric embedding of 3–branching Cayley tree in the hyperbolic unit disc. The proof of this fact is given in Appendix A.
Thus, we can rewrite (2) in the following closed form (recall that Euclidean distance $`r`$ and the corresponding hyperbolic distance $`d`$ are linked by the relation (1))
$$\stackrel{~}{𝒵}(q,d)=\frac{1}{2\pi }\underset{0}{\overset{r(d)}{}}\underset{0}{\overset{2\pi }{}}e^{q\mathrm{ln}\frac{1+r}{1r}}\delta \left(\mathrm{ln}u(r,\alpha )\right)\left|\frac{d\mathrm{ln}u(r,\alpha )}{dr}\right|r𝑑r𝑑\alpha $$
(40)
where for $`\delta [\mathrm{ln}u(x)]`$ we use the standard integral representation $`\delta [\mathrm{ln}u(x)]=\frac{1}{2\pi }\underset{\mathrm{}}{\overset{\mathrm{}}{}}𝑑\xi [u(x)]^{i\xi }`$.
It is noteworthy to pay attention to the difference between the partition functions $`𝒵(q,N)`$ (Eq.(2)) and $`\stackrel{~}{𝒵}(q,d)`$ (Eqs.(27) and (40)). The function $`𝒵(q,N)`$ counts the weighted number of Cayley tree vertices up to the generation $`N`$ for nonfixed maximal radius $`r(d)=\mathrm{tanh}(d/2)`$ in the hyperbolic unit disc, while the function $`\stackrel{~}{𝒵}(q,d)`$ counts the weighted number of Cayley tree vertices within the hyperbolic disc of radius $`r(d)`$ for nonfixed maximal generation $`N`$. The last partition function is in fact related to the number of tree vertices inside the disc of radius $`d`$. This is the content of the famous circle problem first formulated by Gauss for the Euclidean lattice $`\mathrm{𝖹𝖹}^2`$. The extension to the non-Euclidean case is due to Delsarte (see also ).
## III Multifractality in 2D quasi–classical Liouville field theory
Our starting point is the family of normalized wave functions $`\psi _k(𝐱)`$ defined as follows
$$\psi _k(𝐱)=\frac{|𝐱|^ke^{\phi (𝐱)}}{\left[𝑑𝐱|𝐱|^{2k}e^{2\phi (𝐱)}\right]^{1/2}}$$
(41)
where integration extends to a disc of radius $`R`$ and the potential $`\phi (𝐱)`$ is distributed with Gaussian distribution function
$$P[\phi (𝐱)]\mathrm{exp}\left\{\frac{1}{2g}𝑑𝐱\left(\phi (𝐱)\right)^2\right\}$$
(42)
The problem defined in (41)–(42) appears in various models which will be discussed in the next section.
The multifractal exponent for the quenched and annealed distributions of disorder in (41)–(42) can be computed in the standard way
$$\begin{array}{cc}\tau _{\mathrm{qu}}(q)=\underset{R\mathrm{}}{lim}\frac{\mathrm{ln}Q(q,R)}{\mathrm{ln}R}\hfill & \text{for quenched disorder}\hfill \\ \tau _{\mathrm{an}}(q)=\underset{R\mathrm{}}{lim}\frac{\mathrm{ln}Q(q,R)}{\mathrm{ln}R}\hfill & \text{for annealed disorder}\hfill \end{array}$$
(43)
where $`Q(q,R)=𝑑𝐱|\psi _k(𝐱)|^{2q}`$ and the brackets $`\mathrm{}`$ denote averaging with the distribution $`P[\phi (𝐱)]`$.
We pay attention to the case of annealed disorder and our aim is to evaluate the correlation function
$$Q(q,R)=𝑑𝐱|\psi _k(𝐱)|^{2q}$$
(44)
The averaging $`\mathrm{}`$ in (44) means
$$\mathrm{}_{S[\phi ]}=𝒟[\phi (𝐱)]e^{S[\phi (𝐱)]}$$
(45)
where
$$S[\phi (𝐱)]=\frac{1}{2g}𝑑𝐱\left\{\left(\phi (𝐱)\right)^2\right\}$$
(46)
In order to take into account proper normalization of the wave function $`\psi _k(𝐱)`$ it is convenient to use a Lagrange multiplier $`\lambda `$, so that eventually
$$Q(q,R)=𝑑𝐱_0𝒟[\phi (𝐱)]e^{2q\left(\phi (𝐱_0)k\mathrm{ln}|𝐱_0|\right)}e^{S_k[\phi (𝐱)]}$$
(47)
where
$$S_k[\phi (𝐱)]=\frac{1}{2g}𝑑𝐱\left\{\left(\phi (𝐱)\right)^2+2\lambda g\left(|𝐱|^{2k}e^{2\phi (𝐱)}\frac{1}{\pi R^2}\right)\right\}$$
(48)
is the action of 2D Liouville Field Theory (LFT).
The careful treatment of the quantum LFT in the case $`k=0`$ (see for review ) enables one to find the conformal weights $`\mathrm{\Delta }\left(e^{2q\phi }\right)=q(𝒬q)`$, where $`𝒬(g)`$ is the “background charge”, obtained by imposing conformal invariance of $`S_0[\phi ]`$. The authors of work have related the value $`\mathrm{\Delta }\left(e^{2q\phi }\right)`$ to the critical exponent $`\tau _{\mathrm{an}}(q)`$ in the scaling dependence of the average inverse participation ratio (44)
$$Q(q,R)R^{\tau _{\mathrm{an}}(q)}$$
(49)
Despite the multiscaling exponent $`\tau _{\mathrm{an}}(q)`$ has been computed in the general framework of Conformal Field Theory (CFT) few years ago, from our point of view, the geometrical interpretation of the multifractal behavior in the model has not yet been cleared up. A more physical approach put forward in exploits an analogy between this model and the problem of directed polymers on a Cayley tree. This analogy is supported by the fact that in both cases the correlation functions grow logarithmically with the distance. For directed polymers it is the correlation function of the random potential defined on the tree vertices that scales logarithmically with the ultrametric distance (i.e. distance along the tree). The same logarithmic behaviour occurs in $`2D`$ Gaussian Field Theory .
We adopt a different point of view. Let us notice that the tree structure (conjectured by C.Chamon et al) emerges quite naturally from the Liouville field theory treated at a semi-classical level. Our idea is as follows. Indeed there is not just one saddle point solution but a whole orbit of solutions parametrized by $`SL(2,\mathrm{𝖨𝖱})`$. If one further assumes that the integration has to be performed not over the whole group but only over a subgroup (for instance $`\mathrm{\Gamma }_\theta `$), one recovers quite naturally the model defined in section II.
Our starting point is the semi–classical ($`g0`$) treatment of (47) (a similar approach can be found in ). Using a saddle point method, one is led to the classical equation
$$\mathrm{\Delta }\phi 2qg\delta ^2(𝐱𝐱_0)+2\lambda g|𝐱|^{2k}e^{2\phi }=0$$
(50)
which gives, after integration
$$\lambda =qF/2g$$
(51)
where
$$F=𝑑𝐱\mathrm{\Delta }\phi $$
is the magnetic flux. We then introduce the shifted field
$$\stackrel{~}{\phi }(𝐱)=\phi (𝐱)\mathrm{ln}|𝐱|^k$$
(52)
and taking into account that
$$\mathrm{\Delta }\stackrel{~}{\phi }(𝐱)=\mathrm{\Delta }\phi (𝐱)2\pi k\delta ^2(𝐱)$$
(53)
we end up with the following equation
$$\mathrm{\Delta }\stackrel{~}{\phi }+2\pi k\delta ^2(𝐱)2qg\delta ^2(𝐱𝐱_0)+(2qgF)e^{2\stackrel{~}{\phi }}=0$$
(54)
The most general solution of (54) (away from singularities) in Euclidean space of complex coordinate $`z`$ can be written as follows
$$e^{2\stackrel{~}{\phi }}=\frac{4}{|F2qg|}\frac{_zA(z)_{\overline{z}}B(\overline{z})}{\left(1+ϵA(z)B(\overline{z})\right)^2}$$
(55)
where $`A(z)`$ and $`B(\overline{z})`$ are correspondingly holomorphic and anti–holomorphic functions of $`z`$ and $`ϵ=\mathrm{sign}(F2qg)`$. The semi–classical treatment assumes $`g`$ to be small, hence in order to have real $`\stackrel{~}{\phi }`$ in Eq.(55) we should put $`ϵ=1`$, i.e. $`qF/2g`$. The relevant solution of (55), compatible with the singularities, reads
$$e^{2\stackrel{~}{\phi }_{cl}(z)}=\frac{4(k+1)^2}{F}\frac{(z\overline{z})^k}{\left(1+(z\overline{z})^{k+1}\right)^2}$$
(56)
The normalization condition of the wave function is then satisfied only if $`F=4\pi (k+1)`$. We would like to stress that for $`k=0`$ we here recover the critical value $`4\pi `$ of the magnetic flux: uniqueness of the ground state wave function holds only below this value. It also should be mentioned that our analysis does not depend on a peculiar basis of eigenfunctions and the results presented here can be extended to wave functions of the form $`\psi (z)=P_k(z)e^{\phi (z)}`$, $`P_k`$ being a polynomial of degree $`k`$. It is noteworthy that (56) is an algebraically decaying wave function, what is, following , a signature of the existence of prelocalized states.
Using the fact that the Liouville field is not exactly a scalar but varies under holomorphic coordinate transformations $`zw(z)`$ as
$$\stackrel{~}{\phi }(z)\stackrel{~}{\phi }\left(w(z)\right)\mathrm{ln}|w^{}(z)|,$$
(57)
one can check that the set of solutions (56) is invariant under the following family of transformations, parametrized by the group $`PSL(2,|𝖢)`$:
$$zw_k(z;a,b,c,d)=\left(\frac{az^{k+1}+b}{cz^{k+1}+d}\right)^{\frac{1}{k+1}}$$
(58)
The orbit of $`\stackrel{~}{\phi }_{cl}(z)`$ is then given by
$$\stackrel{~}{\phi }_{cl}(z;a,b,c,d)=\frac{1}{2}\mathrm{ln}\left(\frac{\pi \left(|az^{k+1}+b|^2+|cz^{k+1}+d|^2\right)^2}{(z\overline{z})^k}\right)$$
(59)
Up to redefinition of the measure $`d\tau `$ on $`PSL(2,|𝖢)`$, we restrict the domain of integration to $`PSL(2,\mathrm{𝖨𝖱})`$.
Due to the angular symmetry of (56), we take points of the form $`z=i^{\frac{1}{k+1}}\rho `$ ($`\rho \mathrm{𝖨𝖱}`$), and following we rewrite (44)–(46)
$$\rho ^{2qk}e^{2q\phi (\rho )}_{S_k[\phi ]}=R^{\frac{4\pi (k+1)^2}{g}}\mathrm{Det}\left[\frac{\delta ^2S_k}{\delta \phi ^2}\right]^{1/2}_{PSL(2,\mathrm{𝖨𝖱})}e^{2q\stackrel{~}{\phi }_{cl}(\rho ,\tau )}𝑑\tau $$
(60)
Let us denote
$$\tau =\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\mathrm{and}\nu _\rho =\left(\begin{array}{cc}\rho ^{(k+1)/2}\hfill & 0\hfill \\ 0\hfill & \rho ^{(k+1)/2}\hfill \end{array}\right)$$
(61)
then we can rewrite (60) as follows
$$\rho ^{2qk}e^{2q\phi (\rho )}_{S_k[\phi ]}\rho ^{2q}_{PSL(2,\mathrm{𝖨𝖱})}e^{2q\mathrm{ln}\mathrm{Tr}\left[(\tau \nu _\rho )(\tau \nu _\rho )^{}\right]}𝑑\tau =\rho ^{2q}I(\rho ,q)$$
(62)
where we have got rid of irrelevant factors and the function $`I(\rho ,q)`$ reads
$$I(\rho ,q)=_{PSL(2,\mathrm{𝖨𝖱})}\left[2\mathrm{cosh}d(i,\tau \nu _\rho (i))\right]^{2q}𝑑\tau $$
(63)
Instead of summing over the whole group $`PSL(2,\mathrm{𝖨𝖱})`$, we restrict the sum over a discrete subgroup, $`\mathrm{\Gamma }_\theta `$ in our case. Even if this discretization of the saddle manifold has no evident physical justification, we believe that the model obtained yields interesting results. It leads to consider the so-called Poincaré series (see for review) $`H`$, defined as follows
$$I(\rho ,q)=H(i,\nu _\rho (i),q)=\underset{\tau \mathrm{\Gamma }_\theta }{}\left[2\mathrm{cosh}d(i,\tau \nu _\rho (i))\right]^{2q}$$
(64)
As shown in , the series $`H`$ does not converge for $`q<q^{}`$ with $`q^{}`$ depending on $`\mathrm{\Gamma }_\theta `$ only (the analysis of the previous sections show that we roughly may set $`q^{}q_0/2`$). A new dependence on $`\rho `$ occurs only if $`H`$ does not converge, we will therefore consider only this regime. We must introduce in this case a cut-off $`𝒩`$ to regularize the series, and finally study asymptotics of the finite sum
$$I_𝒩(\rho ,q)=H_𝒩(i,\nu _\rho (i),q)=\underset{\tau \mathrm{\Gamma }_\theta /l(\tau )𝒩}{}\left[2\mathrm{cosh}d(i,\tau \nu _\rho (i))\right]^{2q}$$
(65)
This Poincaré series has the same asymptotic properties as the one that defines our model. In particular it will exhibit a multifractal behaviour in $`𝒩`$. However what really matters for a physical system is the multifractal behaviour under transformations parametrized by $`\rho `$. We therefore have to relate the behaviour in the group manifold to the behaviour in the real space. This will be achieved through a renormalization transformation of the form
$$I_𝒩(\rho ,q)=C\rho ^\kappa I_{𝒩^{}(𝒩,\rho )}(1,q)$$
(66)
Appendix B provides a heuristic derivation which gives
$$I_𝒩(\rho ,q)=C\rho ^{(k+1)\mathrm{ln}2/\gamma }I_{𝒩+\frac{(k+1)}{\gamma }\mathrm{ln}\rho }(1,q)$$
(67)
Comparing (27) and (64) gives $`I_𝒩(1,q)𝒵(2q,𝒩)`$. Therefore
$$I_𝒩(\rho ,q)=C\rho ^{(k+1)\mathrm{ln}2/\gamma }𝒵(2q,𝒩+\frac{k+1}{\gamma }\mathrm{ln}\rho )$$
(68)
This relation allows to extract the scale dependence in $`\rho `$ for a given cut-off $`𝒩`$. Using the asymptotics of $`𝒵`$ obtained in previous sections yields
$$\rho ^{2qk}e^{2q\phi (\rho )}_{S_k[\phi ]}\rho ^{2q(k+1)\mathrm{ln}2/\gamma }\rho ^{\frac{k+1}{\gamma }\tau (2q)}$$
(69)
with the notations of previous sections. After integrating over the whole domain we arrive at the final expression for the multifractal exponent $`\tau _{\mathrm{an}}(q)`$:
$$\tau _{\mathrm{an}}(q)=\underset{R\mathrm{}}{lim}\frac{\mathrm{ln}Q(q,R)}{\mathrm{ln}R}=2(q1)\left(1\frac{(k+1)\sigma ^2}{\gamma }q\right)\text{for }q<q_0/2$$
(70)
with $`q_0`$ defined in (33). The regular term $`2(q1)`$ corresponds to the one obtained in for $`g0`$. Multifractality of the wave function is induced by the quadratic term, which is directly related to geometric properties of the saddle point hyperbolic manifold (target space), and holds in absence of any random potential in this target space. The regime $`q>q_0/2`$ is not affected by these geometric properties.
## IV Conclusion
The wave function $`\psi _k(𝐱)`$ introduced in (41) belongs to the general class of exponential functionals of free fields. Such functionals appear in several physical contexts.
1. The square of the wave function (for $`k=0`$) may be interpreted as the equilibrium Gibbs measure in the random potential $`\phi (𝐱)`$. In the 1D case the problem was first studied in . A rather deep and complete analysis of this problem was recently presented in .
2. Exponential functionals of free fields play an important role in the context of one dimensional classical diffusion in a random environment. Their probability distribution controls the anomalous diffusive behaviour of particles at large time . They also arise in the study of disordered samples of finite length . Some mathematical properties are discussed in .
3. In the context of one dimensional localization in a random potential such functionals arise in the study of the Wigner time delay .
4. The function $`\psi _k(𝐱)`$ is the ground state wave function of 2D Dirac fermions in a random magnetic field with $`B=\mathrm{\Delta }\phi `$. The multifractal behaviour first conjectured in has been recently confirmed by an independent investigation based on renormalization group method . The scenario of multifractality which is presented here relies mainly on a geometric approach to a semiclassical quantization scheme of the Liouville field theory. The fact that tree like structure emerges quite naturally in our consideration is an interesting feature which obviously deserves further investigation. The multifractality in our approach appears as a by-product of isometric embedding of a Cayley tree in the hyperbolic plane. The objects which possess multifractal behavior are the moments of the partition function defined as sums over all vertices of a Cayley tree isometrically embedded in the hyperbolic plane where each vertex carries a Botzmanm weight depending on the hyperbolic distance from the root point. No randomness is imposed in the model.
From the mathematical side our work reveals some interesting links between the theory of automorphic functions, invariant measures and spectral theory. We hope to return to these problems in a forthcoming publication.
Acknowledgments
The authors are grateful to Ph. Bougerol, J. Marklof, O. Martin, H. Saleur, and Ch. Texier for valuable discussions and useful comments of different aspects of the problem.
## A
Let us prove that the function $`u(w)=(1w\overline{w})^2f(w)`$ where $`f(w)`$ is defined in (36) has the following properties:
* At all centers $`w=w_c`$ of zero–angled triangles tesselating the Poincaré hyperbolic unit disc $`u(w_c)=1`$;
* The function $`u(w)`$ has local maxima at the points $`w_c`$.
I. The proof of the first statement implies the proof of the fact that the function $`u(w)`$ is invariant with respect to the conformal transform $`w^{(1)}(w)`$ of the unit Poincaré disc to itself where
$$w^{(1)}(w)=\frac{ww_0}{w\overline{w}_01},$$
(A1)
and $`w_0`$ is the coordinate of any center of zero–angled triangle in the hyperbolic Poincaré disc obtained by successive transformations from the initial one.
Hence, it is neccessary and sufficient to show that the values $`u(w=0)`$, $`u\left(w=\frac{i}{2}\right)`$, $`u\left(w=\frac{1}{2}e^{i\pi /6}\right)`$ and $`u\left(w=\frac{1}{2}e^{i5\pi /6}\right)`$ coincide. Then, performing the conformal transform and taking $`w_0=\{\frac{i}{2},\frac{1}{2}e^{i\pi /6},\frac{1}{2}e^{i5\pi /6}\}`$, we move the centers of the first generation of zero–angled triangles to the center of the disc $`w^{(1)}`$. Now we can repeat recursively the contruction, i.e. find the new coordinates of the centers of the second generation of zero–angled triangles in the disc $`w^{(1)}`$ and compute the function $`u(w)`$ at these points, then we perform the conformal transform $`w^{(2)}(w^{(1)})`$ and so on…
We have at the point $`w=0`$:
$$u(w=0)=c\left|\theta _1^{}(0,e^{i\pi [1/2+i\sqrt{3}/2]})\right|^{8/3}=1$$
while at the point $`w=\frac{i}{2}`$ the function $`u(w)`$ can be written in the form
$$u\left(w=\frac{i}{2}\right)=\frac{c}{\left|1\frac{1}{2}e^{i\pi }\right|^4}\left|\theta _1^{}(0,e^{i\pi [1/2+i/(2\sqrt{3})]})\right|^{8/3}\left(1\frac{1}{4}\right)^2$$
(A2)
Let us use the properties of Jacobi $`\theta `$–functions:
$$\{\begin{array}{cc}\theta _1^{}(0,e^{i\pi (w+k)})=\theta _1^{}(0,e^{i\pi w});\hfill & kN\hfill \\ \theta _1^{}(0,e^{i\pi [1/2+i/(2\sqrt{\lambda })]})\lambda ^{3/4}=\theta _1^{}(0,e^{i\pi [1/2+i\sqrt{\lambda }/2]});\hfill & \lambda R\hfill \end{array}$$
(A3)
Taking into account (A3) we can rewrite (A2) in the form
$$\begin{array}{ccc}u\left(w=\frac{i}{2}\right)\hfill & =\hfill & \frac{2^4c}{3^4}(3^{3/4})^{8/3}\left|\theta _1^{}(0,e^{i\pi [1/2+i/(2\sqrt{3})]3})\right|^{8/3}\left(\frac{3}{4}\right)^2\hfill \\ & =\hfill & c\left|\theta _1^{}(0,e^{i\pi (1/2+i\sqrt{3}/2)+i\pi })\right|^{8/3}=1\hfill \end{array}$$
(A4)
Thus, $`u(w=0)=u\left(w=\frac{i}{2}\right)=1`$.
At the point $`u\left(w=\frac{1}{2}e^{i\pi /6}\right)`$ we have
$$\begin{array}{ccc}u\left(w=\frac{1}{2}e^{i\pi /6}\right)\hfill & =\hfill & \frac{2^4c}{3^2}\left|\theta _1^{}(0,e^{i\pi [3/2+i\sqrt{3}/2]})\right|^{8/3}\left(1\frac{1}{4}\right)^2\hfill \\ & =\hfill & c\left|\theta _1^{}(0,e^{i\pi [1/2+i\sqrt{3}/2]+i\pi })\right|^{8/3}=1\hfill \end{array}$$
(A5)
In the same way we can transform the function $`u\left(w=\frac{1}{2}e^{i5\pi /6}\right)`$:
$$\begin{array}{ccc}u\left(w=\frac{1}{2}e^{i5\pi /6}\right)\hfill & =\hfill & \frac{2^4c}{3^2}\left|\theta _1^{}(0,e^{i\pi [1/2+i\sqrt{3}/2]})\right|^{8/3}\left(1\frac{1}{4}\right)^2\hfill \\ & =\hfill & c\left|\theta _1^{}(0,e^{i\pi [1/2+i\sqrt{3}/2]i\pi })\right|^{8/3}=1\hfill \end{array}$$
(A6)
The transforms (A4)–(A6) complete the proof of the part I.
II. Let us prove that the function $`u(w)`$ has local maxima at all centers $`w_c`$ of zero–angled triangles tesselating the Poicaré hyperbolic disc. Actually, the function $`f(w)`$ by construction gives a metric of some discrete subgroup of the group of motions of Poicaré hyperbolic disc. Hence the function $`f(w)`$ cannot grow faster than the isortopic hyperbolic metric $`(1w\overline{w})^2`$ and the following inequality is valid
$$0<u(w)1$$
for all points $`w`$ inside the unit disc. But we have shown that $`u(w)=1`$ at $`w=w_c`$ what means that the function $`u(w)`$ reaches its local maxima at the points $`w_c`$ and at all these maximal points the function $`u(w)`$ has one and the same value $`u(w_c)=1`$. The part II is proved.
## B
Our aim is to extract explicitly the $`\rho `$–dependence of the truncated series (65) and to connect it to $`I(1,q)`$. More precisely we are looking for a renormalization transformation of the form
$$I_𝒩(\rho ,q)=C\rho ^\kappa I_{𝒩^{}(𝒩,\rho )}(1,q)$$
(B1)
Using the correspondance (up to the volume of $`SO(2)`$) between $`PSL(2,\mathrm{𝖨𝖱})`$ and $`^2`$, we interpret the shift $`\tau \tau \nu _\rho `$ as a change of hyperbolic coordinates—see fig 7.
Note that the expression (65) does not depend on the particular representation of the hyperbolic 2-space, since the hyperbolic distance is invariant. The only one requirement is to define a compatible action of $`\mathrm{\Gamma }_\theta `$ in the space under consideration. We use for conveniency the unit disc representation whose center is the image of the point $`\nu _\rho (i)`$ where $`\nu _\rho `$ is defined by Eq.(61) in the $`^2`$ representation. For shortness the generic element of $`\mathrm{\Gamma }_\theta `$ is labelled by $`\tau `$ independent on the representation. We parametrize the point $`\tau \nu _\rho (i)`$ by its hyperbolic polar coordinates $`(d_1^\tau ,\alpha _1^\tau )`$ with the origin at the point $`i`$ and by “shifted” hyperbolic coordinates $`(d_2^\tau ,\alpha _2^\tau )`$ with the origin at the point $`\nu _\rho (i)`$. Note that $`d_\rho =d(i,\nu _\rho (i))(k+1)\mathrm{ln}\rho `$, and the following “triangle equation” in hyperbolic 2-space holds:
$$\mathrm{cosh}d_1^\tau =\mathrm{cosh}d_2^\tau \mathrm{cosh}d_\rho +\mathrm{sinh}d_2^\tau \mathrm{sinh}d_\rho \mathrm{cos}\alpha _2^\tau $$
(B2)
In order to extract the scaling conjectured in (66), we make an approximation which consists in neglecting fluctuations of $`W_n(d)`$. In this approximation we can sum over the generations $`n`$ and the angles $`\alpha _{j_n}`$ of the vertices within each generation ($`1j_n3\times 2^{n1}`$). Namely $`(d_2^\tau ,\alpha _2^\tau )=(\gamma n,\alpha _{j_n})`$. Thus one has
$$I_𝒩(\rho ,q)=\underset{n=1}{\overset{𝒩}{}}\left[2\mathrm{cosh}\gamma n\mathrm{cosh}d_\rho \right]^{2q}\underset{j_n=1}{\overset{3\times 2^{n1}}{}}(1+\mathrm{tanh}\gamma n\mathrm{tanh}d_\rho \mathrm{cos}\alpha _{j_n})^{2q}$$
(B3)
Assuming that $`\alpha _{j_n}`$ are uniformly distributed, we get for $`n1`$ the following expression
$$\underset{j_n=1}{\overset{3\times 2^{n1}}{}}(1+\mathrm{tanh}\gamma n\mathrm{tanh}d_\rho \mathrm{cos}\alpha _{j_n})^{2q}\frac{3\times 2^{n1}}{2\pi }_0^{2\pi }\frac{d\alpha }{(1+\mathrm{tanh}\gamma n\mathrm{tanh}d_\rho \mathrm{cos}\alpha )^{2q}}$$
(B4)
which leads to the asymptotic behavior:
$$\underset{n\mathrm{}}{lim}2^n\underset{j_n=1}{\overset{3\times 2^{n1}}{}}(1+\mathrm{tanh}\gamma n\mathrm{tanh}d_\rho \mathrm{cos}\alpha _{j_n})^{2q}|_\rho \mathrm{}\{\begin{array}{cc}\mathrm{const}\hfill & \text{for }q1/4\hfill \\ e^{2d_\rho (2q\frac{1}{2})}\hfill & \text{for }q>1/4\hfill \end{array}$$
(B5)
As justified hereafter we consider as relevant only the case $`q1/4`$. Using that for $`n1`$ and $`\rho 1`$ one has $`2\mathrm{cosh}\gamma n\mathrm{cosh}d_\rho \mathrm{cosh}(\gamma n+d_\rho )`$ we can rewrite Eq.(B3) as follows
$$I_𝒩(\rho ,q)=C\rho ^{(k+1)\mathrm{ln}2/\gamma }\underset{n=1}{\overset{𝒩}{}}2^{(n+d_\rho /\gamma )}[\mathrm{cosh}(\gamma n+d_\rho )]^{2q}$$
(B6)
Performing the shift $`\stackrel{~}{n}=n+d_\rho /\gamma `$ we get finally
$$I_𝒩(\rho ,q)=C\rho ^{(k+1)\mathrm{ln}2/\gamma }I_{𝒩+\frac{(k+1)}{\gamma }\mathrm{ln}\rho }(1,q)$$
(B7)
This expression fulfills the condition (B1). We assume that this renormalization also holds for the full function $`I_𝒩(\rho ,q)`$.
Let us pay attention to some contradiction between (63) and (B5) raised by the set of successive approximations of (63) which however is irrelevant for our final conclusions about multifractality. The equation (63) shows that if the integral over $`PSL(2,\mathrm{𝖨𝖱})`$ converges, it should not depend on $`\rho `$. Using the Poincaré series, the convergence of (63) occurs for $`q>q^{}`$. For $`q>1/4`$ and $`q>q^{}`$ the $`\rho `$–dependence shown in (B5) should then cancel by summing over all $`n`$. The discrepancy between (63) and (B5) appears for $`q[1/4,q^{}]`$. First of all we should note that the interval $`[\frac{1}{4},q^{}]`$ is numerically small (following previous sections we have $`q^{}q_0/20.4`$) and is nonuniversal, i.e. depends on the particular choose of the subgroup under consideration. Moreover, both the threshold $`q=1/4`$ and the asymptotics (B5) depend on the distribution of $`\alpha _{j_n}`$ and we believe that more careful treatment of angle dependence in (63) would allow disregard the region $`[1/4,q^{}]`$.
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# Hankel transform via double Hecke algebra
## 1 L-operator
We begin with the classical operator
$$=(\frac{}{x})^2+\frac{2k}{x}\frac{}{x}.$$
Upon the conjugation:
$$=|x|^k|x|^k,=(\frac{}{x})^2+\frac{k(1k)}{x^2}.$$
(1)
Here $`k`$ is a complex number. Both operators are symmetric $`=`$ even.
The $`\phi `$-function is introduced as follows:
$$\phi _\lambda (x,k)=4\lambda ^2\phi _\lambda (x,k),\phi _\lambda (x,k)=\phi _\lambda (x.k),\phi _\lambda (0,k)=1.$$
(2)
We will mainly write $`\phi _\lambda (x)`$ instead of $`\phi _\lambda (x,k).`$ Since $``$ is a DO of second order, the eigenvalue problem has a two-dimensional space of solutions. The even ones form a one-dimensional subspace and the normalization condition fixes $`\phi _\lambda `$ uniquely. Indeed, the operator $``$ preserves the space of even functions holomorphic at $`0.`$ The $`\phi _\lambda `$ can be of course constructed explicitly, without any references to the general theory of ODE.
We look for a solution in the form $`\phi _\lambda (x,k)=f(x\lambda ,k)`$. Set $`x\lambda =t`$. The resulting ODE is
$$\frac{d^2f}{dt^2}(t)+2k\frac{1}{t}\frac{df}{dt}(t)4f(t)=0,\text{ a Bessel-type equation. }$$
Its even normalized solution is given by the following series
$$f(t,k)=\underset{m=0}{\overset{\mathrm{}}{}}\frac{t^{2m}}{m!(k+1/2)\mathrm{}(k1/2+m)}=\mathrm{\Gamma }(k+\frac{1}{2})\underset{m=0}{\overset{\mathrm{}}{}}\frac{t^{2m}}{m!\mathrm{\Gamma }(k+1/2+m)}.$$
(3)
So
$$f(t,k)=\mathrm{\Gamma }(k+\frac{1}{2})t^{k+\frac{1}{2}}J_{k\frac{1}{2}}(2it).$$
The existence and convergence is for all $`t𝐂`$ subject to the constraint:
$$k1/2+n,n𝐙_+.$$
(4)
The symmetry $`\phi _\lambda (x,k)=\phi _x(\lambda ,k)`$ plays a very important role in the theory. Here it is immediate. In the multi-dimensional setup, it is a theorem.
Let us discuss other (nonsymmetric) solutions of (3) and (2). Looking for $`f`$ in the form $`t^\alpha (1+ct+\mathrm{})`$ in a neighborhood of $`t=0,`$ we get that the coefficients of the expansion
$$f(t)=t^{12k}\underset{m=0}{\overset{\mathrm{}}{}}c_mt^{2m}\text{ at }t=0$$
can be readily calculated from (3) and are well-defined for all $`k.`$
The convergence is easy to control. Generally speaking, such $`f`$ are neither regular nor even. To be precise, we get even functions $`f`$ regular at $`0`$ when $`k=1/2n`$ for an integer $`n0,`$ i.e. when (4) does not hold. These solutions cannot be normalized as above because they vanish at $`0.`$
Note that we do not need nonsymmetric $`f`$ and the corresponding $`\phi _\lambda (x)=f(x\lambda )`$ in the paper. Only even normalized $`\phi `$ will be considered. The nonsymmetric $`\psi `$-functions discussed in the next sections are of different nature.
###### Lemma 1.1
(a) Let $`^{}`$ be the adjoint operator of $``$ with respect to the $`𝐂`$-valued scalar product $`f,g_0=2_0^+\mathrm{}f(x)g(x)𝑑x`$. Then $`|x|^{2k}^{}|x|^{2k}=`$.
(b) Setting $`f,g=2_0^+\mathrm{}f(x)g(x)x^{2k}𝑑x,`$ the $``$ is self-adjoint with respect to this scalar product, i.e. $`(f),g=f,(g)`$.
Proof. First, the operator multiplication by $`x`$ is self-adjoint. Second, $`(\frac{}{x})^{}=\frac{}{x}`$ via integration by parts. Finally,
$`x^{2k}^{}x^{2k}=x^{2k}(({\displaystyle \frac{d^2}{dx^2}})^{}+({\displaystyle \frac{}{x}})^{}({\displaystyle \frac{2k}{x}}))x^{2k}=x^{2k}(({\displaystyle \frac{}{x}})^2{\displaystyle \frac{}{x}}({\displaystyle \frac{2k}{x}}))x^{2k}`$
$`=x^{2k}(x^{2k}({\displaystyle \frac{}{x}})^2+4kx^{2k1}{\displaystyle \frac{}{x}}+2k(2k1)x^{2k2}2kx^{2k1}{\displaystyle \frac{}{x}}2k(2k1)x^{2k2})`$
$`=({\displaystyle \frac{}{x}})^2+{\displaystyle \frac{2k}{x}}{\displaystyle \frac{}{x}}=.`$ (5)
Therefore, $`(f),g=2_0^{\mathrm{}}(f)gx^{2k}𝑑x=`$
$`2{\displaystyle _{𝐑_+}}f^{}(x^{2k}g)𝑑x=2{\displaystyle _{𝐑_+}}fx^{2k}x^{2k}(x^{2k}g)𝑑x=f,(g).`$ (6)
Actually this calculation is not necessary if (1) is used. Indeed, $`^{}=.`$ $`\mathrm{}`$
## 2 Hankel transform
Let us define the symmetric Hankel transform on the space of continuous functions $`f`$ on $`𝐑`$ such that $`lim_x\mathrm{}f(x)e^{cx}=0`$ for any $`c𝐑`$. Provided (4),
$$(𝔽_kf)(\lambda )=\frac{2}{\mathrm{\Gamma }(k+1/2)}_0^+\mathrm{}\phi _\lambda (x,k)f(x)x^{2k}𝑑x.$$
(7)
The growth condition makes the transform well-defined for all $`\lambda 𝐂,`$ because
$$\phi _\lambda (x,k)\text{Const}(e^{2\lambda x}+e^{2\lambda x})\text{ at }x=\mathrm{}.$$
The latter is standard.
We switch from $`𝔽`$ on functions to the transform of the operators: $`𝔽(A)(𝔽(f)=𝔽(A(f)).`$ Remark that the Hankel transform of the function is very much different from the transform of the corresponding multiplication operator. The key point of the operator technique is the following lemma.
###### Lemma 2.1
Using the upper index to denote the variable ($`x`$ or $`\lambda `$),
$$(a)𝔽(^x)=4\lambda ^2;(b)𝔽(4x^2)=^\lambda ;(c)𝔽(4x\frac{}{x})=4\lambda \frac{d}{d\lambda }48k.$$
Proof. Claim (a) is a direct consequence of Lemma 1.1 (b) with $`g(x)=\phi _\lambda (x):`$
$$𝔽(f)=f,\phi _\lambda =f,\phi _\lambda =4\lambda ^2f,\phi _\lambda =4\lambda ^2𝔽(f)$$
. Claim (b) results directly from the $`x\lambda `$ symmetry of $`\varphi ,`$ namely, from the relation $`^\lambda \phi _\lambda (x)=4x^2\psi _\lambda (x)`$. Concerning (c), there are no reasons, generally speaking, to expect any simple Fourier transforms for the operators different from $`.`$ However in this particular case: $`[^x,x^2]=4x\frac{}{x}+2+4k`$. Appling $`𝔽`$ to both sides and using (a), (b), $`[4\lambda ^2,^\lambda /4]=𝔽(4x\frac{}{x})+2+4k`$. Finally
$$𝔽(4x\frac{}{x})=4\lambda \frac{d}{d\lambda }24k24k=4\lambda \frac{d}{d\lambda }48k$$
.
Note that
$$[x\frac{}{x},x^2]=2x^2,[x\frac{}{x},^x]=2^x,$$
because operators $`x\frac{}{x},`$ are homogeneous of degree $`2`$ and $`2`$. So $`e=x^2,`$ $`f=^x/4,`$ and $`h=x\frac{}{x}+k+1/2=[e,f]`$ generate a representation of the Lie algebra $`sl_2(𝐂).`$ $`\mathrm{}`$
###### Theorem 2.2
(Master Formula) Assuming that $`\mathrm{𝖱𝖾}k>\frac{1}{2},`$
$`2{\displaystyle _0^{\mathrm{}}}\phi _\lambda (x)\phi _\mu (x)e^{x^2}x^{2k}𝑑x=\mathrm{\Gamma }(k+{\displaystyle \frac{1}{2}})e^{\lambda ^2+\mu ^2}\phi _\lambda (\mu ),`$ (8)
$`2{\displaystyle _0^{\mathrm{}}}\phi _\lambda (x)\mathrm{exp}({\displaystyle \frac{}{4}})(f(x))e^{x^2}x^{2k}𝑑x=\mathrm{\Gamma }(k+{\displaystyle \frac{1}{2}})e^{\lambda ^2}f(\lambda ),`$
provided the existence of $`\mathrm{exp}(\frac{}{4})(f(x))`$ and the integral in the second formula.
Proof. The left-hand side of the first formula equals $`\mathrm{\Gamma }(k+1/2)𝔽(e^{x^2}\phi _\mu (x))`$. We set
$$\phi _\mu ^{}(x)=e^{x^2}\phi _\mu (x),\phi _\mu ^+(x)=e^{x^2}\phi _\mu (x).$$
They are eigenfunctions of the operators
$$_{}=e^{x^2}e^{x^2},_+=e^{x^2}e^{x^2}.$$
To be more exact, $`\phi _\mu ^\pm `$ is a unique eigenfunction of $`^\pm `$ with eigenvalue $`2\mu `$, normalized by $`\phi _\mu ^\pm (0)=1`$.
Express $`_{}`$ in terms of the operators from the previous lemma.
$`e^{x^2}({\displaystyle \frac{}{x}})^2e^{x^2}`$ $`=e^{x^2}(e^{x^2}({\displaystyle \frac{}{x}})^2+2(2x)e^{x^2}{\displaystyle \frac{}{x}}+(2+4x^2)e^{x^2})=({\displaystyle \frac{}{x}})^2+4x{\displaystyle \frac{}{x}}+2+4x^2,`$
$`e^{x^2}{\displaystyle \frac{2k}{x}}{\displaystyle \frac{}{x}}e^{x^2}`$ $`=e^{x^2}(e^{x^2}{\displaystyle \frac{2k}{x}}{\displaystyle \frac{}{x}}+2xe^{x^2}{\displaystyle \frac{2k}{x}})={\displaystyle \frac{2k}{x}}{\displaystyle \frac{}{x}}+4k,`$
$`_{}`$ $`=e^{x^2}(({\displaystyle \frac{}{x}})^2+{\displaystyle \frac{2k}{x}}{\displaystyle \frac{}{x}})e^{x^2}=+4x{\displaystyle \frac{}{x}}+2+4k+4x^2.`$ (9)
Analogously, $`_+=4x\frac{}{x}24k+4x^2`$. Now we may use Lemma 2.1:
$`𝔽(_{}^x)=`$ $`𝔽(^x)+𝔽(4x^2)+𝔽(4x{\displaystyle \frac{}{x}})+𝔽(2+4k)=`$
$`=`$ $`4\lambda ^2+^\lambda 4\lambda {\displaystyle \frac{d}{d\lambda }}48k+2+4k=_+^\lambda .`$ (10)
Thus
$$L_+^\lambda (𝔽\phi _\mu ^{})=𝔽(_{}^x)(𝔽\phi _\mu ^{})=𝔽(_{}^x\phi _\mu ^{})=2\mu 𝔽(\phi _\mu ^{}),$$
i.e. $`𝔽\phi _\mu ^{}`$ is an eigenfunction of $`_+`$ with the eigenvalue $`2\mu `$. Using the uniqueness, we conclude that $`𝔽(\phi _\mu ^{})(\lambda )=`$ $`C(\mu )e^{\mu ^2}\phi _\mu ^+(\lambda )`$ for a constant $`C(\mu ).`$ However the left-hand side of the master formula is $`\lambda \mu `$ symmetric as well as $`e^{\mu ^2}\phi _\mu ^+(\lambda )=e^{\lambda ^2+\mu ^2}\phi _\mu (\lambda ).`$ So $`C(\mu )=C(\lambda )=C.`$ Setting $`\lambda =0=\mu ,`$ we get the desired.
The second formula follows from the first for $`f(x)=\phi _\mu (x,k).`$ Move $`\mathrm{exp}(\mu ^2)`$ to the left to see this. It is linear in terms of $`f(x)`$ and holds for finite linear combinations of $`\phi `$ and infinite ones provided the convergence. So it is valid for any reasonable $`f.`$ We skip the detail. $`\mathrm{}`$
## 3 Dunkl operator
The above proof is straightforward. One needs the self-duality of the Hankel transform and the the commutator representation for $`x/x.`$ The self-duality holds in the general multi-dimensional theory. The second property is more special. Also our proof does not clarify why the master formula is so simple. There is a “one-line” proof of this important formula, which can be readily generalized. It involves the Dunkl operator:
$$𝒟=\frac{}{x}\frac{k}{x}(s1),\text{ where }s\text{ is the reflection }s(f(x))=f(x).$$
(11)
The operator $`𝒟`$ is not local anymore, because $`s`$ is a global operator apart from a neighborhood of $`x=0.`$ We are going to find its eigenfunctions. Generally speaking, this may create problems since we cannot use the uniquness theorems from the theory of ODE. However everything is surprisingly smooth.
###### Lemma 3.1
Considering $`x`$ as the multiplication operator,
$$sx=xs,s\frac{}{x}=\frac{}{x}s,$$
(12)
$`(a)𝒟^2=`$ upon the restriction to even functions,
$`(b)s𝒟s=𝒟`$ and $`𝒟^2`$ fixes the space of even functions.
Proof. Indeed, $`(sx)(f(x))=s(xf(x))=xf(x)=(xs)(f(x)).`$ The $`\frac{}{x}`$ is analogous. Then
$`𝒟^2`$ $`=({\displaystyle \frac{}{x}})^2{\displaystyle \frac{k}{x}}(s1){\displaystyle \frac{}{x}}{\displaystyle \frac{}{x}}{\displaystyle \frac{k}{x}}(s1)+{\displaystyle \frac{k}{x}}(s1){\displaystyle \frac{k}{x}}(s1)`$
$`=({\displaystyle \frac{}{x}})^2+{\displaystyle \frac{k}{x}}{\displaystyle \frac{}{x}}(s+1){\displaystyle \frac{}{x}}{\displaystyle \frac{k}{x}}(s1)+{\displaystyle \frac{k}{x}}(s1){\displaystyle \frac{k}{x}}(s1).`$ (13)
It is simple to calculate the final formula but unnecessary. Applying (3) to symmetric (i.e. even) functions $`f(x)`$, the two last terms will vanish, because $`(s1)(f(x))=f(x)f(x)=0.`$ So $`(s+1)(f(x))=f(x)+f(x)=2f(x),`$ and $`𝒟^2|_{even}=(\frac{}{x})^2+2\frac{k}{x}\frac{}{x}=.`$
Claim (b) is obvious. Indeed, $`s^2=1,`$ $`s\frac{}{x}s=\frac{}{x}s^2=dx`$, and $`s(\frac{k}{x}(s1))s=\frac{k}{x}s(s^2s)=\frac{k}{x}(s^3s^2)=\frac{k}{x}(s1)`$. Thus $`s𝒟s=𝒟`$.
By the way, this implies that $`s𝒟^2s=𝒟^2`$, i.e. $`𝒟^2`$ commutes with $`s.`$ So we do not need an explicit formula for $`𝒟^2|_{even}`$ to see that it preserves even functions. $`\mathrm{}`$
Let us consider the standard scalar product $`f,g_0=_{\mathrm{}}^+\mathrm{}f(x)g(x)𝑑x.`$ Here the functions are continuous $`𝐂`$-valued continues on the real line $`𝐑.`$ One may add the complex conjugation to $`g`$ but we will not do this. The scalar product is non-degenerate, so adjoint operators are well-defined. We continue to use the notation $`H^{}`$ for the pairing $`f,g_0.`$ Let us calculate the adjoint of $`𝒟`$ with respect to $`|x|^{2k}.`$
###### Proposition 3.2
Setting $`f,g=_{\mathrm{}}^+\mathrm{}f(x)g(x)|x|^{2k}𝑑x,`$ the Dunkl operator $`𝒟`$ is anti self-adjoint with respect to this scalar product, i.e. $`𝒟(f),g=f,𝒟(g)`$. Equivalently, $`|x|^{2k}𝒟^{}|x|^{2k}=𝒟`$.
Proof. Recall that $`x^{}=x`$ and $`(\frac{}{x})^{}=\frac{}{x},`$ where $`x`$ is considered as the multiplication operator. Then $`s^{}=s:`$
$$s(f),g_0=_{\mathrm{}}^+\mathrm{}f(x)g(x)𝑑x=_+\mathrm{}^{\mathrm{}}f(t)g(t)(dt)=f,s(g)_0$$
for $`t=x.`$ Hence,
$`|x|^{2k}𝒟^{}|x|^{2k}`$ $`=|x|^{2k}({\displaystyle \frac{}{x}}{\displaystyle \frac{k}{x}}(s1))^{}|x|^{2k}=|x|^{2k}(({\displaystyle \frac{}{x}})^{}(s1)^{}({\displaystyle \frac{k}{x}})^{})|x|^{2k}`$
$`=|x|^{2k}({\displaystyle \frac{}{x}}(s1){\displaystyle \frac{k}{x}})|x|^{2k}=|x|^{2k}({\displaystyle \frac{}{x}}+{\displaystyle \frac{k}{x}}(1+s))|x|^{2k}`$
$`=|x|^{2k}|x|^{2k}({\displaystyle \frac{}{x}})+|x|^{2k}({\displaystyle \frac{2k}{x}}|x|^{2k})+|x|^{2k}|x|^{2k}{\displaystyle \frac{k}{x}}(1+s)`$
$`={\displaystyle \frac{}{x}}+{\displaystyle \frac{k}{x}}(s1)=𝒟.`$ (14)
Finally,
$`𝒟(f),g={\displaystyle _{\mathrm{}}^+\mathrm{}}𝒟(f(x))g(x)|x|^{2k}𝑑x={\displaystyle _{\mathrm{}}^+\mathrm{}}f(x)𝒟^{}(|x|^{2k}g(x))𝑑x`$ (15)
$`={\displaystyle _{\mathrm{}}^+\mathrm{}}f(x)|x|^{2k}(|x|^{2k}𝒟^{}|x|^{2k})(g(x))𝑑x`$
$`={\displaystyle _{\mathrm{}}^+\mathrm{}}f(x)(𝒟(g(x)))|x|^{2k}dx=f,𝒟(g).\mathrm{}`$
The proposition readily gives that $`|x|^{2k}^{}|x|^{2k}=`$ on even functions $`f.`$ Indeed,
$$(f),g=𝒟^2(f),g=f,𝒟^2(g)=f,(g),$$
provided that $`g`$ is even too. Recall that it was not difficult to check this relation directly. In the multi-dimensional theory, this calculation is more involved and the usage of the (generalized) Dunkl operators makes perfect sense.
## 4 Nonsymmetric eigenfunctions
Our next step will be a study of the eigenfunctions of the Dunkl operator:
$$𝒟\psi _\lambda (x,k)=2\lambda \psi _\lambda (x,k),\psi _\lambda (0,k)=1.$$
(16)
We will use the shortcut notation $`f^\iota (x)=s(f(x))=f(x).`$
###### Lemma 4.1
There exists a unique solution of the eigenfunction problem (16) for $`\lambda 0`$. It is represented in the form $`\psi _\lambda (x)=g(\lambda x).`$ In the case $`\lambda =0,`$ the solution is given by the formula $`\psi _0=1+Cx|x|^{2k1}`$, where $`C𝐂`$ is an arbitrary constant.
Proof. Assuming that $`\psi _\lambda `$ is a solution of (16), let
$$\psi _\lambda ^0=\frac{1}{2}(\psi _\lambda +\psi _\lambda ^\iota ),\psi _\lambda ^1=\frac{1}{2}(\psi _\lambda \psi _\lambda ^\iota ),$$
be its even and odd parts. By Lemma 3.1 (b), $`𝒟s(\psi _\lambda (x))=s𝒟(\psi _\lambda (x))=2\lambda s(\phi _\lambda (x))`$. Hence, (16) is equivalent to
$`𝒟\psi _\lambda ^0=2\lambda \psi _\lambda ^1;\psi _\lambda ^0(0)=1`$
$`𝒟\psi _\lambda ^1=2\lambda \psi _\lambda ^0\psi _\lambda ^1(0)=0.`$ (17)
Furthermore, $`𝒟^2\psi _\lambda ^0=4\lambda ^2\psi _\lambda ^0`$. Since $`\psi _\lambda ^0`$ is even, $`\psi _\lambda ^0=4\lambda ^2\psi _\lambda `$ due to Lemma 3.1. Therefore $`\psi _\lambda ^0`$ has to coincide with $`\phi _\lambda `$ from the first section. This is true for all $`\lambda .`$ If $`\lambda 0,`$
$$\psi _\lambda ^1=\frac{1}{2\lambda }𝒟\psi _\lambda ^0=\frac{1}{2\lambda }\left(\frac{d\psi _\lambda ^0}{dx}\frac{k}{x}(s1)\psi _\lambda ^0\right)=\frac{1}{2\lambda }\frac{d\phi _\lambda }{dx}.$$
(18)
The last equality holds because $`\psi _\lambda ^0=\phi _\lambda `$ is even. Finally,
$$\psi _\lambda (x)=\phi _\lambda (x)+\frac{1}{2\lambda }\phi _\lambda ^{}(x)=g(\lambda x)\text{ for }g=f+\frac{1}{2}f^{},$$
(19)
where $`\phi _\lambda (x)=f(\lambda x),`$ $`f`$ is from (3), and $`f^{}`$ is the derivative. It is for $`\lambda 0.`$
Let us consider the case $`\lambda =0`$. We have $`𝒟\psi _\lambda ^0=0`$, $`\psi _\lambda ^0(0)=1`$ and $`𝒟\psi _\lambda ^1=0`$, $`\psi _\lambda ^1(0)=0`$. Since $`\psi _\lambda ^0`$ is even, $`𝒟\psi _\lambda ^0=\frac{d\psi _\lambda ^0}{dx}=0`$. Thus $`\psi _\lambda ^0(x)=1`$. The $`\psi _\lambda ^1`$ is odd. So
$$𝒟\psi _\lambda ^1(x)=\frac{d\psi _\lambda ^1(x)}{dx}\frac{k}{x}(s1)\psi _\lambda ^1(x)=\frac{d\psi _\lambda ^1(x)}{dx}+\frac{k}{x}(\psi _\lambda ^1(x)\psi _\lambda ^1(x))=\frac{d\psi _\lambda ^1(x)}{dx}+\frac{2k}{x}\psi _\lambda ^1(1).$$
Solving the resulting ODE, $`\psi _\lambda ^1(x)=Cx|x|^{2k1}`$. $`\mathrm{}`$
In this proof, we used that $`𝒟f(x)=f^{}(x)`$ on even functions and
$$𝒟f(x)=f^{}(x)+\frac{k}{x}(f(x)f(x))=f^{}(x)+\frac{2k}{x}f(x)=(\frac{}{x}+\frac{2k}{x})f(x)$$
on odd functions. By the way, it makes obvious the coincidence of $``$ with $`𝒟^2`$ on even $`f`$. Indeed, $`𝒟^2f(x)=(\frac{}{x}+\frac{2k}{x})(\frac{}{x}f(x)).`$ For odd $`f,`$ it is the other way round: $`𝒟^2f(x)=𝒟(𝒟f(x))=`$ $`\frac{}{x}(𝒟f(x))=\frac{}{x}(\frac{}{x}+\frac{2k}{x})(f(x)).`$ In particular,
$$𝒟^2\psi _\lambda ^1(x)=\frac{}{x}(\frac{}{x}+\frac{2k}{x})\psi _\lambda ^1(x)=((\frac{}{x})^2+\frac{}{x}\frac{2k}{x})\psi _\lambda ^1(x).$$
Hence $`\psi _\lambda ^1(x)`$ is also a solution of a second order differential equation. This equation is different from that for $`\phi ,`$ but not too different. Comparing them we come to the important definition of the shift operator. We show the dependence of $``$ on $`k`$ and set $`\stackrel{~}{_k}=(\frac{}{x})^2+\frac{}{x}\frac{2k}{x}`$.
###### Lemma 4.2
(a) $`x^1\stackrel{~}{_k}x=_{k+1}`$.
(b) $`\stackrel{~}{_k}\psi _\lambda ^1=4\lambda ^2\psi _\lambda ^1`$.
(c) $`(x^1\stackrel{~}{_k}x)(x^1\psi _\lambda ^1)=4\lambda ^2(x^1\psi _\lambda ^1)`$.
Proof. The first claim:
$`x^1({\displaystyle \frac{}{x}})^2x`$ $`=x^1(x({\displaystyle \frac{}{x}})^2+2{\displaystyle \frac{}{x}})=({\displaystyle \frac{}{x}})^2+{\displaystyle \frac{2}{x}}{\displaystyle \frac{}{x}},`$
$`x^1{\displaystyle \frac{}{x}}{\displaystyle \frac{2k}{x}}x`$ $`=x^1{\displaystyle \frac{}{x}}2k={\displaystyle \frac{2k}{x}}{\displaystyle \frac{}{x}},`$
$`x^1\stackrel{~}{_k}x`$ $`=({\displaystyle \frac{}{x}})^2+{\displaystyle \frac{2}{x}}{\displaystyle \frac{}{x}}+{\displaystyle \frac{2k}{x}}{\displaystyle \frac{}{x}}=({\displaystyle \frac{}{x}})^2+{\displaystyle \frac{2(k+1)}{x}}dx=_{k+1}.`$ (20)
Then $`\stackrel{~}{_k}\psi _\lambda ^1=𝒟^2\psi _\lambda ^1=4\lambda ^2\psi _\lambda ^1`$ due to (4). Claim (c) is a combination of (a) and (b). $`\mathrm{}`$
###### Proposition 4.3
(Shift Formula)
$$\frac{1}{x}\psi _\lambda ^1(x,k)=\frac{2\lambda }{1+2k}\psi _\lambda ^0(x,k+1),\text{ i.e. }\frac{1}{x}\frac{d\phi _\lambda }{dx}(x,k)=\frac{4\lambda ^2}{1+2k}\phi _\lambda (x,k+1).$$
(21)
Proof. Lemma 4.2 (c) implies that $`x^1\psi _\lambda ^1(x,k)=C(\lambda ,k)\phi _\lambda (x,k+1)`$, because $`\phi _\lambda (x,k+1)`$ is a unique even normalized solution of (2) for $`k+1.`$ Thanks to (18) $`\psi _\lambda ^1(x,k)=(2\lambda )^1\frac{d\phi _\lambda }{dx}(x,k)`$. Thus $`x^1\phi _\lambda ^{}(x,k)=C(\lambda ,k)\phi _\lambda (x,k+1)`$. The constant $`C`$ readily results from the expantion (3) of $`\phi _\lambda (x,k).`$ Explicitly:
$`0`$ $`=(_k\phi _\lambda 4\lambda ^2\phi _\lambda )(0)`$
$`0`$ $`=(2k+1)(x^1{\displaystyle \frac{d\phi _\lambda }{dx}})(0,k)4\lambda ^2.`$ (22)
The shift formula can be of course checked directly without $`\psi _\lambda ^1,`$ a good exercise. $`\mathrm{}`$
## 5 Master formula
Let us define the nonsymmetric Hankel transform. We consider complex-valued $`C^{\mathrm{}}`$ functions $`f`$ on $`𝐑`$ such that $`lim_x\mathrm{}f(x)e^{cx}=0`$ for any $`c𝐑`$ and set
$$(f)(\lambda )=\frac{1}{\mathrm{\Gamma }(k+1/2)}_{\mathrm{}}^+\mathrm{}\psi _\lambda (x,k)f(x)|x|^{2k}𝑑x$$
(23)
We assume that $`\mathrm{𝖱𝖾}k>\frac{1}{2}`$ and always take $`\psi _0(x,k)=1.`$ Recall that the case $`\lambda =0`$ is exceptional (Lemma 4.1 ): the dimension of the space of eigenfunctions is $`2.`$
Let us compute the transforms of our main operators. Compare it with Lemma 2.1 : it is much more comfortable to deal with the operators of the first order. The upper index denotes the variable.
###### Lemma 5.1
$$(a)(𝒟^x)=2\lambda ;(b)(2x)=𝒟^\lambda ;(c)(s^x)=s^\lambda .$$
Proof. The first formula is an immediate consequence of Proposition 3.2 (a) with $`g(x)=\psi _\lambda (x):`$
$$(𝒟f)=𝒟f,\psi _\lambda =f,𝒟\psi _\lambda =2\lambda f,\psi _\lambda =2\lambda (f).$$
Claim (b) follows from the $`x\lambda `$ symmetry. As to (c), use that $`\psi _\lambda (x)=\psi _\lambda (x).`$ $`\mathrm{}`$
###### Theorem 5.2
(Nonsymmetric Master Formula)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\psi _\lambda (x)\psi _\mu (x)e^{x^2}|x|^{2k}𝑑x=\mathrm{\Gamma }(k+{\displaystyle \frac{1}{2}})e^{\lambda ^2+\mu ^2}\psi _\lambda (\mu ),`$ (24)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\psi _\lambda (x)\mathrm{exp}({\displaystyle \frac{𝒟^2}{4}})(f(x))e^{x^2}|x|^{2k}𝑑x=\mathrm{\Gamma }(k+{\displaystyle \frac{1}{2}})e^{\lambda ^2}f(\lambda ).`$
Proof. In the first formula, the left-hand side equals $`\mathrm{\Gamma }(k+1/2)(e^{x^2}\psi _\mu (x))`$. We set
$$\psi _\mu ^{}(x)=e^{x^2}\psi _\mu (x),\psi _\mu ^+(x)=e^{x^2}\psi _\mu (x),$$
$$𝒟_{}=e^{x^2}𝒟e^{x^2},𝒟_+=e^{x^2}𝒟e^{x^2}.$$
The function $`\psi _\mu ^\pm `$ is an eigenfunction of $`𝒟^\pm `$ with eigenvalue $`2\mu .`$ The normalization fixes it uniquely with the standard reservation about $`\mu =0.`$
One gets:
$$𝒟_{}=e^{x^2}(\frac{}{x}\frac{k}{x}(s1))e^{x^2}=\frac{}{x}+2x\frac{k}{x}(s1)=𝒟+2x.$$
Correspondingly, $`𝒟_+=𝒟2x`$. Using Lemma 5.1 ,
$$(𝒟_{}^x)=(𝒟^x)+(2x)=2\lambda +𝒟^\lambda =𝒟_+^\lambda .$$
Therefore
$$D_+^\lambda (\psi _\mu ^{})=(𝒟_{}^x)(\psi _\mu ^{})=(𝒟_{}^x\psi _\mu ^{})=2\mu (\psi _\mu ^{}),$$
i.e. $`(\psi _\mu ^{})`$ is an eigenfunction of $`𝒟_+`$ with the eigenvalue $`2\mu ,`$ and $`(\psi _\mu ^{})(\lambda )=C(\mu )e^{\mu ^2}\psi _\mu ^+(\lambda )`$. Using the $`\lambda \mu `$-symmetry, $`C(\mu )=C(\lambda )=C`$ and
$$C=_𝐑e^{x^2}|x|^{2k}𝑑x=\mathrm{\Gamma }(k+\frac{1}{2}).$$
Cf. the proof of the symmetric master formula.
The second formula readily follows from the first provided the existence of the function $`\mathrm{exp}(\frac{𝒟^2}{4})(f(x))`$ and the corresponding integral. The latter function has to go to zero at $`x=\mathrm{}`$ faster than $`e^{cx}`$ for any $`c𝐑.`$ $`\mathrm{}`$
The symmetric master theorem is of course a particular case of (24). Indeed, we may replace $`\psi _\lambda (x)`$ by $`2\phi _\lambda (x)=\psi _\lambda (x)+\psi _\lambda (x)=\psi _\lambda (x)+\psi _\lambda (x)`$ on the left-hand side. Then
$$\psi _\lambda (\mu )\psi _\lambda (\mu )+\psi _\lambda (\mu )=\psi _\lambda (\mu )+\psi _\lambda (\mu )=2\phi _\lambda (\mu )$$
on the right-hand side. We use that the factor $`e^{\lambda ^2+\mu ^2}`$ is even. Now we either repeat the same transfer for $`\mu `$ or simply symmetrize the integrand.
It is more surprising that the nonsymmetric theorem can be deduced from the symmetric one. It is a special feature of the one-dimensional case. Generally speaking, there is no reason to expect such an implication. This may clarify why the nonsymmetric Hankel transform and $`\psi `$ were of little importance in the classical theory of Bessel functions. They could be considered as a minor technical improvement of the symmetric theory. Now we have the opposite point of view.
Let us deduce Theorem 5.2 from Theorem 2.2. We may assume that $`\lambda ,\mu 0.`$ Discarding the odd summands in the integrand,
$`{\displaystyle _𝐑}\psi _\lambda \psi _\mu e^{x^2}|x|^{2k}𝑑x={\displaystyle _𝐑}(\psi _\lambda ^0+\psi _\lambda ^1)(\psi _\mu ^0+\psi _\mu ^1)e^{x^2}|x|^{2k}𝑑x`$
$`={\displaystyle _𝐑}(\psi _\lambda ^0\psi _\mu ^0+\psi _\lambda ^1\psi _\mu ^1)e^{x^2}|x|^{2k}𝑑x={\displaystyle _𝐑}(\phi _\lambda \phi _\mu +\psi _\lambda ^1\psi _\mu ^1)e^{x^2}|x|^{2k}𝑑x.`$ (25)
The integral of $`\phi _\lambda \phi _\mu `$ is nothing else but (8). Let us use the shift formula to manage $`\psi _\lambda ^1\psi _\mu ^1.`$ See Proposition 4.3.
We get $`(\psi _\lambda ^1\psi _\mu ^1)(x,k)=\frac{4\lambda \mu x^2}{(1+2k)^2}(\phi _\lambda \phi _\mu )(x,k+1)`$ and
$`{\displaystyle _𝐑}(\psi _\lambda ^1\psi _\mu ^1)(x,k)e^{x^2}|x|^{2k}𝑑x={\displaystyle \frac{4\lambda \mu }{(1+2k)^2}}{\displaystyle _𝐑}(\phi _\lambda \phi _\mu )(x,k+1)e^{x^2}|x|^{2(k+1)}𝑑x`$
$`={\displaystyle \frac{4\lambda \mu }{(1+2k)^2}}\phi _\lambda (\mu ,k+1)e^{\lambda ^2+\mu ^2}\mathrm{\Gamma }(k+{\displaystyle \frac{3}{2}})={\displaystyle \frac{2\lambda \mu }{(1+2k)}}\phi _\lambda (\mu ,k+1)e^{\lambda ^2+\mu ^2}\mathrm{\Gamma }(k+{\displaystyle \frac{1}{2}})`$
$`=\psi _\lambda ^1(\mu ,k)e^{\lambda ^2+\mu ^2}\mathrm{\Gamma }(k+{\displaystyle \frac{1}{2}}).`$ (26)
This concludes the deduction. $`\mathrm{}`$
## 6 Double H double prime
Let $`^{\prime \prime }`$ be the double degeneration of the double affine Hecke algebra:
$$^{\prime \prime }=,x,ssxs=x,ss=,[,x]=1+2ks.$$
(27)
Its polynomial representation $`\rho :^{\prime \prime }\mathrm{End}(𝒫)`$ in $`𝒫=𝐂[x]`$ is as follows:
$$\rho (x)=\text{ multiplication by }x,\rho (s)=s,\rho ()=𝒟,$$
where $`s`$ is the reflection $`ff^\iota ,`$ $`D`$ is the Dunkl operator. The first two of the defining relations of $`^{\prime \prime }`$ are satisfied thanks to Lemma 3.1. As to the last,
$`[𝒟,x]=({\displaystyle \frac{}{x}}{\displaystyle \frac{k}{x}}(s1))xx({\displaystyle \frac{}{x}}{\displaystyle \frac{k}{x}}(s1))`$
$`=`$ $`x{\displaystyle \frac{}{x}}+1x{\displaystyle \frac{k}{x}}(s1)x{\displaystyle \frac{}{x}}+x{\displaystyle \frac{k}{x}}(s1)=1+2ks.`$ (28)
###### Theorem 6.1
(a) Any nonzero finite linear combination $`H={\displaystyle \underset{m,n,ϵ}{}}c_{m,n,ϵ}x^m^ns^ϵ`$, where $`n,m𝐙_+`$, $`ϵ\{0,1\}`$, acts as a nonzero operator in $`𝒫`$.
(b) Any element $`H^{\prime \prime }`$ can be uniquely expressed in the form $`H={\displaystyle \underset{m,n,ϵ}{}}x^m^ns^ϵ.`$ The representation $`\rho `$ is faithful for any $`k.`$
Proof. Let $`H={\displaystyle \underset{n,ϵ}{}}f_n^ns^ϵ={\displaystyle \underset{m=0}{\overset{N}{}}}(g_m^m)(1+s)+{\displaystyle \underset{m=0}{\overset{M}{}}}(h_m^m)(1s)`$ for polynomials $`f_m,g_m,`$ and $`h_m.`$ We may assume that and at least one of the leading coefficients $`g_N(x),`$ $`h_M(x)`$ is nonzero. Then $`\rho (H)=L_+(1+s)+L_{}(1s)`$ for differential operators $`L_+=g_N(x)(\frac{}{x})^N+\mathrm{}`$ and $`L_{}=h_M(x)(\frac{}{x})^M+\mathrm{}`$ modulo differential operators of lower orders. Applying $`\rho (H)`$ to even and odd functions, we get that $`\rho (H)=0`$ implies that both $`L_+`$ and $`L_{}`$ have infinite dimensional spaces of eigenfunctions. This is impossible. Claim (a) is verified.
Concerning (b), any element $`H^{\prime \prime }`$ can be obviously expressed in the desired form. Such expression is unique and the representation $`\rho `$ is faithful thanks to (a). $`\mathrm{}`$
The theorem is the key point of the representation theory of the double H. It is a variant of the so-called PBW theorem. There are not many algebras in mathematics and physics possessing this property. All have important applications. The double Hecke algebra is one of them.
Next, we study the irreducibility of $`\rho `$.
###### Lemma 6.2
The Dunkl operator $`𝒟`$ has only one eigenvalue in $`𝒫`$, namely, $`\lambda =0`$. If $`k1/2n`$ for any $`n𝐙_+`$, then $`𝒟`$ has a unique (up to a constant) eigenfunction in $`𝒫`$, the constant function $`1`$. When $`k=1/2n`$ for $`n𝐙_+`$, the space of $`0`$-eigenfunctions is $`𝐂+𝐂x^{2n+1}`$.
Proof. Let $`p(x)𝒫`$ be an eigenfunction for $`𝒟`$. Since $`𝒟`$ lowers the degree of any polynomial by 1, we have $`𝒟^{m+1}p=0`$, where $`m=\mathrm{deg}p`$. Therefore all eigenvalues of $`𝒟`$ are zero. Representing $`p`$ as the sum $`p(x)=p^0(x)+p^1(x)`$ of even $`p^0`$ and odd $`p^1,`$ $`𝒟p=\frac{}{x}p^0+(\frac{}{x}p^1+\frac{2k}{x}p^1)=0`$. Both the even and the odd parts of this expression have to be zero. Hence $`\frac{}{x}p^0=0`$ and $`\frac{}{x}p^1+\frac{2k}{x}p^1=0`$. Therefore $`p^0=`$Const. Setting $`p^1(x)=a_lx^{2l+1},`$
$$(\frac{}{x}+\frac{2k}{x})p^1(x)=a_l(2l+1+2k)x^{2l}=0.$$
If $`k1/2n`$ for any $`n𝐙_+`$ then $`a_l=0`$ for any $`l`$, i.e $`p^1=0.`$ Otherwise $`k=1/2n`$ for a certain $`n𝐙_+`$ and $`p^1(x)`$ is proportional to $`x^{2n+1}`$. $`\mathrm{}`$
###### Theorem 6.3
(a) The representation $`\rho `$ is irreducible if and only if $`k1/2n`$ for any $`n𝐙_+`$.
(b) If $`k=1/2n`$ for $`n𝐙_+`$, then there exists a unique non-trivial $`^{\prime \prime }`$-submodule of $`𝒫`$, namely, $`(x^{2n+1})=x^{2n+1}𝒫`$.
Proof. Let $`\{0\}V𝒫`$ be a $`^{\prime \prime }`$ submodule of $`𝒫`$. Let $`0vV`$ and $`m=\mathrm{deg}v`$. Set $`𝒫^{(m)}=\{p𝒫|\mathrm{deg}pm\}`$. We have $`V^{(m)}=V𝒫^{(m)}\{0\}`$. Then $`V^{(m)}`$ is $`𝒟`$ invariant. More exactly, $`𝒟(V^{(m)})V^{(m1)}.`$ Thus it contains an eigenfunction $`v_0`$ of $`𝒟`$. If $`k1/2n`$ for $`n𝐙_+`$ then Lemma 6.2 implies that $`v_0=1`$. So $`1V`$ and $`V\rho (x^m)(1)=x^m`$ for any $`m𝐍.`$ This means that $`V=𝒫`$ and completes (a).
If $`k=1/2n`$ for $`n𝐙_+`$ then Lemma 6.2 states that $`v_0=c_1+c_2x^{2n+1}V`$ for constants $`c_1,c_2.`$ If $`c_10`$ then $`s(v_0)+v_0=2c_1V`$ (the latter is $`s`$-invariant and $`V=𝒫`$ as above. If $`v_0=x^{2n+1}V`$ then this argument gives that $`(x^{2n+1})=x^{2n+1}𝒫V.`$ Moreover $`V`$ cannot contain the polynomials of degree less than $`2n+1.`$ Otherwise we can find a $`0`$-eigenvector of $`𝒟`$ in the space of such polynomials, which is impossible. Hence $`V=(x^{2n+1}).`$
The latter is invariant with respect to $`x`$ and $`s.`$ Its $`𝒟`$-invariance readily follows from the formula $`𝒟(x^l)=(l+(1(1)^l)k)x^{l1}`$ considered in the range $`l2n+1.`$ $`\mathrm{}`$
We can reformulate the theorem as follows. The polynomial representation has a nontrivial (proper) $`^{\prime \prime }`$-quotient if and only if $`k=1/2n`$ for $`n𝐙_+.`$ In the latter case, such quotient is unique, namely, $`V_{2n+1}=𝒫/(x^{2n+1}).`$ Its dimension is $`2n+1.`$
Note that the subspace $`V_{2n+1}^0`$ of $`V_{2n+1}`$ generated by even polynomials is invariant with respect to the action of $`h=x\frac{}{x}+k+1/2`$ and $`e=x^2,`$ $`f=/4,`$ satisfying the defining relations of $`sl_2(𝐂)`$ (see Section 2). We get an irreducible representation of $`sl_2(𝐂)`$ of dimension $`n+1.`$
## 7 Algebraization
Let us use $`^{\prime \prime }`$ to formalize the previous considerations and to switch to the standard terminology of the representation theory.
(a) Inner product. We call a representation $`V`$ of $`^{\prime \prime }`$ pseudo-unitary if it possesses a non-degenerate $`𝐂`$-bilinear form $`(u,w)`$ such that $`(Hu,w)=(u,H^{}v)`$ for $`H^{\prime \prime }`$ for the anti-involution
$$^{}=,s^{}=s,x^{}=x.$$
(29)
By anti-involution, we mean a $`𝐂`$-automorphism satisfying $`(AB)^{}=B^{}A^{}.`$ We call such form $``$-invariant. Formulas (29) are compatible with the defining relations (27) of $`^{\prime \prime }`$ and therefore can be extended to the whole $`^{\prime \prime }.`$ This is straightforward. For instance,
$$[x^{},^{}]=[x,]=[,x]=1+2ks=1+2ks^{}=[,x]^{}.$$
We add “pseudo” because the pairing, generally speaking, is not supposed to be positive and the functions can be complex-valued.
The pairing $`f,g=_𝐑f(x)g(x)|x|^{2k}𝑑x`$ gives an example, provided the existence of the integral. Taking real-valued functions, we make this inner product positive (no “pseudo”). Assuming that the functions are $`C^{\mathrm{}},`$ we need to examine the convergence at $`x=0`$ and $`x=\mathrm{}.`$ If $`\mathrm{𝖱𝖾}(k)>1/2`$ then it suffices to take regular $`f`$ at $`x=0.`$ At infinity, $`f(x)|x|^k`$ has to be of type $`L^1(𝐑).`$ Polynomials times the Gaussian $`e^{x^2}`$ are fine.
(b) Gaussians. A homomorphism $`\gamma :VW`$ for two $`^{\prime \prime }`$-modules $`V,W`$ is called a Gaussian if $`\gamma H=\tau (H)\gamma `$ for the following automorphism $`\tau `$ of $`^{\prime \prime }:`$
$$\tau ()=2x,\tau (x)=x,\tau (s)=s.$$
(30)
These formulas can be extended to an automorphism of $`^{\prime \prime }.`$ Indeed, $`\tau (s)\tau ()\tau (s)=s(+2x)s=2x=\tau (),`$ the same holds for $`x`$, and
$$[\tau (),\tau (x)]=[d2x,x]=1+2s=\tau (1+2s).$$
Note that $`2`$ can be replaced by any constant $`\alpha 𝐂`$ in this definition. We get a family of automorphisms $`\tau _\alpha ()=2\alpha x`$ of $`^{\prime \prime }.`$ They lead to the following generalization of the master formula:
$$_𝐑\psi _\lambda (x)\psi _\mu (x)e^{\alpha x^2}|x|^{2k}𝑑x=\psi _\lambda (\frac{\mu }{\alpha })e^{\frac{\lambda ^2+\mu ^2}{\alpha }}\alpha ^k\mathrm{\Gamma }(k+\frac{1}{2}).$$
Here $`\alpha >0`$ to ensure the convergence. The substitution $`u=\sqrt{\alpha }x`$ readily makes it equvalent to (24): use that $`\psi _\lambda (x)`$ is a function of the $`x\lambda .`$ One can also follow the proof of the master formula employing $`e^{\alpha x^2}𝒟e^{\alpha x^2}=𝒟+2\alpha x.`$
If representations $`V,W`$ are algebras of functions on the same set then $`\gamma `$ can be assumed to be a function, to be more exact, the operator of multuplication by a function. For instance, the multiplication by $`e^{x^2}`$ sends the polynomial representation $`𝒫`$ to $`𝒫e^{\pm x^2}.`$ The latter is a $`^{\prime \prime }`$-module too. Adding all integral powers of $`e^{x^2}`$ to $`𝒫`$ we make this multiplication an inner automorphism of the resulting algebra. However it is somewhat artificial. Algebraically, the resulting representation $`𝒫[e^{mx^2},m𝐙]`$ is “too” reducible. Analytically, we mix together $`e^{x^2}`$ and $`e^{x^2},`$ functions with absolutely different behaviour at infinity. There are interesting examples of inner automorphisms $`\tau ,`$ but they are finite-dimensional.
(c) Hankel transform. Following Lemma 5.1, the operator Hankel transform is the following automorphism of $`^{\prime \prime }:`$
$$\sigma (s)=s,\sigma ()=2x,\sigma (2x)=.$$
(31)
These realations are obviously comapatible with the defining relations of $`^{\prime \prime }.`$ Any homomorphism $`:VW`$ of $`^{\prime \prime }`$-modules inducing $`\sigma `$ on $`^{\prime \prime }`$ can be called a Hankel transform. The main example so far is $`:𝒫e^{x^2}𝒫e^{+x^2},`$ where we identify $`x`$ and $`\lambda `$ in 23. Indeed, $`𝒟=2x`$, $`2x=𝒟`$, $`s=s`$ upon this identification.
It is interesting to interpret the master formula from this viewpoint. It is nothing else but the following identities for $`𝐅=e^{x^2}e^{^2/4}e^{x^2}`$:
$$𝐅s=s𝐅,𝐅=2x𝐅,𝐅(2x)=𝐅.$$
(32)
This means that $`𝐅`$ is Hankel transform whenever it is well-defined. Relations (32) can be deduced directly from the defining relations. In the first place, check that $`[,x^2]=2x,[^2,x]=2.`$ Then get that $`[,e^{x^2}]=2xe^{x^2},`$ $`[e^{^2/4},x]=`$ $`e^{d^2}/2`$ and use it as follows:
$`e^{x^2}e^{^2/4}e^{x^2}2x=e^{x^2}e^{^2/4}(2x)e^{x^2}=e^{x^2}(+2x)e^^2e^{x^2}=`$
$`(2x+2x)e^{x^2}e^{d^2}e^{x^2}=e^{x^2}e^^2e^{x^2},`$
$`e^{x^2}e^{^2/4}e^{x^2}2x=e^{x^2}e^{^2/4}(2x)e^{x^2}=e^{x^2}(+2x)e^^2e^{x^2}=`$
$`(2x+2x)e^{x^2}e^{d^2})e^{x^2}=e^{x^2}e^^2e^{x^2}.`$
The commutativity of $`𝐅`$ with $`s`$ is obvious.
Note the following braid identity which can be proved similarly:
$$𝐅=e^{x^2}e^{^2/4}e^{x^2}==e^^2e^{x^2/4}e^^2.$$
Actually we do need calculations from scratch because it suffices to use the nonsymmetric master formula Lemma 5.1 and the fact that $`𝒫`$ is a faithful representation. For instance, $`𝐅(2x)𝐅^1`$ and $``$ coincide in $`𝒫.`$ The previous consideration shows that the former is an element of $`^{\prime \prime }.`$ Hence they must coincide identically, i.e. in $`^{\prime \prime }.`$
A good demonstration of the convenience of such an algebraization will be the case of negative half-integral $`k.`$ Before switching to this case, let us conclude the “analytic” theory calculating the inverse Hankel transform.
## 8 Inverse transform and Plancherel formula
Let $`\mathrm{𝖱𝖾}k>1/2.`$ We use $`\psi _\lambda (x)`$ from (16).
$`(_{re}f)(\lambda )`$ $`={\displaystyle \frac{1}{\mathrm{\Gamma }(k+1/2)}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\psi _\lambda (x)f(x)|x|^{2k}𝑑x,`$ (33)
$`(_{im}g)(x)`$ $`={\displaystyle \frac{1}{\mathrm{\Gamma }(k+1/2)}}{\displaystyle _i\mathrm{}^{+i\mathrm{}}}\psi _x(\lambda )g(\lambda )|\lambda |^{2k}𝑑\lambda .`$ (34)
The first is nothing else but $``$ from (23). We just show explicitely that the integration is real.
Here we may consider $`𝐂`$-valued functions $`f`$ on $`𝐑`$ and $`g`$ on $`𝐂`$ respectively such that
$$f(x)=o(e^{cx})\text{ at }x=\mathrm{}c𝐑\text{ and }fg(\lambda )=o(e^{ci\lambda })\text{ at }\lambda =i\mathrm{}c𝐑.$$
Restricting ourselves with the polynomials times the Gaussian,
$$_{re}:𝐂[x]e^{x^2}𝐂[\lambda ]e^{\lambda ^2},_{im}:𝐂[\lambda ]e^{\lambda ^2}𝐂[x]e^{x^2}.$$
The first map is an isomorphism. Let us discuss the latter.
Let $`p(x)𝐂[x]`$. Applying the master formula to $`p_i(x)=p(ix),`$
$$\frac{1}{\mathrm{\Gamma }(k+1/2)}_{\mathrm{}}^+\mathrm{}\psi _\lambda (x)p(ix)e^{x^2}|x|^{2k}𝑑x=_{re}(e^{x^2}p(ix))=e^{\lambda ^2}\mathrm{exp}((𝒟^\lambda )^2/4)(p(i\lambda )).$$
Since $`𝒟I=iI𝒟`$ for $`I(f)(x)=f(ix),`$
$$\frac{(𝒟^\lambda )^2}{4}p(i\lambda )=\left(\frac{(𝒟^u)^2}{4}p(u)\right)|_{u=i\lambda }.$$
Now we replace $`\lambda i\lambda ,`$ use that $`\psi _{i\lambda }(x)=\psi _\lambda (ix),`$ and then integrate by substitution using $`z=ix.`$ The resulting formula reads:
$`{\displaystyle \frac{1}{\mathrm{\Gamma }(k+1/2)}}{\displaystyle _{\mathrm{}}^+\mathrm{}}\psi _\lambda (ix)p(ix)e^{x^2}|x|^{2k}𝑑x=e^{\lambda ^2}\mathrm{exp}((𝒟^\lambda )^2/4)(p(\lambda ))`$ (35)
$`={\displaystyle \frac{1}{i\mathrm{\Gamma }(k+1/2)}}{\displaystyle _i\mathrm{}^{+i\mathrm{}}}\psi _\lambda (z)p(z)e^{z^2}|z|^{2k}𝑑z.`$ (36)
Switching to $`\lambda ,`$ $`_{im}(e^{\lambda ^2}p(\lambda ))=e^{x^2}\mathrm{exp}((𝒟^x)^2/4)(p(x)).`$
We come to the inversion theorem.
###### Theorem 8.1
(Inversion Formula) $`_{im}_{re}=\mathrm{𝐢𝐝}`$ in $`𝐂[x]e^{x^2},`$ $`_{re}_{im}=\mathrm{𝐢𝐝}`$ in $`𝐂[\lambda ]e^{\lambda ^2}.`$
Proof: $`_{im}_{re}(e^{x^2}p(x))=e^{x^2}\mathrm{exp}(𝒟^2/4)\mathrm{exp}(𝒟^2/4)(p(x))=e^{x^2}p(x).`$ The second formula is analogous. $`\mathrm{}`$
There is a simple “algebraic” proof based on the facts that the transform $`_{im}_{re}`$ sends $`𝒟𝒟`$, $`2x2x`$, $`ss.`$ Thanks to the irreducibility of $`𝐂[x]e^{\pm x^2},`$ we may apply the Schur lemma. However the spaces are infinite dimensional, so a minor additional consideration is necessary. We will skip it because it practically coincides with that from the proof of the Plancherel formula.
Provided $`\mathrm{𝖱𝖾}k>\frac{1}{2},`$ the inner products
$$f,g_{re}=_{\mathrm{}}^+\mathrm{}f(x)g(x)|x|^{2k}𝑑x,f,g_{im}=\frac{1}{i}_i\mathrm{}^{+i\mathrm{}}f(\lambda )g(\lambda )|\lambda |^{2k}𝑑x$$
(37)
are non-degenerate respectively in $`𝐂[x]e^{x^2}`$ and $`𝐂[\lambda ]e^{\lambda ^2}`$.
It is obvious when $`𝐑[x],𝐑[\lambda ]`$ are considered instead of $`𝐂[x],𝐂[\lambda ]`$ and $`k𝐑.`$ Indeed, both forms become positive in this case. Concerning the second, use that $`g(\lambda )=\overline{g(\lambda )}`$ for a real polynomial $`g.`$
For complex-valued functions and $`k𝐂,`$ the claim requires proving. Let us use the irreducibility of $`^{\prime \prime }`$ modules $`𝐂[x]e^{x^2}`$ and $`𝐂[x]e^{x^2},`$ which is equivalent to the irreducibility of the polynomial representation $`𝒫=𝐂[x],`$ which we already know for $`\mathrm{𝖱𝖾}k>\frac{1}{2}.`$ Then the radical of the form $`.,._{re}`$ is a submodule of $`𝐂[x]e^{x^2}.`$ It does not coincide with the whole space, since
$$e^{x^2},e^{x^2}_{re}=_𝐑e^{2x^2}|x|^{2k}𝑑x=(\sqrt{2})^{2k1}\mathrm{\Gamma }(k+\frac{1}{2}).$$
(38)
The same argument works for the imaginary integration.
###### Theorem 8.2
(Plancherel Formula)
$$f,g_{re}=\widehat{f},\widehat{g}_{im}\text{ for all }f,g𝐂[x]e^{x^2},\widehat{f}=_{re}(f),\widehat{g}=_{re}(g).$$
(39)
Proof. Setting $`𝒫_{}=𝐂[x]e^{x^2},`$ we need to check that $`.,.=.,._{re}`$ coincides with $`f,g_1=\widehat{f},\widehat{g}_{im}`$ for all $`f,g𝒫_{}`$. In the first place, $`Hf,g_1=f,H^{}g_1`$ for any $`H^{\prime \prime },`$ i.e. this bilinear form is $``$-invariant. Indeed,
$`𝒟f,g_1=\widehat{𝒟f},\widehat{g}_{im}=2x\widehat{f},\widehat{g}_{im}=\widehat{f},2(x)\widehat{g}_{im}=\widehat{f},\widehat{𝒟g}_{im}=f,𝒟g_1.`$
Similarly, $`2xf,g_1=f,2xg_1`$ and $`s(f),g_1=f,s(g)_1.`$
Setting $`.,._2=.,..,._1,`$ we get $`e^{x^2},e^{x^2}_2=0.`$ Cf. 38. It is a $``$-invariant form as well. Let us demonstrate that it vanishes identically.
First, $`e^{x^2},(𝒟2x)f=(𝒟+2x)e^{x^2},f=0,f=0`$ for any $`f𝒫_{}`$ due to the $``$-invariance. So it is applicable to $`.,._2`$ too. Second, $`𝐂e^{x^2}(𝒟2x)𝒫_{}=\mathrm{}`$ because $`e^{x^2},e^{x^2}0.`$ Third, $`𝒫_{}=𝐂e^{x^2}(𝒟2x)𝒫_{}.`$ Really, the dimension of $`(𝒟2x)𝒫_n`$ for $`𝒫_n=`$ $`𝐂[1,x,\mathrm{},x^n]e^{x^2}`$ is $`n+1`$ since the kernel of the operator $`𝒟2x`$ in $`𝒫_{}`$ is zero. However $`(𝒟2x)𝒫_n𝒫_{(n+1)}.`$ Therefore $`(𝒟2x)𝒫_n=𝒫_{(n+1)}.`$ Finally,
$$e^{x^2},𝒫_{}_2=e^{x^2},e^{x^2}+(𝒟2x)𝒫_{}_2=0$$
and $`e^{x^2}`$ belongs to the radical of $`.,._2.`$ Since the module $`𝒫_{}`$ is irreducible, the radical has to coincide with the whole $`𝒫_{}.`$ $`\mathrm{}`$
Taking real $`k,`$ the forms $`.,._{re}`$ $`.,._{im}`$ are positive on $`𝐑[x]e^{x^2}`$ and $`𝐑[\lambda ]e^{\lambda ^2}.`$ The Plancherel formula allows us to complete the function spaces extending the Fourier transforms $`_{re},_{im}`$ to the spaces of square integrable real-valued functions with respect to the “Bessel measure”:
$$^2(𝐑,|x|^{2k}dx)^2(i𝐑,|\lambda |^{2k}d\lambda )^2(𝐑,|x|^{2k}dx).$$
The inversion and Plancherel formulas remain valid.
Here we assume that $`k>\frac{1}{2}.`$ Let us discuss the case of negative half-integers.
## 9 Finite-dimensional case
Let $`k=n\frac{1}{2}`$ for $`n𝐙_+.`$ Then $`V_{2n+1}=𝒫/(x^{2n+1})`$ is an irreducible representation of $`^{\prime \prime }`$. The elements of $`V_{2n+1}`$ can be identified with polynomials of degree less than $`2n+1`$.
###### Theorem 9.1
Finite-dimensional representations of $`^{\prime \prime }`$ exist only as $`k=n1/2`$ for $`n𝐙_+.`$ Given such $`k,`$ the algebra $`^{\prime \prime }`$ has a unique finite-dimensional irreducible representation up to isomorphisms, namely, $`V_{2n+1}.`$
Proof. We will use that
$`[h,x]=x,[h,]=\text{ for }h=(x+x)/2.`$ (40)
It readily follows from the defining relations of $`^{\prime \prime }.`$ Actually (40) determines a super Lie algebra. One may use a general theory of such Lie algebras. However in this particular case a reduction to $`sl_2`$ is more than sufficient.
Note that $`h`$ is $`x\frac{}{x}+k+1/2`$ in the polynomial representation. Since the latter is faithful, (40) is exactly the claim that $`x,𝒟`$ are homogeneous operators of degree $`\pm 1,`$ which is obvious.
We will employ that $`e=x^2,`$ $`f=^2/4,`$ and $`h`$ satisfy the defining relations of $`sl_2(𝐂).`$ Namely, $`[e,f]=h`$ because
$$[^2,x^2]=[^2,x]x+x[^2,x]=2x+x(2),$$
and the relations $`[h,e]=2e,[h,f]=2f`$ readily result from (40). Cf. Section 6.
Let $`V`$ be a finite-dimensional representation of $`^{\prime \prime }.`$ Then the subspaces $`V^0,V^1`$ of $`V`$ formed respectively by $`s`$-invariant and $`s`$-ani-invariant vectors are preserved by $`h,`$ $`e,`$ and $`f.`$ So they are $`sl_2(𝐂)`$-modules. One gets
$$x=h+k+1/2,x=hk1/2\text{ in }V^0,$$
and the other way round in $`V^1.`$
Let us check that $`k1/2𝐙_+.`$ All $`h`$-eigenvalues in $`V`$ are integers thanks to the general theory of finite-dimensional representations of $`sl_2(𝐂)`$. We pick a nonzero $`h`$-eigenvector $`vV`$ with the maximal possible eigenvalue $`m.`$ Then $`m𝐙_+`$ (the theory of $`sl_2`$) and $`x(v)=0`$ because the latter is an $`h`$\- eigenvector with the eigenvalue $`m+1.`$ Hence $`x(v)=0,`$ $`m+k+1/2=0,`$ and $`k=1/2m.`$
Let $`U^0`$ be a nonzero irreducible $`sl_2(𝐂)`$submodule of $`V^0.`$ The spectrum of $`h`$ in $`U^0`$ is $`\{n,n+2,\mathrm{},n2,n\}`$ for an integer $`n0.`$ Let $`v_l0`$ be an $`h`$-eigenvector with the eigenvalue $`l.`$ If $`e(v)=0`$ then $`v=cv_n`$ for a constant $`c,`$ and if $`f(v)=0`$ then $`v=cv_n.`$
Let us check that $`x(v_n)=0,`$ $`x(v_n)=0,`$ and
$$x(v_l)0\text{ for }ln,x(v_l)0\text{ for }ln.$$
Both operators, $`x`$ and $`x,`$ obviously preserve $`U^0:`$
$$x(v_l)=(l+k+1/2)v_l,x(v_l)=(lk1/2)v_l.$$
Hence,
$$^2x^2(v_l)=((x)^2+(12k)(x))(v_l)=(l+k+1/2)(lk+3/2)v_l.$$
Setting $`l=n,`$ we get that $`(n+k+1/2)(nk+3/2)=0`$ and $`k=1/2n,`$ because $`k<0`$ and $`nk+3/2>0.`$ Thus $`x(v_n)=0.`$ The case of $`x`$ is analogous.
The next claim is that $`x(v_n)=0,(v_n)=0.`$ Indeed, $`x(v^{})=0`$ and $`(v^{})=0`$ for $`v^{}=x(v_n).`$ Therefore
$$0=[,x](v^{})=(1+2ks)(v^{})=(11+n)v^{}=nv^{}.$$
This means that either $`v^{}=0`$ or $`n=0.`$ In the latter case, $`v^{}`$ is proportional to $`v_0`$ and therefore $`v^{}=x(v_0)=0`$ as well. Similarly, $`(v_n)=0.`$
Now we use the formula
$$(x^2(v_l))=x(2+x)(v_l)=(2+lk1/2)x(v_l)=(2+l+n)x(v_l),$$
and get that $`x(v_l)(U^0)`$ for any $`nln.`$ Hence $`U=U^0+(U^0)`$ is $`x`$-invariant. It is obviously $``$-invariant and $`s`$-invariant ($``$ $`(V^0)V^1`$). Also the sum is direct.
Finally, $`U`$ is a $`^{\prime \prime }`$-module and has to coincide with $`V`$ because the latter was assumed to be irreducible. The above formulas are sufficient to establish a $`^{\prime \prime }`$\- isomorphism $`UV_{2n+1}.`$ Explicitly, the $`h`$-eigenvectors $`x^i(v_n)U`$ will be identified with the monomials $`x^iV_{2n+1}.`$ $`\mathrm{}`$
Let us discuss the Hankel transform and related structures in the case of $`V_{2n+1}.`$ We follow Section 7.
(a) Form. To make $``$ “inner” we have to construct a non-degenerate bilinear pairing $`(u,v)`$ on $`V_{2n+1}`$ such that $`(Hu,v)=(u,H^{}v).`$ Here it is:
$$f,gV_{2n+1}\text{ set }(f,g)=\mathrm{Res}(f(x)g(x)x^{2n1}),\text{ where }\mathrm{Res}(a_ix^i)=a_1.$$
(41)
The pairing is non-degenerate, because if $`f=ax^l+`$ lower order terms, where $`a0`$ and $`0l2n,`$ then $`(f,x^{2nl})=a`$.
We introduce a scalar product $`(f,g)_0=\mathrm{Res}(fg)`$ for polynomials in terms of $`x`$ and $`x^1.`$ Denoting the conjugate of an operator $`A`$ with respect to this pairing by $`A^{},`$
$$s^{}=s,x^{}=x,\frac{}{x}^{}=\frac{}{x},\text{ and }x^{2n+1}𝒟^{}x^{2n1}=𝒟.$$
The relation $`x^{}=x`$ is obvious. Concerning $`\frac{}{x}^{}=\frac{}{x},`$ it follows from the property $`\mathrm{Res}(df/dx)`$=0 for any polynomial $`f(x,x^1).`$ The formula $`s^{}=s`$ results from $`\mathrm{Res}(s(f(x))=\mathrm{Res}(f(x)).`$
Switching from to $`^{},`$ we have
$$(sf,g)=\mathrm{Res}(s(f)gx^{2n1})=\mathrm{Res}(fs(g)s(x^{2n1}))=(1)^{2n+1}(fs(g)x^{2n1})=(f,s(g)).$$
Finally,
$`𝒟^{}`$ $`=({\displaystyle \frac{}{x}}+{\displaystyle \frac{k}{x}}(1s))^{}={\displaystyle \frac{}{x}}+(1+s){\displaystyle \frac{k}{x}},`$
$`x^{2k}𝒟^{}x^{2k}`$ $`={\displaystyle \frac{}{x}}{\displaystyle \frac{2k}{x}}+{\displaystyle \frac{k}{x}}(1+s)={\displaystyle \frac{}{x}}+{\displaystyle \frac{k}{x}}(1+s)=𝒟.`$ (42)
The first equality on the second line holds because $`\frac{k}{x}x^{2k}=kx^{2n2}`$ is an even function, and thus it commutes with the action of $`s`$. Finally
$$(𝒟f,g)=(𝒟f,x^{2k}g)_0=(f,x^{2k}x^{2k}𝒟^{}x^{2k}(g))_0=(f,x^{2k}𝒟g)_0=(f,𝒟g).\mathrm{}$$
(b) Gaussian. The Gaussian does not exist in polynomials but of course can be introduced as a power series $`e^{x^2}=_{m=0}^{\mathrm{}}(x^2)^m/m!`$ in the algebra of formal series $`𝐂[[x]],`$ a completion of the polynomial representation. The conjugation by this series induces $`\tau `$ on $`^{\prime \prime }.`$ Its inverse is $`e^{x^2}=_{m=0}^{\mathrm{}}(x^2)^m/m!.`$ The multiplication by the Gaussian does not preserve the space of polynomials but is well-defined on $`V_{2n+1}`$ because $`fV_{2n+1}`$ we have $`x^mf=0`$ for $`m2n+1.`$ Finally,
$$\gamma ^\pm =\underset{m=0}{\overset{2n}{}}(\pm x^2)^m/m!.$$
(c) Hankel transform. The operator $`𝒟`$ is nilpotent in $`V_{2n+1}`$ because it lowers the degree of $`fV_{2n+1}`$ by one. Therefore the operators $`\mathrm{exp}(\pm 𝒟^2/4)𝐂[[𝒟]]`$ are well-defined in this representation as well as the Gaussians. It suffices to take $`_{m=0}^{2n}(\pm (D/2)^2)^m/m!`$. Thus we may set
$$𝐅=e^{x^2}e^{\frac{𝒟^2}{4}}e^{x^2}\text{ in }V_{2n+1}.$$
(43)
###### Proposition 9.2
The map $``$ is a Hankel transform on $`V_{2n+1}`$, i.e. $`𝐅𝒟=𝐅2x`$, $`𝐅2x=𝒟𝐅`$, $`𝐅s=𝐅s`$. These relations fix it uniquely up to proportionality.
Proof. We already know that $`𝐅`$ is a Hankel transform (the previous section). If $`\stackrel{~}{𝐅}`$ is another one then the ratio $`\stackrel{~}{𝐅}𝐅^1`$ commutes with $`x,𝒟`$, and $`s`$ because of the very definition. Since $`V_{2n+1}`$ is irreducible (and finite dimensional) we get that $`\stackrel{~}{𝐅}`$ is proportional to $`𝐅.`$ $`\mathrm{}`$
## 10 Truncated Bessel functions
Recall that $`𝒟`$ has only one eigenvalue in $`V_{2n+1}`$, namely, $`0.`$ Therefore we cannot define the $`\psi _\lambda `$ as an eigenfunction of $`𝒟`$ in $`V_{2n+1}`$ any longer. Instead, it will be introduced as the kernel of the Hankel transform.
Any linear operator $`A:V_{2n+1}^xV_{2n+1}^\lambda `$ (the upper index indicates the variable) is a matrix. It means that
$`A(f)(\lambda )=(f(x),\alpha (x,\lambda ))=\mathrm{Res}(f(x)\alpha (x,\lambda )x^{2n1}),\text{ where}`$
$`\alpha (x,\lambda )={\displaystyle \underset{l,m=0}{\overset{2n}{}}}c_{l,m}x^l\lambda ^m={\displaystyle \underset{l=0}{\overset{2n}{}}}x^{2nl}A(x^l).`$ (44)
So here the kernel $`\alpha (x,\alpha )`$ is uniquely defined by $`A`$ and vice versa.
The truncated $`\psi `$-function is the kernel of $`𝐅:`$
$$𝐅(f)(\lambda )=(f(x),\psi _\lambda (x))=\mathrm{Res}(f(x)\psi _\lambda (x)x^{2n1}).$$
(45)
There is a somewhat different approach. Let us use that the relations from Lemma 9.2 determine $`𝐅`$ uniquely up to proportionality. These relations are equivalent to the following properties of $`\psi _\lambda (x):`$
$$𝒟\psi _\lambda (x)=2\lambda \psi _\lambda (x)mod(x^{2n+1},\lambda ^{2n+1}),\psi _\lambda (x)=\psi _x(\lambda ),\psi _\lambda (s(x))=\psi _{s(\lambda )}(x).$$
(46)
Let us solve the first equation. Setting $`\psi _\lambda (x)=_{l,m=0}^{2n}c_{l,m}x^l\lambda ^m,`$
$`{\displaystyle \underset{l=1,m=0}{\overset{l=2n,m=2n}{}}}c_{l,m}(l+(1(1)^l)(n{\displaystyle \frac{1}{2}}))x^{l1}\lambda ^m`$
$`=`$ $`{\displaystyle \underset{l=0,m=0}{\overset{l=2n,m=2n1}{}}}2c_{l,m}x^l\lambda ^{m+1}mod(x^{2n+1},\lambda ^{2n+1}),`$
$`c_{l,m}={\displaystyle \frac{2}{l+(1(1)^l)(n1/2)}}c_{l1,m1}\text{ for }2nl>0,2nm>0,`$
$`\text{where }c_{l,0}=0=c_{2n,m}\text{ for }l>0,m<2n.`$ (47)
Using the $`x\lambda `$ symmetry, we conclude that $`c_{l,0}=0=c_{0,l}`$ for nonzero $`l`$ and $`c_{l,m}=0`$ for $`lm.`$ Thus
$$\psi _\lambda (x)=g_n(\lambda x)\text{ for }g_n=\underset{l=0}{\overset{2n}{}}c_lt^l,c_l=c_{l,l},$$
where the coefficients are given by (47).
Finally, $`g_n(t)=f_n(t)+(1/2)df_n/dt`$ for the truncated Bessel function $`f_n(t)=_{m=0}^nc_{2m}t^{2m}`$ which is an even solution of the truncated Bessel equation (cf. Section 1):
$$\frac{d^2f}{dt^2}(t)+2k\frac{1}{t}\frac{df}{dt}(t)4f(t)=0mod(t^{2n}),k=n1/2.$$
(48)
This equation is sufficient to determine the coefficients of $`f_n`$ uniquely for any constant term $`c_0=c_{0,0}.`$ They are given by the same formula (3) till $`c_{2n}`$ up to proportionality. This can be checked directly using explicit formulas which will be discussed next.
Still $`c_0`$ remains arbitrary. Recall that $`\psi `$ was initially introduced as the kernel of $`𝐅.`$ So it comes with its own normalization. Let us calculate its $`c_0.`$ One gets:
$`𝐅(e^{x^2})=e^{\lambda ^2}\mathrm{exp}(𝒟^2/4)(e^{\lambda ^2}e^{\lambda ^2})=e^{\lambda ^2}\mathrm{exp}(𝒟^2/4)(1)=e^{\lambda ^2},\text{ so}`$
$`𝐅(1{\displaystyle \frac{x^2}{1!}}+{\displaystyle \frac{x^4}{2!}}+\mathrm{}+(1)^n{\displaystyle \frac{x^{2n}}{n!}})=1+{\displaystyle \frac{\lambda ^2}{1!}}+\mathrm{}{\displaystyle \frac{\lambda ^{2n}}{n!}}.`$ (49)
Here the transform of $`1`$ is proportional to $`\lambda ^{2n}`$ since the latter has to be an eigenfunction of $`\lambda ,`$ i.e. the solution of the equation $`\lambda 𝐅(1)=0`$ in $`V_{2n+1}^\lambda .`$ Similarly, $`𝐅(x^l)=(𝒟^\lambda /2)^l𝐅(1)`$ is proportional to $`\lambda ^{2nl}`$ for $`0l2n.`$ Thus (49) leads to the relations
$`𝐅(x^{2m})=(1)^m{\displaystyle \frac{m!}{(nm)!}}\lambda ^{2n2m}.`$ (50)
For instance, $`𝐅(x^{2n})=(1)^nn!.`$ This is exactly the coefficient $`c_0`$ above.
We obtain that the normalization serving the truncated Hankel transform is
$$\psi _\lambda (0)=n!,c_0=g_n(0)=f_n(0)=n!.$$
(51)
Formula (50) also results in
$`𝐅(x^{2m+1})=𝐅(x(x^{2m}))=(𝒟/2)𝐅(x^{2m})=(1)^m{\displaystyle \frac{m!}{(nm1)!}}\lambda ^{2n2m1}.`$ (52)
Substituting,
$`\psi _\lambda (x)={\displaystyle \underset{m=0}{\overset{n}{}}}x^{2n2m}𝐅(x^{2m})+{\displaystyle \underset{m=0}{\overset{n1}{}}}x^{2n2m1}𝐅(x^{2m+1})=f_n(x\lambda )+{\displaystyle \frac{1}{2}}f_n^{}(x\lambda )\text{ for}`$
$`f_n(t)={\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \frac{(1)^mm!}{(nm)!}}t^{2n2m}={\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \frac{(1)^{nm}(nm)!}{m!}}t^{2m}.`$ (53)
It is exactly the solution of (48) with the truncated normalization $`f_n(0)=(1)^nn!.`$
Truncated inversion. Concluding the consideration of the case $`k=n\frac{1}{2}`$ for $`n𝐙_+,`$ let us discuss the inversion. We have the following transformations and scalar products:
$`𝐅_+:𝐂[x]/(x^{2n+1})𝐂[\lambda ]/(\lambda ^{2n+1}),`$ $`𝐅_+(f)=\mathrm{Res}(f(x)\psi _\lambda (x)x^{2n1}),`$
$`𝐅_{}:𝐂[\lambda ]/(\lambda ^{2n+1})𝐂[x]/(x^{2n+1}),`$ $`𝐅_{}(f)=\mathrm{Res}(f(\lambda )\psi _x(\lambda )\lambda ^{2n1}).`$
$`f,g_+=\mathrm{Res}(f(x)g(x)x^{2n1}),`$ $`f,g𝐂[x]/(x^{2n+1});`$
$`f,g_{}=\mathrm{Res}(f(\lambda )g(\lambda )\lambda ^{2n1}),`$ $`f,g𝐂[\lambda ]/(\lambda ^{2n+1}).`$ (54)
Here $`𝐅_+(f)=𝐅(f)=\widehat{f}`$ in the notation above. The transform $`𝐅_{}(f)`$ coincides with $`𝐅_+^\lambda (f)`$ for even $`f(\lambda )`$ and with $`𝐅_+^\lambda (f)`$ for odd $`f(\lambda ).`$
We can follow the “analytic” case and check that $`𝐅_{}𝐅_+`$ commutes with $`𝒟,x,s.`$ Hence it is the multiplication by a constant thanks to the irreducibility of $`V_{2n+1}.`$ The constant is $`𝐅_{}𝐅_+(1)`$ and can be readily calculated. It is equally simple to calculate all $`𝐅_{}𝐅_+(x^l)`$ using (50) and (52). For instance,
$`𝐅_{}𝐅_+(x^{2m})=𝐅_{}\left({\displaystyle \frac{(1)^mm!}{(nm)!}}\lambda ^{2n2m}\right)=`$
$`={\displaystyle \frac{(1)^mm!}{(nm)!}}{\displaystyle \frac{(1)^{nm}(nm)!}{m!}}x^{2m}=(1)^nx^{2m}.`$
Thus the truncated inversion reads:
$$𝐅_{}𝐅_+=(1)^n\text{id}=𝐅_+𝐅_{}.$$
Concerning the Plancherel formula, we may use the proportionality of the forms $`f,g_+`$ and $`\widehat{f},\widehat{g}_{}`$ for $`f,gV_{2n+1}`$ and their transforms $`\widehat{f}=𝐅(f),\widehat{g}=𝐅(g).`$ It results from the irreducibility of $`V_{2n+1}.`$ A direct calculation is simple as well. Let
$`f,f_+=f,f={\displaystyle \underset{l=0}{\overset{2n}{}}}a_la_{2nl}\text{ for }f={\displaystyle \underset{l=0}{\overset{2n}{}}}a_lx^l,`$
$`g,g_{}={\displaystyle \underset{l=0}{\overset{2n}{}}}(1)^lb_lb_{2nl}\text{ for }g=𝐅(f)={\displaystyle \underset{l=0}{\overset{2n}{}}}b_l\lambda ^l.`$ (55)
It is easy to check that
$$b_lb_{2nl}=(1)^{l+n}a_la_{2n1}.$$
Indeed, using (50) and (52):
$`b_{2m}b_{2n2m}=(1)^ma_{2m}{\displaystyle \frac{m!}{(nm)!}}(1)^{nm}a_{2n2m}{\displaystyle \frac{(nm)!}{m!}}`$
$`=(1)^na_{2m}a_{2n2m},`$
$`b_{2m+1}b_{2n2m1}=(1)^ma_{2m+1}{\displaystyle \frac{m!}{(nm1)!}}(1)^{nm1}a_{2n2m1}{\displaystyle \frac{(nm1)!}{m!}}`$
$`=(1)^{n1}a_{2m+1}a_{2n2m1}.`$
We get the truncated Plancherel formula:
$$\widehat{f},\widehat{g}_{}=(1)^nf,g_+.$$
The above consideration proves the coincidence for $`f=g,`$ i.e. for the corresponding quadratic forms. It is of course sufficient.
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# BOUND OF AUTOMORPHISMS OF PROJECTIVE VARIETIES OF GENERAL TYPE
## 1 Introduction
The automorphim group of a projective variety of general type is known to be finite. For every curve $`C`$ of genus $`g2`$, we have the estimate :
$$\mathrm{}\text{Aut}(C)84(g1)$$
by well known Hurwitz’s theorem.
In the case of surfaces, G. Xiao proved that for every smooth minimal surface of general type
$$\mathrm{}\text{Aut}(S)1764K_S^2$$
holds . The main purpose of this article is to prove the following theorem.
###### Theorem 1.1
There exists a positive number $`C_n`$ which depends only on $`n`$ such that for every smooth projective $`n`$-fold $`X`$ of general type defined over complex numbers, the automorphism group $`\text{Aut}(X)`$ of $`X`$ satisfies the estimate:
$$\mathrm{}\text{Aut}(X)C_n\mu (X,K_X),$$
where $`\mu (X,K_X)`$ is the volume of $`X`$ with respect to $`K_X`$ (cf. Definition 2.3).
The method of the proof of Theorem 1.1 is a combination of the ideas in and . Let $`X`$ be a projective $`n`$-fold of general type and let $`G`$ denotes the automorphism group of $`X`$. Since $`G`$ acts on the canonical ring $`R(X,K_X)`$ of $`X`$, by we may assume that $`X`$ is a canonical model, i.e. $`X`$ has only canonical singularity and $`K_X`$ is ample (our proofs of Theorem 1.1 and Therem 1.2 below depend on the finite generation of canonical rings of varieties of general type in which has not yet published. For the safe side, one may restrict oneselves to the case of $`dimX3`$ (cf. )) The quotient $`X/G`$ is a projective variety. Let $`K_{X/G,orb}`$ be the orbifold canonical divisor of $`X/G`$. Then we see that
$$G=K_X^n/K_{X/G,orb}^n$$
holds, where $`G`$ denotes the order of $`G`$. Since $`\mu (X,K_X)=K_X^n`$ holds in this case, we see that Theorem 1.1 follows from the following theorem.
###### Theorem 1.2
Let $`X`$, $`G`$ be as above. There exists a positive constant $`c_n`$ depending only on $`n`$ such that
$$K_{X/G,orb}^nc_n$$
holds.
It is easy to see $`c_1`$ can be taken to be $`1/42`$. This leads to Hurwicz’s theorem. G. Xiao proved that $`c_2`$ can be taken as $`1/1764`$ ().
The key ingredient of the proof of Theorem 1.2 is the subadjunction formula in which relates the canonical divisor of the minimal center of logcanonical singularities and the canonical divisor of the ambient space. Using this we see that $`X/G`$ with $`\mu (X/G,K_{X/G,orb})=K_{X/G,orb}^n1`$ is birationally bounded by the inductive procedure in . Then Theorem 1.1 and Theorem 1.2 follows from a Diophantine consideration.
Therem 1.1 and Theorem 1.2 are not effective in the sense that there exist no explicit estimates of $`C_n`$ and $`c_n`$.
## 2 Preliminaries
### 2.1 Orbifold canonical divisors
Let $`X`$ be a projective variety of general type with only canonical singularities. Let $`G`$ denote the automorphism group of $`X`$. Iis well known that $`G`$ is a finite group. The quotient $`X/G`$ is a projective variety. Let $`\stackrel{~}{X}`$ be the equivalent resolution of $`X`$ with respect to $`G`$ such that $`\stackrel{~}{X}/G`$ is also smooth. We may take $`\stackrel{~}{X}`$ such that the ramification divisor $`R`$ of
$$\stackrel{~}{\pi }:\stackrel{~}{X}\stackrel{~}{X}/G$$
and the branch locus $`B=(\stackrel{~}{\pi }_{}(R))_{red}`$ is a divisor with normal crossings. Let $`B=_iB_i`$ be the irreducible decomposition of $`B`$. Then there exists a set of positive integers $`m_i`$ such that
$$K_{\stackrel{~}{X}}=\stackrel{~}{\pi }^{}(K_{\stackrel{~}{X}/G}+\underset{i}{}\frac{m_i1}{m_i}B_i)$$
Let
$$\varpi :\stackrel{~}{X}/GX/G$$
be the natural morphism. We set
$$K_{X/G,orb}:=\varpi _{}(K_{\stackrel{~}{X}/G}+\underset{i}{}\frac{m_i1}{m_i}B_i)$$
and call it the obifold canonical divisor of $`X/G`$. Let
$$\pi :XX/G$$
be the natural morphism. Then
$$K_X=\pi ^{}K_{X/G,orb}$$
holds. The orbifold canonical ring is defined by
$$R(X/G,K_{X/G,orb}):=R(X,K_X)^G.$$
And the linear system $`[mK_{X/G,orb}]`$ is given by
$$[mK_{X/G,orb}]=mK_X^G.$$
Hence we have that
$$R(X/G,K_{X/G,orb})=_{m0}\mathrm{\Gamma }(X/G,𝒪_{X/G}([mK_{X/G,orb}]))$$
holds.
### 2.2 Multiplier ideal sheaves
In this section, we shall review the basic definitions and properties of multiplier ideal sheaves.
###### Definition 2.1
Let $`L`$ be a line bundle on a complex manifold $`M`$. A singular hermitian metric $`h`$ is given by
$$h=e^\phi h_0,$$
where $`h_0`$ is a $`C^{\mathrm{}}`$-hermitian metric on $`L`$ and $`\phi L_{loc}^1(M)`$ is an arbitrary function on $`M`$.
The curvature current $`\mathrm{\Theta }_h`$ of the singular hermitian line bundle $`(L,h)`$ is defined by
$$\mathrm{\Theta }_h:=\mathrm{\Theta }_{h_0}+\sqrt{1}\overline{}\phi ,$$
where $`\overline{}`$ is taken in the sense of a current. The $`L^2`$-sheaf $`^2(L,h)`$ of the singular hermitian line bundle $`(L,h)`$ is defined by
$$^2(L,h):=\{\sigma \mathrm{\Gamma }(U,𝒪_M(L))h(\sigma ,\sigma )L_{loc}^1(U)\},$$
where $`U`$ runs opens subsets of $`M`$. In this case there exists an ideal sheaf $`(h)`$ such that
$$^2(L,h)=𝒪_M(L)(h)$$
holds. We call $`(h)`$ the multiplier ideal sheaf of $`(L,h)`$. If we write $`h`$ as
$$h=e^\phi h_0,$$
where $`h_0`$ is a $`C^{\mathrm{}}`$ hermitian metric on $`L`$ and $`\phi L_{loc}^1(M)`$ is the weight function, we see that
$$(h)=^2(𝒪_M,e^\phi )$$
holds. We have the following vanishing theorem.
###### Theorem 2.1
(Nadel’s vanishing theorem \[9, p.561\]) Let $`(L,h)`$ be a singular hermitian line bundle on a compact Kähler manifold $`M`$ and let $`\omega `$ be a Kähler form on $`M`$. Suppose that $`\mathrm{\Theta }_h`$ is strictly positive, i.e., there exists a positive constant $`\epsilon `$ such that
$$\mathrm{\Theta }_h\epsilon \omega $$
holds. Then $`(h)`$ is a coherent sheaf of $`𝒪_M`$-ideal and for every $`q1`$
$$H^q(M,𝒪_M(K_M+L)(h))=0$$
holds.
### 2.3 Analytic Zariski decomposition
To study a big line bundle we introduce the notion of analytic Zariski decompositions. By using analytic Zariski decompositions, we can handle big line bundles like a nef and big line bundles.
###### Definition 2.2
Let $`M`$ be a compact complex manifold and let $`L`$ be a line bundle on $`M`$. A singular hermitian metric $`h`$ on $`L`$ is said to be an analytic Zariski decomposition, if the followings hold.
1. $`\mathrm{\Theta }_h`$ is a closed positive current,
2. for every $`m0`$, the natural inclusion
$$H^0(M,𝒪_M(mL)(h^m))H^0(M,𝒪_M(mL))$$
is isomorphim.
###### Remark 2.1
If an AZD exists on a line bundle $`L`$ on a smooth projective variety $`M`$, $`L`$ is pseudoeffective by the condition 1 above.
###### Theorem 2.2
() Let $`L`$ be a big line bundle on a smooth projective variety $`M`$. Then $`L`$ has an AZD.
### 2.4 Volume of projective varieties
To measure the positivity of big line bundles on a projective variety we shall introduce a volume of a projective variety with respect to a line bundle.
###### Definition 2.3
Let $`L`$ be a line bundle on a compact complex manifold $`M`$ of dimension $`n`$. We define the $`L`$-volume of $`M`$ by
$$\mu (M,L):=n!\underset{m\mathrm{}}{lim\; sup}m^ndimH^0(M,𝒪_M(mL)).$$
###### Definition 2.4
() Let $`L`$ be a big line bundle on a smooth projective variety $`X`$. Let $`Y`$ be a subvariety of $`X`$ of dimension $`r`$. We define the volume $`\mu (Y,L)`$ of $`Y`$ with respect to $`L`$ by
$$\mu (Y,L):=r!\underset{m\mathrm{}}{lim\; sup}m^rdimH^0(Y,𝒪_Y(mL)(h^m)/tor),$$
where $`h`$ is an AZD of $`L`$ and $`tor`$ denotes the torsion part of the sheaf $`𝒪_Y(mL)(h^m)`$. This definition can be easily generalized to the case that $`L`$ is a Q-line bundle.
## 3 Stratification of varieties by multiplier ideal sheaves
Let $`X`$ be a smooth projective $`n`$-fold of general type. Then the canonical ring $`R(X,K_X)`$ is finitely generated by . Let $`X_{can}`$ be the canonical model of $`X`$. $`K_{X_{can}}`$ is an ample $`𝐐`$-Cartier divisor on $`X_{can}`$. We assume that the natural rational map
$$\phi :X\mathrm{}X_{can}$$
is a morphism. Let $`h_{can}`$ be a $`C^{\mathrm{}}`$-hermitian metric on $`K_{X_{can}}`$ induced from the Fubini-Study metric on the hyperplane bundle of a projective space by a projective embedding of $`X_{can}`$ associated with $`rK_{X_{can}}`$ where $`r`$ is a sufficiently large positive integer such that $`rK_{X_{can}}`$ is Cartier. Then $`h_{can}`$ has strictly positive curvature on $`X_{can}`$. $`h_{can}`$ induces a singular hermitian metric $`h`$ on $`K_X`$ in a natural manner. By the definition, $`h`$ is an AZD of $`K_X`$. To prove Theorem 1.1, we may replace $`X`$ by any birational model of $`X`$, we may assume that there exists an effective Q-divisor $`N`$ such that
$$(h^m)=𝒪_X([mN])$$
holds for every $`m0`$. In particular we may and do assume that $`(h^m)`$ is locally free for every $`m0`$. Let us denote $`\mu (X/G,K_{X/G,orb})`$ by $`\mu _0`$. We set
$$X^{}=\{xX\phi \text{ is a local isomorphism around }x\}.$$
Let $`G`$ be the group of the birational automorphism of $`X`$. To prove Theorem 1.1, we may assume that $`G`$ acts $`X`$ regularly and $`X/G`$ is also smooth. Let
$$\pi :XX/G$$
be the natural morphism. We set
$$(X/G)^{}=\pi (X^{}).$$
###### Lemma 3.1
Let $`x,y`$ be distinct points on $`(X/G)^{}`$. We set
$$_{x,y}=_x_y$$
Let $`\epsilon `$ be a sufficiently small positive number. Then
$$H^0(X/G,𝒪_{X/G}(mK_{X/G,orb})_{x,y}^{\sqrt[n]{\mu _0}(1\epsilon )\frac{m}{\sqrt[n]{2}}})0$$
for every sufficiently large $`m`$, where $`_x,_y`$ denote the maximal ideal sheaf of the points $`x,y`$ respectively.
Proof of Lemma 3.1. Let us consider the exact sequence:
$$0H^0(X/G,𝒪_{X/G}(mK_{X/G,orb})_{x,y}^{\sqrt[n]{\mu _0}(1\epsilon )\frac{m}{\sqrt[n]{2}}})H^0(X/G,𝒪_{X/G}(mK_{X/G,orb}))$$
$$H^0(X/G,𝒪_{X/G}(mK_{X/G,orb})/_{x,y}^{\sqrt[n]{\mu _0}(1\epsilon )\frac{m}{\sqrt[n]{2}}}).$$
Since
$$n!\underset{m\mathrm{}}{lim\; sup}m^ndimH^0(X/G,𝒪_{X/G}(mK_{X/G,orb})/_{x,y}^{\sqrt[n]{\mu _0}(1\epsilon )\frac{m}{\sqrt[n]{2}}})=\mu _0(1\epsilon )^n<\mu _0$$
hold, we see that Lemma 3.1 holds. Q.E.D.
Let us take a sufficiently large positive integer $`m_0`$ and let $`\sigma `$ be a general (nonzero) element of $`H^0(X/G,𝒪_{X/G}(m_0K_{X/G,orb})_{x,y}^{\sqrt[n]{\mu _0}(1\epsilon )\frac{m_0}{\sqrt[n]{2}}})`$. We define a singular hermitian metric $`h_0`$ on $`K_{X/G,orb}`$ by
$$h_0(\tau ,\tau ):=\frac{\tau ^2}{\sigma ^{2/m_0}}.$$
Then
$$\mathrm{\Theta }_{h_0}=\frac{2\pi }{m_0}(\sigma )$$
holds, where $`(\sigma )`$ denotes the closed positive current defined by the divisor $`(\sigma )`$. Hence $`\mathrm{\Theta }_{h_0}`$ is a closed positive current. Let $`\alpha `$ be a positive number and let $`(\alpha )`$ denote the multiplier ideal sheaf of $`h_0^\alpha `$, i.e.,
$$(\alpha )=^2(𝒪_{X/G},(\frac{h_0}{h_{X/G}})^\alpha ),$$
where $`h_{X/G}`$ is an arbitrary $`C^{\mathrm{}}`$-hermitian metric on $`K_{X/G,orb}`$. Let us define a positive number $`\alpha _0(=\alpha _0(x,y))`$ by
$$\alpha _0:=inf\{\alpha >0(𝒪_{X/G}/(\alpha ))_x0\text{and}(𝒪_{X/G}/(\alpha ))_y0\}.$$
Since $`(_{i=1}^nz_i^2)^n`$ is not locally integrable around $`O\text{C}^n`$, by the construction of $`h_0`$, we see that
$$\alpha _0\frac{n\sqrt[n]{2}}{\sqrt[n]{\mu _0}(1\epsilon )}$$
holds. Then one of the following two cases occurs.
Case 1.1: For every small positive number $`\delta `$, $`𝒪_{X/G}/(\alpha _0\delta )`$ has $`0`$-stalk at both $`x`$ and $`y`$.
Case 1.2: For every small positive number $`\delta `$, $`𝒪_{X/G}/(\alpha _0\delta )`$ has nonzero-stalk at one of $`x`$ or $`y`$ say $`y`$.
First we consider Case 1.1. Let $`\delta `$ be a sufficiently small positive number and let $`V_1`$ be the germ of subscheme at $`x`$ defined by the ideal sheaf $`(\alpha _0+\delta )`$. By the coherence of $`(\alpha )(\alpha >0)`$, we see that if we take $`\delta `$ sufficiently small, then $`V_1`$ is independent of $`\delta `$. It is also easy to verify that $`V_1`$ is reduced if we take $`\delta `$ sufficiently small. In fact if we take a log resolution of $`(X/G,\frac{\alpha _0}{m_0}(\sigma ))`$, $`V_1`$ is the image of the divisor with discrepancy $`1`$ (for example cf. \[4, p.207\]). Let $`(X/G)_1`$ be a subvariety of $`X/G`$ which defines a branch of $`V_1`$ at $`x`$. We consider the following two cases.
Case 2.1: $`(X/G)_1`$ passes through both $`x`$ and $`y`$,
Case 2.2: Otherwise
For the first we consider Case 2.1. Suppose that $`(X/G)_1`$ is not isolated at $`x`$. Let $`n_1`$ denote the dimension of $`(X/G)_1`$. Let us define the volume $`\mu _1`$ of $`(X/G)_1`$ with respect to $`K_{X/G,orb}`$ by
$$\mu _1:=\mu ((X/G)_1,K_{X/G,orb}).$$
Since $`xX/G^{}`$, we see that $`\mu _1>0`$ holds.
###### Lemma 3.2
Let $`\epsilon `$ be a sufficiently small positive number and let $`x_1,x_2`$ be distinct regular points on $`(X/G)_1X/G^{}`$. Then for a sufficiently large $`m>1`$ divisible by $`G`$,
$$H^0((X/G)_1,𝒪_{(X/G)_1}(mK_{X/G,orb})(h^m)_{x_1,x_2}^{\sqrt[n_1]{\mu _1}(1\epsilon )\frac{m}{\sqrt[n_1]{2}}})0$$
holds.
The proof of Lemma 3.2 is identical as that of Lemma 3.1, since
$$(h^m)_{x_i}=𝒪_{X/G,x_i}(i=1,2)$$
hold for every $`m`$.
By Kodaira’s lemma there is an effective $`𝐐`$-divisor $`E`$ such that $`K_{X/G,orb}E`$ is ample. Let $`\mathrm{}`$ be a sufficiently large positive integer such that
$$L:=\mathrm{}(K_{X/G,orb}E)$$
is a line bundle and statisfies the property in Lemma 3.3.
###### Lemma 3.3
If we take $`\mathrm{}`$ sufficiently large, then
$$\varphi _m:H^0(X/G,𝒪_{X/G}(mK_{X/G,orb}+L)(h^m))$$
$$H^0((X/G)_1,𝒪_{(X/G)_1}(mK_{X/G,orb}+L)(h^m))$$
is surjective for every $`m0`$ divisible by $`G`$.
Proof. Let us take a locally free resolution of the ideal sheaf $`_{(X/G)_1}`$ of $`(X/G)_1`$.
$$0_{(X/G)_1}_1_2\mathrm{}_k0.$$
Then by the trivial extention of the case of vector bundles, if $`r`$ is sufficiently large, we see that
$$H^q(X/G,𝒪_{X/G}(mK_{X/G,orb}+L)(h^m)_j)=0$$
holds for every $`m1`$, $`q1`$ and $`1jk`$. In fact if we take $`\mathrm{}`$ sufficiently large, we see that for every $`j`$, $`𝒪_{X/G}(LK_{X/G})_j`$ admits a $`C^{\mathrm{}}`$-hermitian metric $`g_j`$ such that
$$\mathrm{\Theta }_{g_j}\text{Id}_{E_j}\omega $$
holds, where $`\omega `$ is a Kähler form on $`X/G`$. By \[2, Theorem 4.1.2 and Lemma 4.2.2\] we have the desired vanishing.
Hence
$$H^1(X/G,𝒪_{X/G}(mK_{X/G,orb}+L)(h^m)_{(X/G)_1})=0$$
holds. This completes the proof of Lemma 3.3. Q.E.D.
Let $`\tau `$ be a general section in $`H^0(X/G,𝒪_{X/G}(L))`$.
Let $`m_1`$ be a sufficiently large positive integer divisible by $`G`$ and let $`\sigma _1^{}`$ be a general element of
$$H^0((X/G)_1,𝒪_{(X/G)_1}(m_1K_{X/G,orb})(h^{m_1})_{x_1,x_2}^{\sqrt[n_1]{\mu _1}(1\epsilon )\frac{m_1}{\sqrt[n_1]{2}}}),$$
where $`x_1,x_2(X/G)_1`$ are distinct nonsingular points on $`(X/G)_1`$.
By Lemma 3.2, we may assume that $`\sigma _1^{}`$ is nonzero. Then by Lemma 3.3 we see that
$$\sigma _1^{}\tau H^0((X/G)_1,𝒪_{(X/G)_1}(m_1K_{X/G,orb}+L)(h^{m_1})_{x_1,x_2}^{\sqrt[n_1]{\mu _1}(1\epsilon )\frac{m_1}{\sqrt[n_1]{2}}})$$
extends to a section
$$\sigma _1H^0(X/G,𝒪_{X/G}((m+\mathrm{})K_{X/G,orb})(h^{m+\mathrm{}}))$$
We may assume that there exists a neighbourhood $`U_{x,y}`$ of $`\{x,y\}`$ such that the divisor $`(\sigma _1)`$ is smooth on $`U_{x,y}(X/G)_1`$ by Bertini’s theorem, if we take $`\mathrm{}`$ sufficiently large, since by Theorem 2.1,
$$H^0(X/G,𝒪_{X/G}(mK_{X/G,orb}+L)(h^m))$$
$$H^0(X/G,𝒪_{X/G}(mK_{X/G,orb}+L)(h^m))/𝒪_{X/G}((X/G)_1)_y)$$
is surjective for every $`yX/G`$ and $`m0`$ divisible by $`G`$, where $`𝒪_{X/G}((X/G)_1)`$ is the ideal sheaf of $`(X/G)_1`$. We define a singular hermitian metric $`h_1`$ on $`K_{X/G,orb}`$ by
$$h_1=\frac{1}{\sigma _1^{\frac{2}{m_1+\mathrm{}}}}.$$
Let $`\epsilon _0`$ be a sufficiently small positive number and let $`_1(\alpha )`$ be the multiplier ideal sheaf of $`h_0^{\alpha _0\epsilon _0}h_1^\alpha `$,i.e.,
$$_1(\alpha )=^2(𝒪_{X/G},h_0^{\alpha _0\epsilon _0}h_1^\alpha /h_{X/G}^{(\alpha _0+\alpha \epsilon _0)}).$$
Suppose that $`x,y`$ are nonsingular points on $`(X/G)_1`$. Then we set $`x_1=x,x_2=y`$ and define $`\alpha _1(=\alpha _1(x,y))>0`$ by
$$\alpha _1:=inf\{\alpha (𝒪_{X/G}/_1(\alpha ))_x0\text{and}(𝒪_{X/G}/_1(\alpha ))_y0\}.$$
By Lemma 3.3 we may assume that we have taken $`m_1`$ so that
$$\frac{\mathrm{}}{m_1}\epsilon _0\frac{\sqrt[n_1]{\mu _1}}{n_1\sqrt[n_1]{2}}$$
holds.
###### Lemma 3.4
$$\alpha _1\frac{n_1\sqrt[n_1]{2}}{\sqrt[n_1]{\mu _1}}+O(\epsilon _0)$$
holds.
To prove Lemma 3.4, we need the following elementary lemma.
###### Lemma 3.5
(\[16, p.12, Lemma 6\]) Let $`a,b`$ be positive numbers. Then
$$_0^1\frac{r_2^{2n_11}}{(r_1^2+r_2^{2a})^b}𝑑r_2=r_1^{\frac{2n_1}{a}2b}_0^{r_1^{2a}}\frac{r_3^{2n_11}}{(1+r_3^{2a})^b}𝑑r_3$$
holds, where
$$r_3=r_2/r_1^{1/a}.$$
Proof of Lemma 3.3. Let $`(z_1,\mathrm{},z_n)`$ be a local coordinate on a neighbourhood $`U`$ of $`x`$ in $`X/G`$ such that
$$U(X/G)_1=\{qUz_{n_1+1}(q)=\mathrm{}=z_n(q)=0\}.$$
We set $`r_1=(_{i=n_1+1}^nz_1^2)^{1/2}`$ and $`r_2=(_{i=1}^{n_1}z_i^2)^{1/2}`$. Then there exists a positive constant $`C`$ such that
$$\sigma _1^2C(r_1^2+r_2^{2\sqrt[n_1]{\mu _1}(1\epsilon )\frac{m_1}{\sqrt[n_1]{2}}})$$
holds on a neighbourhood of $`x`$, where $``$ denotes the norm with respect to $`h_{X/G}^{m_1+\mathrm{}}`$. We note that there exists a positive integer $`M`$ such that
$$\sigma ^2=O(r_1^M)$$
holds on a neighbourhood of the generic point of $`U(X/G)_1`$, where $``$ denotes the norm with respect to $`h_{X/G}^{m_0}`$. Then by Lemma 3.5, we have the inequality
$$\alpha _1(\frac{m_1+\mathrm{}}{m_1})\frac{n_1\sqrt[n_1]{2}}{\sqrt[n_1]{\mu _1}}+O(\epsilon _0)$$
holds. By using the fact that
$$\frac{\mathrm{}}{m_1}\epsilon _0\frac{\sqrt[n_1]{\mu _1}}{n_1\sqrt[n_1]{2}}$$
we obtain that
$$\alpha _1\frac{n_1\sqrt[n_1]{2}}{\sqrt[n_1]{\mu _1}}+O(\epsilon _0)$$
holds. Q.E.D.
If $`x`$ or $`y`$ is a singular point on $`(X/G)_1`$, we need the following lemma.
###### Lemma 3.6
Let $`\phi `$ be a plurisubharmonic function on $`\mathrm{\Delta }^n\times \mathrm{\Delta }`$. Let $`\phi _t(t\mathrm{\Delta })`$ be the restriction of $`\phi `$ on $`\mathrm{\Delta }^n\times \{t\}`$. Assume that $`e^{\phi _t}`$ does not belong to $`L_{loc}^1(\mathrm{\Delta }^n,O)`$ for every $`t\mathrm{\Delta }^{}`$.
Then $`e^{\phi _0}`$ is not locally integrable at $`O\mathrm{\Delta }^n`$.
Lemma 3.6 is an immediate consequence of . Using Lemma 3.6 and Lemma 3.5, we see that Lemma 3.4 holds by letting $`x_1x`$ and $`x_2y`$.
For the next we consider Case 1.2 and Case 2.2. We note that in Case 2.2 by modifying $`\sigma `$ a little bit , if necessary we may assume that $`(𝒪_{X/G}/(\alpha _0\epsilon ))_y0`$ and $`(𝒪_{X/G}/(\alpha _0\epsilon ^{}))_x=0`$ hold for a sufficiently small positive number $`\epsilon ^{}`$. For example it is sufficient to replace $`\sigma `$ by the following $`\sigma ^{}`$ constructed below.
Let $`X/G_1^{}`$ be a subvariety which defines a branch of
$$\text{Spec}(𝒪_{X/G}/(\alpha +\delta ))$$
at $`y`$. By the assumption (changing $`(X/G)_1`$, if necessary) we may assume that $`(X/G)_1^{}`$ does not contain $`x`$. Let $`m^{}`$ be a sufficiently large positive integer divisible by $`G`$ such that $`m^{}/m_0`$ is sufficiently small (we can take $`m_0`$ arbitrary large).
Let $`\tau _y`$ be a general element of
$$H^0(X/G,𝒪_{X/G}(m^{}K_{X/G,orb})_{(X/G)_1^{}}),$$
where $`_{(X/G)_1^{}}`$ is the ideal sheaf of $`(X/G)_1^{}`$. If we take $`m^{}`$ sufficiently large, $`\tau _y`$ is not identically zero. We set
$$\sigma ^{}=\sigma \tau _y.$$
Then we see that the new singular hermitian metric $`h_0^{}`$ defined by $`\sigma ^{}`$ satisfies the desired property.
In these cases, instead of Lemma 3.2, we use the following simpler lemma.
###### Lemma 3.7
Let $`\epsilon `$ be a sufficiently small positive number and let $`x_1`$ be a smooth point on $`(X/G)_1`$. Then for a sufficiently large $`m>1`$ divisible by $`G`$,
$$H^0((X/G)_1,𝒪_{(X/G)_1}(mK_{X/G,orb})(h^m)_{x_1}^{\sqrt[n_1]{\mu _1}(1\epsilon )m})0$$
holds.
Then taking a general $`\sigma _1^{}`$ in
$$H^0((X/G)_1,𝒪_{(X/G)_1}(m_1K_{X/G,orb})(h^{m_1})_{x_1}^{\sqrt[n_1]{\mu _1}(1\epsilon )m_1}),$$
for a sufficiently large $`m_1`$. As in Case 1.1 and Case 2.1 we obtain a proper subvariety $`(X/G)_2`$ in $`(X/G)_1`$ also in this case.
Inductively for distinct points $`x,yX/G^{}`$, we construct a strictly decreasing sequence of subvarieties
$$X/G=(X/G)_0(x,y)(X/G)_1(x,y)\mathrm{}$$
$$(X/G)_r(x,y)(X/G)_{r+1}(x,y)=\{x\}\text{or}\{x,y\},$$
where $`R_y`$ (or $`R_x`$) is a subvariety such that $`x`$ deos not belong to $`R_y`$ and $`y`$ belongs to $`R_y`$. and invariants :
$$\alpha _0(x,y),\alpha _1(x,y),\mathrm{},\alpha _r(x,y),$$
$$\mu _0,\mu _1(x,y),\mathrm{},\mu _r(x,y)$$
and
$$n>n_1>\mathrm{}>n_r.$$
By Nadel’s vanishing theorem (Theorem 2.1) we have the following lemma.
###### Lemma 3.8
Let $`x,y`$ be two distinct points on $`X/G^{}`$. Then for every $`m_{i=0}^r\alpha _i(x,y)+1`$, $`\mathrm{\Phi }_{mK_{X/G,orb}}`$ separates $`x`$ and $`y`$.
Proof. For simplicity let us denote $`\alpha _i(x,y)`$ by $`\alpha _i`$. Let us define the singular hermitian metric $`h_{x,y}`$ of the Q-line bundle $`(m1)K_{X/G,orb}`$ defined by
$$h_{x,y}=(\underset{i=0}{\overset{r1}{}}h_i^{\alpha _i\epsilon _i})h_r^{\alpha _r+\epsilon _r}h^{(m1(_{i=0}^{r1}(\alpha _i\epsilon _i))(\alpha _r+\epsilon _r)\mathrm{}\delta _L)}h_L^{\delta _L},$$
where $`h_L`$ is a $`C^{\mathrm{}}`$-hermitian metric on $`L`$ with strictly positive curvature and $`\delta _L`$ be a sufficiently small positive number. Then we see that $`(h_{x,y})`$ defines a subscheme of $`X/G`$ with isolated support around $`x`$ or $`y`$ by the definition of the invariants $`\{\alpha _i\}`$’s. By the construction the curvature current $`\mathrm{\Theta }_{h_{x,y}}`$ is strictly positive on $`X/G`$. Then by Nadel’s vanishing theorem (Theorem 2.1) we see that
$$H^1(X/G,𝒪_{X/G}(K_{X/G}+(m1)K_{X/G,orb})(h_{x,y}))=0.$$
Hence
$$H^0(X/G,𝒪_{X/G}(K_{X/G}+(m1)K_{X/G,orb}))$$
separates $`x`$ and $`y`$. We note that
$$H^0(X/G,𝒪_{X/G}(K_{X/G}+(m1)K_{X/G,orb}))$$
is a subspace of
$$H^0(X,𝒪_X(mK_X))^G$$
by the definition of $`K_{X/G,orb}`$. This implies that $`\mathrm{\Phi }_{[mK_{X/G,orb}]}`$ separates $`x`$ and $`y`$. Q.E.D.
We note that for a fixed $`x`$, $`_{i=0}^r\alpha _i(x,y)`$ depends on $`y`$. We set
$$\alpha (x)=\underset{yU_0}{sup}\underset{i=0}{\overset{r}{}}\alpha _i$$
and let
$$X/G=(X/G)_0(X/G)_1(X/G)_2\mathrm{}$$
$$(X/G)_r(X/G)_{r+1}=\{x\}\text{or}\{x,y\}$$
be the stratification which attains $`\alpha (x)`$. In this case we call it the maximal stratification at $`x`$. We see that there exists a nonempty open subset $`U`$ in countable Zariski topology of $`X/G`$ such that on $`U`$ the function $`\alpha (x)`$ is constant and there exists an irreducible family of stratification which attains $`\alpha (x)`$ for every $`xU`$.
In fact this can be verified as follows. We note that the cardinarity of
$$\{(X/G)_i(x,y)x,yX/G,xy(i=0,1,\mathrm{})\}$$
is uncontably many, while the cardinarity of the irreducible components of Hilbert scheme of $`X/G`$ is countably many. We see that for fixed $`i`$ and very general $`x`$, $`\{(X/G)_i(x,y)\}`$ should form a family on $`X/G`$. Similary we see that for very general $`x`$, we may assume that the maximal stratification $`\{(X/G)_i(x)\}`$ forms a family. This implies the existence of $`U`$.
And we may also assume that the corresponding invariants $`\{\alpha _0,\mathrm{},\alpha _r\}`$, $`\{\mu _0,\mathrm{},\mu _r\}`$, $`\{n=n_0\mathrm{},n_r\}`$ are constant on $`U`$. Hereafter we denote these invariants again by the same notations for simplicity. The proof of the following lemma is parallel to that of Lemma 3.4.
###### Lemma 3.9
$$\alpha _i\frac{n_i\sqrt[n_i]{2}}{\sqrt[n_i]{\mu _i}}+O(\epsilon _{i1})$$
hold for $`1ir`$.
###### Proposition 3.1
For every
$$m>\underset{i=0}{\overset{r}{}}\alpha _i+1$$
$`[mK_{X/G,orb}]`$ gives a birational rational map from $`X/G`$ into a projective space.
###### Lemma 3.10
If $`\mathrm{\Phi }_m_{(X/G)_i}`$ is birational rational map onto its image, then
$$\mathrm{deg}\mathrm{\Phi }_m((X/G)_i)m^{n_i}\mu _i$$
holds.
Proof. Let $`p:\stackrel{~}{X/G}X/G`$ be the resolution of the base locus of $`mK_{X/G,orb}`$ and let
$$p^{}[mK_{X/G,orb}]=P_m+F_m$$
be the decomposition into the free part $`P_m`$ and the fixed component $`F_m`$. Let $`p_i:\stackrel{~}{X/G}_i(X/G)_i`$ be the resolution of the base locus of $`\mathrm{\Phi }_{mK_{X/G,orb}}_{(X/G)_i}`$ obtained by the restriction of $`p`$ on $`p^1((X/G)_i)`$. Let
$$p_i^{}(mK_{X/G,orb}_{(X/G)_i})=P_{m,i}+F_{m,i}$$
be the decomposition into the free part $`P_{m,i}`$ and the fixed part $`F_{m,i}`$. We have
$$\mathrm{deg}\mathrm{\Phi }_{[mK_{X/G,orb}]}((X/G)_i)=P_{m,i}^{n_i}$$
holds. Then by the ring structure of $`R(X/G,K_{X/G,orb})`$, we have that there exists a natural injection
$$H^0(X/G,𝒪_{X/G}(\nu P_m))H^0(X/G,𝒪_{X/G}([m\nu K_{X/G,orb}])(h^{m\nu }))$$
for every $`\nu 1`$. Hence there exists a natural morphism
$$H^0((X/G)_i,𝒪_{(X/G)_i}(\nu P_{m,i}))H^0((X/G)_i,𝒪_{(X/G)_i}([m\nu K_{X/G,orb}])(h^{m\nu }))$$
for every $`\nu 1`$. This morphism is clearly injective. This implies that
$$\mu _im^{n_i}\mu ((X/G)_i,P_{m,i})$$
holds. Since $`P_{m,i}`$ is nef and big on $`(X/G)_i`$ we see that
$$\mu ((X/G)_i,P_{m,i})=P_{m,i}^{n_i}$$
holds. Hence
$$\mu _im^{n_i}P_{m,i}^{n_i}$$
holds. This implies that
$$\mathrm{deg}\mathrm{\Phi }_{mK_{X/G,orb}}((X/G)_i)\mu _im^{n_i}$$
holds. Q.E.D.
## 4 Proof of Theorem 1.1
To prove Theorem 1.1 we use the following subadjunction formula.
###### Theorem 4.1
() Let $`X/G`$ be a normal projective variety. Let $`D^{}`$ and $`D`$ be effective Q-divisor on $`X`$ such that $`D^{}<D`$, $`(X,D^{})`$ is logterminal and $`(X,D)`$ is logcanonical. Let $`W`$ be a minimal center of logcanonical singularities for $`(X,D)`$. Let $`H`$ be an ample Cartier divisor on $`X`$ and $`ϵ`$ a positive rational number. Then there exists an effective Q-divisor $`D_W`$ on $`D`$ such that
$$(K_X+D+ϵH)_W_\text{Q}K_W+D_W$$
and $`(W,D_W)`$ is logterminal. In particular $`W`$ has only rational singularities.
Let us start the proof of Theorem 1.1. We prove Theorem 1.1 by induction on $`n=dimX`$. Suppose that Theorem 1.1 holds for varities of general type of dimension $`<n`$. Then there exists a positive constant $`C(m)(m<n)`$ depending only on $`m`$ such that for every smooth projective varietiey $`Y`$ of general type of dimension $`m`$
$$\mu (Y,K_Y)/\mathrm{}\text{Aut}(Y)C(m)$$
holds. Let $`X`$ be a smooth projective variety of general type as in Section 3. We use the same notations as in Section 3. Let $`x,y`$ be distinct points on $`(X/G)^{}`$ and let
$$X/G=(X/G)_0(X/G)_1\mathrm{}(X/G)_r(X/G)_{r+1}=\{x\}\text{or}\{x,y\}$$
be the stratification constructed as in Section 3 and let
$$\mu _0,\mathrm{},\mu _r$$
$$n_1,\mathrm{},n_r$$
be the invariants as in Section 3. Let
$$X=X_0X_1\mathrm{}X_rX_{r+1}$$
be the corresponding stratification of $`X`$. If we take $`x,y`$ general, $`X_i(0ir)`$ are projective varieties of general type. Let
$$X_{can}:=\text{Proj}R(X,K_X)$$
be the canonical model of $`X`$.
We have the corresponding stratification
$$X_{can}=X_{0,can}X_{1,can}\mathrm{}X_{r,can}X_{r+1,can}$$
on $`X_{can}`$ (here we note that $`X_{i,can}`$ does not denote the canonical model of $`X_i`$ for $`i1`$).
Then we see that
$$\mu _i=\frac{1}{G}\mu (X_i,K_X)=\frac{1}{G}(K_{X_{can}})^{n_i}X_{i,can}$$
holds. Let $`H`$ be an ample divisor on $`X`$. By the subadjunction formula, we see that for every positive rational number $`ϵ`$
$$K_{X_{i,can}}<_\text{Q}(1+\underset{j=0}{\overset{i1}{}}\alpha _j)K_{X_{can}}+ϵH$$
holds, where $`<_\text{Q}`$ means that the righthandside minus the lefthandside is Q-linear equivalent to an effective divisor and $`K_{X_i,can}`$ denotes the pushforward of the canonical divisor of a nonsingular model of $`K_{X_{i,can}}`$. This can be verified as follows. Let
$$\pi :XX/G$$
be the natural morphism. Let $`D_i`$ be the divisor on $`X`$ which corresponds to the singular hermitian metric
$$\pi ^{}(h_0^{\alpha _0\epsilon _0}\mathrm{}h_{i1}^{\alpha _{i1}\epsilon _{i1}}h_i^{\alpha _i}).$$
$`D_i`$ is a positive linear combinations of $`\{\pi ^{}(\sigma _0),\mathrm{},(\sigma _j)\}`$ by the constructions of $`h_0,\mathrm{},h_i`$. Also we may assume that $`D_i`$ is a $`𝐐`$-divisor by perturbations of $`\epsilon _0,\mathrm{},\epsilon _{i1}`$. $`X_{i,can}`$ may not be the minimal center of $`(X,D_i)`$ and $`(X,D_i)`$ may not be logcanonical. But if we take a suitable modification
$$\pi _i:Y_iX_{i,can},$$
we may assume that there exists an effective Q-divisor $`E_i`$ such that
1. $`\pi _i^{}D_iE_i`$ is effective,
2. $`(Y_i,\pi _i^{}D_iE_i)`$ is logcanonical and the proper transform of $`X_{i,can}`$ is the minimal center of $`(Y_i,\pi _i^{}D_iE_i)`$.
Then by Theorem 4.1, we have that for every positive rational number $`ϵ`$
$$K_{X_{i,can}}<_\text{Q}(1+\underset{j=0}{\overset{i1}{}}\alpha _j)K_{X_{can}}+ϵH$$
holds. By the inductive assumption this implies that
$$(1+\underset{j=0}{\overset{i1}{}}\alpha _j)^{n_i}\mu _iC(n_i)$$
holds. Since
$$\alpha _i\frac{\sqrt[n_i]{2}n_i}{\sqrt[n_i]{\mu _i}}+O(\epsilon _{i1})$$
holds by Lemma 3.9, we see that
$$()\frac{1}{\sqrt[n_i]{\mu _i}}(1+\underset{j=0}{\overset{i1}{}}\frac{\sqrt[n_j]{2}n_j}{\sqrt[n_j]{\mu _j}})C(n_i)^1$$
holds for every $`i1`$. Inductively we see that if $`\mu _01`$ holds,
$$\frac{1}{\sqrt[n_i]{\mu _i}}\frac{1}{\sqrt[n]{\mu _0}}C(C(1),\mathrm{},C(n1))$$
holds where $`C(C(1),\mathrm{}C(n1))`$ is a positive constant depending only on $`C(1),\mathrm{},C(n1)`$. Hence if $`\mu _0<1`$ holds then we see that
$$\mathrm{deg}\mathrm{\Phi }_{(1+_{i=0}^r\alpha _i)K_{X/G,orb}}(X)C(C(1),\mathrm{},C(n1))^n$$
holds. This implies that $`X/G`$ is birationally bounded, if
$$\mu _0(=\frac{1}{G}\mu (X,K_X))1$$
holds. We set
$$\alpha :=\underset{i=0}{\overset{r}{}}\alpha _i+1.$$
Then using Lemma 3.10, we have the following lemma.
###### Lemma 4.1
If $`\mu _01`$ holds, then there exists a positive constant $`A(n)`$ depending only on $`n`$ such that
$$1\alpha ^n\mu _0A(n)$$
holds.
Let
$$\alpha K_X^G=P+F$$
be the decomposition of $`\alpha K_X^G`$ into the movable part $`P`$ and the fixed component $`F`$. Taking a suitable successive $`G`$-equivariant blowing ups, we may assume that $`P`$ is base point free. And also we may assume that the canonical birational map
$$f:XX_{can}$$
is a morphism.
###### Lemma 4.2
There exists a positive constant $`c_n`$ depending only on $`n`$ such that
$$f^{}K_{X_{can}}P^{n1}c_nG$$
holds. In particular
$$\alpha ^{n1}K_{X_{can}/G,orb}^nc_n$$
holds.
Proof. Let
$$f_G:X/GX_{can}/G$$
be the natural morphism. Let us write
$$K_{X/G}=f_G^{}(K_{X_{can}/G})+a_iE_i$$
where $`\{E_i\}`$ are irreducible exceptional divisor of $`f_G`$. We set
$$Y:=\mathrm{\Phi }_{\alpha K_X^G}(X).$$
and we set
$$\varphi :=\mathrm{\Phi }_P:XY.$$
Let
$$\varphi _G:X/GY$$
be the birational morphism induced by $`\varphi `$. Then
$$f^{}K_{X_{can}}P^{n1}=\varphi _{}f^{}K_{X_{can}}H^{n1}$$
holds, where $`H`$ denotes the hyperplane section of $`Y`$. Also
$$\varphi _{}f^{}K_{X_{can}}H^{n1}=G(\varphi _G)_{}f_G^{}K_{X_{can}/G,orb}H^{n1}$$
holds. On the other hand
$`()(\varphi _G)_{}f_G^{}K_{X_{can}/G}H^{n1}`$ $`=`$ $`(\varphi _G)_{}(K_{X/G}{\displaystyle a_iE_i})H^{n1}`$
$`=`$ $`K_YH^{n1}{\displaystyle \underset{i}{}}a_i(\varphi _G)_{}E_iH^{n1}`$
holds, where $`K_Y`$ denotes the pushforward of the canonical divisor of the normalization of $`Y`$ to $`Y`$. We note that $`K_YH^{n1}(=K_{X/G}P^{n1}`$)is an integer. Since $`E_i`$’s appear as fixed components of $`[\alpha K_{X_{can}/G,orb}]^G`$, we see that
$$\underset{i}{}(\varphi _G)_{}E_iH^{n1}\alpha ^n\mu _0C(n)$$
hold. Hence $`_i(\varphi _G)_{}E_i`$ is bounded.
Since $`_i(\varphi _G)_{}E_i`$ is an exceptional divisor of the birational rational map
$$f_G\varphi _G^1:Y\mathrm{}X_{can}/G,$$
$`\{a_i\}`$ is of finitely many possibilities. Hence there exists a positive constant $`K_n`$ depending only on $`n`$ such that
$$(\mathrm{})(\varphi _G)_{}f_G^{}(K_{X_{can}/G})H^{n1}K_n$$
holds. Let $`\{D_j\}`$ be the irreducible divisors such that
$$K_{X_{can}/G,orb}=K_{X_{can}/G}+\underset{j}{}\frac{m_j1}{m_j}D_j$$
for some positive integers $`\{m_j\}`$. Then we see that
$`(\mathrm{})(f_G^{}K_{X_{can}/G,orb})\varphi _G^{}H^{n1}`$ $`=`$ $`f_G^{}K_{X_{can}/G}\varphi _G^{}H^{n1}+{\displaystyle \underset{j}{}}{\displaystyle \frac{m_j1}{m_j}}f_G^{}D_j\varphi _G^{}H^{n1}`$
$``$ $`\alpha ^n\mu _0`$
$``$ $`A(n)`$
hold. By $`(\mathrm{})`$ this implies that $`_j(\varphi _G)_{}f_G^{}D_j`$ is bounded and
$$\mathrm{}\{j(\varphi _G)_{}f_G^{}D_j0\}$$
is uniformly bounded by a positive integer, say $`N`$ depending only on $`n`$.
###### Lemma 4.3
Let $`N`$ and $`B`$ are fixed positive integers. Then
$$\{\{\underset{j=1}{\overset{N}{}}\frac{b_j}{a_j}\},a_j,b_j\text{are integers such that}b_jB\}\{0\}$$
is bounded below by a positive constant, where for a rational number $`c`$ $`\{c\}`$ denotes the fractional part of $`c`$. i.e.
$$\{c\}:=c[c].$$
Proof. Suppose not. Then there exists a sequence of positive integers
$$\{a_{j,k}\},\{b_{j,k}\}1jN,k=1,2,\mathrm{}$$
such that
$$b_{j,k}B,$$
$$\{\underset{j=1}{\overset{N}{}}\frac{b_{j,k}}{a_{j,k}}\}0,$$
$$\underset{k\mathrm{}}{lim}\frac{b_{j,k}}{a_{j,k}}$$
exists for every $`j`$ and
$$\underset{k\mathrm{}}{lim}\{\underset{j=1}{\overset{N}{}}\frac{b_{j,k}}{a_{j,k}}\}=0$$
hold. We note that if
$$\underset{k\mathrm{}}{lim}\frac{b_{j,k}}{a_{j,k}}0$$
then by the boundedness of $`b_{j,k}`$ the sequence is constant for every sufficiently large $`k`$ and if
$$\underset{k\mathrm{}}{lim}\frac{b_{j,k}}{a_{j,k}}=0$$
then $`a_{j,k}`$ tends to infinity $`k`$ goes to infinity. Since
$$\underset{k\mathrm{}}{lim}\{\underset{j=1}{\overset{N}{}}\frac{b_{j,k}}{a_{j,k}}\}=0$$
holds, there is no $`j`$ such that
$$\underset{k\mathrm{}}{lim}\frac{b_{j,k}}{a_{j,k}}=0$$
holds. Hence by the above observation we see that for every $`j`$ the sequence $`\{b_{j,k}/a_{j,k}\}_{k=1}^{\mathrm{}}`$ is constant for every sufficiently large $`k`$ and $`j`$. This contradicts to the fact that
$$\{\underset{j=1}{\overset{N}{}}\frac{b_{j,k}}{a_{j,k}}\}0$$
holds for every $`k`$. This completes the proof of Lemma 4.3 . Q.E.D.
We note that by $`()`$, the finiteness properties of $`\{a_i\}`$ and the boundedness of $`_i(\varphi _G)_{}E_i`$, we see that the rational number $`f_G^{}K_{X_{can}/G}H^{n1}`$ is of finitely many possibilities. By $`(\mathrm{})`$, the boundedness of $`_j(\varphi _G)_{}f_G^{}D_j`$ and Lemma 4.3, we see that there exists a positive constant $`c_n`$ depending only on $`n`$ such that
$$f^{}K_{X_{can}}P^{n1}c_nG$$
holds. Since $`R(X_{can}/G,K_{X_{can}/G,orb})`$ is a ring,
$$\alpha ^{n1}K_{X_{can}/G,orb}^nc_n$$
holds. This completes the proof of Lemma 4.1. Q.E.D.
By Lemma 4.1 and Lemma 4.2 we see that
$$\alpha \frac{A(n)}{c_n}$$
holds. By Lemma 4.1, we see that
$$\mu _0\frac{1}{\alpha ^n}$$
holds. Hence we have that
$$\mu _0(\frac{c_n}{A(n)})^n$$
holds. This completes the proof of Theorem 1.2. Since
$$\mu _0=\frac{1}{G}\mu (X,K_X)$$
holds, we have that
$$G(\frac{A(n)}{c_n})^n\mu (X,K_X)$$
holds. This completes the proof of Theorem 1.1.
Author’s address
Hajime Tsuji
Department of Mathematics
Tokyo Institute of Technology
2-12-1 Ohokayama, Megro 152-8551
Japan
e-mail address: tsuji@math.titech.ac.jp
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# Planetary Transits Toward the Galactic Bulge
## 1. Introduction
The search for extrasolar planets has garnered enormous attention in recent years, due primarily to the successful implementation of radial velocity searches (Mayor & Queloz (1995), Marcy & Butler (1996)). These searches have led to the discovery of a population of massive, close-in planets with orbital separations of $`a\begin{array}{c}<\hfill \\ \hfill \end{array}0.1\mathrm{AU}`$. Recently, it was discovered that one such planet, the companion to HD 209458, also transits its parent star (Charbonneau et al. (2000); Henry et al. (2000)), yielding a measurement of the mass, radius, and density of the companion.
Clearly, transit observations can be used to extract additional information about known companions. The discovery of an extrasolar planet using transits, however, has remained elusive. There are two primary difficulties with detecting planets with transits. First, the photometric requirements are quite stringent: a planet of radius $`R_\mathrm{p}R_\mathrm{J}`$ (where $`R_\mathrm{J}`$ is the radius of Jupiter) transiting an primary of radius $`R_{}=R_{}`$ would produce a fractional deviation of $`\begin{array}{c}<\hfill \\ \hfill \end{array}1\%`$ during the course of the transit. Second, the probability that a planet will transit its parent is small: for a planet with separation $`0.05\mathrm{AU}`$ orbiting a star with $`R_{}=R_{}`$, the probability is $`\begin{array}{c}<\hfill \\ \hfill \end{array}10\%`$. Several methods for dealing will the small probability have been proposed. For instance, one can monitor eclipsing binary stars, where the orbital plane is known to be (nearly) perpendicular to the sky (Deeg et al. (1998)).
Another way of overcoming this small probability is to simply monitor many stars simultaneously. This can be done by employing a camera with a large field-of-view, or by monitoring very dense stellar fields. Here I focus on the latter possibility. Specifically I determine the number of planets that might be detected in a campaign monitoring stars toward the Galactic bulge.
## 2. Formalism
The flux of a star being occulted by a planet is given by,
$$F(t)=F_0[\delta (t)+1]+F_b,$$
(1)
where $`F_0`$ is the unocculted flux of the star, $`F_b`$ is the total flux from any unrelated sources, and $`\delta (t)`$ is the fractional deviation of the flux due to the transit, which depends on the radius of the planet relative to the star, the inclination angle, $`i`$, and the limb-darkening of the star (Sackett (1999)). For a small planet ($`R_\mathrm{p}R_{}`$) and no limb-darkening, $`\delta =(R_\mathrm{p}/R_{})^2\mathrm{\Theta }(1\tau )`$, where $`\mathrm{\Theta }(x)`$ is the step function, and $`\tau `$ is a normalized time, $`\tau (tt_0)/t_\mathrm{T}`$. Here $`t_0`$ is the time of the midpoint of the transit, and $`t_\mathrm{T}`$ is one-half the transit duration, which for circular orbits is,
$$t_\mathrm{T}=\frac{P}{2\pi }\mathrm{arcsin}\left(\sqrt{\left(\frac{R_{}+R_\mathrm{p}}{a}\right)^2\mathrm{cos}^2i}\right).$$
(2)
In reality, $`\delta `$ depends very sensitively on $`R_\mathrm{p}`$ and $`\mathrm{cos}i`$, and less so on the limb-darkening. I will therefore use the explicit form for $`\delta `$ given in Sackett (1999), but assume no limb-darkening.
Since the proposed search for planets will be carried out in dense stellar fields, and transits produce time-dependent variations in the flux of the stars, the data will likely be reduced with image-subtraction techniques (Tomaney & Crotts (1996), Alard & Lupton (1998)). With image-subtraction, one measures only the time variable portion of the flux, $`\stackrel{~}{F}(t)=F_0[\delta (t)]`$.
## 3. Detection Probability
There are three requirements to detect a planet of separation $`a`$ and radius $`R_\mathrm{p}`$ around a star of mass $`M`$, radius $`R_{}`$ and flux $`F_0`$. These are: (1) the planet must transit the star, (2) at least two transits must occur during the time when observations are made, and (3) the transit must cause a detectable deviation in the light curve. If the duration of the transit is much smaller than the window of observations, than these three requirements can be considered independent.
For a planet to transit its parent star, it must have an inclination angle $`\mathrm{cos}i\mathrm{cos}i_{\mathrm{min}}(R_{}+R_\mathrm{p})/a`$. The probability $`P_\mathrm{T}`$ that a planet will transit its parent star is therefore,
$$P_\mathrm{T}=\frac{_0^{\mathrm{cos}i_{\mathrm{min}}}\mathrm{d}(\mathrm{cos}i)}{_0^1\mathrm{d}(\mathrm{cos}i)}=\frac{R_{}+R_\mathrm{p}}{a}.$$
(3)
Consider a campaign lasting $`N`$ nights with $`T_\mathrm{W}`$ hours per night. Defining $`t=0`$ as the beginning of the first night, then the times when observations are possible on (integer) night $`n`$ satisfy $`𝒯(t)=|tn\lambda T_\mathrm{W}/2|T_\mathrm{W}/20`$, where $`\lambda =1\mathrm{day}`$. Both the first transit occurring at time $`t_1`$ and the second transit at time $`t_2`$ must satisfy this relation on some (integer) nights $`n_1`$ and $`n_2`$. Note $`n_1n_2`$. The time $`t_2`$ is given by $`t_2=t_1+n_\mathrm{P}P`$, where $`n_\mathrm{P}`$ is the number of periods between $`t_1`$ and $`t_2`$. Since $`t_1`$ can occur anywhere in the time span $`0t_1P`$, then the probability that both transits will occur during the window(s) of observations is,
$$P_\mathrm{W}=\frac{1}{P}_0^Pdt_1\mathrm{\Theta }[𝒯(t_1)]\mathrm{\Theta }[𝒯(t_1+n_\mathrm{P}P)],$$
(4)
for any combination of $`n_1=0,1,\mathrm{},N1`$, $`n_2=n_1,n_1+1,\mathrm{}.,N1`$, and $`n_\mathrm{P}=1,2,\mathrm{},N\lambda /P`$.
Finally, consider a transit of duration $`2t_\mathrm{T}`$ that occurs well within the observing window. Assuming that the transit is monitored continuously with a telescope that records $`\alpha `$ electrons per second per unit flux, the total signal-to-noise of the transit is,
$$Q=(2\alpha t_\mathrm{T})^{1/2}\frac{F_0}{[S_{\mathrm{tot}}\mathrm{\Omega }_{\mathrm{PSF}}+F_0]^{1/2}}\left(\frac{R_\mathrm{p}}{R_{}}\right)^2G,$$
(5)
where $`S_{\mathrm{tot}}=S_{\mathrm{sky}}+S_{\mathrm{back}}`$ is the total surface brightness (sky + unresolved background), $`\mathrm{\Omega }_{\mathrm{PSF}}`$ is the area of the PSF, and the function $`G`$ is defined as,
$$G\left(\frac{R_{}}{R_\mathrm{p}}\right)^2\left[\frac{1}{2}_1^1d\tau [\delta (\tau )]^2\right]^{1/2},$$
(6)
and depends on $`R_\mathrm{p}/R_{}`$, $`\mathrm{cos}i`$, and the limb darkening of the star. For $`R_\mathrm{p}R_{}`$ and no limb-darkening, $`G=1`$. Note that an implicit assumption in equation (5) is that $`\delta 1`$. For a transit to be detectable, $`Q`$ must exceed some minimum threshold, $`Q_{\mathrm{min}}`$. Integration over $`\mathrm{cos}i`$ then defines the probability that a transit will satisfy the signal-to-noise requirement,
$$P_{\mathrm{S}/\mathrm{N}}=(\mathrm{cos}i_{\mathrm{min}})^1_0^{\mathrm{cos}i_{\mathrm{min}}}\mathrm{d}(\mathrm{cos}i)\mathrm{\Theta }(QQ_{\mathrm{min}}).$$
(7)
The total detection probability is then just $`P_{\mathrm{tot}}=P_\mathrm{T}\times P_\mathrm{W}\times P_{\mathrm{S}/\mathrm{N}}`$.
Consider a population of stars with luminosity function $`\mathrm{\Phi }(F_0)`$ (in units of number per area), mass-flux relation $`M(F_0)`$ and radius-flux relation $`R(F_0)`$. Assuming a fraction $`f`$ of these stars have planets of radius $`R_\mathrm{p}`$ distributed according to $`(a)\mathrm{d}a`$ (which I will assume is independent of $`F_0`$), then the number of planets detected in a field of view of area $`\mathrm{\Omega }_{\mathrm{CCD}}`$ is,
$$N_{\mathrm{det}}=f\mathrm{\Omega }_{\mathrm{CCD}}da(a)dF_0P_{\mathrm{tot}}\mathrm{\Phi }(F_0).$$
(8)
## 4. Application to the Galactic Bulge
Before applying the results of § 2 to the Galactic bulge, several assumptions must be made about the population being monitored, and also the telescope and observational setup. I will consider observations in the $`I`$-band, which provides a good compromise between dust extinction and high sky background. I construct an $`I`$-band luminosity function by combining the determination toward Baade’s window by Holtzman et al. (1998) on the bright ($`M_I9`$) end with the local M-dwarf luminosity function as determined by Gould, Bahcall & Flynn (1997) for the faint end. I normalize the latter to agree with Holtzman et al. (1998) at $`M_I=7.25`$. I adopt a distance modulus of $`14.52`$ and an extinction of $`A_I=1.0`$, appropriate for Baade’s window. For $`M(F_0)`$ and $`R(F_0)`$, I use the 10 Gyr, solar metallicity isochrone of Girardi et al. (2000). These relations are shown in Figure 1. Varying the age and/or metallicity of the population within a reasonable range does not affect the results substantially (see Fig. 1).
FIG. 1 (a) The adopted luminosity function (LF) for the Galactic bulge, constructed from the Holtzman et al. (1997) LF for $`M_I9`$, and the Gould et al. (1997) LF for $`M_I9`$. (b) The solid line shows the adopted mass-luminosity relation, as determined from the 10 Gyr, solar-metallicity isochrones of Giraldi et al. (2000). Also shown are the mass-luminosity relations for a 5 Gyr and solar-metalliticity population (dashed line) and a 10 Gyr and half-solar metallicity population (dotted line). (c) The radius-luminosity relation, as determined from the same isochrones as in panel (b). The bottom axis shows $`M_I`$, while the top axis shows $`I`$, assuming a distance modulus of 14.52 and $`A_I=1.0`$.
I assume that $`S_{\mathrm{sky}}=19.5\mathrm{mag}\mathrm{arcsec}^2`$, $`\mathrm{\Omega }_{\mathrm{PSF}}=\pi \theta ^2`$, where $`\theta `$ is the seeing, and $`\alpha =600(D/10\mathrm{m})^2e^{}s^1`$ at $`I=20`$, where $`D`$ is the telescope diameter. In the crowded fields toward the Galactic bulge, the surface brightness $`S_{\mathrm{back}}`$ due to unresolved sources will depend strongly on the seeing. To estimate $`S_{\mathrm{back}}`$, I first use the LF to determine the magnitude at which the sources become unresolved, i.e., I determine the $`I_{\mathrm{min}}`$ such that all source brighter than $`I_{\mathrm{min}}`$ contribute on average one star per seeing disk. Then $`S_{\mathrm{back}}`$ is just the total surface brightness due to all stars fainter than $`I_{\mathrm{min}}`$.
These assumptions can now be used, along with the results of § 2 to determine the number of planets that may be detected in a monitoring campaign toward the Galactic bulge as a function of the radius and separation of the planet, and to explore the effects of the diameter $`D`$ of the telescope, the seeing, $`\theta `$, and the total number of nights $`N`$ the field is monitored, on the number of detections.
FIG. 2 The detection probability as a function of $`I`$ magnitude for a planet with radius $`R_\mathrm{p}=1.0R_\mathrm{J}`$ and separation $`a=0.05\mathrm{AU}`$, or period $`P=4.08[M(I)/M_{}]^{1/2}\mathrm{days}`$, assuming 10 nights of 8 hours per night on a 10m telescope with 0.75” seeing. The solid curve shows the total detection probability, which is the given by $`P_{\mathrm{tot}}=P_\mathrm{T}\times P_\mathrm{W}\times P_{\mathrm{S}/\mathrm{N}}`$, where $`P_\mathrm{T}`$ is the probability that the planet will transit its parent star, $`P_\mathrm{W}`$ is the probability that two transits will occur during observation windows, and $`P_{\mathrm{S}/\mathrm{N}}`$ is the probability that the transit will have total signal-to-noise $`>10`$.
Figure 2 shows the total detection probability $`P_{\mathrm{tot}}`$ for a planet of radius $`R_\mathrm{p}=1R_\mathrm{J}`$ and separation $`a=0.05\mathrm{AU}`$ as a function of $`I`$ magnitude, under the assumptions of 10 nights of 8 hours per night on a 10m telescope with $`\theta =0.75^{\prime \prime }`$, and a minimum signal-to-noise of $`Q_{\mathrm{min}}=10`$ for a detection. For planets orbiting stars slightly fainter than the turn-off, $`19I21`$, almost all transits occurring during the windows of observation create significant ($`QQ_{\mathrm{min}}=10`$) transits, i.e. $`P_{\mathrm{S}/\mathrm{N}}1`$. For $`I19`$, the radii of the sources rapidly increases, rendering the transits undetectable. For $`I21`$, the sources produce too few photons to pass the signal-to-noise criterion. For this particular separation, $`a=0.05\mathrm{AU}`$, the probability $`P_\mathrm{W}`$ that the planet will transit twice during the windows of observation drops precipitously for $`20I21`$, since the period of the planet, $`P=4.08[M(I)/M_{}]^{1/2}\mathrm{days}`$, moves into “anti-resonance” with the observation window, $`T_\mathrm{W}=8\mathrm{hours}`$. However, such effects will approximately average out when a range of separations is considered.
The number of planets detected during a monitoring campaign can now be determined by integrating over the luminosity function of the sources and the separation of the companions (c.f. Eq. 8). This requires knowledge of the frequency and distribution of planetary companions to the bulge sources. Obviously, little is known about planetary companions to bulge stars. However, radial velocity surveys do provide information on the frequency and distribution of planetary companions to solar-type stars in the local neighborhood. Cumming, Marcy & Butler (1999) performed a statistical study of 74 solar-type stars from the Lick radial velocity survey. Of these, 2 had confirmed planetary ($`M_\mathrm{p}10M_\mathrm{J}`$) companions with separations $`0.1\mathrm{AU}`$, i.e. $`3\pm 2\%`$ of the sample. Furthermore, they note that the distribution in orbital radius shows a “piling-up” toward small orbital radii, but that this trend is not statistically significant. It does hint, however, that the distribution in $`a`$ may not be uniform. I will therefore assume that $`f=1\%`$ of all stars have planetary companions distributed uniformly in $`\mathrm{log}a`$ between $`0.01`$ and $`0.1\mathrm{AU}`$. While this may not accurately reflect the frequency and distribution of planets in either the solar neighborhood or the bulge, it is at least consistent with available observations.
FIG. 3 The number of planets detected per AU per square arcminute monitored as a function of orbital separation, assuming that 1% of all stars have planets distributed uniformly in $`\mathrm{log}a`$ between $`a=0.01`$ and $`0.1`$. All parameters are held constant at $`R_\mathrm{p}=1.0R_\mathrm{J}`$, $`a=0.05\mathrm{AU}`$, $`N=10\mathrm{nights}`$, $`T_\mathrm{W}=8\mathrm{hours}`$, $`D=10\mathrm{m}`$, and $`\theta =0.75^{\prime \prime }`$, unless otherwise stated. (a) The effect of varying the total number of nights. (b) The effect of varying the seeing, $`\theta `$. (c) The effect of varying the radius of the planet, $`R_\mathrm{p}`$.
Figure 3 shows the differential distribution of the number of detected planets per unit area, $`\mathrm{d}(N_{\mathrm{det}}\mathrm{\Omega }_{CCD}^1)/\mathrm{d}a`$, as a function of $`a`$ for $`R_\mathrm{p}=1R_\mathrm{J}`$, assuming $`f=1\%`$, $`(a)1/a`$, and the fiducial campaign with parameters $`N=10`$, $`T_\mathrm{W}=8`$ hours, $`D=10`$m, and $`\theta =0.75^{\prime \prime }`$. Each panel shows the result of varying $`N`$, $`\theta `$, and $`R_\mathrm{p}`$. Decreasing the duration of the campaign to $`N=5`$ nights will not substantially decrease the number of detections: most of the planets lost will be at large orbital separations, where the detection efficiency and frequency are already low. Similarly, doubling the number of nights will not substantially enhance the number of detections, although it enables the detection of planets at orbital separations larger than $`0.1\mathrm{AU}`$. The number of detected planets depends quite crucially on the seeing: increasing the seeing increases the number of unresolved sources, and therefore adds to the background flux. As $`\theta `$ increases, the signal-to-noise degrades, and transits quickly fall below the minimum detectable threshold. Thus detections are lost, and preferentially so for smaller separations (where the duration of the transits are shorter). Conversely, improving the seeing dramatically increases the number of detections. Therefore, transit searches toward the Galactic bulge should be carried out at good sites with seeing better than 1”.
FIG. 4 The number of planets detected square arcminute as a function of the radius of the planet for three different telescope apertures, assuming 10 nights of 8 hours per night. I have assumed a seeing of $`\theta =0.75^{\prime \prime }`$, and that 1% of all stars have planets distributed uniformly in $`\mathrm{log}a`$ between $`a=0.01`$ and $`0.1`$. The thick errorbar on the point for $`D=10`$m and $`R_\mathrm{p}=1R_\mathrm{J}`$ corresponds to changing the number of nights by $`N=10_5^{+10}`$. The thin errorbar corresponds to changing the seeing by $`\theta =[0.75\pm 0.25]^{\prime \prime }`$.
Figure 4 shows the number of planets detected per unit area $`N_{\mathrm{det}}\mathrm{\Omega }_{CCD}^1`$, as a function of $`R_\mathrm{p}`$, for telescope apertures of $`D=10`$m, 4m, and 1m, and the fiducial parameters, $`N=10`$, $`T_\mathrm{W}=8`$ hours, and $`\theta =0.75^{\prime \prime }`$. For a $`5^{}\times 5^{}`$ field-of-view, a 10m telescope would detect $`100`$ planets of radius $`R_\mathrm{p}=1.5R_\mathrm{J}`$, and $`30`$ planets if $`R_\mathrm{p}=1.0R_\mathrm{J}`$. Most of these planets will be discovered at $`a0.02\mathrm{AU}`$ around stars at or slightly below the turn off: the number-weighted $`I`$-magnitude of the sources for which detections are made is $`\overline{I}=19.4`$ and the number-weighted orbital separation is $`\overline{a}=0.021`$ for $`R_\mathrm{p}=1.0R_\mathrm{J}`$. These values approximately constant for $`1.0R_\mathrm{p}/R_\mathrm{J}2.0`$. For small planetary radii, $`R_\mathrm{p}\begin{array}{c}<\hfill \\ \hfill \end{array}0.8`$, $`N_{\mathrm{det}}<1`$. Thus if most planets have radii less than that of Jupiter, it will be quite difficult to detect them around stars in the Galactic bulge, unless the seeing is excellent, $`\theta 0.5^{\prime \prime }`$ (see Fig. 4). For 1m and 4m class telescope, the number of detected event is almost negligible below $`1.5R_\mathrm{J}`$. Therefore, such monitoring campaigns are unlikely to detect any planets, unless there exists a substantial population of companions in the Galactic bulge with radii $`R_\mathrm{p}>1.5R_\mathrm{J}`$.
## 5. Conclusions
In this Letter, I have estimated the number of planets that may be detected by transits in a monitoring campaign toward the Galactic bulge. An investment of a relatively modest amount of telescope resources, 10 clear nights of 8 hours per night on a 10m telescope at a site with excellent ($`0.75^{\prime \prime }`$) median seeing, would result in the detection of $`30`$ planets of Jupiter size, if the frequency and distribution of planetary companions to stars in the Galactic bulge is similar to those of G-dwarfs in the solar neighborhood. Most of these planets will be found at orbital separations of $`a0.02\mathrm{AU}`$ around stars slightly fainter than the turn-off, i.e. evolved G or early K dwarfs. Modifications to the observing strategy, such as decreasing the number of nights to 5 instead of 10, will not result in substantially fewer detections. However, if the seeing is substantially worse than $`0.75^{\prime \prime }`$, the number of detections will be considerably smaller. Therefore an excellent site is required. Similar campaigns involving 1m- or 4m-class telescopes are unlikely to result in any detections toward the bulge. Thus, collaborations currently monitoring the Galactic bulge for microlensing events are unlikely to serendipitously detect any transits.
## Acknowledgements
I would like to thank Andreas Berlind, Andrew Gould, and Paul Martini for helpful discussions, and Alberto Conti and Penny Sackett for reading the manuscript. The original idea for this paper is due to Penny Sackett. This work was supported by an Ohio State University Presidential Fellowship
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# Phase-ordering and persistence: relative effects of space-discretization, chaos, and anisotropy
## 1 Introduction
Following a quench at low temperature, bistable “ferromagnetic” systems usually exhibit domain coarsening dynamics. This phase separation process has been observed in various experimental setups, as well as in numerous model systems . The Ising model and its continuous counterpart, the so-called time-dependent Ginzburg-Landau (TDGL) equation are usually taken as paradigms of non-conservative coarsening dynamics . But many other systems, e.g. the diffusion equation, also exhibit phase-ordering dynamics once suitable “phases” are defined (for the diffusion equation, one can for instance consider the sign of a zero-mean field). Spatially-extended chaotic systems such as coupled map lattices (CMLs) can also exhibit coarsening transients, after which one chaotic phase dominates another, leaving the system in a long-range ordered state that is usually accompanied by a non-trivial evolution of spatially-averaged quantities .
The standard theory of domain coarsening predicts that, in general, the correlation length $`L(t)`$ grows algebraically in time, $`L(t)t^{1/z}`$ . Moreover, the exponent usually takes the value $`z=2`$ in systems with a non-conserved, scalar order parameter when the domain growth is driven by curvature. This universality of domain growth processes is now well established. Another important quantifier of phase ordering dynamics is persistence , defined, e.g., as the fraction of space that has remained in the same phase since some given initial time $`t_0`$. Persistence is seen to decay algebraically $`p(t)t^\theta `$ with exponent $`\theta `$ in systems with algebraic domain growth, reflecting the stationarity of two-time correlations expressed in logarithmic time. As a matter of fact, the correlation length $`L(t)`$ provides a better, “natural” measure of time, by which persistence decays as $`p(t)L(t)^{\overline{\theta }}`$ (with $`\overline{\theta }=z\theta `$) even in systems with, say, logarithmic growth of domains. The (degree of) universality of persistence exponents is still largely an open question today. Even models as close to each other as the zero-temperature Ising model and the TDGL equation —they share the same exponent $`z=2`$, the same structure function, the same Fisher-Huse exponent— seem to possess different persistence exponents (currently available measurements give $`\theta =0.20`$ for TDGL and $`0.22`$ for Ising ).
In fact, since it depends on the whole history of each point in the space, persistence is a rather complex, non-local quantity. Exact or approximate values of persistence exponents are usually not simple numbers. Understanding the extent to which persistence properties are universal is one of the challenging problems of modern statistical physics. It is also important from an experimental point of view, since persistence is an easily measurable quantifier of phase ordering. One can hope to bring new light to this problem by studying non-conventional models like CMLs. In a sense, they are not constrained as are traditional models: their local dynamics cannot be rigorously reduced to that of standard models: there is no Hamiltonian, no detailed balance, and they have been shown (numerically) to exhibit Ising-like phase transitions with critical exponents that are significantly different from those of the Ising model .
In , domain growth has been studied in simple CMLs. Numerical simulations showed the expected scaling behavior, but with exponents $`z`$ and $`\theta `$ continuously varying with the strength of the diffusive coupling between chaotic units, although $`\overline{\theta }`$ was found to be universal. It was also found that “normal”, TDGL, values were recovered in the continuous space limit of CMLs, suggesting that a role is played by lattice effects.
In this article, we come back to these somewhat surprising results and study in some detail the dynamics of the interfaces delimiting domains, in order to unravel the origin of the peculiar behavior observed on more global quantities, such as $`L(t)`$.
## 2 Domain growth and interface dynamics in chaotic CMLs
We consider $`d=2`$ dimensional square lattices $``$ of diffusively-coupled identical maps $`S_\mu `$ acting on real variables $`(X_\stackrel{}{r})_\stackrel{}{r}`$:
$$X_\stackrel{}{r}^{t+1}=(12dg)S_\mu (X_\stackrel{}{r}^t)+g\underset{\stackrel{}{e}𝒱}{}S_\mu (X_{\stackrel{}{r}+\stackrel{}{e}}^t),$$
(1)
where $`𝒱`$ is the set of $`2d`$ nearest neighbors $`\stackrel{}{e}`$ of site $`\stackrel{}{0}`$. For simplicity, we present results obtained for the piecewise linear, odd, local map $`S_\mu `$ defined by:
$$S_\mu (X)=\{\begin{array}{ccc}\mu X\hfill & \mathrm{if}\hfill & X[1/3,1/3]\hfill \\ 2\mu /3\mu X\hfill & \mathrm{if}\hfill & X[1/3,1]\hfill \\ 2\mu /3\mu X\hfill & \mathrm{if}\hfill & X[1,1/3].\hfill \end{array}$$
(2)
Choosing $`\mu [1,2]`$ guarantees that the local map has two symmetric invariant intervals, allowing us to define “spins” as $`\sigma _\stackrel{}{r}=\mathrm{sign}(X_\stackrel{}{r})`$. The deterministic system thus defined is similar to the Ising model at zero temperature in the sense that local variables can flip only when crossed by an interface.
Take $`\mu =1.9`$ and consider different choices for coupling strength $`g`$. Two cases are illustrated in Fig. 1: at small $`g`$, the evolution leads to blocked clusters of the two “phases” corresponding to the two invariant intervals (which “shrink” but are preserved by the linear diffusive coupling). Domain walls are pinned between lattice sites; the system is multistable. On the contrary, for larger values of the coupling, domain coarsening never stops, leading to the emergence of long-range order.
Thus, there exists a threshold value $`g_\mathrm{e}`$ separating these regimes: domain growth is expected to slow down as $`g`$ goes to $`g_\mathrm{e}`$ from above, corresponding to a gradual “freezing” of interface dynamics.
### 2.1 Previous results
We first review some of the results presented in for the sake of completeness.
Starting from random initial conditions with exactly zero magnetization (for the “spin” variables $`\sigma _\stackrel{}{r}=\mathrm{sign}(X_\stackrel{}{r})`$), we measured $`L(t)`$ defined by the two-point correlation function at mid-height, as well as $`p(t)`$, the persistent fraction of sites since $`t_0=0`$. Fig. 2 shows the results obtained from single runs on lattices of $`2048^2`$ sites simulated up to time $`t=10^4`$ (at later times, finite-size fluctuations become too important). Algebraic growth of $`L(t)`$ and decay of $`p(t)`$ is observed, but with effective exponents $`z`$ and $`\theta `$ which vary with $`g`$, the coupling strength. The decay of $`\theta `$ with $`g`$ is best fitted with some power-law dependence which defines the threshold coupling $`g_\mathrm{e}0.169(1)`$.
Normal, $`z=2`$ coarsening is only recovered when one approaches the continuous-space limit of CMLs which, in practice, consists in applying the coupling operator many times. In this limit, the persistence exponent $`\theta `$ approaches the value known for the TDGL equation, i.e., $`\theta 0.20`$. The underlying lattice thus seems to influence the propagation of fronts in the system, by discretization and/or anisotropy effects.
### 2.2 Interface dynamics
Anisotropic curvature-driven domain growth can usually be defined by the following expression for the normal velocity of the interface :
$$v_\mathrm{n}=M(\phi )[\sigma (\phi )+\sigma ^{\prime \prime }(\phi ))\kappa +h]$$
(3)
where $`M(\phi )`$ is the mobility, $`\sigma (\phi )`$ the surface tension, $`\kappa `$ the local curvature, $`h`$ an “external field”, and $`\phi `$ an angle specifying the orientation of the vector normal to the interface in space.
It is thus important to know whether these quantities can be defined in the CMLs of interest here, and, in the affirmative case, to try to estimate them in order to investigate possible discrepancies with Eq.(3).
The mobility is usually measured as the response of a flat interface ($`\kappa =0`$) to the driving field $`h`$. Previous work has indicated that the mobility may not be defined in deterministic lattice systems . This is also the case of the CMLs studied here. Replacing the evolution equation (1) by:
$$X_\stackrel{}{r}^{t+1}=h+(1h)((12dg)S_\mu (X_\stackrel{}{r}^t)+g\underset{\stackrel{}{e}𝒱}{}S_\mu (X_{\stackrel{}{r}+\stackrel{}{e}}^t)),$$
(4)
one observes moving interfaces for large-enough $`h`$, but the velocity may not be constant, and the $`h`$ dependence of the average velocity is not always linear. Moreover, there generally exists a finite critical field value below which the interface does not move (Fig. 3). In fact, the mobility seems to be ill-defined for all rational values of $`\mathrm{tan}\phi `$. Thus the mobility of the interfaces of the CML studied here is not well defined, a further indication of the special character of domain growth.
Another typical experiment to try to assess the validity of Eq. (3) is to study the fate of droplets of one phase immersed in the other phase. In anisotropic situations such as those described by Eq. (3), one expects, from some circular initial droplet, a transient stage during which the droplet takes its asymptotic shape and then shrinks self-similarly, its volume decreasing linearly with time; its lifetime is therefore proportional to its initial volume . This is roughly what is observed with the CML studied here, except that the expected scaling is recorded only for droplets of large initial radius (Fig.4). For smaller size, one finds that the lifetime of droplets scales like $`r^{2\alpha }`$, with $`1/2\alpha 1/z`$, where $`z`$ is the growth exponent measured in Fig.2. Interestingly, the crossover size is approximately of the order of the typical size of domains at the end of the runs shown in Fig.2.
This suggests that the non-trivial scalings laws reported in Fig.2 might only be transient, ultimately leaving place to normal (curvature-driven-like) phase ordering.
### 2.3 Large-scale domain growth
We thus went back to domain growth transients following random initial conditions, but with larger system sizes and longer simulation times than before. Fig. 5 shows the growth of $`L(t)`$ for a system of $`8192^2`$ sites during $`10^5`$ timesteps (runs performed on a parallel machine for a total of approximately 2000 cpu hours). A log-log plot of $`L(t)`$ reveals an increase of the local slope at long times, which, however, does not reach the “normal” value $`1/z=1/2`$. On the other hand, the emergence of the normal growth behavior is clearly seen in a plot of $`L`$ vs $`t^{1/2}`$, which becomes linear after approximately $`10^4`$ iterations (Fig. 5a). The approach of the freezing transition at $`g=g_\mathrm{e}`$ is therefore contained in the prefactor of the growth law $`L(t)A(g)t^{1/2}`$. The effective “mobility” $`A(g)`$ seems to decrease continuously to zero. An excellent fit of the data is $`A(g)(gg_\mathrm{e})^w`$ with $`w=0.52(3)`$ (Fig. 5b), yielding the new estimate $`g_\mathrm{e}0.171(1)`$, but we cannot exclude a linear behavior near threshold which would yield a $`g_\mathrm{e}`$ value closer to our previous estimate in .
To sum up: domain coarsening in the CML studied here does show discrepancies with Eq. (3) (such as the ill-defined character of the mobility), but this seems to quantitatively influence the dynamics at early times only, so that normal growth behavior is recovered asymptotically at long times. It must be noted, though, that the effect of these features, specific to lattice determinitic systems, decreases rather slowly in time, as testified by the long crossover times recorded in our simulations.
## 3 Persistence issues
### 3.1 Persistence at late times
We now turn to the other set of results presented in , those related to the decay of $`p(t)`$, the fraction of persistent sites. Fig. 6a shows the decay of $`p(t)`$ for the same run as in Fig. 5. The behavior of the local exponent clearly changes at late times, in coincidence with the crossover observed in the growth of $`L(t)`$. The persistence exponents estimated in thus reflect only the short-time behavior of the coarsening process.
Given that $`L(t)`$ reaches its asymptotic behavior rather late, one would ideally like to start measuring persistence from an initial time $`t_0`$ larger than the crossover time, so that the whole history of the system coming into account in the decay of $`p(t)`$ be situated in the asymptotic scaling regime. Unfortunately, this is basically out of reach of current computers’ power, since it would require very large systems and their simulation over times much larger than $`t_0`$.
On the other hand, it was noticed in that the decay of the persistence probability recorded as a function of $`L(t)`$ is in a sense “more universal”, as it seems to be roughly independent of $`g`$. Those measurements were performed in the intermediate scaling regime where the growth of $`L(t)`$ is not normal. For later times and larger system sizes, the variation of $`p`$ with $`L`$ is seen to be better behaved than the simple decay of $`p(t)`$ (Fig. 6b). The local slope is approximately constant, and especially so after the crossover time for the growth of $`L(t)`$. We also checked that changing $`t_0`$ does not significantly influence the results. We can thus estimate a reliable persistence exponent, whose variation with $`g`$ is not significant given our numerical accuracy, and which takes approximately the value known for the TDGL equation, i.e. $`\theta 0.20(2)`$.
### 3.2 Generalized persistence
The behavior of persistence seems independent of $`g`$, or rather we cannot resolve the possible variation of exponent $`\theta `$ as we go from the threshold (near $`g_\mathrm{e}`$) to the continuous-space limit. Moreover, although the estimated $`\theta `$ is close to the TDGL value, the error bar does not allow us to rule out the currently accepted value for the Ising model.
That $`\theta `$ is close to the TDGL value is not surprising if one considers that both our system and the TDGL equation are governed by deterministic dynamics. In the strong coupling limit of CMLs, in particular, interfaces evolve smoothly on a continuous space, and their motion is expected to resemble that of interfaces in TDGL. On the other hand, near $`g_e`$, strong lattice/anisotropy effects arise, yielding somewhat jittery interface motion, similarly to the discrete walls observed in Ising model.
Although persistence seems to depend only weakly on these lattice effects, it was recently proposed that these discrepancies become obvious when generalized persistence is considered . In order to define generalized persistence, let us consider the time-average of the coarse-grained variable $`\sigma _\stackrel{}{r}^t=\mathrm{sign}(X_\stackrel{}{r}^t)`$. At each point $`\stackrel{}{r}`$ in space and starting from some initial time $`t_0`$, it reads,
$$M_\stackrel{}{r}^t=\frac{1}{tt_0}\underset{t^{}=t_0}{\overset{t}{}}\sigma _\stackrel{}{r}^t^{}.$$
Generalized persistence $`P(t;x)`$ can then be defined as the probability that $`M_\stackrel{}{r}^t`$ has always remained above some threshold $`x[1,1]`$ :
$$P(t;x)=\mathrm{Prob}\left(M^t^{}x;t^{}t\right).$$
Generalized persistence is nothing but the “standard” persistence for the process $`\mathrm{sign}(M^tx)`$, and is thus also expected to decay algebraically in time with exponent $`\theta (x)`$. This provides a spectrum $`\theta (x)`$ of generalized persistence exponents, which contains, in particular, the “standard” persistence exponent $`\theta (1)=\theta `$.
As they discriminate in deeper details the possible paths followed by the variable $`\sigma _\stackrel{}{r}^t`$, generalized persistence exponents are expected to encode more information on the properties of interface motion than simple persistence and, in particular, to be sensitive to the influence of the jittery motion of interfaces .
Here we estimate $`\theta (x)`$ near and away from $`g_\mathrm{e}`$. The measurement of generalized persistence has been carried out on the same runs as those presented previously. Spectra of generalized persistence exponents are displayed in Fig. 7. It should first be noted that these spectra are difficult to resolve numerically for small $`x`$ values, since the corresponding data relies on the relatively rare spins which have spent most of their time in the phase opposite to their original phase. Thus, the data presented in Fig. 7 for this region can only have an indicative value (it is believed on general grounds, however, that $`\theta (x)`$ goes to zero as $`(1+x)^\theta `$ when $`x1`$).
On the other hand, the normalized exponents $`\theta (x)/\theta `$ are rather well-resolved numerically near $`x=1`$, and their behavior in this region is expected to characterize the short-scale motion of interfaces. The insert in Figure 7 shows $`\theta (x)/\theta `$ for our CML as well as for the Ising model and the TDGL equation. The normalized exponent spectra of our CML tend to the TDGL behavior as $`g`$ goes away from $`g_\mathrm{e}`$ and approaches the continuous-space limit. One can notice, moreover, that our CML exhibits, under identical experimental conditions, the same value $`\theta 0.204`$ for both values of $`g`$ presented here, despite the important change in the strength of discretisation effects. Thus, these effects are not felt at the level of simple persistence but are well captured, in a sense, by generalized persistence. Space discretisation involves at least two factors: the special behavior of interfaces (as seen in Section 2.2) and anisotropy. The results reported here do not allow one to separate them but imply that they seem to have only a marginal effect on $`\theta `$. We note, further, that this is in agreement with recent studies of persistence in anisotropic partial differential equations in which no detectable change of $`\theta `$ with anisotropy could be recorded , but in disagreement with the claims of .
## 4 Conclusion
The results presented here have elucidated the surprising coarsening behavior of chaotic CMLs previously reported in , which we have shown to be representative only of the long, intermediate scaling region present in such deterministic systems. “Normal”, $`z=2`$ phase-ordering is recovered at large time/length scales for all coupling strengths $`g>g_\mathrm{e}`$.
In fact, both the growth law and the persistence exponent $`\theta `$ seem to be independent of $`g`$ in this limit. Although our numerical data does not allow for a definitive statement about this, it provides an interesting case for the current debate considering which factors —space-discretization, anisotropy, chaos or probabilistic motion— are possibly influencing persistence exponents .
The general picture emerging for the phase-ordering properties of our CML is as follows: while TDGL-like, smooth interface motion is recovered in the continuous-space limit, near $`g_\mathrm{e}`$, important lattice/anisotropy effects yield interfacial properties somewhat similar to that of the Ising model; but these are only felt in generalized persistence scaling, confirming that this quantity captures details of interface dynamics. Generalized persistence spectra $`\theta (x)`$ show a significant qualitative change as $`g`$ goes from $`g_\mathrm{e}`$ to the continuous-space limit, similar to the recently observed difference for $`\theta (x)`$ between the Ising and TDGL. Thus our CML, to some extent, interpolates between these two models.
In this regard, the ultimate question of the “true” asymptotic behavior of the phase-ordering properties of our CML must depend on further knowledge about the status of the commonly observed numerical differences in persistence properties between the TDGL equation and the Ising model. Future work should focus on resolving the question of observed discrepancy’s origin, possibly with the help of “intermediate” systems such as the CML studied here.
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# ENTANGLEMENTS AND COMPOUND STATES IN QUANTUM INFORMATION THEORY
## 1. Introduction
Recently, the specifically quantum correlations, called in quantum physics entanglements, are used to study quantum information processes, in particular, quantum computation, quantum teleportation, quantum cryptography . There have been mathematical studies of the entanglements in , in which the entangled state is defined by a compound state which can not be written as a convex combination $`_n\mu \left(n\right)\varsigma _n\varrho _n`$ with any states $`\varrho _n`$ and $`\varsigma _n.`$ However it is obvious that there exist several important applications with correlated states written as separable forms above. Such correlated, or entangled states have been also discussed in several contexts in quantum probability such as quantum measurement and filtering , quantum compound state and lifting . In this paper, we study the mathematical structure of quantum entangled states to provide a finer classification of quantum sates, and we discuss the informational degree of entanglement and entangled quantum mutual entropy.
We show that the pure entangled states can be treated as generalized compound states, the nonseparable states of quantum compound systems which are not representable by convex combinations of the product states.
The mixed compound states, defined as convex combinations by orthogonal decompositions of their input marginal states $`\varrho _0`$, have been introduced in for studying the information in a quantum channel with the general output C\*-algebra $`𝒜`$. This o-entangled compound state is a particular case of so called separable state of a compound system, the convex combination of the arbitrary product states which we call c-entangled. We shall prove that the o-entangled compound states are most informative among c-entangled states in the sense that the maximum of mutual information over all c-entanglements to the quantum system $`(𝒜,\varrho )`$ is achieved on the extreme o-entangled states, defined by a Schatten decomposition of a given state $`\varrho `$ on $`𝒜`$. This maximum coincides with von Neumann entropy $`S\left(\varrho \right)`$ of the state $`\varrho `$, and it can also be achieved as the maximum of the mutual information over all couplings with classical probe systems described by a maximal Abelian subalgebra $`𝒜^{}𝒜`$. Thus the couplings described by c-entanglements of (quantum) probe systems $``$ to a given system $`𝒜`$ don’t give an advantage in maximizing the mutual information in comparison with the quantum-classical couplings, corresponding to the Abelian $`=𝒜^{}`$. The achieved maximal information $`S\left(\varrho \right)`$ coincides with the classical entropy on the Abelian subalgebra $`𝒜^{}`$ of a Schatten decomposition for $`\varrho `$, and is bounded by $`\mathrm{ln}\mathrm{rank}𝒜=\mathrm{ln}dim𝒜^{}`$, where $`\mathrm{rank}𝒜`$ is the rank of the von Neumann algebra $`𝒜`$ defined as the dimensionality of a maximal Abelian subalgebra. Due to $`dim𝒜\left(\mathrm{rank}𝒜\right)^2`$, it is achieved on the normal central $`\rho =\left(\mathrm{rank}𝒜\right)^1I`$ only in the case of finite dimensional $`𝒜`$.
More general than o-entangled states, the d-entangled states, are defined as c-entangled states by orthogonal decomposition of only one marginal state on the probe algebra $``$. They can give bigger mutual entropy for a quantum noisy channel than the o-entangled state which gains the same information as d-entangled extreme states in the case of a deterministic channel.
We prove that the truly (strongest) entangled states are most informative in the sense that the maximum of mutual entropy over all entanglements to the quantum system $`𝒜`$ is achieved on the quasi-compound state, given by an extreme entanglement of the probe system $`=𝒜`$ with coinciding marginals, called standard for a given $`\varrho `$. The standard entangled state is o-entangled only in the case of Abelian $`𝒜`$ or pure marginal state $`\varrho `$. The gained information for such extreme q-compound state defines another type of entropy, the quasi-entropy $`S_q\left(\varrho \right)`$ which is bigger than the von Neumann entropy $`S\left(\varrho \right)`$ in the case of non-Abelian $`𝒜`$ (and mixed $`\varrho `$.) The maximum of mutual entropy over all quantum couplings, described by true quantum entanglements of probe systems $``$ to the system $`𝒜`$ is bounded by $`\mathrm{ln}dim𝒜`$, the logarithm of the dimensionality of the von Neumann algebra $`𝒜`$, which is achieved on a normal tracial $`\rho `$ in the case of finite dimensional $`𝒜`$. Thus the q-entropy $`S_q\left(\varrho \right)`$, which can be called the dimensional entropy, is the true quantum entropy, in contrast to the von Neumann rank entropy $`S\left(\varrho \right)`$, which is semi-classical entropy as it can be achieved as a supremum over all couplings with the classical probe systems $``$. These entropies coincide in the classical case of Abelian $`𝒜`$ when $`\mathrm{rank}𝒜=dim𝒜`$. In the case of non-Abelian finite-dimensional $`𝒜`$ the q-capacity $`C_q=\mathrm{ln}dim𝒜`$ is achieved as the supremum of mutual entropy over all q-encodings (correspondences), described by entanglements. It is strictly bigger then the semi-classical capacity $`C=\mathrm{ln}\mathrm{rank}𝒜`$ of the identity channel, which is achieved as the supremum over usual encodings, described by the classical-quantum correspondences $`𝒜^{}𝒜`$.
In this paper we consider the case of a discrete decomposable C\*-algebra $`𝒜`$ for which the results are achieved by relatively simple proofs. The purely quantum case of a simple algebra $`𝒜=\left(\right)`$, for which some proofs are rather obvious was considered in a short paper . The general case of decomposable C\*-algebra $`𝒜`$ to include the continuous systems, and will be published elsewhere.
## 2. Compound States and Entanglements
Let $``$ denote the (separable) Hilbert space of a quantum system, and $`\left(\right)`$ be the algebra of all linear bounded operators on $``$. In order to include the classical discrete systems as a particular quantum case, we shall fix a decomposable subalgebra $`𝒜\left(\right)`$ of bounded observables $`A𝒜`$ of the form $`A=\left[A\left(i\right)\delta _i^k\right]`$, where $`A\left(i\right)\left(_i\right)`$ are arbitrary bounded operators in Hilbert subspaces $`_i`$, corresponding to an orthogonal decomposition $`=_i_i`$. A bounded linear functional $`\varrho :𝒜𝐂`$ is called a state on $`𝒜`$ if it is positive (i.e., $`\varrho \left(A\right)0`$ for any positive operator $`A`$ in $`𝒜`$) and normalized $`\varrho (I)=1`$ for the identity operator $`I`$ in $`𝒜`$ . A normal state can be expressed as
(1)
$$\varrho \left(A\right)=\mathrm{tr}_𝒢\kappa ^{}A\kappa =\mathrm{tr}A\rho ,\text{ }A𝒜$$
where $`𝒢`$ is another separable Hilbert space, $`\kappa `$ is a linear Hilbert-Schmidt operator from $`𝒢`$ to $``$ and $`\kappa ^{}`$ is the adjoint operator of $`\kappa `$ from $``$ to $`𝒢`$. This $`\kappa `$ is called the amplitude operator, and it is called just the amplitude if $`𝒢`$ is one dimensional space $``$, corresponding to the pure state $`\varrho \left(A\right)=\kappa ^{}A\kappa `$ for a $`\kappa `$ with $`\kappa ^{}\kappa =\kappa ^2=1,`$ in which case $`\kappa ^{}`$ is the adjoint functional from $``$ to $``$. The density operator $`\rho =\kappa \kappa ^{}`$ is uniquely defined by the condition $`\rho 𝒜`$ as a decomposable trace class operator $`\mathrm{P}_𝒜=\mathrm{P}_𝒜\left(i\right)`$ with P$`{}_{𝒜}{}^{}\left(i\right)\left(_i\right)`$,
$$\nu \left(i\right)=\mathrm{tr}__i\mathrm{P}_𝒜\left(i\right)0,\underset{i}{}\nu \left(i\right)=1.$$
Thus the predual space $`𝒜_{}`$ can be identified with the direct sum $`𝒯\left(_i\right)`$ of the Banach spaces $`𝒯\left(_i\right)`$ of trace class operators in $`_i`$ (the density operators $`\mathrm{P}_𝒜𝒜_{}`$, $`\mathrm{P}_{}_{}`$ of the states $`\varrho `$, $`\varsigma `$ on different algebras $`𝒜`$, $``$ will be usually denoted by different letters $`\rho ,\sigma `$ corresponding to their Greek variations $`\varrho `$, $`\varsigma `$.)
In general, $`𝒢`$ is not one dimensional, the dimensionality $`dim𝒢`$ must be not less than $`\mathrm{rank}\rho `$, the dimensionality of the range $`\mathrm{ran}\rho `$ of the density operator $`\rho .`$ We shall equip it with an isometric involution $`J=J^{}`$, $`J^2=I`$, having the properties of complex conjugation on $`𝒢`$,
$$J\lambda _j\zeta _j=\overline{\lambda _j}J\zeta _j,\lambda _j𝐂,\zeta _j𝒢$$
with respect to which $`J\sigma =\sigma J`$ for the positive and so self-adjoint operator $`\sigma =\kappa ^{}\kappa =\sigma ^{}`$ on $`𝒢`$. The latter can also be expressed as the symmetricity property $`\stackrel{~}{\varsigma }=\varsigma `$ of the state $`\varsigma \left(B\right)=`$ $`\mathrm{tr}B\sigma `$ given by the real and so symmetric density operator $`\overline{\sigma }=\sigma =\stackrel{~}{\sigma }`$ on $`𝒢`$ with respect to the complex conjugation $`\overline{B}=JBJ`$ and the tilda operation ($`𝒢`$-transponation) $`\stackrel{~}{B}=JB^{}J`$ on the algebra $`\left(𝒢\right)`$, and thus on any tilda invariant decomposable subalgebra $`\left(𝒢\right)`$ containing $`\kappa ^{}𝒜\kappa \sigma `$.
For example, $`𝒢`$ can be realized as a subspace of $`l^2(𝐍)`$ of complex sequences $`𝐍n\zeta \left(n\right)`$, with $`_n\left|\zeta \left(n\right)\right|^2<+\mathrm{}`$ in the diagonal representation $`\sigma =\left[\mu \left(n\right)\delta _n^m\right]`$. The involution $`J`$ can be identified with the complex conjugation $`C\zeta \left(n\right)=\overline{\zeta }\left(n\right)`$, i.e.,
$$C:\zeta =\underset{n}{}|n\zeta \left(n\right)C\zeta =\underset{n}{}|n\overline{\zeta }\left(n\right)$$
in the standard basis $`\left\{|n\right\}𝒢`$ of $`l^2(𝐍)`$. In this case $`\kappa =\kappa _nn|`$ is given by orthogonal eigen-amplitudes $`\kappa _n`$, $`\kappa _m^{}\kappa _n=0`$, $`mn`$, normalized to the eigen-values $`\lambda \left(n\right)=\kappa _n^{}\kappa _n=\mu \left(n\right)`$ of the density operator $`\rho `$ such that $`\rho =\kappa _n\kappa _n^{}`$ is a Schatten decomposition, i.e. the spectral decomposition of $`\rho `$ into one-dimensional orthogonal projectors. In any other basis the operator $`J`$ is defined then by $`J=U^{}CU`$, where $`U`$ is the corresponding unitary transformation. One can also identify $`𝒢`$ with $``$ by $`U\kappa _n=\lambda \left(n\right)^{1/2}|n`$ such that the operator $`\rho `$ is real and symmetric, $`J\rho J=\rho =J\rho ^{}J`$ in $`𝒢=`$ with respect to the involution $`J`$ defined in $``$ by $`J\kappa _n=\kappa _n`$. Here $`U`$ is an isometric operator $`l^2\left(\right)`$ diagonalizing the operator $`\rho `$: $`U\rho U^{}=|n\lambda \left(n\right)n|`$. The amplitude operator $`\kappa =\rho ^{1/2}`$ corresponding to $`=𝒜`$, $`\sigma =\rho `$ is called standard.
Given the amplitude operator $`\kappa `$, one can define not only the states $`\varrho `$ and $`\varsigma `$ on the algebras $`𝒜=\left(\right)`$ and $`=\left(𝒢\right)`$ but also a pure entanglement state $`\varpi `$ on the algebra $`𝒜`$ of all bounded operators on the tensor product Hilbert space $`𝒢`$ by
$$\varpi \left(BA\right)=\mathrm{tr}_𝒢\stackrel{~}{B}\kappa ^{}A\kappa =\mathrm{tr}_{}A\kappa \stackrel{~}{B}\kappa ^{}.$$
Indeed, thus defined $`\varpi `$ is uniquely extended by linearity to a normal state on the algebra $`𝒜`$ generated by all linear combinations $`C=\lambda _jB_jA_j`$ due to $`\varpi \left(II\right)=\mathrm{tr}\kappa ^{}\kappa =1`$ and
$`\varpi \left(C^{}C\right)`$ $`=`$ $`{\displaystyle \underset{i,k}{}}\overline{\lambda }_i\lambda _k\mathrm{tr}_𝒢\stackrel{~}{B}_k\stackrel{~}{B}_i^{}\kappa ^{}A_i^{}A_k\kappa `$
$`=`$ $`{\displaystyle \underset{i,k}{}}\overline{\lambda }_i\lambda _k\mathrm{tr}_𝒢\stackrel{~}{B}_i^{}\kappa ^{}A_i^{}A_k\kappa \stackrel{~}{B}_k=\mathrm{tr}_𝒢\chi ^{}\chi 0,`$
where $`\chi =_jA_j\kappa \stackrel{~}{B}_j`$. This state is pure on $`\left(𝒢\right)`$ as it is given by an amplitude $`\vartheta 𝒢`$ defined as
$$\left(\zeta \eta \right)^{}\vartheta =\eta ^{}\kappa J\zeta ,\zeta 𝒢,\eta ,$$
and it has the states $`\varrho `$ and $`\varsigma `$ as the marginals of $`\varpi `$:
(2)
$$\varpi \left(IA\right)=\mathrm{tr}_{}A\rho ,\varpi \left(BI\right)=\mathrm{tr}_𝒢B\sigma .$$
As follows from the next theorem for the case $`=`$ , any pure state
$$\varpi \left(BA\right)=\vartheta ^{}\left(BA\right)\vartheta ,B,A𝒜$$
given on $`\left(𝒢\right)`$ by an amplitude $`\vartheta 𝒢`$ with $`\vartheta ^{}\vartheta =1`$, can be achieved by a unique entanglement of its marginal states $`\varsigma `$ and $`\varrho `$.
Theorem 2.1. Let $`\varpi :𝒜`$ be a compound state
$$\varpi \left(BA\right)=\mathrm{tr}_{}\upsilon ^{}\left(BA\right)\upsilon ,$$
defined by an amplitude operator $`\upsilon :𝒢`$ on a separable Hilbert space $``$ into the tensor product Hilbert space $`𝒢`$ with
$$\upsilon \upsilon ^{}𝒜,\mathrm{tr}_{}\upsilon ^{}\upsilon =1.$$
Then this state can be achieved as an entanglement
(3)
$$\varpi \left(BA\right)=\mathrm{tr}_𝒢\stackrel{~}{B}\kappa ^{}\left(IA\right)\kappa =\mathrm{tr}_{}\left(IA\right)\kappa \stackrel{~}{B}\kappa ^{}$$
of the states (2) with $`\sigma =\kappa ^{}\kappa `$ and $`\rho =\mathrm{tr}_{}\kappa \kappa ^{}`$, where $`\kappa `$ is an amplitude operator $`𝒢`$ with
(4)
$$\kappa ^{}\left(I𝒜\right)\kappa ,\mathrm{tr}_{}\kappa \kappa ^{}𝒜.$$
The entangling operator $`\kappa `$ is uniquely defined by $`\stackrel{~}{\kappa }U=\upsilon `$ up to a unitary transformation $`U`$ of the minimal domain $`=\mathrm{dom}\upsilon `$.
###### Proof.
Without loss of generality we can assume that the space $``$ is equipped with an isometric involution $`J`$ as well as the space $`𝒢`$ is equipped with $`J`$ . The entangling operator $`\kappa `$ can be defined then as $`\kappa =\left(UI\right)\stackrel{~}{\upsilon }`$ by
$$(U\xi \eta )^{}\kappa \zeta =(\xi \eta )^{}\stackrel{~}{\upsilon }\zeta :=\left(J\zeta \eta \right)^{}\upsilon J\xi ,\xi ,\zeta 𝒢,\eta ,$$
where $`U`$ is arbitrary linear isometry in $``$. Indeed, let $`\left\{\xi _k\right\}`$ be an orthonormal basis of $``$, in which, say (but not necessary,) the density operator $`\upsilon ^{}\upsilon `$ is diagonal, and $`J:`$ be the complex conjugation in this basis, $`J\xi _k=\xi _k`$, defining an isometric involution in $``$. In general $`J`$ is different from the complex conjugation $`C`$, given by $`C|n=|n`$ in the standard basis $`\left\{|n;n\right\}`$ if $``$ is identified with a subspace $`l^2\left(𝐍\right)`$ for the diagonal representation of $`\upsilon ^{}\upsilon `$. Note that although the isometric transformation $`U=_k|k\xi _k^{}`$ of the arbitrary basis $`\left\{\xi _k\right\}`$ into $`\left\{|k\right\}l^2\left(𝐍\right)`$ is also arbitrary, it can be always considered as real with respect to $`C`$ and $`J=U^{}CU`$, in the sense $`\overline{U}:=CUJ=U`$, and so $`\stackrel{~}{U}:=CU^{}J=U^{}`$. Defining $`\kappa =\kappa _nn|`$ in the standard basises of $``$ and $`𝒢`$ as the block-matrix $`_{kn}|k\psi _k\left(n\right)n|`$ transposed to $`_{kn}|n\psi _k\left(n\right)k|`$, where the amplitudes $`\psi _k\left(n\right)`$ are given by the matrix elements $`\eta ^{}\psi _k\left(n\right)=\left(n|\eta ^{}\right)\upsilon \xi _k`$, we obtain
$`\mathrm{tr}_𝒢\stackrel{~}{B}\kappa ^{}\left(IA\right)\kappa `$ $`=`$ $`{\displaystyle \underset{n,m}{}}n\left|\stackrel{~}{B}\right|m\psi _k^{}\left(m\right)A\psi _k\left(n\right)`$
$`=`$ $`{\displaystyle \underset{n,m}{}}\psi _k^{}\left(m\right)m\left|B\right|nA\psi _k\left(n\right)=\mathrm{tr}_{}\upsilon ^{}\left(BA\right)\upsilon \text{ .}`$
Hence $`\kappa :𝒢`$ , defined by $`\kappa _n=|k\psi _k\left(n\right)`$ as the transposed to $`\upsilon U^{}=\upsilon \stackrel{~}{U}\stackrel{~}{\kappa }`$, is the required entangling operator of the form $`\kappa =\left(UI\right)\stackrel{~}{\upsilon }`$ with $`\kappa ^{}\kappa `$$`=\sigma =\mathrm{tr}_{}\upsilon \upsilon ^{}`$ and $`\mathrm{tr}_{}\kappa \kappa ^{}=\rho =\mathrm{tr}_𝒢\upsilon \upsilon ^{}`$. Moreover, it satisfies the conditions (4) as $`\omega =\upsilon \upsilon ^{}𝒜`$ and thus
$$\kappa ^{}\left(IA\right)\kappa =\mathrm{tr}_{}\left(IA\right)\omega ,\mathrm{tr}_{}\kappa \stackrel{~}{B}\kappa ^{}=\mathrm{tr}_𝒢\left(BI\right)\omega 𝒜.$$
The uniqueness follows from the obvious isometricity of the families
$$\{\underset{k}{}|k\eta ^{}\psi _k\left(n\right):n,\eta \},\{\underset{k}{}\eta ^{}\psi _k\left(n\right)\xi _k^{}:n,\eta \}$$
of vectors $`\left(I\eta ^{}\right)\kappa |n`$ in $`l^2\left(𝐍\right)`$ and of $`\left(n|\eta ^{}\right)\upsilon `$ in $`^{}`$ which follows from
$$\mathrm{tr}_𝒢|nm|\kappa ^{}\left(I\eta \eta ^{}\right)\kappa =\mathrm{tr}_{}\upsilon ^{}\left(|mn|\eta \eta ^{}\right)\upsilon .$$
Thus they are unitary equivalent in the minimal space $``$. So the entangling operator $`\kappa `$ is defined in the minimal $``$ up to the unitary equivalence, corresponding to the arbitrary of the unitary operator $`U`$ in $``$, intertwining the involutions $`C`$ and $`J`$.∎
Note that the entangled state (3) is written as
$$\varpi \left(BA\right)=\mathrm{tr}_𝒢\stackrel{~}{B}\pi \left(A\right)=\mathrm{tr}_{}A\pi _{}\left(\stackrel{~}{B}\right),$$
where $`\pi \left(A\right)=\kappa ^{}\left(IA\right)\kappa `$, bounded by $`A\sigma _{}`$ for any $`A\left(\right)`$, is in the predual space $`_{}`$ of all trace-class operators in $`𝒢`$, and $`\pi _{}\left(B\right)=\mathrm{tr}_{}\kappa B\kappa ^{}`$, bounded by $`B\rho 𝒜_{}`$, is in $`𝒜_{}𝒜`$. The map $`\pi `$ is the Steinspring form of the general completely positive map $`𝒜_{}`$, written in the eigen-basis $`\left\{|k\right\}`$ of the density operator $`\upsilon ^{}\upsilon `$ as
(5)
$$\pi \left(A\right)=\underset{m,n}{}|m\kappa _m^{}\left(IA\right)\kappa _nn|,A𝒜$$
while the dual operation $`\pi _{}`$ is the Kraus form of the general completely positive map $`\mathrm{A}𝒜_{}`$, given in this basis as
(6)
$$\pi _{}\left(B\right)=\underset{n,m}{}n\left|B\right|m\mathrm{tr}_{}\kappa _n\kappa _m^{}=\mathrm{tr}_𝒢\stackrel{~}{B}\omega .$$
It corresponds to the general form
(7)
$$\omega =\underset{m,n}{}|nm|\mathrm{tr}_{}\kappa _n\kappa _m^{}$$
of the density operator $`\omega =\upsilon \upsilon ^{}`$ for the entangled state $`\varpi \left(BA\right)=\mathrm{tr}\left(BA\right)\omega `$ in this basis, characterized by the weak orthogonality property
(8)
$$\mathrm{tr}_{}\psi \left(m\right)^{}\psi \left(n\right)=\mu \left(n\right)\delta _n^m$$
in terms of the amplitude operators $`\psi \left(n\right)=\left(In|\right)\stackrel{~}{\kappa }=\stackrel{~}{\kappa }_n`$.
Definition 2.1. The dual map $`\pi _{}:𝒜_{}`$ to a completely positive map $`\pi :𝒜_{}`$, normalized as $`\mathrm{tr}_𝒢\pi \left(I\right)=1`$, is called the quantum entanglement of the state $`\varsigma =\pi \left(I\right)`$ on $``$ to the state $`\varrho =\pi _{}\left(I\right)`$ on $`𝒜`$. The entanglement by
(9)
$$\pi _{}^{}\left(A\right)=\rho ^{1/2}A\rho ^{1/2}=\pi ^{}\left(A\right)$$
of the state $`\varsigma =\varrho `$ on the algebra $`=𝒜`$ is called standard for the system $`(𝒜,\varrho )`$.
The standard entanglement defines the standard compound state
$$\varpi _0\left(BA\right)=\mathrm{tr}_{}\stackrel{~}{B}\rho ^{1/2}A\rho ^{1/2}=\mathrm{tr}_{}A\rho ^{1/2}\stackrel{~}{B}\rho ^{1/2}$$
on the algebra $`𝒜𝒜`$, which is pure, given by the amplitude $`\vartheta _0=\stackrel{~}{\kappa }_0`$, where $`\kappa _0=\rho ^{1/2}`$ in the case of the simple algebra $`𝒜=\left(\right)`$. In the general case of decomposable $`𝒜=\left(_i\right)`$ with the density operator $`\rho =\rho \left(i\right)`$ having more than one components $`\rho \left(i\right)=\rho _i\nu \left(i\right)`$ with $`\nu \left(i\right)=\mathrm{tr}\rho \left(i\right)0`$ and positive $`\rho _i\left(_i\right)`$, the standard state $`\varpi _0`$ is a mixture
(10)
$$\varpi _0\left(BA\right)=\underset{i}{}\vartheta _0^i\left(B\left(i\right)A\left(i\right)\right)\vartheta _0^i\nu \left(i\right),A,B𝒜$$
of such pure compound states given by the amplitudes $`\vartheta _0^i_i_i`$ with $`\stackrel{~}{\vartheta }_0^i\stackrel{~}{\vartheta }_0^i=\rho _i`$. The standard amplitudes $`\vartheta _0^i_i_i`$ for an orthogonal decomposition $`\upsilon _0=_i\vartheta _0^i\xi _i^{}\nu \left(i\right)^{1/2}`$ of the standard amplitude operator $`\upsilon _0:_0`$ are defined as $`\stackrel{~}{\kappa }_0\left(i\right)/\stackrel{~}{\kappa }_0\left(i\right)`$ by the entangling components $`\kappa _0\left(i\right)=\rho \left(i\right)^{1/2}`$ with
$$\left(\zeta _i\eta _i\right)^{}\stackrel{~}{\kappa }_0\left(i\right)=\eta _i^{}\kappa _0\left(i\right)J\zeta _i,\eta _i,\zeta _i_i.$$
\[Example\] In quantum physics the entangled states are usually obtained by a unitary transformation $`U`$ of an initial disentangled state, described by the density operator $`\sigma _0\rho _0\tau _0`$ on the tensor product Hilbert space $`𝒢𝒦`$ , that is,
$$\varpi \left(BA\right)=\mathrm{tr}U_1^{}\left(BAI\right)U_1\left(\sigma _0\rho _0\tau _0\right).$$
In the simple case, when $`𝒦=`$, $`\tau _0=1`$, the joint amplitude operator $`\upsilon `$ is defined on the tensor product $`=𝒢_0`$ with $`_0=\mathrm{ran}\rho _0`$ as $`\upsilon =U_1\left(\sigma _0\rho _0\right)^{1/2}`$. The entangling operator $`\kappa `$, describing the entangled state $`\varpi `$, is constructed as it was done in the proof of Theorem 1 by transponation of the operator $`\upsilon U^{}`$, where $`U`$ is arbitrary isometric operator $`𝒢_0`$. The dynamical procedure of such entanglement in terms of the completely positive map $`\pi _{}:𝒜_{}`$ is the subject of Belavkin quantum filtering theory . The quantum filtering dilation theorem proves that any entanglement $`\pi `$ can be obtained the unitary entanglement as the result of quantum filtering by tracing out some degrees of freedom of a quantum environment, described by the density operator $`\tau _0`$ on the Hilbert space $`𝒦`$, even in the continuous time case.
## 3. C- and D-Entanglements and Encodings
The compound states play the role of joint input-output probability measures in classical information channels, and can be pure in quantum case even if the marginal states are mixed. The pure compound states achieved by an entanglement of mixed input and output states exhibit new, non-classical type of correlations which are responsible for the EPR type paradoxes in the interpretation of quantum theory. The mixed compound states on $`𝒜`$ which are given as the convex combinations
$$\varpi =\underset{n}{}\varsigma _n\varrho _n\mu \left(n\right),\mu \left(n\right)0,\underset{n}{}\mu \left(n\right)=1$$
of tensor products of pure or mixed normalized states $`\varrho _n𝒜_{}`$, $`\varsigma _n_{}`$ as in classical case, do not exhibit such paradoxical behavior, and are usually considered as the proper candidates for the input-output states in the communication channels. Such separable compound states are achieved by c-entanglements, the convex combinations of the primitive entanglements $`B\mathrm{tr}_𝒢B\omega _n`$, given by the density operators $`\omega _n=\sigma _n\rho _n`$ of the product states $`\varpi _n=\varsigma _n\varrho _n`$:
(11)
$$\pi _{}\left(B\right)=\underset{n}{}\rho _n\mathrm{tr}_𝒢B\sigma _n\mu \left(n\right),$$
A compound state of this sort was introduced by Ohya in order to define the quantum mutual entropy expressing the amount of information transmitted from an input quantum system to an output quantum system through a quantum channel, using a Schatten decomposition $`\sigma =_n\sigma _n\mu \left(n\right)`$, $`\sigma _n=|nn|`$ of the input density operator $`\sigma `$. It corresponds to a particular, diagonal type
(12)
$$\pi \left(A\right)=\underset{n}{}|n\kappa _n^{}\left(IA\right)\kappa _nn|$$
of the entangling map (5) in an eigen-basis $`\left\{|n\right\}𝒢`$ of the density operator $`\sigma `$, and is discussed in this section.
Let us consider a finite or infinite input system indexed by the natural numbers $`n𝐍`$. The associated space $`𝒢l^2\left(𝐍\right)`$ is the Hilbert space of the input system described by a quantum projection-valued measure $`n|nn|`$ on $``$, given an orthogonal partition of unity $`I=|nn|`$ $``$ of the finite or infinite dimensional input Hilbert space $`𝒢`$. Each input pure state, identified with the one-dimensional density operator $`|nn|`$ corresponding to the elementary symbol $`n`$, defines the elementary output state $`\varrho _n`$ on $`𝒜`$. If the elementary states $`\varrho _n`$ are pure, they are described by output amplitudes $`\eta _n`$ satisfying $`\eta _n^{}\eta _n=1=\mathrm{tr}\rho _n`$, where $`\rho _n=`$ $`\eta _n\eta _n^{}`$ are the corresponding output one-dimensional density operators. If these amplitudes are non-orthogonal $`\eta _m^{}\eta _n\delta _n^m`$, they cannot be identified with the input amplitudes $`|n`$.
The elementary joint input-output states are given by the density operators $`|nn|\rho _n`$ in $`𝒢`$. Their mixtures
(13)
$$\omega =\underset{n}{}\mu \left(n\right)|nn|\rho _n,$$
define the compound states on $`𝒜`$, given by the quantum correspondences $`n|nn|`$ with the probabilities $`\mu \left(n\right)`$. Here we note that the quantum correspondence is described by a classical-quantum channel, and the general d-compound state for a quantum-quantum channel in quantum communication can be obtained in this way due to the orthogonality of the decomposition (13), corresponding to the orthogonality of the Schatten decomposition $`\sigma =_n|n\mu \left(n\right)n|`$ for $`\sigma =\mathrm{tr}_{}\omega `$.
The comparison of the general compound state (7) with (13) suggests that the quantum correspondences are described as the diagonal entanglements
(14)
$$\pi _{}\left(B\right)=\underset{n}{}\mu \left(n\right)n|B|n\rho _n,$$
They are dual to the orthogonal decompositions (12):
$$\pi \left(A\right)=\underset{n}{}\mu \left(n\right)|n\eta _n^{}A\eta _nn|=\underset{n}{}|n\eta \left(n\right)^{}A\eta \left(n\right)n|,$$
where $`\eta \left(n\right)=\mu \left(n\right)^{1/2}\eta _n`$. These are the entanglements with the stronger orthogonality
(15)
$$\psi \left(m\right)\psi \left(n\right)^{}=\rho \left(n\right)\delta _n^m,$$
for the amplitude operators $`\psi \left(n\right):`$ of the decomposition $`\upsilon =_n|n\psi \left(n\right)`$ in comparison with the orthogonality (8). The orthogonality (15) can be achieved in the following manner: Take in (5) $`\kappa _n=|n\eta \left(n\right)`$ with $`m|n=\delta _n^m`$ so that
$$\kappa _m^{}\left(IA\right)\kappa _n=\mu \left(n\right)\eta _n^{}A\eta _n\delta _n^m$$
for any $`A𝒜`$. Then the strong orthogonality condition (15) is fulfilled by the amplitude operators $`\psi \left(n\right)=\eta \left(n\right)n|=\stackrel{~}{\kappa }_n`$, and
$$\kappa ^{}\kappa =\underset{n}{}\mu \left(n\right)|nn|=\sigma ,\kappa \kappa ^{}=\underset{n}{}\eta \left(n\right)\eta \left(n\right)^{}=\rho .$$
It corresponds to the amplitude operator for the compound state (13) of the form
(16)
$$\upsilon =\underset{n}{}|n\psi \left(n\right)U,$$
where $`U`$ is arbitrary unitary operator from $``$ onto $`𝒢`$, i.e. $`\upsilon `$ is unitary equivalent to the diagonal amplitude operator
$$\kappa =\underset{n}{}|nn|\eta \left(n\right)$$
on $`=𝒢`$ into $`𝒢`$. Thus, we have proved the following theorem in the case of pure output states $`\rho _n=\eta _n\eta _n^{}`$.
Theorem 3.1. Let $`\pi `$ be the operator (13), defining a d-compound state of the form
(17)
$$\varpi \left(BA\right)=\underset{n}{}n|B|n\mathrm{tr}__n\psi _n^{}A\psi _n\mu \left(n\right)$$
Then it corresponds to the entanglement by the orthogonal decomposition (12) mapping the algebra $`𝒜`$ into a diagonal subalgebra of $``$.
###### Proof.
Let $`_n_n`$ be the Hilbert orthogonal sum of the domains $`_n`$ for the amplitude operators $`\psi _n`$ in (17) with an isometric involution $`C_n`$. In the case $`_n=`$ of the amplitudes $`\psi _n`$ corresponding to pure states $`\rho _n`$ the involution $`_nC_n`$ is the componentwise complex conjugation in $`_nl^2\left(\right)`$; in the general case it is given by some isometric involutions $`C_n`$ in the Hilbert spaces $`_n`$, which are equivalent to the ranges $`_n=\rho _n`$ of the density operators $`\rho _n=\psi _n\psi _n^{}`$ with the standard involutions in their eigen-representations, or contain these ranges. We can define the global output amplitude operator $`\psi \left(n\right)`$ on $`=_n_n`$ by
$$\psi \left(n\right)=\mu \left(n\right)^{1/2}\psi _nϵ_n^{},$$
where $`ϵ_n:_n`$ are the canonical orthogonal isometries, $`ϵ_m^{}ϵ_n=I_n\delta _n^m`$, and by (16) an amplitude operator $`\upsilon :𝒢`$ of the compound state (17), defining its density operator $`\omega =\upsilon \upsilon ^{}`$ independently of the unitary transformation $`U`$ of the Hilbert space onto $`_n_n`$.
The entangling operator $`\kappa =_n\kappa _nn|`$ is then defined by its components $`\kappa _n`$ of the form
$$\kappa _n=\left(ϵ_nI\right)\stackrel{~}{\psi }_n\mu \left(n\right)^{1/2}=\stackrel{~}{\psi }\left(n\right),$$
Here $`\stackrel{~}{\psi }_n`$ are the amplitudes in $`_n`$ obtained from the operators $`\psi _n:_n`$ by
$$\left(\xi _n\eta \right)^{}\stackrel{~}{\psi }_n=\eta ^{}\psi _nC_n\xi _n,\eta ,\xi _n_n$$
In particular $`\kappa `$ is the diagonal amplitude operator with the components $`\kappa _n=_m\delta _n^m\stackrel{~}{\psi }\left(n\right)`$ in $`_m_m`$:
(18)
$$\kappa =\underset{n}{}\kappa _nn|=_m\stackrel{~}{\psi }\left(m\right)m|.$$
Thus the entanglement (6) corresponding to (17) is given by the dual to (12) diagonal map (14) with the density operators $`\rho \left(n\right)=\psi \left(n\right)\psi \left(n\right)^{}=\mathrm{tr}_{}\kappa _n\kappa _n^{}`$ normalized to the probabilities $`\mu \left(n\right)=\kappa _n^{}\kappa _n`$.∎
Note that (2.9) defines the general form of a positive map on $`𝒜`$ with values in the simultaneously diagonal trace-class operators in $`\mathrm{A}`$.
Definition 3.1. A convex combination (11) of the primitive CP maps $`\rho _n\varsigma _n`$ is called c-entanglement, and is called d-entanglement, or quantum encoding if it has the diagonal form (14) on $``$. The d-entanglement is called o-entanglement and compound state is called o-compound if all density operators $`\rho _n`$ are orthogonal: $`\rho _m\rho _n=\rho _n\rho _m`$ for all $`m`$ and $`n`$.
Note that due to the commutativity of the operators $`BI`$ with $`IA`$ on $`𝒢`$, one can treat the correspondences as the nondemolition measurements in $``$ with respect to $`𝒜`$. So, the compound state is the state prepared for such measurements on the input $`𝒢`$. It coincides with the mixture of the states, corresponding to those after the measurement without reading the sent message. The set of all d-entanglements corresponding to a given Schatten decomposition of the input state $`\sigma `$ on $``$ is obviously convex with the extreme points given by the pure output states $`\rho _n`$ on $`𝒜`$, corresponding to a not necessarily orthogonal decompositions $`\rho =_n\rho \left(n\right)`$ into one-dimensional density operators $`\rho \left(n\right)=\mu \left(n\right)\rho _n.`$
The Schatten decompositions $`\rho =_n\lambda \left(n\right)\rho _n`$ correspond to the extreme d-entanglements, $`\rho _n=\eta _n\eta _n^{}`$, $`\mu \left(n\right)=\lambda \left(n\right)`$, characterized by orthogonality $`\rho _m\rho _n=0`$, $`mn`$ . They form a convex set of d-entanglements with mixed commuting $`\rho _n`$ for each Schatten decomposition of $`\rho `$. The orthogonal d-entanglements were used in to construct a particular type of Accardi’s transitional expectations and to define the entropy in a quantum dynamical system via such transitional expectations.
The established structure of the general q-compound states suggests also the general form
$$\mathrm{\Phi }_{}(B,\varrho _0)=\mathrm{tr}__1X^{}\left(B\rho _0\right)X=\mathrm{tr}_𝒢\left(\stackrel{~}{B}I\right)Y\left(I\rho _0\right)Y^{}$$
of transitional expectations $`\mathrm{\Phi }_{}:\times 𝒜_{}^{}𝒜_{}`$, describing the entanglements $`\pi _{}=\mathrm{\Phi }_{}\left(\varrho _0\right)`$ of the states $`\varsigma =\pi \left(I\right)`$ to $`\varrho =\pi _{}\left(I\right)`$ for each initial state $`\varrho _0𝒜_{}^{}`$ with the density operator $`\rho _0𝒜^{}\left(_0\right)`$ by $`\pi _{}\left(B\right)=\mathrm{tr}_{}\kappa \left(BI\right)\kappa ^{}`$, where $`\kappa =X^{}\left(I\rho _0\right)^{1/2}`$. It is given by an entangling transition operator $`X:𝒢_0`$, which is defined by a transitional amplitude operator $`Y:_0𝒢`$ up to a unitary operator $`U`$ in $``$ as
$$\left(\zeta \eta _0\right)^{}X\left(U\xi \eta \right)=\left(\eta _0J\xi \right)^{}Y^{}\left(J\zeta \eta \right)\text{.}$$
The dual map $`\mathrm{\Phi }:𝒜_{}𝒜^{}`$ is obviously normal and completely positive,
(19)
$$\mathrm{\Phi }\left(A\right)=X\left(IA\right)X^{}_{}𝒜^{},A𝒜,$$
with $`\mathrm{tr}_𝒢\mathrm{\Phi }\left(I\right)=I^{}`$, and is called filtering map with the output states
$$\varsigma =\mathrm{tr}__0\mathrm{\Phi }\left(I\right)\left(I\rho _0\right)$$
in the theory of CP flows over $`𝒜=𝒜^{}`$. The operators $`Y`$ normalized as $`\mathrm{tr}_{}Y^{}Y=I^{}`$ describe $`𝒜`$-valued q-compound states
$$\mathrm{E}\left(BA\right)=\mathrm{tr}_{}Y^{}\left(BA\right)Y=\mathrm{tr}_𝒢\left(\stackrel{~}{B}I\right)\mathrm{\Phi }\left(A\right),$$
defined as the normal completely positive maps $`𝒜𝒜^{}`$ with $`\mathrm{E}\left(II\right)=I^{}`$ .
If the $`𝒜`$-valued compound state has the diagonal form given by the orthogonal decomposition
(20)
$$\mathrm{\Phi }\left(A\right)=\underset{n}{}|n\mathrm{tr}_{}\mathrm{\Psi }\left(n\right)^{}A\mathrm{\Psi }\left(n\right)n|,$$
corresponding to $`Y`$ $`=_n|n\mathrm{\Psi }\left(n\right)`$, where $`\mathrm{\Psi }\left(n\right):_0`$, it is achieved by the d-transitional expectations
$$\mathrm{\Phi }_{}(B,\varrho _0)=\underset{n}{}n|B|n\mathrm{\Psi }\left(n\right)\left(\rho _0I\right)\mathrm{\Psi }\left(n\right)^{}.$$
The d-transitional expectations correspond to the instruments of the dynamical theory of quantum measurements. The elementary filters
$$\mathrm{\Theta }_n\left(A\right)=\frac{1}{\mu \left(n\right)}\mathrm{tr}_{}\mathrm{\Psi }^{}\left(n\right)A\mathrm{\Psi }\left(n\right),\mu \left(n\right)=\mathrm{tr}\mathrm{\Psi }\left(n\right)\left(\rho _0I\right)\mathrm{\Psi }^{}\left(n\right)$$
define posterior states $`\varrho _n=\varrho _0\mathrm{\Theta }_n`$ on $`𝒜`$ for quantum nondemolition measurements in $``$, which are called indirect if the corresponding density operators $`\rho _n`$ are non-orthogonal. They describe the posterior states with orthogonal
$$\rho _n=\mathrm{\Psi }_n\left(\rho _0I\right)\mathrm{\Psi }_n^{},\mathrm{\Psi }_n=\mathrm{\Psi }\left(n\right)/\mu \left(n\right)^{1/2}$$
for all $`\rho _0`$ iff $`\mathrm{\Psi }\left(n\right)^{}\mathrm{\Psi }\left(n\right)=\delta _n^mM\left(n\right)`$.
## 4. Quantum Entropy via Entanglements
As it was shown in the previous section, the diagonal entanglements describe the classical-quantum encodings $`\varkappa :𝒜_{}`$, i.e. correspondences of classical symbols to quantum, in general not orthogonal and pure, states. As we have seen in contrast to the classical case, not every entanglement can be achieved in this way. The general entangled states $`\varpi `$ are described by the density operators $`\omega =\upsilon \upsilon ^{}`$ of the form (7) which are not necessarily block-diagonal in the eigen-representation of the density operator $`\sigma `$, and they cannot be achieved even by a more general c-entanglement (11). Such nonseparable entangled states are called in the quasicompound (q-compound) states, so we can call also the quantum nonseparable correspondences the quasi-encodings (q-encodings) in contrast to the d-correspondences, described by the diagonal entanglements.
As we shall prove in this section, the most informative for a quantum system $`(𝒜,\varrho )`$ is the standard entanglement $`\pi _{}^{}=\pi _0`$ of the probe system $`(^{},\varsigma _0)=(𝒜,\varrho )`$, described in (9). The other extreme cases of the self-dual input entanglements
$$\pi _{}\left(A\right)=\underset{n}{}\rho \left(n\right)^{1/2}A\rho \left(n\right)^{1/2}=\pi \left(A\right),$$
are the pure c-entanglements, given by the decompositions $`\rho =\rho \left(n\right)`$ into pure states $`\rho \left(n\right)=\eta _n\eta _n^{}\mu \left(n\right)`$. We shall see that these c-entanglements, corresponding to the separable states
(21)
$$\omega =\underset{n}{}\eta _n\eta _n^{}\eta _n\eta _n^{}\mu \left(n\right),$$
are in general less informative then the pure d-entanglements, given in an orthonormal basis $`\left\{\eta _n^{}\right\}`$ by
$$\pi ^{}\left(A\right)=\underset{n}{}\eta _n^{}\eta _n^{}A\eta _n\eta _n^{}\mu \left(n\right)\pi _{}^{}\left(A\right).$$
Now, let us consider the entangled mutual information and quantum entropies of states by means of the above three types of compound states. To define the quantum mutual entropy, we need the relative entropy of the compound state $`\varpi `$ with respect to a reference state $`\phi `$ on the algebra $`𝒜`$. In our discrete case of the decomposable algebras it is defined by the density operators $`\omega ,\varphi 𝒜`$ of these states as
(22)
$$S(\varpi ,\phi )=\mathrm{tr}\omega \left(\mathrm{ln}\omega \mathrm{ln}\varphi \right).$$
It has a positive value $`S(\varpi ,\phi )[0,\mathrm{}]`$ if the states are equally normalized, say (as usually) $`\mathrm{tr}\omega =1=\mathrm{tr}\varphi `$, and it can be finite only if the state $`\varpi `$ is absolutely continuous with respect to the reference state $`\phi `$, i.e. iff $`\varpi \left(E\right)=0`$ for the maximal null-orthoprojector $`E\varphi =0`$.
The mutual information $`I_{𝒜,}\left(\varpi \right)`$ of a compound state $`\varpi `$ achieved by an entanglement $`\pi _{}:`$ $`𝒜_{}`$ with the marginals
$$\varsigma \left(B\right)=\varpi \left(BI\right)=\mathrm{tr}_𝒢B\sigma ,\varrho \left(A\right)=\varpi \left(IA\right)=\mathrm{tr}_{}A\rho $$
is defined as the relative entropy (22) with respect to the product state $`\phi =\varsigma \varrho `$:
(23)
$$I_{𝒜,}\left(\varpi \right)=\mathrm{tr}\omega \left(\mathrm{ln}\omega \mathrm{ln}\left(\sigma I\right)\mathrm{ln}\left(I\rho \right)\right).$$
Here the operator $`\omega `$ is uniquely defined by the entanglement $`\pi _{}`$ as its density in (6), or the $`𝒢`$-transposed to the operator $`\stackrel{~}{\omega }`$ in
$$\pi \left(A\right)=\kappa ^{}\left(IA\right)\kappa =\mathrm{tr}_{}A\stackrel{~}{\omega }.$$
This quantity describes an information gain in a quantum system $`(𝒜,\varrho )`$ via an entanglement $`\pi _{}`$ of another system $`(,\varsigma ).`$ It is naturally treated as a measure of the strength of an entanglement, having zero value only for completely disentangled states, corresponding to $`\varpi =\varsigma \varrho `$.
Proposition 4.1. Let $`\pi _{}^{}:^{}𝒜_{}`$ be an entanglement $`\pi _{}^{}`$ of a state $`\varsigma _0=\pi ^{}\left(I\right)`$ on a discrete decomposable algebra $`^{}\left(𝒢_0\right)`$ to the state $`\varrho =\pi _{}^{}\left(I\right)`$ on $`𝒜`$, and $`\pi _{}=\pi _{}^{}\mathrm{K}`$ be an entanglement defined as the composition with a normal completely positive unital map $`\mathrm{K}:^{}`$. Then $`I_{𝒜,}\left(\varpi \right)I_{𝒜,^{}}\left(\varpi _0\right)`$, where $`\varpi ,\varpi _0`$ are the compound states achieved by $`\pi _{}^{}`$ , $`\pi _{}`$ respectively. In particular, for any c-entanglement $`\pi _{}`$ to $`(𝒜,\varsigma )`$ there exists a not less informative d-entanglement $`\pi _{}^{}=\varkappa `$ with an Abelian $`^{}`$, and the standard entanglement $`\pi _0\left(A\right)=\rho ^{1/2}A\rho ^{1/2}`$ of $`\varsigma _0=\varrho `$ on $`^{}=𝒜`$ is the maximal one in this sense.
###### Proof.
The first follows from the monotonicity property
(24)
$$\varpi =\mathrm{K}_{}\varpi _0,\phi =\mathrm{K}_{}\phi _0S(\varpi ,\phi )S(\varpi _0,\phi _0)$$
of the general relative entropy on a von Neuman algebra $``$ with respect to the predual $`\mathrm{K}_{}`$ to any normal completely positive unital map $`\mathrm{K}:^{}`$. It should be applied to the ampliation $`\mathrm{K}\left(BA\right)=\mathrm{K}\left(B\right)A`$ of the CP map $`\mathrm{K}`$ from $`^{}`$ to $`𝒜^{}𝒜`$, with the compound state $`\mathrm{K}_{}\varpi _0=\varpi _0\left(\mathrm{K}\mathrm{I}\right)`$ ($`\mathrm{I}`$ denotes the identity map $`𝒜𝒜`$) corresponding to the entanglement $`\pi _{}=\pi _{}^{}\mathrm{K}`$ and $`\mathrm{K}_{}\phi _0=\varsigma \varrho `$ , $`\varsigma =\varsigma _0\mathrm{K}`$ corresponding to $`\phi _0=\varsigma _0\varrho `$.
This monotonicity property proves in particular that for any separable compound state on $`𝒜`$, which is prepared by a c-entanglement (11), there exists a diagonal entanglement $`\pi _{}^{}`$ to the system $`(𝒜,\varrho )`$with the same, or even bigger information gain (23). One can take even a classical system $`(^{},\varsigma _0)`$, say the diagonal sublagebra $`^{}`$ on $`𝒢_0=`$ $`l^2\left(𝐍\right)`$ with the state $`\varsigma _0`$, induced by the measure $`\nu `$, and consider the classical-quantum correspondence (encoding)
$$\pi _{}^{}\left(B^{}\right)=\underset{n}{}\beta \left(n\right)\rho _n\nu \left(n\right),B^{}=\underset{n}{}|n\beta \left(n\right)n|,$$
prescribing the states $`\varrho _n\left(A\right)=\mathrm{tr}A\rho _n`$ to the letters $`n`$ with the probabilities $`\nu \left(n\right)`$. The information gain
$$I_{𝒜,^{}}\left(\varpi _0\right)=\underset{n}{}\mu \left(n\right)\mathrm{tr}\rho _n\left(\mathrm{ln}\rho _n\mathrm{ln}\rho \right).$$
is equal or bigger then $`I_{𝒜,}\left(\varpi \right)`$ corresponding to $`\omega =_n\sigma _n\rho _n\nu \left(n\right)`$ because the entanglement (11) is represented as the composition $`\pi _{}^{}\mathrm{K}`$ with the CP map
$$\mathrm{K}\left(B\right)=\underset{n}{}|n\varsigma _n\left(B\right)n|,B$$
into the diagonal algebra $`^{}`$.
The inequality (24) can be also applied to the standard entanglement, corresponding to the compound state (10) on $`𝒜𝒜=_{i,k}𝒜\left(i\right)𝒜\left(k\right)`$, where $`𝒜\left(i\right)=\left(_i\right)`$. It is described by the density operator
(25)
$$\omega _0=_{i,k}\mathrm{P}_{𝒜𝒜}(i,k)=_i\vartheta _0^i\vartheta _0^i\nu \left(i\right)\text{,}$$
with $`\mathrm{P}_{𝒜𝒜}(i,k)=0`$, $`ik`$ concentrated on the diagonal $`_i𝒜\left(i\right)𝒜\left(i\right)`$ of $`𝒜𝒜`$. The amplitudes $`\vartheta _0^i_i_i`$ are defined in (10) by orthogonal components $`\kappa _0\left(i\right)=\rho \left(i\right)^{1/2}`$ of the central decomposition $`\kappa _0=_n|n\kappa _0\left(n\right)`$ for the standard entangling operator $`\kappa _0:l^2(𝐍)`$. Indeed, any entanglement $`\pi _{}\left(B\right)=\mathrm{tr}_{}\kappa B\kappa ^{}`$ as a normal CP map $`𝒜`$ normalized to the density operator $`\rho =\mathrm{tr}_{}\kappa \kappa ^{}`$ can be represented as the composition $`\pi _{}^{}\mathrm{K}`$ of the standard entanglement $`\pi _{}^{}=\pi _0`$ on $`(^{},\varsigma _0)=(𝒜,\varrho )`$ and a normal unital CP map $`\mathrm{K}:𝒜`$. The CP map $`\mathrm{K}`$ is defined by $`\rho ^{1/2}\mathrm{K}\left(B\right)\rho ^{1/2}=\pi _{}\left(B\right)`$ as
$$\mathrm{K}\left(B\right)=\mathrm{tr}_{_{}}X^{}BX,B$$
where $`X`$ is an operator $`_{}𝒢`$, $`\mathrm{tr}_{_{}}X^{}X=I`$ such that $`\kappa =\left(I^{}\kappa _0\right)X^{}`$ is an entangling operator for $`\pi `$. Thus the standard entanglement $`\pi _{}^{}`$ corresponds to the maximal mutual information.∎
Note that any extreme d-entanglement
$$\pi _{}^{}\left(B\right)=\underset{n}{}n|B|n\rho _n^{}\mu \left(n\right),B^{},$$
with $`\rho =_n\rho _n^{}\mu \left(n\right)`$ decomposed into pure normalized states $`\rho _n^{}=\eta _n\eta _n^{}`$, is maximal among all c-entanglements in the sense $`I_{𝒜,}\left(\varpi _0\right)I_{𝒜,}\left(\varpi \right)`$. This is because $`\mathrm{tr}\rho _n^{}\mathrm{ln}\rho _n^{}=0`$, and therefore the information gain
$$I_{𝒜,}\left(\varpi \right)=\underset{n}{}\mu \left(n\right)\mathrm{tr}\rho _n\left(\mathrm{ln}\rho _n\mathrm{ln}\rho \right).$$
with a fixed $`\pi _{}\left(I\right)=\rho `$ achieves its supremum $`\mathrm{tr}_{}\rho \mathrm{ln}\rho `$ at any such extreme d-entanglement $`\pi _{}^{}`$. Thus the supremum of the information gain (23) over all c-entanglements to the system $`(𝒜,\varrho )`$ is the von Neumann entropy
(26)
$$S_𝒜\left(\varrho \right)=\mathrm{tr}_{}\rho \mathrm{ln}\rho .$$
It is achieved on any extreme $`\pi _{}^{}`$, for example given by the maximal Abelian subalgebra $`^{}𝒜`$, with the measure $`\mu =\lambda `$, corresponding to a Schatten decomposition $`\rho =_n\eta _n^{}\eta _n^{}\lambda \left(n\right)`$, $`\eta _m^{}\eta _n^{}=\delta _n^m`$. The maximal value $`\mathrm{ln}\mathrm{rank}𝒜`$ of the von Neumann entropy is defined by the dimensionality $`\mathrm{rank}𝒜=dim^{}`$ of the maximal Abelian subalgebra of the decomposable algebra $`𝒜`$, i.e. by $`dim`$.
Definition 4.1. The maximal mutual information
(27)
$$H_𝒜\left(\varrho \right)=\underset{\pi _{}\left(I\right)=\rho }{sup}I_{𝒜,}\left(\varpi \right)=I_{𝒜,^{}}\left(\varpi _0\right),$$
achieved on $`^{}=𝒜`$ by the standard q-entanglement $`\pi _{}^{}\left(A\right)=\rho ^{1/2}A\rho ^{1/2}`$ for a fixed state $`\varrho \left(A\right)=\mathrm{tr}_{}A\rho `$, is called q-entropy of the state $`\varrho `$. The differences
$$H_{|𝒜}\left(\varpi \right)=H_{}\left(\varsigma \right)I_{𝒜,}\left(\varpi \right)$$
$$S_{|𝒜}\left(\varpi \right)=S_{}\left(\varsigma \right)I_{𝒜,}\left(\varpi \right)$$
are respectively called the q-conditional entropy on $``$ with respect to $`𝒜`$ and the degree of disentanglement for the compound state $`\varpi `$.
Obviously, $`H_{|𝒜}\left(\varpi \right)`$ is positive in contrast to the disentanglement $`S_{|𝒜}\left(\varpi \right)`$, having the positive maximal value $`S_{|𝒜}\left(\varpi \right)=S_{}\left(\varsigma \right)`$ in the case $`\varpi =\varsigma \varrho `$ of complete disentanglement, but which can achieve also a negative value
(28)
$$\underset{\pi _{}\left(I\right)=\rho }{inf}S_{|𝒜}\left(\varpi \right)=S_𝒜\left(\varsigma \right)H_𝒜\left(\varrho \right)=\underset{i}{}\nu \left(i\right)\mathrm{tr}__i\rho _i\mathrm{ln}\rho _i$$
for the entangled states as the following theorem states. Here $`\rho _i\left(_i\right)`$ are the density operators of normalized factor-states $`\varrho _i=\nu \left(i\right)^1\varrho |\left(_i\right)`$with $`\nu \left(i\right)=\varrho \left(I^i\right)`$, where $`I^i`$ are the orthoprojectors onto $`_i`$. Obviously $`H_𝒜\left(\varrho \right)=S_𝒜\left(\varrho \right)`$ if the algebra $`𝒜`$ is completely decomposable, i.e. Abelian, and the maximal value $`\mathrm{ln}\mathrm{rank}𝒜`$ of $`S_𝒜\left(\varrho \right)`$ can be written as $`\mathrm{ln}dim𝒜`$ in this case. The disentanglement $`S_{|𝒜}\left(\varpi \right)`$ is always positive in this case, as well as in the case of Abelian $``$ when $`H_{|𝒜}\left(\varpi \right)=S_{|𝒜}\left(\varpi \right)`$.
Theorem 4.2. Let $`𝒜`$ be the discrete decomposable algebra on $`=_i_i`$, with a normal state given by the density operator $`\rho =\rho \left(i\right)`$, and $`𝒞𝒜`$ be its center with the state $`\nu =\varrho |𝒞`$ induced by the probability distribution $`\nu \left(i\right)=\mathrm{tr}\rho \left(i\right).`$ Then the q-entropy is given by the formula
(29)
$$H_𝒜\left(\varrho \right)=\underset{i}{}\left(\nu \left(i\right)\mathrm{ln}\nu \left(i\right)2\mathrm{t}\mathrm{r}__i\rho \left(i\right)\mathrm{ln}\rho \left(i\right)\right),$$
i.e. $`H_𝒜\left(\varrho \right)=H_{𝒜|𝒞}\left(\varrho \right)+H_𝒞\left(\nu \right)`$, where $`H_𝒞\left(\nu \right)=_i\nu \left(i\right)\mathrm{ln}\nu \left(i\right)=S_𝒞\left(\nu \right)`$, and
$$H_{𝒜|𝒞}\left(\varrho \right)=2\underset{i}{}\nu \left(i\right)\mathrm{tr}__i\rho _i\mathrm{ln}\rho _i=2S_{𝒜|𝒞}\left(\varrho \right),$$
with $`\rho _i=\rho \left(i\right)/\nu \left(i\right)`$. It is positive, $`H_𝒜\left(\varrho \right)[0,\mathrm{}]`$, and if $`𝒜`$ is finite dimensional, it is bounded, with the maximal value $`H_𝒜\left(\varrho ^{}\right)=\mathrm{ln}dim𝒜`$ which is achieved on the tracial $`\rho ^{}=\rho _i^{}\nu ^{}\left(i\right)`$,
$$\rho _i^{}=\left(dim_i\right)^1I^i,\nu ^{}\left(i\right)=dim𝒜\left(i\right)/dim𝒜,$$
where $`dim𝒜\left(i\right)=\left(dim_i\right)^2`$, $`dim𝒜=_idim𝒜\left(i\right)`$.
###### Proof.
The q-entropy $`H_𝒜\left(\varrho \right)`$ is the supremum (27) of the mutual information (23) which is achieved on the standard entanglement, corresponding to the density operator (25) of the standard compound state (10) with $`=𝒜`$, $`\sigma =\rho `$. Thus $`H_𝒜\left(\rho \right)=I_{𝒜,𝒜}\left(\varpi _0\right)`$, where
$`I_{𝒜,𝒜}\left(\varpi _0\right)`$ $`=`$ $`\mathrm{tr}_{}\omega _0\left(\mathrm{ln}\omega _0\mathrm{ln}\left(\rho I\right)\mathrm{ln}\left(I\rho \right)\right)`$
$`=`$ $`{\displaystyle \underset{i}{}}\nu \left(i\right)\mathrm{ln}\nu \left(i\right)2\mathrm{t}\mathrm{r}\rho \mathrm{ln}\rho ={\displaystyle \underset{i}{}}\nu \left(i\right)\left(\mathrm{ln}\nu \left(i\right)+2\mathrm{t}\mathrm{r}__i\rho _i\mathrm{ln}\rho _i\right).`$
Here we used that $`\mathrm{tr}\omega _0\mathrm{ln}\omega _0=_i\nu \left(i\right)\mathrm{ln}\nu \left(i\right)`$ due to
$$\omega _0\mathrm{ln}\omega _0=_{i,k}\mathrm{P}_{𝒜𝒜}(i,k)\mathrm{ln}\mathrm{P}_{𝒜𝒜}(i,k)=_i\nu \left(i\right)\vartheta _0^i\vartheta _0^i\mathrm{ln}\nu \left(i\right)$$
for the orthogonal diagonal decomposition (25) of $`\omega _0`$ into one-dimensional orthoprojectors $`\vartheta _0^i\vartheta _0^i=\mathrm{P}_{𝒜𝒜}(i,i)/\nu \left(i\right)`$, and that $`\mathrm{tr}\rho \mathrm{ln}\rho =_i\nu \left(i\right)\left(\mathrm{ln}\nu \left(i\right)S_{𝒜_i}\left(\varrho _i\right)\right)`$ due to
$$\rho \mathrm{ln}\rho =_i\mathrm{P}_𝒜\left(i\right)\mathrm{ln}\mathrm{P}_𝒜\left(i\right)=_i\nu \left(i\right)\rho _i\left(\mathrm{ln}\nu \left(i\right)+\mathrm{ln}\rho _i\right)$$
for the orthogonal decomposition $`\rho =_i\nu \left(i\right)\mathrm{P}_{𝒜\left(i\right)}`$, where $`\mathrm{P}_{𝒜\left(i\right)}=\mathrm{P}_𝒜\left(i\right)/\nu \left(i\right)=\rho _i`$, $`\nu \left(i\right)=\mathrm{trP}_𝒜\left(i\right)`$, $`\mathrm{P}_𝒜\left(i\right)=_k\mathrm{tr}_{}\mathrm{P}_{𝒜𝒜}(i,k)=\rho \left(i\right)`$.
Thus $`H_𝒜\left(\varrho \right)=H_{𝒜|𝒞}\left(\varrho \right)+H_𝒞\left(\nu \right)=2S_{𝒜|𝒞}\left(\varrho \right)+S_𝒞\left(\nu \right)`$ is positive, and it is bounded by
$`C_𝒜`$ $`=`$ $`\underset{\nu }{sup}{\displaystyle \underset{i}{}}\nu \left(i\right)\left(2\underset{\varrho _i}{sup}S_{𝒜\left(i\right)}\left(\varrho _i\right)\mathrm{ln}\nu \left(i\right)\right)`$
$`=`$ $`\underset{\nu }{inf}{\displaystyle \underset{i}{}}\nu \left(i\right)\left(\mathrm{ln}\nu \left(i\right)2\mathrm{ln}dim_i\right)=\mathrm{ln}dim𝒜.`$
Here we used the fact that the supremum of von Neumann entropies
$$S_{𝒜\left(i\right)}\left(\varrho _i\right)=\underset{i}{}\mathrm{tr}__i\rho _i\mathrm{ln}\rho _i$$
for the simple algebras $`𝒜\left(i\right)=\left(_i\right)`$ with $`dim𝒜\left(i\right)=\left(dim_i\right)^2<\mathrm{}`$ is achieved on the tracial density operators $`\rho _i=\left(dim_i\right)^1I^i\rho _i^{}`$, and the infimum of the relative entropy
$$S(\nu ,\nu ^{})=\underset{i}{}\nu \left(i\right)\left(\mathrm{ln}\nu \left(i\right)\mathrm{ln}\nu ^{}\left(i\right)\right),$$
where $`\nu ^{}\left(i\right)=dim𝒜\left(i\right)/dim𝒜`$, is zero, achieved at $`\nu =\nu ^{}`$.∎
## 5. Quantum Channel and its Q-Capacity
Let $`_0`$ be a Hilbert space describing a quantum input system and $``$ describe its output Hilbert space. A quantum channel is an affine operation sending each input state defined on $`_0`$ to an output state defined on $``$ such that the mixtures of states are preserved. A deterministic quantum channel is given by a linear isometry $`Y:_0`$ with $`Y^{}Y=I^{}`$ ($`I^{}`$ is the identify operator in $`_0`$) such that each input state vector $`\eta _0`$, $`\eta =1`$ is transmitted into an output state vector $`Y\eta `$, $`Y\eta =1`$. The orthogonal mixtures $`\rho _0=_n\mu \left(n\right)\rho _n^{}`$ of the pure input states $`\rho _n^{}=\eta _n^{}\eta _n^{}`$ are sent into the orthogonal mixtures $`\rho =_n\mu \left(n\right)\rho _n`$ of the corresponding pure states $`\rho _n=Y\rho _n^{}Y^{}`$.
A noisy quantum channel sends pure input states $`\varrho _0`$ into mixed ones $`\varrho =\mathrm{\Lambda }^{}\left(\varrho _0\right)`$ given by the dual $`\mathrm{\Lambda }^{}`$ to a normal completely positive unital map $`\mathrm{\Lambda }:𝒜𝒜_0`$,
$$\mathrm{\Lambda }\left(A\right)=\mathrm{tr}__1Y^{}AY,A𝒜$$
where $`Y`$ is a linear operator from $`_0_+`$ to $``$ with $`\mathrm{tr}__+Y^{}Y=I^{}`$, and $`_+`$ is a separable Hilbert space of quantum noise in the channel. Each input mixed state $`\varrho _0`$ on $`𝒜^{}\left(_0\right)`$ is transmitted into an output state $`\varrho =\varrho _0\mathrm{\Lambda }`$ given by the density operator
$$\mathrm{\Lambda }_{}\left(\rho _0\right)=Y\left(\rho _0I^+\right)Y^{}𝒜_{}$$
for each density operator $`\rho _0𝒜_{}^{}`$, where $`I^+`$ is the identity operator in $`_+`$. Without loss of generality we can assume that the input algebra $`𝒜^{}`$ is the smallest decomposable algebra, generated by the range $`\mathrm{\Lambda }\left(𝒜\right)`$ of the given map $`\mathrm{\Lambda }`$.
The input entanglements $`\varkappa :𝒜_{}^{}`$ described as normal CP maps with $`\varkappa \left(I\right)=\varrho _0`$, define the quantum correspondences (q-encodings) of probe systems $`(,\varsigma )`$, $`\varsigma =\varkappa ^{}\left(I\right)`$, to $`(𝒜^{},\varrho _0)`$. As it was proven in the previous section, the most informative is the standard entanglement $`\varkappa =\pi _{}^{}`$, at least in the case of the trivial channel $`\mathrm{\Lambda }=\mathrm{I}`$. This extreme input q-entanglement
$$\pi ^{}\left(A^{}\right)=\rho _0^{1/2}A^{}\rho _0^{1/2}=\pi _{}^{}\left(A^{}\right),A^{}𝒜^{},$$
corresponding to the choice $`(,\varsigma )=(𝒜^{},\varrho _0)`$, defines the following density operator
(30)
$$\omega =\left(\mathrm{I}\mathrm{\Lambda }\right)_{}\left(\omega _q^{}\right),\omega _q^{}=\underset{i}{}\left(\vartheta _0^\iota \vartheta _0^\iota \right)\nu _0\left(i\right)$$
of the input-output compound state $`\varpi _q^{}\mathrm{\Lambda }`$ on $`𝒜^{}𝒜`$. It is given by the central decomposition $`\rho _0=\rho _{0i}\nu _0\left(i\right)`$ of the density operator $`\rho _0𝒜_{}^{}=𝒯\left(_{0i}\right)`$, with the amplitudes $`\vartheta _0^i_{0i}^2`$ defined by $`\stackrel{~}{\vartheta }_0^\iota =\rho _{0i}^{1/2}`$. The other extreme cases of the self-dual input entanglements, the pure c-entanglements corresponding to (21), can be less informative then the d-entanglements, given by the decompositions $`\rho _0=\rho _0\left(n\right)`$ into pure states $`\rho _0\left(n\right)=\eta _n\eta _n^{}\mu \left(n\right)`$. They define the density operators
(31)
$$\omega =\left(\mathrm{I}\mathrm{\Lambda }\right)_{}\left(\omega _d^{}\right),\omega _d^{}=\underset{n}{}\eta _n^{}\eta _n^{}\eta _n\eta _n^{}\mu _0\left(n\right),$$
of the $`𝒜^{}𝒜`$-compound state $`\varpi _d^{}\mathrm{\Lambda }`$, which are known as the Ohya compound states $`\varpi _o^{}\mathrm{\Lambda }`$ in the case
$$\rho _0\left(n\right)=\eta _n^{}\eta _n^{}\lambda _0\left(n\right),\eta _m^{}\eta _n^{}=\delta _n^m,$$
of orthogonality of the density operators $`\rho _0\left(n\right)`$ normalized to the eigen-values $`\lambda _0\left(n\right)`$ of $`\rho _0`$. They are described by the input-output density operators
(32)
$$\omega =\left(\mathrm{I}\mathrm{\Lambda }\right)_{}\left(\omega _o^{}\right),\omega _o^{}=\underset{n}{}\eta _n^{}\eta _n^{}\eta _n^{}\eta _n^{}\lambda _0\left(n\right),$$
coinciding with (30) in the case of Abelian $`𝒜^{}`$. These input-output compound states $`\varpi `$ are achieved by compositions $`\lambda =\pi ^{}\mathrm{\Lambda }`$, describing the entanglements $`\lambda ^{}`$ of the extreme probe system $`(^{},\varsigma _0)=(𝒜^{},\varrho _0)`$ to the output $`(𝒜,\varrho )`$ of the channel.
If $`\mathrm{K}:^{}`$ is a normal completely positive unital map
$$\mathrm{K}\left(B\right)=\mathrm{tr}_{_{}}X^{}BX,B,$$
where $`X`$ is a bounded operator $`_{}𝒢_0𝒢`$ with $`\mathrm{tr}_{_{}}X^{}X=I^{}`$, the compositions $`\varkappa =\pi _{}^{}\mathrm{K}`$, $`\pi _{}=\mathrm{\Lambda }_{}\varkappa `$ are the entanglements of the probe system $`(,\varsigma )`$ to the channel input $`(𝒜^{},\varrho _0)`$ and to the output $`(𝒜,\varrho )`$ via this channel. The state $`\varsigma =\varsigma _0\mathrm{K}`$ is given by
$$\mathrm{K}_{}\left(\sigma _0\right)=X\left(I^{}\sigma _0\right)X^{}_{}$$
for each density operator $`\sigma _0_{}^{}`$, where $`I^{}`$ is the identity operator in $`_{}`$. The resulting entanglement $`\pi _{}=\lambda _{}\mathrm{K}`$ defines the compound state $`\varpi =\varpi _0\left(\mathrm{K}\mathrm{\Lambda }\right)`$ on $`𝒜`$ with
$$\varpi _0\left(B^{}A^{}\right)=\mathrm{tr}\stackrel{~}{B}^{}\pi ^{}\left(A^{}\right)=\mathrm{tr}\upsilon _0^{}\left(B^{}A^{}\right)\upsilon _0.$$
on $`^{}𝒜^{}`$. Here $`\upsilon _0:_0𝒢_0_0`$ is the amplitude operator, uniquely defined by the input compound state $`\varpi _0_{}^{}𝒜_{}^{}`$ up to a unitary operator $`U^{}`$ on $`_0`$, and the effect of the input entanglement $`\varkappa `$ and the output channel $`\mathrm{\Lambda }`$ can be written in terms of the amplitude operator of the state $`\varpi `$ as
$$\upsilon =\left(XY\right)\left(I^{}\upsilon _0I^+\right)U$$
up to a unitary operator $`U`$ in $`=_{}_0_+`$. Thus the density operator $`\omega =\upsilon \upsilon ^{}`$ of the input-output compound state $`\varpi `$ is given by $`\varpi _0\left(\mathrm{K}\mathrm{\Lambda }\right)`$ with the density
(33)
$$\left(\mathrm{K}\mathrm{\Lambda }\right)_{}\left(\omega _0\right)=\left(XY\right)\omega _0\left(XY\right)^{},$$
where $`\omega _0=\upsilon _0\upsilon _0^{}`$.
Let $`𝒦_q`$ be the convex set of normal completely positive maps $`\varkappa :𝒜_{}^{}`$ normalized as $`\mathrm{tr}\varkappa \left(I\right)=1`$, and $`𝒦_q^{}`$ be the convex subset $`\{\varkappa 𝒦_q:\varkappa \left(I\right)=\varrho _0\}`$. Each $`\varkappa 𝒦_q^{}`$ can be decomposed as $`\pi _{}^{}\mathrm{K}`$, where $`\pi _{}^{}=\pi ^{}`$ is the standard entanglement on $`(𝒜^{},\varrho _0)`$, and $`\mathrm{K}`$ is a normal unital CP map $`𝒜^{}`$. Further let $`𝒦_c`$ be the convex set of the maps $`\varkappa `$, dual to the input maps of the form (11), described by the combinations
(34)
$$\varkappa \left(B\right)=\underset{n}{}\varsigma \left(B\right)\rho _0\left(n\right).$$
of the primitive maps $`\varkappa _n:B\varsigma _n\left(B\right)\rho _0\left(n\right)`$, and $`𝒦_d`$ be the subset of the diagonal decompositions
(35)
$$\varkappa \left(B\right)=\underset{n}{}n|B|n\rho _0\left(n\right).$$
As in the first case $`𝒦_c^{}`$ and $`𝒦_d^{}`$ denote the convex subsets corresponding to a fixed $`\varkappa \left(I\right)=\varrho _0`$, and each $`\varkappa 𝒦_c^{}`$ can be represented as $`\pi _{}^{}\mathrm{K}`$, where $`\pi _{}^{}`$ is a d-entanglement, which can be always be made pure by a proper choice of the CP map $`\mathrm{K}:𝒜^{}`$. Furthermore let $`𝒦_o`$ ($`𝒦_o^{}`$) be the subset of all decompositions (34) with orthogonal $`\rho _0\left(n\right)`$ (and fixed $`_n\rho _0\left(n\right)=\rho _0`$):
$$\rho _0\left(m\right)\rho _0\left(n\right)=0,mn.$$
Each $`\varkappa 𝒦_o^{}`$ can be also represented as $`\pi _{}^{}\mathrm{K}`$, where $`\pi _{}^{}`$ is a diagonal pure o-entanglement $`𝒜^{}`$.
Now, let us maximize the entangled mutual information for a given quantum channel $`\mathrm{\Lambda }`$ and a fixed input state $`\varrho _0`$ by means of the above four types of compound states. The mutual information (23) was defined in the previous section by the density operators of the compound state $`\varpi `$ on $`𝒜`$, and the product-state $`\phi =\varsigma \varrho `$ of the marginals $`\varsigma ,\varrho `$ for $`\varpi `$. In each case
$$\varpi =\varpi _0\left(\mathrm{K}\mathrm{\Lambda }\right),\phi =\phi _0\left(\mathrm{K}\mathrm{\Lambda }\right),$$
where $`\mathrm{K}`$ is a CP map $`^{}`$, $`\varpi _0`$ is one of the corresponding extreme compound states $`\varpi _q^{}`$, $`\varpi _c^{}=\varpi _d^{}`$, $`\varpi _o^{}`$ on $`𝒜^{}𝒜^{}`$, and $`\phi _0=\varrho _0\varrho _0`$. The density operator $`\omega =\left(\mathrm{K}\mathrm{\Lambda }\right)_{}\left(\omega _0\right)`$ is written in (33), and $`\varphi =\sigma \rho `$ can be written as
$$\varphi =\varkappa _{}\left(I\right)\lambda _{}\left(I\right),$$
where $`\lambda _{}=\mathrm{\Lambda }_{}\pi _{}^{}`$.
Proposition 5.1. The entangled mutual informations achieve the following maximal values
(36)
$$\underset{\varkappa 𝒦_q^{}}{sup}I_{𝒜,}\left(\varpi \right)=I_q(\varrho _0,\mathrm{\Lambda }):=I_{𝒜,𝒜^{}}\left(\varpi _q^{}\mathrm{\Lambda }\right),$$
$$I_c(\varrho _0,\mathrm{\Lambda })=\underset{\varkappa 𝒦_c^{}}{sup}I_{𝒜,}\left(\varpi \right)=\underset{\varpi _d^{}}{sup}I_{𝒜,𝒜^{}}\left(\varpi _d^{}\mathrm{\Lambda }\right)=I_d(\varrho _0,\mathrm{\Lambda }),$$
(37)
$$\underset{\varkappa 𝒦_o^{}}{sup}I_{𝒜,}\left(\varpi \right)=I_o(\varrho _0,\mathrm{\Lambda }):=\underset{\varpi _o^{}}{sup}I_{𝒜,𝒜^{}}\left(\varpi _o^{}\mathrm{\Lambda }\right),$$
where $`\varpi _{}^{}`$ are the corresponding extremal input entangled states on $`𝒜^{}𝒜^{}`$ with marginals $`\varrho _0`$. They are ordered as
(38)
$$I_q(\varrho _0,\mathrm{\Lambda })I_c(\varrho _0,\mathrm{\Lambda })=I_d(\varrho _0,\mathrm{\Lambda })I_o(\varrho _0,\mathrm{\Lambda }).$$
###### Proof.
Due to the monotonicity
$$I_{𝒜,}\left(\varpi _d^{}\left(\mathrm{K}\mathrm{\Lambda }\right)\right)I_{𝒜,𝒜^{}}\left(\varpi _d^{}\left(\mathrm{I}\mathrm{\Lambda }\right)\right)$$
the supremum over all c-entanglements $`\varkappa 𝒦_c^{}`$ coinsides with the supremum over $`𝒦_d^{}𝒦_c^{}`$ which is achieved on the pure d-entanglements on $`(𝒜^{},\varrho _0)`$ corresponding to the extreme compound states $`\varpi _d^{}`$. By the same monotonicity arguments we can get the equalities (36) and (37). The entanglements $`\varkappa 𝒦_q^{}`$ can be written as
$$\varkappa \left(B\right)=\underset{m,n}{}m|B|n\chi \left(m\right)\chi \left(n\right)^{}$$
in a basis $`\left\{|n\right\}𝒢`$ for the Schatten decompositions $`\sigma =_n|n\mu \left(n\right)n|`$ corresponding to weakly orthogonal amplitude operators $`\chi \left(n\right)=\left(n|XI\right)\left(I^{}\upsilon _0\right):`$
$$\mathrm{tr}\chi \left(m\right)\chi \left(n\right)^{}=\mu \left(n\right)\delta _n^m.$$
The maps $`\varkappa 𝒦_d^{}`$ can be written as
$$\varkappa \left(B\right)=\underset{n}{}n|B|n\chi \left(n\right)\chi \left(n\right)^{}$$
corresponding to stronger orthogonal amplitude operators
$$\chi \left(m\right)\chi \left(n\right)^{}=\rho _0\left(n\right)\delta _n^m,$$
defining not necessarily orthogonal decompositions $`\rho _0=_n\rho _0\left(n\right)`$. The extreme maps $`\varkappa 𝒦_o^{}`$ can be written as
$$\varkappa \left(B\right)=\underset{n}{}n|B|n\chi \left(n\right)\chi \left(n\right)^{}$$
with amplitude operators $`\chi \left(n\right)`$, satisfying the second orthogonality condition
$$\chi \left(n\right)^{}\chi \left(m\right)=\mu \left(n\right)\tau _n^{}\delta _n^m,$$
where $`\tau _n^{}`$ are density operators in $`_0`$ with the traces $`\mathrm{tr}\tau _n^{}=1`$. Thus, the inequalities in (38) follow from $`𝒦_q𝒦_c𝒦_d𝒦_o`$.∎
We shall denote the maximal informations $`I_c(\varrho _0,\mathrm{\Lambda })=I_d(\varrho _0,\mathrm{\Lambda })`$ simply as $`I(\varrho _0,\mathrm{\Lambda })`$.
Definition 5.1. The supremums
$$C_q\left(\mathrm{\Lambda }\right)=\underset{\varkappa 𝒦_q}{sup}I_{𝒜,}\left(\varpi \right)=\underset{\varrho _0}{sup}I_q(\varrho _0,\mathrm{\Lambda }),$$
(39)
$$\underset{\varkappa 𝒦_c}{sup}I_{𝒜,}\left(\varpi \right)=C\left(\mathrm{\Lambda }\right):=\underset{\varrho _0}{sup}I(\varrho _0,\mathrm{\Lambda }),$$
$$C_o\left(\mathrm{\Lambda }\right)=\underset{\varkappa 𝒦_o}{sup}I_{𝒜,}\left(\varpi \right)=\underset{\varrho _0}{sup}I_o(\varrho _0,\mathrm{\Lambda }),$$
are called the q-, c- or d-, and o-capacities respectively for the quantum channel defined by a normal unital CP map $`\mathrm{\Lambda }:𝒜𝒜^{}`$.
Obviously the capacities (39) satisfy the inequalities
$$C_o\left(\mathrm{\Lambda }\right)C\left(\mathrm{\Lambda }\right)C_q\left(\mathrm{\Lambda }\right).$$
Theorem 5.2. Let $`\mathrm{\Lambda }\left(A\right)=Y^{}AY`$ be a unital CP map $`𝒜𝒜^{}`$ describing a quantum deterministic channel. Then
$$I(\varrho _0,\mathrm{\Lambda })=I_o(\varrho _0,\mathrm{\Lambda })=S\left(\varrho _0\right),I_q(\varrho _0,\mathrm{\Lambda })=S_q\left(\varrho _0\right),$$
where $`S_q\left(\varrho _0\right)=H_𝒜^{}\left(\varrho _0\right)`$, and thus in this case
$$C\left(\mathrm{\Lambda }\right)=C_o\left(\mathrm{\Lambda }\right)=\mathrm{ln}\mathrm{rank}𝒜^{},C_q\left(\mathrm{\Lambda }\right)=\mathrm{ln}dim𝒜^{}.$$
###### Proof.
It was proved in the previous section for the case of the identity channel $`\mathrm{\Lambda }=\mathrm{I}`$, and thus it is also valied for any isomorphism $`\mathrm{\Lambda }`$ described by a unitary operator $`Y`$. In the case of non-unitary $`Y`$ we can use the identity
$$\mathrm{tr}Y\left(\rho _0I^+\right)Y^{}\mathrm{ln}Y\left(\rho _0I^+\right)Y^{}=\mathrm{tr}R\left(\omega _0I^+\right)\mathrm{ln}R\left(\omega _0I^+\right),$$
where $`R=Y^{}Y`$. Due to this $`S\left(\varrho _0\mathrm{\Lambda }\right)=\mathrm{tr}R\left(\rho _0I^+\right)\mathrm{ln}R\left(\rho _0I^+\right)`$, and $`S\left(\varpi _0\left(\mathrm{K}\mathrm{\Lambda }\right)\right)=`$
$$\mathrm{tr}\left(SR\right)\left(I^{}\omega _0I^+\right)\mathrm{ln}\left(SR\right)\left(I^{}\omega _0I^+\right),$$
where $`S=X^{}X`$. Thus $`S\left(\varrho _0\mathrm{\Lambda }\right)=S\left(\varrho _0\right)`$, $`S\left(\varpi _0\left(\mathrm{K}\mathrm{\Lambda }\right)\right)=S\left(\varpi _0\left(\mathrm{K}\mathrm{I}\right)\right)`$ if $`Y^{}Y=I`$, and
$`I_{𝒜,}\left(\varpi _0\left(\mathrm{K}\mathrm{\Lambda }\right)\right)`$ $`=`$ $`S\left(\varsigma _0\mathrm{K}\right)+S\left(\varrho _0\right)S\left(\varpi _0\left(\mathrm{K}\mathrm{I}\right)\right)`$
$``$ $`S\left(\varsigma _0\right)+S\left(\varrho _0\right)S\left(\varpi _0\right)=I_{𝒜^{},^{}}\left(\varpi _0\right)`$
for any normal unital CP map $`\mathrm{K}:^{}`$ and a compound state $`\varpi _0`$ on $`^{}𝒜^{}`$. The supremum (36), which is achieved at the standard entanglement, corresponding to $`\varpi _0=\varpi _q`$, coincides with q-entropy $`H_𝒜^{}\left(\varrho _0\right)`$, and the supremum (37), coinciding with $`S_𝒜^{}\left(\varrho _0\right)`$, is achieved for a pure o-entanglement, corresponding to $`\varpi _0=\varpi _o`$ given by any Schatten decomposition for $`\rho _0`$. Moreover, the entropy $`H_𝒜^{}\left(\varrho _0\right)`$ is also achieved by any pure d-entanglement, corresponding to $`\varpi _0=\varpi _d`$ given by any extreme decomposition for $`\rho _0`$, and thus is the maximal mutual information $`I(\varrho _0,\mathrm{\Lambda })`$ in the case of deterministic $`\mathrm{\Lambda }`$. Thus the capacity $`C\left(\mathrm{\Lambda }\right)`$ of the deterministic channel is given by the maximum $`C_o=\mathrm{ln}dim_0`$ of the von Neumann entropy $`S_𝒜^{}`$, and the q-capacity $`C_q\left(\mathrm{\Lambda }\right)`$ is equal $`C_𝒜^{}=\mathrm{ln}dim𝒜^{}`$.∎
In the general case d-entanglements can be more informative than o-entanglements as it can be shown on an example of a quantum noisy channel for which
$$I(\varrho _0,\mathrm{\Lambda })>I_o(\varrho _0,\mathrm{\Lambda }),C\left(\mathrm{\Lambda }\right)>C_o\left(\mathrm{\Lambda }\right).$$
The last equalities of the above theorem will be related to the work on entropy by Voiculescu .
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# Untitled Document
Schubert Calculus and Threshold Polynomials of Affine Fusion
S.E. Irvine supported in part by a NSERC Undergraduate Research Award M.A. Walton supported in part by NSERC. E-mail: walton@uleth.ca
Physics Department, University of Lethbridge
Lethbridge, Alberta, Canada T1K 3M4
Abstract
We show how the threshold level of affine fusion, the fusion of Wess-Zumino-Witten (WZW) conformal field theories, fits into the Schubert calculus introduced by Gepner. The Pieri rule can be modified in a simple way to include the threshold level, so that calculations may be done for all (non-negative integer) levels at once. With the usual Giambelli formula, the modified Pieri formula deforms the tensor product coefficients (and the fusion coefficients) into what we call threshold polynomials. We compare them with the $`q`$-deformed tensor product coefficients and fusion coefficients that are related to $`q`$-deformed weight multiplicities. We also discuss the meaning of the threshold level in the context of paths on graphs.
1. Introduction
Gepner found geometrical and topological interpretations of the fusion rings of Wess-Zumino-Witten (WZW) conformal field theories . He described them using a Schubert calculus, a “quantum version” of the classical Schubert calculus that is fundamental in the geometry and topology of complex manifolds (see , e.g.).
Gepner also pointed out a correspondence between the WZW fusion rings and the chiral rings of $`N=2`$ superconformal theories. These two observations have been seminal. For example, their relation was clarified in , where the new Schubert calculus was shown to describe the quantum cohomology of Grassmannians. Also, the $`N=2`$ interpretation led to new realisations of WZW fusion rings in topological theories .
We study the Schubert calculus of WZW fusion rings. Our initial motivation was computational. In Gepner’s approach, a fusion potential is introduced whose derivatives give the fusion constraints to be implemented. The fusion potential, and so the fusion constraints, are level dependent. Therefore a significant part of any computation must be re-done whenever the level is changed. By the depth rule , however, the results are simpler than this procedure indicates. A threshold level exists for each coupling ; in any fusion product of two fixed fields, a third appears in the decomposition for all integer levels greater than or equal to a characteristic one To the best of our knowledge, however, a completely rigorous demonstration of the existence of threshold levels is still lacking.. Therefore, finding the threshold levels for a fixed product amounts to finding the fusion rules for all levels at once.
We show how to incorporate the notion of threshold level into Gepner’s Schubert calculus for WZW fusion. This is done in section 3, after a review is given in section 2, where the notation is also established.
Another motivation for this work emerges in section 2: it is convenient to encode the threshold levels in generating polynomials, dubbed threshold polynomials. These are then polynomial deformations of tensor product coefficients and fusion coefficients. Similar objects, the quantum group ($`q`$-)deformations of tensor product coefficients and fusion coefficients have been studied previously. Most importantly to us, the $`q`$-deformed coefficients are related to the $`q`$-deformed weight multiplicities defined by Lusztig . In section 4 we compare our deformations with the $`q`$-deformations. We show that the new deformations are related in a similar way to deformations of the weight multiplicities, that are natural from the point of view of a conjectured refinement of the Gepner-Witten depth rule .
As we have argued, the threshold level has computational advantages over the use of fusion potentials, and the relations derived from them. But again, the connection with geometry, topology and $`N=2`$ superconformal theories was the point of , not computation. There the fusion potential played a central role. But one can’t have it both ways: we indicate at the end of section 3 that a deformed fusion potential that incorporates the threshold levels cannot be written. Nevertheless, one might hope to give the threshold level a somewhat deeper motivation, perhaps through its meaning in the many different realisations of WZW fusion rings. In section 5 we make a very small start in this direction; we discuss the meaning of the threshold level in the context of paths on graphs (see , and references therein).
Section 6 is a short conclusion.
2. WZW fusion, threshold level, and threshold polynomials
Let us first establish notation. For the most part, we restrict attention to the simple Lie algebras $`A_r`$ and the affine algebras $`A_r^{(1)}`$ that are the untwisted central extensions of their loop algebras. When the level $`k`$ is fixed, we denote the affine algebra by $`A_{r,k}`$. However, we use a notation that is easily adapted to any untwisted affine Kac-Moody algebra $`X_r^{(1)}`$ (or $`X_{r,k}`$) based on a simple Lie algebra $`X_r`$, and expect that such generalisation is straightforward.
The set of roots of $`X_r`$ will be written as $`R`$, and the set of positive (negative) roots as $`R_>`$ ($`R_<`$). If $`\alpha R`$ is a root, then the corresponding co-root is defined as $`\alpha ^{}:=2\alpha /(\alpha ,\alpha )`$.
Let $`F=\{\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_r\}`$ denote the set of fundamental weights of $`X_r`$, and
$$P:=\{\lambda =\underset{i=1}{\overset{r}{}}\lambda _i\mathrm{\Lambda }_i|\lambda _iZ\}$$
the set of integral weights. The set of dominant integral weights,
$$P_{}:=\{\lambda =\underset{i=1}{\overset{r}{}}\lambda _i\mathrm{\Lambda }_i|\lambda _iZ_{}\},$$
is the set of highest weights for irreducible integrable modules of $`X_r`$.
Let $`M(\lambda )`$ denote an irreducible module of $`X_r`$, of highest weight $`\lambda P_{}`$. The set of weights of $`M(\lambda )`$ will be denoted $`P(\lambda )`$.
The irreducible integrable modules of $`X_{r,k}`$ have highest weights that project to the following set of dominant weights of $`X_r`$:
$$P_{}^k:=\{\lambda =\underset{i=1}{\overset{r}{}}\lambda _i\mathrm{\Lambda }_i|\lambda _iZ_{},\underset{i=1}{\overset{r}{}}\lambda _ia_i^{}k\}.$$
The $`a_i^{}`$ are the co-marks, defined by $`a_0^{}=1`$, and
$$\theta ^{}=\theta =\underset{i=1}{\overset{r}{}}a_i^{}\alpha _i^{},$$
where $`\theta `$ ($`\theta ^{}`$) denotes the highest (co-)root of $`X_r`$. We normalise $`(\theta ,\theta )=2`$.
The Weyl group of $`X_r`$ will be denoted by $`W`$, and the shifted action of $`wW`$ on a weight $`\lambda `$ by $`w.\lambda =w(\lambda +\rho )\rho `$, where $`\rho =_{i=1}^r\mathrm{\Lambda }_i=_{\alpha R_>}\alpha /2`$ is the Weyl vector. $`W^k`$ will indicate the projection of the affine Weyl group, the Weyl group of $`X_{r,k}`$, onto the horizontal weight space, the weight space of $`X_r`$. $`W`$ is generated by the primitive reflections $`r_i`$, $`i=1,\mathrm{},r`$, with action
$$r_i\lambda =\lambda (\lambda ,\alpha _i^{})\alpha _i$$
on any weight $`\lambda `$. In order to enlarge $`W`$ to $`W^k`$, we adjoin $`r_0`$ to the generating set. Its shifted action is
$$r_0.\lambda =r_\theta .\lambda +(k+x)\theta ,$$
where $`x`$ is the dual Coxeter number of $`X_r`$. Notice the $`k`$-dependence of the action of $`W^k`$ on $`P`$, coming from that of $`r_0`$.
We write the decomposition of the tensor product of two irreducible integrable $`X_r`$-modules as
$$M(\lambda )M(\mu )=\underset{\nu P_{}}{}T_{\lambda ,\mu }^\nu M(\nu ).$$
We will call the $`T_{\lambda ,\mu }^\nu Z_{}`$ tensor product coefficients. We indicate the affine fusion of two modules of $`X_{r,k}`$ by writing the truncated tensor product of the corresponding modules $`M(\lambda )`$ and $`M(\mu )`$ of $`X_r`$:
$$M(\lambda )_kM(\mu )=\underset{\nu P_{}^k}{}{}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }M(\nu ).$$
The fusion coefficients obey
$${}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }{}_{}{}^{(k+1)}T_{\lambda ,\mu }^{\nu },$$
and furthermore
$$\underset{k\mathrm{}}{lim}{}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }=T_{\lambda ,\mu }^\nu .$$
We can encode the fusion products for all levels by including the threshold levels $`t`$ as subscripts in the tensor product decomposition . If we denote by $`𝒮_t`$ the operator that includes these subscripts, we can write
$$𝒮_t\left[M(\lambda )M(\mu )\right]=\underset{\nu P_{}}{}\underset{tZ_{}}{}T_{\lambda ,\mu }^{\nu (t)}M(\nu )_t.$$
Then
$${}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }=\underset{t=0}{\overset{k}{}}T_{\lambda ,\mu }^{\nu (t)}.$$
We call the fixed-threshold-level coefficients $`T_{\lambda ,\mu }^{\nu (t)}`$, the threshold coefficients. For example, we can modify the $`A_2`$ tensor product to
$$𝒮_t\left[M(1,1)^2\right]=M(2,2)_4M(3,0)_3M(0,3)_32M(1,1)_{2,3}M(0,0)_2,$$
encoding the corresponding fusions at all levels. Here $`M(a,b):=M(\lambda )`$, with $`\lambda =a\mathrm{\Lambda }_1+b\mathrm{\Lambda }_2`$, and $`pM(a,b)_{t_1,\mathrm{},t_p}:=M(a,b)_{t_1}\mathrm{}M(a,b)_{t_p}`$.
From the notational point of view, the threshold levels are unnecessarily large numbers, since
$$T_{\lambda ,\mu }^{\nu (t)}0t(\nu ,\theta ).$$
Consequently, one could also define the threshold delay $`d`$ by
$$d:=t(\nu ,\theta ).$$
Writing the delays as superscripts, the right-hand side of (2.1) is replaced by
$$M(2,2)^2M(3,0)^0M(0,3)^02M(1,1)^{0,1}M(0,0)^2.$$
This is a minor point, so we’ll stick to using the threshold levels. In section 4, however, (2.1) will reappear.
As (2.1) makes clear, we need to consider $`N`$-tuples of threshold levels. It is convenient to encode them in threshold polynomials, defined by
$$T_{\lambda ,\mu }^\nu [\mathrm{}]:=\underset{tZ_{}}{}\mathrm{}^tT_{\lambda ,\mu }^{\nu (t)}.$$
Then
$$^t\left[M(\lambda )M(\mu )\right]=\underset{\nu P_{}}{}T_{\lambda ,\mu }^\nu [\mathrm{}]M(\nu )$$
is equivalent to (2.1). For example, the $`A_2`$ tensor product (2.1) is rewritten as
$$^t\left[M(1,1)^2\right]=\mathrm{}^4M(2,2)\mathrm{}^3M(3,0)\mathrm{}^3M(0,3)(\mathrm{}^2+\mathrm{}^3)M(1,1)\mathrm{}^2M(0,0).$$
The threshold polynomials can be regarded as deformations of the tensor product coefficients, since
$$T_{\lambda ,\mu }^\nu [1]=T_{\lambda ,\mu }^\nu ,$$
so that the tensor product (2.1) is recovered when $`\mathrm{}=1`$. Furthermore, we define the deformation of the fusion coefficient as
$${}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }[\mathrm{}]:=\underset{t=0}{\overset{k}{}}\mathrm{}^tT_{\lambda ,\mu }^{\nu (t)}.$$
So $`{}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }[\mathrm{}]`$ is the degree$`k`$ part of the polynomial $`T_{\lambda ,\mu }^\nu [\mathrm{}]`$, and
$${}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }[1]={}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }.$$
(2.1) is deformed to
$$\underset{k\mathrm{}}{lim}{}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }[\mathrm{}]=T_{\lambda ,\mu }^\nu [\mathrm{}],$$
by construction.
From (2.1), we see that the threshold polynomials are the generating functions for the threshold coefficients. Consequently, they are related to the generating functions for fusion rules studied in , where the threshold level was first introduced (but named later in ). For completeness, we indicate the relation here.
The generating function for fusion rules is defined as
$$G(L,M,N;d):=\underset{\lambda ,\mu ,\nu P_{}}{}\underset{k=0}{\overset{\mathrm{}}{}}{}_{}{}^{(k)}T_{\lambda ,\mu ,\nu }^{}d^kL^\lambda M^\mu N^\nu ,$$
where the dummy variables $`L,M,N`$ satisfy $`L^\lambda L^\lambda ^{}=L^{\lambda +\lambda ^{}}`$, etc. Here $`{}_{}{}^{(k)}T_{\lambda ,\mu ,\nu }^{}:={}_{}{}^{(k)}T_{\lambda ,\mu }^{C\nu }`$, where $`C\nu `$ is the highest weight of the module conjugate to $`M(\nu )`$. Using (2.1) and switching the order of summations, we arrive at
$$G(L,M,N;d)=(1d)^1\underset{\lambda ,\mu ,\nu P_{}}{}L^\lambda M^\mu N^\nu T_{\lambda ,\mu ,\nu }[d].$$
Here $`T_{\lambda ,\mu ,\nu }[d]=_{t=0}^{\mathrm{}}d^tT_{\lambda ,\mu }^{C\nu }[d]`$; see (2.1) and (2.1). Hence the only difference between the generating functions for deformed tensor product coefficients and fusion rules is $`(1d)^1`$, a factor characteristic of the existence of a threshold level .
For successive fusions, we need a memory of the threshold levels. For example, suppose we need to calculate $`𝒮_t\left[M(\varphi )(M(1,1)^2)\right]`$. Then using (2.1), we would encounter products like $`𝒮_t\left[M(\varphi )M(1,1)_3\right]`$. So (2.1) is only a special case of what we need. Abusing notation slightly, we attach threshold levels to the factor modules in the tensor products, and write
$$M(\lambda )_rM(\mu )_s=\underset{\nu P_{}}{}\underset{tZ_{}}{}T_{\lambda (r),\mu (s)}^{\nu (t)}M(\nu )_t.$$
For example, we have
$$M(1,1)_2M(1,1)_3=M(2,2)_4M(3,0)_3M(0,3)_32M(1,1)_{3,3}M(0,0)_3.$$
(2.1) is recovered from (2.1) by setting $`r=(\lambda ,\theta )`$ and $`s=(\mu ,\theta )`$.
Again using polynomials to carry the $`N`$-tuples of threshold levels, we write
$$^t\left[M(\lambda )_rM(\mu )_s\right]=\underset{\nu P_{}}{}T_{\lambda (r),\mu (s)}^\nu [\mathrm{}]M(\nu ).$$
Comparing (2.1) with (2.1), for example, shows that there is a simple relation between the coefficients $`T_{\lambda (r),\mu (s)}^{\nu (t)}`$ and $`T_{\lambda ,\mu }^{\nu (t)}`$. To write it in polynomial form, we define
$$\mathrm{}^a\mathrm{}^b:=\mathrm{}^{\mathrm{max}\{a,b\}}.$$
and extend bilinearly (so that two polynomials can be multiplied). Then we have
$$T_{\lambda (r),\mu (s)}^\nu [\mathrm{}]=\mathrm{}^r\mathrm{}^sT_{\lambda ,\mu }^\nu [\mathrm{}].$$
We see that the polynomials $`T_{\lambda ,\mu }^\nu [\mathrm{}]`$ are fundamental, and so we will concentrate on them henceforth.
The definition (2.1), however, is natural from the point of view of threshold polynomials. With it, we can generalise the crossing symmetry of the tensor product coefficients,
$$\underset{\zeta P_{}}{}T_{\lambda ,\mu }^\zeta T_{\zeta ,\phi }^\nu =\underset{\zeta P_{}}{}T_{\lambda ,\zeta }^\nu T_{\mu ,\phi }^\zeta ,$$
to
$$\underset{\zeta P_{}}{}T_{\lambda ,\mu }^\zeta [\mathrm{}]T_{\zeta ,\phi }^\nu [\mathrm{}]=\underset{\zeta P_{}}{}T_{\lambda ,\zeta }^\nu [\mathrm{}]T_{\mu ,\phi }^\zeta [\mathrm{}].$$
Furthermore, the crossing symmetry for the fusion coefficients
$$\underset{\zeta P_{}^k}{}{}_{}{}^{(k)}T_{\lambda ,\mu }^{\zeta }{}_{}{}^{(k)}T_{\zeta ,\phi }^{\nu }=\underset{\zeta P_{}^k}{}{}_{}{}^{(k)}T_{\lambda ,\zeta }^{\nu }{}_{}{}^{(k)}T_{\mu ,\phi }^{\zeta }$$
deforms to
$$\underset{\zeta P_{}^k}{}{}_{}{}^{(k)}T_{\lambda ,\mu }^{\zeta }[\mathrm{}]{}_{}{}^{(k)}T_{\zeta ,\phi }^{\nu }[\mathrm{}]=\underset{\zeta P_{}^k}{}{}_{}{}^{(k)}T_{\lambda ,\zeta }^{\nu }[\mathrm{}]{}_{}{}^{(k)}T_{\mu ,\phi }^{\zeta }[\mathrm{}].$$
3. Schubert calculus, threshold level, and threshold polynomials
The Schubert calculus is based on the Pieri and Giambelli formulas. For discussions of them emphasising the geometric context see . More relevant to us is their use in representation theory; consult , e.g.
The Pieri formula is
$$T_{\lambda ,\mathrm{\Lambda }}^\nu =\{\begin{array}{cc}1,& \mathrm{if}\nu \lambda P(\mathrm{\Lambda });\\ 0,& \mathrm{otherwise},\end{array}$$
where $`\mathrm{\Lambda }`$ is a fundamental weight, i.e. $`\mathrm{\Lambda }F`$. Here we are specialising to the algebras $`A_r`$, although the formulas for other algebras are only slightly more complicated.
Adapted to include threshold polynomials, the Pieri formula simply becomes
$$T_{\lambda ,\mathrm{\Lambda }}^\nu [\mathrm{}]=\{\begin{array}{cc}\mathrm{}^{(\lambda ,\theta )}\mathrm{}^{(\nu ,\theta )},& \mathrm{if}\nu \lambda P(\mathrm{\Lambda });\\ 0,& \mathrm{otherwise}.\end{array}$$
Fundamental monomials
$$M(\mathrm{\Lambda }^\mu ):=M(\mathrm{\Lambda }_1)^{\mu _1}M(\mathrm{\Lambda }_2)^{\mu _2}\mathrm{}M(\mathrm{\Lambda }_r)^{\mu _r}=\underset{i=1}{\overset{r}{}}M(\mathrm{\Lambda }_i)^{\mu _i},$$
with all $`\mathrm{\Lambda }_iF`$, are easily decomposed using the Pieri formula (3.1). The decompositions are triangular in the irreducible highest-weight modules $`M(\sigma )`$, $`\sigma P_{}`$:
$$M(\mathrm{\Lambda }^\lambda )=\underset{P_{}\sigma \lambda }{}\mathrm{\Omega }_{\lambda ,\sigma }M(\sigma ).$$
Here $`\sigma \lambda `$ means that $`\lambda \sigma `$ is a non-negative integer linear combination of positive roots, i.e. $`\lambda \sigma Z_{}R_>`$. So $`\mathrm{\Omega }_{\lambda ,\sigma }`$ is a triangular matrix.
The polynomial deformation of this, encoding the threshold levels, is just
$$^t\left[M(\mathrm{\Lambda }^\lambda )\right]=\underset{P_{}\sigma \lambda }{}\mathrm{\Omega }_{\lambda ,\sigma }[\mathrm{}]M(\sigma ).$$
Some $`A_2`$ examples (to be used shortly) will make this clear. We find:
$$𝒮_t\left[M(\mathrm{\Lambda }^{(2,2)})\right]=M(2,2)_4M(3,0)_34M(1,1)_{2,2,2,3}M(0,3)_32M(0,0)_{1,2},$$
or
$$\begin{array}{cc}\hfill ^t\left[M(\mathrm{\Lambda }^{(2,2)})\right]=\mathrm{}^4M(2,2)& \mathrm{}^3M(3,0)(\mathrm{}^3+3\mathrm{}^2)M(1,1)\hfill \\ \hfill & \mathrm{}^3M(0,3)(\mathrm{}+\mathrm{}^2)M(0,0);\hfill \end{array}$$
and
$$𝒮_t\left[M(\mathrm{\Lambda }^{(1,1)})\right]=M(1,1)_2M(0,0)_1,𝒮_t\left[M(\mathrm{\Lambda }^{(0,0)})\right]=M(0,0)_0,$$
or
$$^t\left[M(\mathrm{\Lambda }^{(1,1)})\right]=\mathrm{}^2M(1,1)\mathrm{}^1M(0,0),^t\left[M(\mathrm{\Lambda }^{(0,0)})\right]=M(0,0).$$
From (3.1), the highest-weight modules can be expressed as polynomials in the fundamental weights:
$$M(\sigma )=\underset{P_{}\mu \sigma }{}(\mathrm{\Omega }^1)_{\sigma ,\mu }M(\mathrm{\Lambda }^\mu ).$$
That is, $`M(\sigma )`$ can be written as a direct sum of fundamental monomials $`M(\mathrm{\Lambda }^\mu )`$. This is the Giambelli formula, in the non-determinantal form that can be applied to all simple Lie algebras, not just $`A_r`$. Notice that the inversion of $`\mathrm{\Omega }`$ is greatly simplified by its triangularity, and its inverse is also triangular. A simple $`A_2`$ example of (3.1) is
$$M(1,1)=M(\mathrm{\Lambda }^{(1,1)})M(\mathrm{\Lambda }^{(0,0)}).$$
The characters of $`X_r`$ form an algebra with structure constants equal to the tensor product coefficients. The Giambelli formula (3.1) gives rise to a polynomial realisation of this character algebra. One simply replaces $`M(\mathrm{\Lambda }^\mu )`$ with $`_{i=1}^rx_i^{\mu _i}=:x^\mu `$. The resulting polynomial
$$S_\sigma (x):=\underset{P_{}\mu \sigma }{}(\mathrm{\Omega }^1)_{\sigma ,\mu }x^\mu ,$$
is known as a Schur polynomial, a type of Schubert polynomial . For example, $`x_1x_21`$ is the Schur polynomial of the $`A_2`$ module $`M(1,1)`$, by (3.1). With the addition and subtraction of polynomials, the character algebra extends to a ring.
Is there a useful threshold-level version of the Giambelli formula? The inverse $`\mathrm{\Omega }^1[\mathrm{}]`$ of the matrix $`\mathrm{\Omega }[\mathrm{}]`$ in (3.1) has entries that are negative powers of $`\mathrm{}`$. These are difficult to interpret in the context of threshold level. We conclude that the normal, $`\mathrm{}`$-independent matrix $`\mathrm{\Omega }^1`$ should be used. We can write useful formulas for the threshold polynomials in terms of $`\mathrm{\Omega }^1`$, and its deformed inverse $`\mathrm{\Omega }[\mathrm{}]`$:
$$T_{\lambda ,\mu }^\nu [\mathrm{}]=\mathrm{}^{(\lambda ,\theta )}\mathrm{}^{(\mu ,\theta )}\underset{\alpha ,\beta P_{}}{}(\mathrm{\Omega }^1)_{\lambda ,\alpha }(\mathrm{\Omega }^1)_{\mu ,\beta }(\mathrm{\Omega })_{\alpha +\beta ,\nu }[\mathrm{}].$$
We’ll illustrate this formula on the $`A_2`$ example with $`M(\lambda )=M(\mu )=M(1,1)`$, using the subscript notation. First, the required matrix elements of $`\mathrm{\Omega }^1`$ are provided by
$$\begin{array}{cc}\hfill M(1,1)^2=& \left[M(\mathrm{\Lambda }^{(1,1)})M(\mathrm{\Lambda }^{(0,0)})\right]^2\hfill \\ \hfill =& M(\mathrm{\Lambda }^{(2,2)})2M(\mathrm{\Lambda }^{(1,1)})M(\mathrm{\Lambda }^{(0,0)}).\hfill \end{array}$$
Substituting the fundamental monomials (3.1) and (3.1), described by $`\mathrm{\Omega }[\mathrm{}]`$, we get
$$\begin{array}{cc}\hfill ^t\left[M(1,1)^2\right]=& \mathrm{}^2\{\mathrm{}^4M(2,2)\mathrm{}^3M(3,0)(3\mathrm{}^2+\mathrm{}^3)M(1,1)\mathrm{}^3M(0,3)\hfill \\ & (\mathrm{}+\mathrm{}^2)M(0,0)2[\mathrm{}^2M(1,1)\mathrm{}M(0,0)]M(0,0)\}\hfill \\ \hfill =& \mathrm{}^4M(2,2)\mathrm{}^3M(3,0)\mathrm{}^3M(0,3)(\mathrm{}^2+\mathrm{}^3)M(1,1)\mathrm{}^2M(0,0),\hfill \end{array}$$
the correct result.
The deformed Pieri formula (3.1) makes straightforward the calculation of decompositions involving fundamental monomials, like $`M(\mathrm{\Lambda }^\beta )`$. We write
$$^t\left[M(\lambda )M(\mathrm{\Lambda }^\beta )\right]=\underset{\nu P_{}}{}T_{\lambda ,\mathrm{\Lambda }^\beta }^\nu [\mathrm{}]M(\nu ).$$
Then the threshold polynomials may also be calculated from the simpler polynomials $`T_{\lambda ,\mathrm{\Lambda }^\beta }^\nu [\mathrm{}]`$:
$$T_{\lambda ,\mu }^\nu [\mathrm{}]=\mathrm{}^{(\mu ,\theta )}\underset{\beta P_{}}{}(\mathrm{\Omega }^1)_{\mu ,\beta }T_{\lambda ,\mathrm{\Lambda }^\beta }^\nu [\mathrm{}].$$
Using (3.1), an $`A_2`$ example is
$$\begin{array}{cc}\hfill ^t\left[M(1,1)^2\right]=& \mathrm{}^2^t\left[M(1,1)\left(M(\mathrm{\Lambda }^{(1,1)})M(\mathrm{\Lambda }^{(0,0)})\right)\right]\hfill \\ \hfill =& \mathrm{}^4M(2,2)\mathrm{}^3M(3,0)(2\mathrm{}^2+\mathrm{}^3)M(1,1)\mathrm{}^3M(0,3)\mathrm{}^2M(1,1)\hfill \\ & \mathrm{}^2M(1,1).\hfill \end{array}$$
This is again the correct result (see (3.1)).
Finally, we can also multiply two Schur polynomials for $`M(\lambda )`$ and $`M(\mu `$) together, using the coefficients $`T_{\mathrm{\Lambda }^\alpha ,\mathrm{\Lambda }^\beta }^\nu `$, defined in the obvious way:
$$T_{\lambda ,\mu }^\nu [\mathrm{}]=\mathrm{}^{(\lambda ,\theta )}\mathrm{}^{(\mu ,\theta )}\underset{\alpha ,\beta P_{}}{}(\mathrm{\Omega }^1)_{\lambda ,\alpha }(\mathrm{\Omega }^1)_{\mu ,\beta }T_{\mathrm{\Lambda }^\alpha ,\mathrm{\Lambda }^\beta }^\nu [\mathrm{}].$$
To conclude this section, we note that in our deformed Schubert calculus, there is no analogue of the fusion potential that was so important in . We argue that a deformed potential that incorporates the threshold levels cannot be written. Gepner could write a fusion potential because at fixed level $`k`$, the fusion rules are truncations of the tensor product rules. The truncated parts can be set to zero by fusion constraints, that can be derived from the potential. On the other hand, when the threshold level is incorporated into a tensor product, as in (3.1) vs. (3.1), there is no truncation. Instead of constraints, one could only hope to find replacements that would change the right-hand side of (3.1) into that of (3.1), for example. But that is exactly what we do: (3.1) is obtained from (3.1) by replacing $`\mathrm{\Omega }`$ with $`\mathrm{\Omega }[\mathrm{}]`$. A minimal set of such replacements would be those obtained by replacing the right-hand side of the undeformed Pieri rule (3.1) with that of the deformed one (3.1).
Incidentally, we have seen that the Pieri rule with threshold level (3.1) contains the same information as the fusion potentials of Gepner, for all (non-negative integer) levels. So does the generating function for the fusion potentials . It might be interesting to make this more precise.
4. Deformed tensor product coefficients and weight multiplicities
The threshold polynomials (2.1) and (2.1) are deformations of the tensor product coefficients and affine (WZW) fusion coefficients, respectively. It is interesting to compare them with the $`q`$-deformations of these objects studied previously.
WZW fusion coefficients are alternating affine-Weyl ($`W^k`$) sums of tensor product coefficients . In (see also ), the corresponding $`q`$-fusion coefficients (for affine $`A_r`$) are defined in similar fashion in terms of the $`q`$-tensor product coefficients . Since the ordinary (undeformed) tensor product coefficients are also alternating Weyl sums of the weight multiplicities of Lie algebras, the fusion coefficients can also be expressed in that way. In the $`q`$-deformed case, the tensor product and fusion coefficients are related to Lusztig’s $`q`$-deformed weight multiplicities , in turn related to the famous Kazhdan-Lusztig polynomials .
Let us start with an example, taken from . They find, for the $`q`$-deformation of the $`A_3`$ tensor product $`M(1,1,0)^3`$, the following decomposition:
$$\begin{array}{cc}\hfill (q^3+& q^6)M(0,0,3)(2q^4+3q^5+2q^6+q^7)M(0,1,1)(q^2+2q^3+q^4)M(0,3,1)\hfill \\ & (q^5+2q^6+q^7)M(1,0,0)(q^2+2q^3+3q^4+2q^5)M(1,1,2)\hfill \\ & (2q^3+3q^4+3q^5+q^6)M(1,2,0)(q+q^2)M(1,4,0)\hfill \\ & (q^3+3q^4+3q^5+2q^6)M(2,0,1)(q+2q^2+2q^3+q^4)M(2,2,1)\hfill \\ & (q^2+2q^3+q^4)M(3,0,2)(q^2+2q^3+2q^4+q^5)M(3,1,0)\hfill \\ & M(3,3,0)(q+q^2)M(4,1,1)(q^3)M(5,0,0).\hfill \end{array}$$
This should be compared with the threshold-level version of the same tensor product:
$$\begin{array}{cc}\hfill ^t\left[M(1,1,0)^3\right]=& (2\mathrm{}^3)M(0,0,3)(\mathrm{}^2+7\mathrm{}^3)M(0,1,1)(4\mathrm{}^4)M(0,3,1)\hfill \\ & (2\mathrm{}^2+2\mathrm{}^3)M(1,0,0)(8\mathrm{}^4)M(1,1,2)\hfill \\ & (4\mathrm{}^3+5\mathrm{}^5)M(1,2,0)(2\mathrm{}^5)M(1,4,0)\hfill \\ & (5\mathrm{}^3+4\mathrm{}^4)M(2,0,1)(6\mathrm{}^5)M(2,2,1)\hfill \\ & (4\mathrm{}^5)M(3,0,2)(4\mathrm{}^4+2\mathrm{}^5)M(3,1,0)\hfill \\ & (\mathrm{}^6)M(3,3,0)(2\mathrm{}^6)M(4,1,1)(\mathrm{}^5)M(5,0,0).\hfill \end{array}$$
From this example, we see no clear relation between the $`q`$-deformations and the $`\mathrm{}`$-deformations, except that they coincide at $`q=\mathrm{}=1`$.
In order to define the $`q`$-tensor product coefficients, one introduces the $`q`$-deformed Kostant partition function $`K(\beta ;q)`$:
$$\underset{\alpha R_>}{}(1qe^\alpha )^1=:\underset{\beta Z_{}R_>}{}K(\beta ;q)e^\beta .$$
From this we see that the powers of $`q`$ count the number of positive roots in a decomposition of an element of $`Z_{}R_>`$. The $`q`$-deformed weight multiplicities are
$$\mathrm{mult}_\lambda (\mu ;q):=\underset{wW}{}(detw)K(w.\lambda \mu ;q),$$
and we get the $`q`$-deformed tensor product coefficients as
$$T_{\lambda ,\mu }^\nu (q):=\underset{wW}{}(detw)\mathrm{mult}_\mu (w.\nu \lambda ;q).$$
Notice we use different brackets to distinguish the different deformations: $`T_{\lambda ,\mu }^\nu (q)`$ vs. $`T_{\lambda ,\mu }^\nu [\mathrm{}]`$. Finally, the $`q`$-fusion coefficients can be found from
$${}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }(q):=\underset{wW^k}{}(detw)T_{\lambda ,\mu }^{w.\nu }(q).$$
It would be interesting to define the threshold polynomial versions of the $`q`$-Kostant partition function, and the $`q`$-multiplicities. The relation between the $`q`$-tensor product coefficients and the threshold polynomials might then be extracted. We have not succeeded in finding the “$`\mathrm{}`$-Kostant partition function”. But the $`\mathrm{}`$-multiplicities may be defined using a conjectural refinement of the Gepner-Witten depth rule :
$${}_{}{}^{(k)}T_{\lambda ,\mu }^{\nu }=dim\left\{vM(\mu ;\nu \lambda )\right|(f_i)^{\nu _i+1}v=0,i\{0,1,\mathrm{},r\}\}.$$
Here $`M(\mu ;\nu \lambda )`$ is the subspace of weight $`\nu \lambda `$ of the module $`M(\mu )`$, so that $`dimM(\mu ;\sigma )=\mathrm{mult}_\mu (\sigma )`$. The $`f_i`$ are the lowering operators corresponding to the simple roots of the simple Lie algebra $`X_rX_{r,k}`$, for $`i=1,\mathrm{},r`$. $`f_0`$ is identified with $`e_\theta `$, the raising operator corresponding to the highest root $`\theta `$ of $`X_r`$. Recall that the depth of a vector $`v`$ (in a module $`M(\mu )`$, say) is defined as the non-negative integer $`d`$ such that $`(e_\theta )^dv0`$, while $`(e_\theta )^{d+1}v=0`$. The relation between (4.1) and the Gepner-Witten depth rule is then clear.
By (4.1), we see that the subspaces
$$M(\mu ;\sigma |d):=\{vM(\mu ;\sigma )|(e_\theta )^{\mathrm{\hspace{0.17em}1}+d}v=0,(e_\theta )^dv0\}$$
of the spaces $`M(\mu ;\sigma )`$ are relevant. Because of their relation to the depth, the multiplicities that are the dimensions of these spaces,
$$\mathrm{mult}_\mu (\sigma |d):=dimM(\mu ;\sigma |d),$$
were dubbed “profundities” in . They have the same relation to the threshold coefficients (see (2.1)) that the the usual multiplicities have to the tensor product coefficients:
$$T_{\lambda ,\mu }^{\nu (t)}=\underset{wW}{}(detw)\mathrm{mult}_\mu (w.\nu \lambda |t(\nu ,\theta )),$$
where $`t(\nu ,\theta )`$. Substituting this into (2.1) gives
$$T_{\lambda ,\mu }^\nu [\mathrm{}]:=\mathrm{}^{(\nu ,\theta )}\underset{wW}{}(detw)\mathrm{mult}_\mu (w.\nu \lambda ;\mathrm{}),$$
with $`\mathrm{}`$-deformed multiplicities
$$\mathrm{mult}_\mu (\sigma ;\mathrm{}):=\underset{dZ_{}}{}\mathrm{}^d\mathrm{mult}_\mu (\sigma |d).$$
So the only complication is the overall factor $`\mathrm{}^{(\nu ,\theta )}`$, and the $`\mathrm{}`$-deformed multiplicities are generating functions for the profundities.
Incidentally, in deriving this last result, we used the relation
$$d=t(\nu ,\theta )$$
between the depth $`d`$ and the threshold level $`t`$ of a fixed “coupling”. Notice that this is identical to (2.1). Hence the threshold delay of a coupling equals the depth of the corresponding vector in (4.1).
5. Fusion paths and the threshold level
Paths on fusion graphs are important in certain integrable lattice models that are related to conformal field theory (see , and references therein). These graphs may also have a more fundamental significance, indicated by the correspondence between $`A_1`$ modular invariants and the $`ADE`$ graphs .
If we restrict to the case $`X_{r,k}=A_{r,k}`$, then the points of the relevant graphs correspond to the weights of $`P_{}^k`$. There is a distinct directed graph $`{}_{}{}^{(k)}𝒢_{i}^{}`$ for each fundamental weight $`\mathrm{\Lambda }_iF`$. The edges of the graph $`{}_{}{}^{(k)}𝒢_{i}^{}`$ are encoded in its incidence matrix $`{}_{}{}^{(k)}G_{i}^{}`$, which is not necessarily symmetric. $`({}_{}{}^{(k)}G_{i}^{})_{\lambda ,\mu }`$ is the number of edges joining node $`\lambda `$ with node $`\mu `$. The fusion graph is defined by $`({}_{}{}^{(k)}G_{i}^{})_{\lambda ,\mu }={}_{}{}^{(k)}T_{\mathrm{\Lambda }_i,\lambda }^{\nu }`$, or $`{}_{}{}^{(k)}G_{i}^{}={}_{}{}^{(k)}T_{\mathrm{\Lambda }_i}^{}`$, hence the name. One can also define a graph $`{}_{}{}^{(k)}𝒢`$ with incidence matrix $`{}_{}{}^{(k)}G:=_{i=1}^r{}_{}{}^{(k)}G_{i}^{}`$.
A fusion path is a path on a fusion graph. Such paths parametrise the Hilbert space of certain integrable two-dimensional lattice models. The basic construction is a representation of a (quotient of a) Hecke algebra on this space. It guarantees that the models’ Boltzmann weights satisfy the Yang-Baxter equation, ensuring integrability.
Due to (2.1),(2.1), and since $`P_{}^kP_{}^{k+1}\mathrm{}P_{}`$, we can think of the graphs $`{}_{}{}^{(k)}𝒢_{i}^{}`$ and $`{}_{}{}^{(k)}𝒢`$ in the infinite-level limit as tensor product graphs. (Here we restrict consideration to paths involving weights that do not increase with the level.) Such paths on $`P_{}`$ will also be paths on all $`{}_{}{}^{(k)}𝒢`$, for all levels $`k`$ greater than a certain threshold level $`t`$. This threshold level, is just the maximum height $`ht(\lambda ):=(\lambda ,\theta )`$ of the weights $`\lambda P_{}`$ on the path.
Key to the modified Schubert calculus described above were the fundamental monomials, and their decompositions (3.1). But the fundamental monomials $`M(\mathrm{\Lambda }^\mu )`$ generate paths in $`P_{}`$: to every module $`M(\sigma )`$ in the decomposition (3.1) there corresponds a path on $`P_{}`$ that begins at the weight $`0`$ and ends at $`\lambda `$. To each factor $`M(\mathrm{\Lambda })`$, $`\mathrm{\Lambda }F`$, in the monomial corresponds a segment of the path that connects nodes of the graph that differ by some $`\phi P(\mathrm{\Lambda })`$.
The threshold level is included by modifying (3.1) to (3.1), using the $`\mathrm{}`$-Pieri rule. From (3.1), we see that the threshold level is the maximum height of a path on the (infinite) tensor product graph of $`A_r`$. This is the main point of this section.
We should emphasise, however, that the correspondence between the fundamental monomials and tensor product paths is not one-to-one. The polynomial realisation (3.1) of the fusion ring is possible because $`M(\mathrm{\Lambda })M(\mathrm{\Lambda }^{})=M(\mathrm{\Lambda }^{})M(\mathrm{\Lambda })`$, for all $`\mathrm{\Lambda },\mathrm{\Lambda }^{}F`$. But the order of tensor product factors $`M(\mathrm{\Lambda })`$ changes the path. This can be made clear by writing a generating matrix for fundamental monomials:
$$\mathrm{\Phi }:=\underset{\mu P_{}}{}e^\mu T_{\mathrm{\Lambda }_1}^{\mu _1}\mathrm{}T_{\mathrm{\Lambda }_r}^{\mu _r}=\underset{i=1}{\overset{r}{}}\left[1e^{\mathrm{\Lambda }_i}T_{\mathrm{\Lambda }_i}\right]^1.$$
Here $`e^\mu `$ denotes a formal exponential, satisfying $`e^\mu e^\nu =e^{\mu +\nu }`$. $`\mathrm{\Phi }_{\lambda ,\mu }`$ will equal the sum over all fundamental monomials that when tensored with $`M(\lambda )`$, include $`M(\mu )`$ in the decomposition, multiplied by the formal exponential of the monomial weight of each. Putting $`e^\mu 1`$ then gives the number of such monomials. On the other hand, to generate all paths connecting nodes $`\lambda `$ and $`\mu `$, we need $`\mathrm{\Theta }_{\lambda ,\mu }`$ instead, where
$$\mathrm{\Theta }:=\left[\mathrm{\hspace{0.17em}1}\underset{i=1}{\overset{r}{}}e^{\mathrm{\Lambda }_i}T_{\mathrm{\Lambda }_i}\right]^1.$$
The deformations of these two generating matrices are simple to write. One only needs to replace the tensor product matrices $`T_{\mathrm{\Lambda }_i}`$ with their $`\mathrm{}`$-deformations, and insist that they are multiplied in the manner of (2.1). So we get
$$\mathrm{\Phi }[\mathrm{}]=\underset{i=1}{\overset{r}{}}\left\{1e^{\mathrm{\Lambda }_i}T_{\mathrm{\Lambda }_i}[\mathrm{}]\right\}^{(1)},$$
where the notation (we hope) is clear, and a similar formula analogous to (5.1).
6. Conclusion
Our main result is a Schubert-type calculus for affine fusion that incorporates the threshold level. At fixed level, fusion constraints are natural because a fusion rule is a truncation of a tensor product decomposition. Thus fusion potentials that generate constraints are possible, if not necessary. On the other hand, in the threshold level formalism one doesn’t truncate a tensor product, but rather replaces it with a deformed version. So, instead of using a fusion potential to generate constraints, one just deforms the tensor products and then all (non-negative integer) levels are treated on equal footing. The deformations are generated by the deformed version of the Pieri rule, (3.1).
In summary then, to include the threshold levels in a calculus of Schubert type, use the undeformed Giambelli formula (3.1), and the deformed Pieri formula (3.1). Then the threshold polynomials can be calculated by (3.1),(3.1), or (3.1).
Another result is the comparison in section 4 of the threshold polynomials with the $`q`$-deformed tensor product and fusion coefficients. In particular, we found the analogue of the $`q`$-deformed weight multiplicities in our $`\mathrm{}`$-deformation.
We also discussed the interpretation of the threshold level for the decomposition of fundamental monomials, as in (3.1). In the corresponding path on a tensor product graph, the threshold level is just the maximum height of weights on that path.
To close, let us mention a few possible directions from this work.
One could hope to make the connection between the $`q`$-deformations and $`\mathrm{}`$-deformations more precise, extending our section 4. We also expect that one could define a $`q`$-Schubert calculus for the $`q`$-tensor product coefficients (4.1), in a straightforward way. In contrast with the $`\mathrm{}`$-deformed case, both $`\mathrm{\Omega }(q)`$ and $`\mathrm{\Omega }^1(q)`$ should be important. It might be of interest to introduce $`q`$-analogues of the fusion constraints and potentials of Gepner’s calculus, for the $`q`$-fusion-coefficients.
The $`\mathrm{}`$-deformed Schubert calculus is relevant to the search for a Littlewood-Richardson rule for affine fusion . In the present context, the usual Littlewood-Richardson rule for tensor products is related to (3.1), at $`\mathrm{}=1`$. This formula involves the tensor product of two modules $`M(\lambda ),M(\mu )`$, where one is expressed in terms of fundamental monomials by (3.1): $`M(\mu )=_\beta (\mathrm{\Omega })_{\mu ,\beta }^1M(\mathrm{\Lambda }^\beta )`$. The rule gives a way of avoiding the cancellations inherent in (3.1) (see (3.1), e.g.). It identifies a choice of a part of the decompositions of the $`M(\mathrm{\Lambda }^\beta )`$ in (3.1) that leads directly to the result. Unfortunately, the deformed Pieri rule applied to that choice gives incorrect threshold levels (one finds $`2M(1,1)_{2,2}`$ instead of the $`2M(1,1)_{2,3}`$ of (3.1), e.g.). Calculations of the type (3.1), however, show us all the parts. And so we can hope that more in-depth analysis will reveal the appropriate modification.
Finally, if it exists, a motivation other than computational for the threshold level should be found. It might be revealed by finding the meaning of the threshold level in the many different physical and mathematical realisations of affine fusion.
Acknowledgements
We thank L. Bégin, T. Gannon, P. Mathieu and J. Rasmussen for comments, and D. Sénéchal for the use of some computer programs.
References
relax D. Gepner, Commun. Math. Phys. 141 (1990) 381 relax H. Hiller, Geometry of Coxeter groups (Pitman, 1982); W. Fulton, Young tableaux (Cambridge University Press, 1997), Part III; P. Griffiths, J. Harris, Principles of algebraic geometry (Wiley, 1978) relax E. Witten, in: Geometry, Topology and Physics; Conf. Proc. and Lecture Notes in Geom. Topol. VI (1995) 357 relax K. Intriligator, Mod. Phys. Lett. A6 (1991) 3543 relax C. Vafa, in: Essays on Mirror Manifolds, ed. S.-T. Yau (International Press, 1992) relax D. Nemeschansky, N. Warner, Nucl. Phys. B380 (1992) 241 relax D. Gepner, E. Witten, Nucl. Phys. B278 (1986) 493 relax C.J. Cummins, P. Mathieu, M.A. Walton, Phys. Lett. 254B (1991) 386 relax A.N. Kirillov, P. Mathieu, D. Sénéchal, M.A. Walton, Nucl. Phys. B391 (1993) 651 relax A. Schilling, S.O. Warnaar, Commun. Math. Phys. 202 (1999) 359 relax A.N. Kirillov, M. Shimozono, A generalization of the Kostka-Foulkes polynomials, math.QA/9803062 relax M. Shimozono, J. Weyman, Graded characters of modules supported by the closure of a nilpotent conjugacy class, math.QA/9804036 relax B. Leclerc, J-Y. Thibon, Littlewood-Richardson coefficients and Kazhdan-Lusztig polynomials, math.QA/9809122 relax O. Foda, B. Leclerc, M. Okado, J-Y. Thibon, Ribbon tableaux and $`q`$-analogues of fusion rules in WZW conformal field theories, math.QA/9810008 relax A. Schilling, M. Shimozono, Bosonic formula for level-restricted paths, math.QA/9812106 relax G. Lusztig, Astérisque 101- 102 (1983) 208-229 relax A.N. Kirillov, P. Mathieu, D. Sénéchal, M.A. Walton, in: Group-Theoretical Methods in Physics, Proceedings of the XIXth International Colloquium, Salamanca, Spain, 1992, Vol. 1 (CIEMAT, Madrid, 1993) relax M.A. Walton, Can. J. Phys. 72 (1994) 527 relax J.-B. Zuber, Commun. Math. Phys. 179 (1996) 265 relax V.B. Petkova, J.-B. Zuber, Nucl. Phys. B463 (1996) 161 relax W. Fulton, J. Harris, Representation theory: a first course (Springer-Verlag, 1991) relax I. G. Macdonald, Notes on Schubert polynomials (Laboratoire de combinatoire et d’informatique mathématique (LACIM), Université du Québec à Montréal, 1991) relax V.G. Kac, Infinite dimensional Lie algebras, 3rd ed. (Cambridge U. Press, 1990) relax M.A. Walton, Phys. Lett. 241B (1990) 365; Nucl. Phys. B340 (1990) 777 relax P. Furlan, A. Ganchev, V. Petkova, Nucl. Phys. B343 (1990) 205 relax F. Goodman, H. Wenzl, Adv. Math. 82 (1990) 244 relax D. Kazhdan, G. Lusztig, Invent. Math. 53 (1979) 165 relax A. Cappelli, C. Itzykson, J.-B. Zuber, Commun. Math. Phys. 113 (1987) 1 relax M.A. Walton, J. Math. Phys. 39 (1998) 665
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# Acknowledgements
## Acknowledgements
We are grateful to A.L. Malvezzi and O.Techernyshyov for valuable discussions. This work was supported by Lampadia Foundation, by the Brazilian agencies CNPq and Fapesp and by the Netherlands Foundation FOM.
Table 1: Finite size and extrapolated results for the ground state energy and the triplet gap for cases (I) and (II). For comparison we also exhibit the exact results.
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# 1 The Weak-Field Case
## 1 The Weak-Field Case
### 1.1 Introduction
As a test of the applicability of the Weizsäcker-Williams method we first consider the weak-field case in which only a single initial photon is involved:
$$\omega _0+ee^{}+e^+e^{}.$$
(1)
A complete calculation for this process for unpolarized electrons and photons is available , but apparently the analytic form is too complex to be enlightening. A useful summary is given in sec. 11-4 of ref. .
Most discussions of reaction (1) use the frame in which the initial electron is at rest. In our recent experiment in which trident production is a background, the electron was ultrarelativistic in the lab frame. To be able to discuss the process in either frame it is useful to emphasize the relativistic invariants of the problem. In particular, we use
$$s=(\omega _0+e)^2.$$
(2)
where $`s`$ is the square of the center of mass energy of reaction (1) and in eq. (2) $`\omega _0`$ and $`e`$ represent the 4-momenta of the initial electron and photon. The threshold for reaction (1) is $`s_{\mathrm{min}}=9m^2`$, corresponding to the case when all final-state particles are at rest in the c.m. frame.
Just above threshold the cross section for the trident process (1) varies as
$$\sigma _T=9.2\times 10^4\alpha r_0^2\left(\frac{s9m^2}{m^2}\right)^2,(s9m^2m^2),$$
(3)
where $`\alpha =e^2/\mathrm{}c`$ is the fine structure constant, $`r_0=e^2/mc^2`$ is the classical electron radius, $`m`$ is the electron rest mass and $`c`$ is the speed of light. Far above threshold the cross section varies as
$$\sigma _T=\alpha r_0^2\left(\frac{28}{9}\mathrm{ln}\frac{s}{m^2}\frac{100}{9}\right),(s9m^2),$$
(4)
In our experiment, we were near threshold for trident production, but for $`ss_{\mathrm{min}}`$ extended up to $`𝒪(m^2)`$. Hence, we were between the regions of applicability of the asymptotic relations (3) and (4). A numerical tabulation of the trident cross section has been given by Mork based on the analytic calculation of Vortruba . Figure 1 compares theory and experiment for reaction (1).
We will compare the Weizsäcker-Williams approximation to the result of ref. below.
Mork also reported numerical results from a simplified calculation by Borsellino of trident production in which diagrams (b) and (d) of Fig. 2 are neglected. These diagrams are referred to as $`\gamma `$-$`e`$ or as Compton diagrams in that the initial-state photon couples directly to the initial electron rather than the $`e^+e^{}`$ pair. Well above threshold the Compton diagrams contribute little to the cross section, while near threshold they interfere to reduce the cross section, as summarized in Fig. 3 below. Thus the neglect of the Compton diagrams results in an overestimate of the cross section.
### 1.2 The Weizsäcker-Williams Approximation
In a frame in which the initial electron is ultrarelativistic its electric and magnetic fields are nearly transverse and of nearly equal magnitude. That is, the fields appear to be almost identical to those of a packet of real photons, which we label by $`\omega _1`$. A Fourier transform of the time dependence of the electron’s field integrated over observers at impact parameters $`b>b_{\mathrm{min}}`$ to the electron’s trajectory yields the photon number spectrum
$$N(\omega _1)\frac{2\alpha }{\pi \omega _1}\mathrm{ln}\frac{\gamma }{\omega _1b_{\mathrm{min}}},$$
(5)
where $`\gamma `$ is the Lorentz factor for the initial electron . We then suppose that a (virtual) photon $`\omega _1`$ from this spectrum combines with the incident photon $`\omega _0`$ to produce an $`e^+e^{}`$ pair via the Breit-Wheeler process :
$$\omega _0+\omega _1e^+e^{}.$$
(6)
The Weizsäcker-Williams approximation to the trident cross section is then
$$\sigma _T=\frac{2\alpha }{\pi }_{\omega _{1,\mathrm{min}}}^{\omega _{1,\mathrm{max}}}\frac{d\omega _1}{\omega _1}\mathrm{ln}\frac{\gamma }{\omega _1b_{\mathrm{min}}}\sigma _{\mathrm{BW}}(\omega _0,\omega _1),$$
(7)
where $`\sigma _{\mathrm{BW}}`$ is the Breit-Wheeler cross section.
The Breit-Wheeler cross section can be expressed in terms of $`s^{}`$, the square of the center of mass energy of the photon-photon system. For a frame in which the two photons collide head on, $`s^{}=4\omega _0\omega _1`$, where in this expression $`\omega `$ stands for the photon energy. Then
$$\sigma _{\mathrm{BW}}=4\pi r_0^2\frac{m^2}{s^{}}\left[\frac{3\beta ^4}{2}\mathrm{ln}\frac{1+\beta }{1\beta }\beta (2\beta ^2)\right],$$
(8)
where
$$\beta =\sqrt{1\frac{4m^2}{s^{}}}$$
(9)
is $`v/c`$ of the positron (or partner electron) in the pair rest frame. The threshold condition is, of course, $`s_{\mathrm{min}}^{}=4m^2`$. The asymptotic forms are
$$\sigma _{\mathrm{BW}}\pi r_0^2\beta ,(\beta 1s^{}4m^2m^2),$$
(10)
and
$$\sigma _{\mathrm{BW}}2\pi r_0^2\frac{m^2}{s^{}}\left(\mathrm{ln}\frac{s^{}}{m^2}1\right),(s^{}4m^2).$$
(11)
See, for example, sec. 13-3 of ref. .
Before inserting eq. (8) into (7) the latter should be put into a form that is more manifestly covariant. Our approach is to replace $`\omega _1`$ in (7) by $`s^{}=4\omega _0\omega _1`$. Immediately $`d\omega _1/\omega _1=ds^{}/s^{}`$.
Then the lower limit of integration, originally $`\omega _{1,\mathrm{min}}`$, becomes $`s_{\mathrm{min}}^{}=4m^2`$.
The upper limit of integration becomes $`s_{\mathrm{max}}^{}`$ for the Breit-Wheeler process embedded in the trident reaction (1). Another interpretation of $`s^{}`$ is as the square of the invariant mass of the $`e^+e^{}`$ pair: $`s^{}=m_{e^+e^{}}^2`$. Then $`s_{\mathrm{max}}^{}`$ occurs when as much energy as possible goes into the mass of the $`e^+e^{}`$ pair. This occurs when both the pair and the scattered electron $`e^{}`$ are at rest in the c.m. frame of reaction (1). In this case
$$s=(m+m_{e^+e^{}})^2=(m+\sqrt{s_{\mathrm{max}}^{}})^2,\text{and hence}s_{\mathrm{max}}^{}=(\sqrt{s}m)^2,$$
(12)
where $`s`$ is the square of the c.m. energy of reaction (1).
Finally, we need to reinterpret the argument of the logarithm in eq. (7). It should be an invariant, should be greater than 1, and should have $`\omega _1`$ in the denominator. The simplest form is then $`s_{\mathrm{max}}^{}/s^{}`$. This could be multiplied by a number of order 1, which is the usual ambiguity of the Weizsäcker-Williams method.
Altogether, the proposed invariant combination of (7) and (8) is
$`\sigma _T`$ $`=`$ $`{\displaystyle \frac{2\alpha }{\pi }}{\displaystyle _{4m^2}^{s_{\mathrm{max}}^{}}}{\displaystyle \frac{ds^{}}{s^{}}}\mathrm{ln}{\displaystyle \frac{s_{\mathrm{max}}^{}}{s^{}}}\sigma _{\mathrm{BW}}(s^{})`$ (13)
$`=`$ $`8\alpha r_0^2{\displaystyle _{4m^2}^{s_{\mathrm{max}}^{}}}{\displaystyle \frac{m^2ds^{}}{s^{}_{}{}^{}2}}\mathrm{ln}{\displaystyle \frac{s_{\mathrm{max}}^{}}{s^{}}}\left[{\displaystyle \frac{3\beta ^4}{2}}\mathrm{ln}{\displaystyle \frac{1+\beta }{1\beta }}\beta (2\beta ^2)\right],`$
where $`\beta `$ and $`s_{\mathrm{max}}^{}`$ are given by eqs. (9) and (12). This form can be evaluated in a frame in which the initial electron is at rest even though it is unclear that the field of the electron is equivalent to a collection of real photons in this frame.
Figures 3-5 show results of numerical calculations of eq. (13) along with the “exact” cross section as tabulated by Mork and the cross sections calculated by Borsellino by ignoring diagrams (b) and (d) of Fig. 2. The agreement of the Weizsäcker-Williams approximation and the exact calculation is fairly good, although worst near threshold where the rate is very low.
The results are plotted as a function of the invariant $`(ss_{\mathrm{min}})/2m^2`$ which is a measure of how far the reaction is above threshold. For the initial electron at rest this invariant is $`E_\gamma /m`$ where $`E_\gamma `$ is the energy of the initial photon.
In the Weizsäcker-Williams approximation the initial photon $`\omega _0`$ interacts only with the $`e^+e^{}`$ pair, not directly with the initial electron. Thus the approximation neglects diagrams (b) and (d) of Fig. 2. This feature is shared with the approximation of Borsellino , and indeed Figs. 3-5 show that the Weizsäcker-Williams approximation tracks Borsellino’s results more closely than the “exact” results.
The Weizsäcker-Williams approximation corresponds to the use of transverse but otherwise unpolarized virtual photons. This feature is shared with the exact calculation of diagrams (a) and (c) of Fig. 2. Only for the neglected diagrams (b) and (d) could there be polarization of the virtual photons in case the initial photon is polarized. For the Compton diagrams (b) and (d) the virtual photons would take on the polarization of the initial photon for energies of the virtual photon near the maximum.
## 2 The Strong-Field Case
Multiphoton trident production can occur in a strong field of initial-state photons:
$$e+n\omega _0e^{}+e^+e^{}.$$
(14)
Extrapolating from eq. (13), we propose the Weizsäcker-Williams approximation for reaction (14) be formulated as
$$Rate_T=\frac{2\alpha }{\pi }\underset{n}{}_{4\overline{m}^2}^{s_{n,\mathrm{max}}^{}}\frac{ds_n^{}}{s_n^{}}\mathrm{ln}\frac{s_{n,\mathrm{max}}^{}}{s_n^{}}Rate_{\mathrm{BW}}(s_n^{},\eta ),$$
(15)
where $`\overline{m}=m\sqrt{1+\eta ^2}`$ is the shifted mass of the electron in the strong field and $`\eta =e_{\mathrm{rms}}/m\omega c`$ is the field-strength parameter. In this we calculate a rate rather than a cross section, using the results of Nikishov and Ritus . For $`n`$ initial-state (laser) photons, the sub-process is the multiphoton Breit-Wheeler reaction
$$n\omega _0+\omega _1e^+e^{},$$
(16)
for which the square of the c.m. energy is
$$s_n^{}=(n\omega _0+\omega _1)^2,$$
(17)
where in this expression $`\omega `$ stands for the 4-momentum of a photon. Equation (12) becomes
$$s_{n,\mathrm{max}}^{}=(\sqrt{s_n}\overline{m})^2,\text{with}s_n=(e+n\omega _0)^2=\overline{m}^2+2n(e\omega _0),$$
(18)
where $`e`$ and $`\omega _0`$ are the initial-state 4-momenta (including mass-shift effects, hence $`e^2=\overline{m}^2`$).
If $`s_{n,\mathrm{max}}^{}<4\overline{m}^2`$ there is no contribution at order $`n`$. This condition can be stated another way. The threshold condition is that the final-state electron and the $`e^+e^{}`$ pair are both at rest in the c.m. frame of the $`e+n\omega _0`$ system, and that $`m_{e^+e^{}}=2\overline{m}`$. That is,
$$s_{n,\mathrm{min}}=(e+n\omega _0)^2(\overline{m}+2\overline{m})^2=9\overline{m}^2.$$
(19)
For a head-on collision between a relativistic electron and the initial-state photons this becomes
$$\overline{m}^2+4nE_0\omega _09\overline{m}^2,\text{or}n\frac{2\overline{m}^2}{E_0\omega _0},$$
(20)
where $`E_0`$ is the energy of the initial electron. Strictly speaking, $`E_0=q_0`$, the quasienergy of the initial electron which is related by
$$q=e+\frac{\eta ^2m^2}{2(e\omega _0)}\omega _0,$$
(21)
where $`e`$, $`q`$ and $`\omega _0`$ are 4-momenta in this expression. For our example, $`q_0E_0+\eta ^2m^2/4E_0E_0`$.
For $`E_0`$ = 46.6 GeV and $`\omega _0`$ = 2.3 eV we must have $`n5`$. As noted in sec. 1.2, this corresponds to all initial-state photons coupling only to the $`e^+e^{}`$ pair, and ignores the strong-field generalizations of diagrams (b) and (d) of Fig. 2.
In the Weizsäcker-Williams approximation the virtual photon $`\omega _1`$ is unpolarized even if the initial-state photon $`\omega _0`$ is polarized. Hence the Breit-Wheeler rate used in eq. (15) should be for unpolarized $`\omega _1`$ but with whatever polarization holds for the initial-state photons $`\omega _0`$.
### 2.1 Numerical Results
The above procedures have been implemented in a numerical simulation .
The requirements (18) and (19) that energy be conserved during pair creation has a striking effect on the calculated rate. First, the minimum number of laser photons is $`n=5`$ (and it turns out that there is no significant rate unless $`n>5)`$. Second, the maximum energy of the virtual photon, $`\omega _{1,\mathrm{max}}`$, that can contribute is much less than the electron beam energy $`E_0`$. The latter point can be anticipated by the approximation of head-on collisions for which eq. (18) tells us
$$\omega _{1,\mathrm{max}}=E_0\frac{\overline{m}^2}{2n\omega _0}\left(\sqrt{1+\frac{4nE_0\omega _0}{\overline{m}^2}}1\right).$$
(22)
Some representative values are given in the Table.
Figure 6 shows results of the trident-rate calculation for various numbers of laser photons. The solid curves are the proper results while the dashed curves show the effect of setting $`\omega _{1,\mathrm{max}}`$ to $`E_0`$.
Figure 7 shows the contribution to the rate as a function of the invariant measure of energy above threshold. Only the case $`n=6`$ is close enough to threshold that the Weizsäcker-Williams approximation is significantly in error, and the sign of the error is to overestimate the rate.
Finally, Figure 8 shows the total rate of trident production (16) as a function of laser intensity parameter $`\eta ^2`$ for typical conditions of our experiment . Also shown is the calculated rate for pair creation by the two-step process
$$e+n\omega _0e^{}+\omega _1,\text{followed by}\omega _1+m\omega _0e^+e^{}.$$
(23)
The trident process is only a 1% correction to the two-step production process.
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# Contents
## 1 Introduction
### Artificial Intelligence:
The science of Artificial Intelligence (AI) might be defined as the construction of intelligent systems and their analysis. A natural definition of systems is anything which has an input and an output stream. Intelligence is more complicated. It can have many faces like creativity, solving problems, pattern recognition, classification, learning, induction, deduction, building analogies, optimization, surviving in an environment, language processing, knowledge and many more. A formal definition incorporating every aspect of intelligence, however, seems difficult. Further, intelligence is graded, there is a smooth transition between systems, which everyone would agree to be not intelligent and truely intelligent systems. One simply has to look in nature, starting with, for instance, inanimate crystals, then come amino-acids, then some RNA fragments, then viruses, bacteria, plants, animals, apes, followed by the truly intelligent homo sapiens, and possibly continued by AI systems or ET’s. So the best we can expect to find is a partial or total order relation on the set of systems, which orders them w.r.t. their degree of intelligence (like intelligence tests do for human systems, but for a limited class of problems). Having this order we are, of course, are interested in large elements, i.e. highly intelligent systems. If a largest element exists, it would correspond to the most intelligent system which could exist.
Most, if not all known facets of intelligence can be formulated as goal driven or, more precisely, as maximizing some utility function. It is, therefore, sufficient to study goal driven AI. E.g. the (biological) goal of animals and humans is to survive and spread. The goal of AI systems should be to be useful to humans. The problem is that, except for special cases, we know neither the utility function, nor the environment in which the system will operate, in advance.
### Main idea:
We propose a theory which formally<sup>2</sup><sup>2</sup>2With a formal solution we mean a rigorous mathematically definition, uniquely specifying the solution. In the following, a solution is always meant in this formal sense. solves the problem of unknown goal and environment. It might be viewed as a unification of the ideas of universal induction, probabilistic planning and reinforcement learning or as a unification of sequential decision theory with algorithmic information theory. We apply this model to some of the facets of intelligence, including induction, game playing, optimization, reinforcement and supervised learning, and show how it solves these problem classes. This, together with general convergence theorems motivates us to believe that the constructed universal AI system is the best one in a sense to be clarified in the sequel, i.e. that it is the most intelligent environmental independent system possible. The intention of this work is to introduce the universal AI model and give an in breadth analysis. Most arguments and proofs are succinct and require slow reading or some additional pencil work.
### Contents:
Section 2: The general framework for AI might be viewed as the design and study of intelligent agents . An agent is a cybernetic system with some internal state, which acts with output $`y_k`$ to some environment in cycle $`k`$, perceives some input $`x_k`$ from the environment and updates its internal state. Then the next cycle follows. It operates according to some function $`p`$. We split the input $`x_k`$ into a regular part $`x_k^{}`$ and a credit $`c_k`$, often called reinforcement feedback. From time to time the environment provides non-zero credit to the system. The task of the system is to maximize its utility, defined as the sum of future credits. A probabilistic environment is a probability distribution $`\mu (q)`$ over deterministic environments $`q`$. Most, if not all environments are of this type. We give a formal expression for the function $`p^{}`$, which maximizes in every cycle the total $`\mu `$ expected future credit. This model is called the AI$`\mu `$ model. As every AI problem can be brought into this form, the problem of maximizing utility is hence being formally solved, if $`\mu `$ is known. There is nothing remarkable or new here, it is the essence of sequential decision theory . Notation and formulas needed in later sections are simply developed. There are two major remaining problems. The problem of the unknown true prior probability $`\mu `$ is solved in section 4. Computational aspects are addressed in section 10.
Section 3: Instead of talking about probability distributions $`\mu (q)`$ over functions, one could describe the environment by the conditional probability of providing inputs $`x_1\mathrm{}x_n`$ to the system under the condition that the system outputs $`y_1\mathrm{}y_n`$. The definition of the optimal $`p^{}`$ system in this iterative form is shown to be equivalent to the previous functional form. The functional form is more elegant and will be used to define an intelligence order relation and the time-bounded model in section 10. The iterative form is more index intensive but more suitable for explicit calculations and is used in most of the other sections. Further, we introduce factorizable probability distributions.
Section 4: A special topic is the theory of induction. In which sense prediction of the future is possible at all, is best summarized by the theory of Solomonoff. Given the initial binary sequence $`x_1\mathrm{}x_k`$, what is the probability of the next bit being $`1`$? It can be fairly well predicted by using a universal probability distribution $`\xi `$ invented and shown to converge to the true prior probability $`\mu `$ by Solomonoff as long as $`\mu `$ (which needs not be known!) is computable. The problem of unknown $`\mu `$ is hence solved for induction problems. All AI problems where the systems’ output does not influence the environment, i.e. all passive systems are of this inductive form. Besides sequence prediction (SP), classification(CF) is also of this type. Active systems, like game playing (SG) and optimization (FM), can not be reduced to induction systems. The main idea of this work is to generalize universal induction to the general cybernetic model described in sections 2 and 3. For this, we generalize $`\xi `$ to include conditions and replace $`\mu `$ by $`\xi `$ in the rational agent model. In this way the problem that the true prior probability $`\mu `$ is usually unknown is solved. Universality of $`\xi `$ and convergence of $`\xi \mu `$ will be shown. These are strong arguments for the optimality of the resulting AI$`\xi `$ model. There are certain difficulties in proving rigorously that and in which sense it is optimal, i.e. the most intelligent system. Further, we introduce a universal order relation for intelligence.
Sections 59 show how a number of AI problem classes fit into the general AI$`\xi `$ model. All these problems are formally solved by the AI$`\xi `$ model. The solution is, however, only formal because the AI$`\xi `$ model developed thus far is uncomputable or, at best, approximable. These sections should support the claim that every AI problem can be formulated (and hence solved) within the AI$`\xi `$ model. For some classes we give concrete examples to illuminate the scope of the problem class. We first formulate each problem class in its natural way (when $`\mu ^{\text{problem}}`$ is known) and then construct a formulation within the AI$`\mu `$ model and prove its equivalence. We then consider the consequences of replacing $`\mu `$ by $`\xi `$. The main goal is to understand why and how the problems are solved by AI$`\xi `$. We only highlight special aspects of each problem class. Sections 59 together should give a better picture of the AI$`\xi `$ model. We do not study every aspect for every problem class. The sections might be read selectively. They are not necessary to understand the remaining sections.
Section 5: Using the AI$`\mu `$ model for sequence prediction (SP) is identical to Baysian sequence prediction SP$`\mathrm{\Theta }_\mu `$. One might expect, when using the AI$`\xi `$ model for sequence prediction, one would recover exactly the universal sequence prediction scheme SP$`\mathrm{\Theta }_\xi `$, as AI$`\xi `$ was a unification of the AI$`\mu `$ model and the idea of universal probability $`\xi `$. Unfortunately this is not the case. One reason is that $`\xi `$ is only a probability distribution in the inputs $`x`$ and not in the outputs $`y`$. This is also one of the origins of the difficulty of proving error/credit bounds for AI$`\xi `$. Nevertheless, we argue that AI$`\xi `$ is equally well suited for sequence prediction as SP$`\mathrm{\Theta }_\xi `$ is. In a very limited setting we prove a (weak) error bound for AI$`\xi `$ which gives hope that a general proof is attainable.
Section 6: A very important class of problems are strategic games (SG). We restrict ourselves to deterministic strictly competitive strategic games like chess. If the environment is a minimax player, the AI$`\mu `$ model itself reduces to a minimax strategy. Repeated games of fixed lengths are a special case for factorizable $`\mu `$. The consequences of variable game length is sketched. The AI$`\xi `$ model has to learn the rules of the game under consideration, as it has no prior information about these rules. We describe how AI$`\xi `$ actually learns these rules.
Section 7: There are many problems that fall into the category ’resource bounded function minimization’ (FM). They include the Traveling Salesman Problem, minimizing production costs, inventing new materials or even producing, e.g. nice paintings, which are (subjectively) judged by a human. The task is to (approximately) minimize some function $`f:YZ`$ within minimal number of function calls. We will see that a greedy model trying to minimize $`f`$ in every cycle fails. Although the greedy model has nothing to do with downhill or gradient techniques (there is nothing like a gradient or direction for functions over $`Y`$) which are known to fail, we discover the same difficulties. FM has already nearly the full complexity of general AI. The reason being that FM can actively influence the information gathering process by its trials $`y_k`$ (whereas SP and CF cannot). We discuss in detail the optimal FM$`\mu `$ model and its inventiveness in choosing the $`yY`$. A discussion of the subtleties when using AI$`\xi `$ for function minimization, follows.
Section 8: Reinforcement learning, as the AI$`\xi `$ model does, is an important learning technique but not the only one. To improve the speed of learning, supervised learning, i.e. learning by acquiring knowledge, or learning from a constructive teacher is necessary. We show, how AI$`\xi `$ learns to learn supervised. It actually establishes supervised learning very quickly within $`O(1)`$ cycles.
Section 9 gives a brief survey of other general aspects, ideas and methods in AI, and their connection to the AI$`\xi `$ model. Some aspects are directly included, others are or should be emergent.
Section 10: Up to now we have shown the universal character of the AI$`\xi `$ model but have completely ignored computational aspects. Let us assume that there exists some algorithm $`\stackrel{~}{p}`$ of size $`\stackrel{~}{l}`$ with computation time per cycle $`\stackrel{~}{t}`$, which behaves in a sufficiently intelligent way (this assumption is the very basis of AI). The algorithm $`p^{}`$ should run all algorithms of length $`\stackrel{~}{l}`$ for $`\stackrel{~}{t}`$ time steps in every cycle and select the best output among them. So we have an algorithm which runs in time $`\stackrel{~}{l}2^{\stackrel{~}{t}}`$ and is at least as good as $`\stackrel{~}{p}`$, i.e. it also serves our needs apart from the (very large but) constant multiplicative factor in computation time. This idea of the ’typing monkeys’, one of them eventually producing ’Shakespeare’, is well known and widely used in theoretical computer science. The difficult part is the selection of the algorithm with the best output. A further complication is that the selection process itself must have only limited computation time. We present a suitable modification of the AI$`\xi `$ model which solves these difficult problems. The solution is somewhat involved from an implementational aspect. An implementation would include first order logic, the definition of a Universal Turing machine within it and proof theory. The assumptions behind this construction are discussed at the end.
Section 11 contains some discussion of otherwise unmentioned topics and some (personal) remarks. It also serves as an outlook to further research.
Section 12 contains the conclusions.
### History & References:
Kolmogorov65 suggested to define the information content of an object as the length of the shortest program computing a representation of it. Solomonoff64 invented the closely related universal prior probability distribution and used it for binary sequence prediction and function inversion and minimization . Together with Chaitin66&75 this was the invention of what is now called Algorithmic Information theory. For further literature and many applications see . Other interesting ’applications’ can be found in . Related topics are the Weighted Majority Algorithm invented by Littlestone and Warmuth89 , universal forecasting by Vovk92 , Levin search73 , pac-learning introduced by Valiant84 and Minimum Description Length . Resource bounded complexity is discussed in , resource bounded universal probability in . Implementations are rare . Excellent reviews with a philosophical touch are . For an older, but general review of inductive inference see Angluin83 . For an excellent introduction into algorithmic information theory, further literature and many applications one should consult the book of Li and Vitányi97 . The survey or the chapters 4 and 5 of should be sufficient to follow the arguments and proofs in this paper. The other ingredient in our AI$`\xi `$ model is sequential decision theory. We do not need much more than the maximum expected utility principle and the expecimax algorithm . The book of von Neumann and Morgenstern44 might be seen as the initiation of game theory, which already contains the expectimax algorithm as a special case. The literature on decision theory is vast and we only give two possibly interesting references with regard to this paper. Cheeseman85&88 is a defense of the use of probability theory in AI. Pearl88 is a good introduction and overview of probabilistic reasoning.
## 2 The AI$`\mu `$ Model in Functional Form
### The cybernetic or agent model:
A good way to start thinking about intelligent systems is to consider more generally cybernetic systems, in AI usually called agents. This avoids having to struggle with the meaning of intelligence from the very beginning. A cybernetic system is a control circuit with input $`y`$ and output $`x`$ and an internal state. From an external input and the internal state the system calculates deterministically or stochastically an output. This output (action) modifies the environment and leads to a new input (reception). This continues ad infinitum or for a finite number of cycles. As explained in the last section, we need some credit assignment to the cybernetic system. The input $`x`$ is divided into two parts, the standard input $`x^{}`$ and some credit input $`c`$. If input and output are represented by strings, a deterministic cybernetic system can be modeled by a Turing machine $`p`$. $`p`$ is called the policy of the agent, which determines the action to a receipt. If the environment is also computable it might be modeled by a Turing machine $`q`$ as well. The interaction of the agent with the environment can be illustrated as follows:
$`p`$ as well as $`q`$ have unidirectional input and output tapes and bidirectional working tapes. What entangles the agent with the environment, is the fact that the upper tape serves as input tape for $`p`$, as well as output tape for $`q`$, and that the lower tape serves as output tape for $`p`$ as well as input tape for $`q`$. Further, the reading head must always be left of the writing head, i.e. the symbols must first be written, before they are read. $`p`$ and $`q`$ have their own mutually inaccessible working tapes containing their own ’secrets’. The heads move in the following way. In the k<sup>th</sup> cycle $`p`$ writes $`y_k`$, $`q`$ reads $`y_k`$, $`q`$ writes $`x_kc_kx_k^{}`$, $`p`$ reads $`x_kc_kx_k^{}`$, followed by the $`(k+1)^{th}`$ cycle and so on. The whole process starts with the first cycle, all heads on tape start and working tapes being empty. We want to call Turing machines behaving in this way, chronological Turing machines, for obvious reasons. Before continuing, some notations on strings are appropriate.
### Strings:
We will denote strings over the alphabet $`X`$ by $`s=x_1x_2\mathrm{}x_n`$, with $`x_kX`$, where $`X`$ is alternatively interpreted as a non-empty subset of $`IN`$ or itself as a prefix free set of binary strings. $`l(s)=l(x_1)+\mathrm{}+l(x_n)`$ is the length of s. Analogous definitions hold for $`y_kY`$. We call $`x_k`$ the $`k^{th}`$ input word and $`y_k`$ the $`k^{th}`$ output word (rather than letter). The string $`s=y_1x_1\mathrm{}y_nx_n`$ represents the input/output in chronological order. Due to the prefix property of the $`x_k`$ and $`y_k`$, $`s`$ can be uniquely separated into its words. The words appearing in strings are always in chronological order. We further introduce the following abbreviations: $`ϵ`$ is the empty string, $`x_{n:m}:=x_nx_{n+1}\mathrm{}x_{m1}x_m`$ for $`nm`$ and $`ϵ`$ for $`n>m`$. $`x_{<n}:=x_1\mathrm{}x_{n1}`$. Analog for $`y`$. Further, $`yx_n:=y_nx_n`$, $`yx_{n:m}:=y_nx_n\mathrm{}y_mx_m`$, and so on.
### AI model for known deterministic environment:
Let us define for the chronological Turing machine $`p`$ a partial function also named $`p:X^{}Y^{}`$ with $`y_{1:k}=p(x_{<k})`$ where $`y_{1:k}`$ is the output of Turing machine $`p`$ on input $`x_{<k}`$ in cycle k, i.e. where $`p`$ has read up to $`x_{k1}`$ but no further. In an analogous way, we define $`q:Y^{}X^{}`$ with $`x_{1:k}=q(y_{1:k})`$. Conversely, for every partial recursive chronological function we can define a corresponding chronological Turing machine. Each (system,environment) pair $`(p,q)`$ produces a unique I/O sequence $`\omega (p,q):=y_1^{pq}x_1^{pq}y_2^{pq}x_2^{pq}\mathrm{}`$. When we look at the definition of $`p`$ and $`q`$ we see a nice symmetry between the cybernetic system and the environment. Until now, not much intelligence is in our system. Now the credit assignment comes into the game and removes the symmetry somewhat. We split the input $`x_kX:=C\times X^{}`$ into a regular part $`x_k^{}X^{}`$ and a credit $`c_kCIR`$. We define $`x_kc_kx_k^{}`$ and $`c_kc(x_k)`$. The goal of the system should be to maximize received credits. This is called reinforcement learning. The reason for the asymmetry is, that eventually we (humans) will be the environment with which the system will communicate and we want to dictate what is good and what is wrong, not the other way round. This one way learning, the system learns from the environment, and not conversely, neither prevents the system from becoming more intelligent than the environment, nor does it prevent the environment learning from the system because the environment can itself interpret the outputs $`y_k`$ as a regular and a credit part. The environment is just not forced to learn, whereas the system is. In cases where we restrict the credit to two values $`cC=IB:=\{0,1\}`$, $`c=1`$ is interpreted as a positive feedback, called good or correct and $`c=0`$ a negative feedback, called bad or error in the following. Further, let us restrict for a while the lifetime (number of cycles) $`T`$ of the system to a large, but finite value. Let $`C_{km}(p,q):=_{i=k}^mc(x_i)`$ be the total credit, the system $`p`$ receives from the environment $`q`$ in the cycles $`k`$ to $`m`$. It is now natural to call the system, which maximizes the total credit $`C_{1T}`$, called utility, the best or most intelligent one<sup>3</sup><sup>3</sup>3$`maxarg_pC(p)`$ is the $`p`$ which maximizes $`C()`$. If there is more than one maximum we might choose the lexicographically smallest one for definiteness..
$$p^{,T,q}=\underset{p}{maxarg}C_{1T}(p,q)C_{kT}(p^{,T,q},q)C_{kT}(p,q)p$$
For $`k=1`$ this is obvious and for $`k>1`$ easy to see. If $`T`$, $`Y`$ and $`X`$ are finite, the number of different behaviours of the system, i.e. the search space is finite. Therefore, because we have assumed that $`q`$ is known, $`p^{,T,q}`$ can effectively be determined (by pre-analyzing all behaviours). The main reason for restricting to finite $`T`$ was not to ensure computability of $`p^{,T,q}`$ but that the limit $`T\mathrm{}`$ might not exist. This is nothing special, the (unrealistic) assumption of a completely known deterministic environment $`q`$ has simply trivialized everything.
### AI model for known prior probability:
Let us now weaken our assumptions by replacing the environment $`q`$ with a probability distribution $`\mu (q)`$ over chronological functions. $`\mu `$ might be interpreted in two ways. Either the environment itself behaves in a probabilistic way defined by $`\mu `$ or the true environment is deterministic, but we only have probabilistic information, of which environment being the true environment. Combinations of both cases are also possible. The interpretation does not matter in the following. We just assume that we know $`\mu `$ but no more about the environment whatever the interpretation may be.
Let us assume we are in cycle $`k`$ with history $`\dot{y}\dot{x}_1\mathrm{}\dot{y}\dot{x}_{k1}`$ and ask for the best output $`y_k`$. Further, let $`\dot{Q}_k:=\{q:q(\dot{y}_{<k})=\dot{x}_{<k}\}`$ be the set of all environments producing the above history. The expected credit for the next $`mk+1`$ cycles (given the above history) is given by a conditional probability:
$$C_{km}^\mu (p|\dot{y}\dot{x}_{<k}):=\frac{\underset{q\dot{Q}_k}{}\mu (q)C_{km}(p,q)}{_{q\dot{Q}_k}\mu (q)}.$$
(1)
We cannot simply determine $`maxarg_p(C_{1T})`$ unlike the deterministic case because the history is no longer deterministically determined by $`p`$ and $`q`$, but depends on $`p`$ and $`\mu `$ and on the outcome of a stochastic process. Every new cycle adds new information ($`\dot{x}_i`$) to the system. This is indicated by the dots over the symbols. In cycle $`k`$ we have to maximize the expected future credit, taking into account the information in the history $`\dot{y}\dot{x}_{<k}`$. This information is not already present in $`p`$ and $`q/\mu `$ at the system’s start unlike in the deterministic case.
Further, we want to generalize the finite lifetime $`T`$ to a dynamical (computable) farsightedness $`h_km_kk+11`$, called horizon in the following. For $`m_k=T`$ we have our original finite lifetime, for $`m_k=k+m1`$ the system maximizes in every cycle the next $`m`$ expected credits. A discussion of the choices $`m_k`$ is delayed to section 4.
The next $`h_k`$ credits are maximized by
$$p_k^{}:=\underset{p\dot{P}_k}{maxarg}C_{km_k}^\mu (p|\dot{y}\dot{x}_{<k}),$$
where $`\dot{P}_k:=\{p:p(\dot{x}_{<k})=\dot{y}_{<k}\}`$ is the set of systems consistent with the current history. $`p_k^{}`$ depends on $`k`$ and is used only in step $`k`$ to determine $`\dot{y}_k`$ by $`p_k^{}(\dot{x}_{<k};\dot{y}_{<k})=\dot{y}_{<k}\dot{y}_k`$. After writing $`\dot{y}_k`$ the environment replies with $`\dot{x}_k`$ with (conditional) probability $`\mu (\dot{Q}_{k+1})/\mu (\dot{Q}_k)`$. This probabilistic outcome provides new information to the system. The cycle $`k+1`$ starts with determining $`\dot{y}_{k+1}`$ from $`p_{k+1}^{}`$ (which differs from $`p_k`$ as $`\dot{x}_k`$ is now fixed) and so on. Note that $`p_k^{}`$ depends also on $`\dot{y}_{<k}`$ because $`\dot{P}_k`$ and $`\dot{Q}_k`$ do so. But recursively inserting $`p_{k1}^{}`$ and so on, we can define
$$p^{}(\dot{x}_{<k}):=p_k^{}(\dot{x}_{<k};p_{k1}^{}(\dot{x}_{<k1}\mathrm{}p_1^{})))$$
(2)
It is a chronological function and computable if $`X`$, $`Y`$ and $`m_k`$ are finite. The policy $`p^{}`$ defines our AI$`\mu `$ model. For deterministic<sup>4</sup><sup>4</sup>4We call a probability distribution deterministic if it is 1 for exactly one argument and 0 for all others. $`\mu `$ this model reduces to the deterministic case.
It is important to maximize the sum of future credits and not, for instance, to be greedy and only maximize the next credit, as is done e.g. in sequence prediction. For example, let the environment be a sequence of chess games and each cycle corresponds to one move. Only at the end of each game a positive credit $`c=1`$ is given to the system if it won the game (and made no illegal move). For the system, maximizing all future credits means trying to win as many games in as short as possible time (and avoiding illegal moves). The same performance is reached, if we choose $`m_k=k+m`$ with $`m`$ much larger than the typical game lengths. Maximization of only the next credit would be a very bad chess playing system. Even if we would make our credit $`c`$ finer, e.g. by evaluating the number of chessmen, the system would play very bad chess for $`m=1`$, indeed.
The AI$`\mu `$ model still depends on $`\mu `$ and $`m_k`$. $`m_k`$ is addressed in section 4. To get our final universal AI model the idea is to replace $`\mu `$ by the universal probability $`\xi `$, defined later. This is motivated by the fact that $`\xi \mu `$ in a certain sense for any $`\mu `$. With $`\xi `$ instead of $`\mu `$ our model no longer depends on any parameters, so it is truly universal. It remains to show that it produces intelligent outputs. But let us continue step by step. In the next section we develop an alternative but equivalent formulation of the AI model given above. Whereas the functional form is more suitable for theoretical considerations, especially for the development of a timebounded version in section 10, the iterative formulation of the next section will be more appropriate for the explicit calculations in most of the other sections.
## 3 The AI$`\mu `$ Model in Recursive and Iterative Form
### Probability distributions:
Throughout the paper we deal with sequences/strings and conditional probability distributions on strings. Some notations are therefore appropriate.
We use Greek letters for probability distributions and underline their arguments to indicate that they are probability arguments. Let $`\rho _n(\underset{¯}{x}_1\mathrm{}\underset{¯}{x}_n)`$ be the probability that a string starts with $`x_1\mathrm{}x_n`$. We only consider sufficiently long strings, so the $`\rho _n`$ are normalized to 1. Moreover, we drop the index on $`\rho `$ if it is clear from its arguments:
$$\underset{x_nX}{}\rho (\underset{¯}{x}_{1:n})\underset{x_n}{}\rho _n(\underset{¯}{x}_{1:n})=\rho _{n1}(\underset{¯}{x}_{<n})\rho (\underset{¯}{x}_{<n}),\rho (ϵ)\rho _0(ϵ)=1.$$
(3)
We also need conditional probabilities derived from Bayes’ rule. We prefer a notation which preserves the chronological order of the words, in contrast to the standard notation $`\rho (|)`$ which flips it. We extend the definition of $`\rho `$ to the conditional case with the following convention for its arguments: An underlined argument $`\underset{¯}{x}_k`$ is a probability variable and other non-underlined arguments $`x_k`$ represent conditions. With this convention, Bayes’ rule has the form $`\rho (x_{<n}\underset{¯}{x}_n)=\rho (\underset{¯}{x}_{1:n})/\rho (\underset{¯}{x}_{<n})`$. The equation states that the probability that a string $`x_1\mathrm{}x_{n1}`$ is followed by $`x_n`$ is equal to the probability of $`x_1\mathrm{}x_n`$ divided by the probability of $`x_1\mathrm{}x_{n1}`$. We use $`x`$ as a shortcut for ’strings starting with $`x`$’.
The introduced notation is also suitable for defining the conditional probability $`\rho (y_1\underset{¯}{x}_1\mathrm{}y_n\underset{¯}{x}_n)`$ that the environment reacts with $`x_1\mathrm{}x_n`$ under the condition that the output of the system is $`y_1\mathrm{}y_n`$. The environment is chronological, i.e. input $`x_i`$ depends on $`yx_{<i}y_i`$ only. In the probabilistic case this means that $`\rho (y\underset{¯}{x}_{<k}y_k):=_{x_k}\rho (y\underset{¯}{x}_{1:k})`$ is independent of $`y_k`$, hence a tailing $`y_k`$ in the arguments of $`\rho `$ can be dropped. Probability distributions with this property will be called chronological. The $`y`$ are always conditions, i.e. never underlined, whereas additional conditioning for the $`x`$ can be obtained with Bayes’ rule
$$\begin{array}{c}\rho (yx_{<n}y\underset{¯}{x}_n)=\rho (y\underset{¯}{x}_{1:n})/\rho (y\underset{¯}{x}_{<n})\text{and}\\ \rho (y\underset{¯}{x}_{1:n})=\rho (y\underset{¯}{x}_1)\rho (yx_1y\underset{¯}{x}_2)\mathrm{}\rho (yx_{<n}y\underset{¯}{x}_n)\end{array}$$
(4)
The second equation is the first equation applied $`n`$ times.
### Alternative Formulation of the AI$`\mu `$ Model:
Let us define the AI$`\mu `$ model $`p^{}`$ in a different way. In the next subsection we will show that the $`p^{}`$ model defined here is identical to the functional definition of $`p^{}`$ given in the last section.
Let $`\mu (y\underset{¯}{x}_{1:k})`$ be the true chronological prior probability that the environment reacts with $`x_{1:k}`$ if provided with actions $`y_{1:k}`$ from the system. We assume the cybernetic model depicted on page 2 to be valid. Next we define $`C_{k+1,m}^{}(yx_{1:k})`$ to be the $`\mu `$ expected credit sum in cycles $`k+1`$ to $`m`$ with outputs $`y_i`$ generated by system $`p^{}`$ and past responses $`x_i`$ from the environment. Adding $`c(x_k)`$ we get the credit including cycle $`k`$. The probability of $`x_k`$, given $`yx_{<k}y_k`$, is given by the condition probability $`\mu (yx_{<k}y\underset{¯}{x}_k)`$. So the expected credit sum in cycles $`k`$ to $`m`$ given $`yx_{<k}y_k`$ is
$$C_{km}^{}(yx_{<k}y_k):=\underset{x_k}{}[c(x_k)+C_{k+1,m}^{}(yx_{1:k})]\mu (yx_{<k}y\underset{¯}{x}_k)$$
(5)
Now we ask about how $`p^{}`$ chooses $`y_k`$. It should choose $`y_k`$ as to maximize the future credit. So the expected number of errors in cycles $`k`$ to $`m`$ given $`yx_{<k}`$ and $`y_k`$ chosen by $`p^{}`$ is $`C_{km}^{}(yx_{<k}):=\mathrm{max}_{y_k}C_{km}^{}(yx_{<k}y_k)`$. Together with the induction start
$$C_{m+1,m}^{}(yx_{1:m}):=\mathrm{\hspace{0.33em}0}$$
(6)
$`C_{km}`$ is completely defined. We might summarize one cycle into the formula
$$C_{km}^{}(yx_{<k})=\underset{y_k}{\mathrm{max}}\underset{x_k}{}[c(x_k)+C_{k+1,m}^{}(yx_{1:k})]\mu (yx_{<k}y\underset{¯}{x}_k)$$
(7)
If $`m_k`$ is our horizon function of $`p^{}`$ and $`\dot{y}\dot{x}_{<k}`$ is the actual history in cycle $`k`$, the output $`\dot{y}_k`$ of the system is explicitly given by
$$\dot{y}_k=\underset{y_k}{maxarg}C_{km_k}^{}(\dot{y}\dot{x}_{<k}y_k)=:p^{}(\dot{y}\dot{x}_{<k})$$
(8)
Then the environment responds $`\dot{x}_k`$ with probability $`\mu (\dot{y}\dot{x}_{<k}\dot{y}\underset{¯}{\overset{\dot{}}{x}}_k)`$. Then cycle $`k+1`$ starts. We might unfold the recursion (7) further and give $`\dot{y}_k`$ non-recursive as
$$\dot{y}_k=\underset{y_k}{maxarg}\underset{x_k}{}\underset{y_{k+1}}{\mathrm{max}}\underset{x_{k+1}}{}\mathrm{}\underset{y_{m_k}}{\mathrm{max}}\underset{x_{m_k}}{}(c(x_k)+\mathrm{}+c(x_{m_k}))\mu (\dot{y}\dot{x}_{<k}y\underset{¯}{x}_{k:m_k})$$
(9)
This has a direct interpretation: the probability of inputs $`x_{k:m_k}`$ in cycle $`k`$ when the system outputs $`y_{k:m_k}`$ and the actual history is $`\dot{y}\dot{x}_{<k}`$ is $`\mu (\dot{y}\dot{x}_{<k}y\underset{¯}{x}_{k:m_k})`$. The future credit in this case is $`c(x_k)+\mathrm{}+c(x_{m_k})`$. The best expected credit is obtained by averaging over the $`x_i`$ ($`sum_{x_i}`$) and maximizing over the $`y_i`$. This has to be done in chronological order to correctly incorporate the dependency of $`x_i`$ and $`y_i`$ on the history. This is essentially the expectimax algorithm/sequence . The AI$`\mu `$ model is optimal in the sense that no other policy leads to higher expected credit.
These explicit as well as recursive definitions of the AI$`\mu `$ model are more index intensive as compared to the functional form but are more suitable for explicit calculations.
### Equivalence of Functional and Iterative AI model:
The iterative environmental probability $`\mu `$ is given by the functional form in the following way,
$$\mu (y\underset{¯}{x}_{1:k})=\underset{q:q(y_{1:k})=x_{1:k}}{}\mu (q)$$
(10)
as is easy to see. We will prove the equivalence of (2) and (8) only for $`k=2`$ and $`m_2=3`$. The proof of the general case is completely analog except that the notation becomes quite messy.
Let us first evaluate (1) for fixed $`\dot{y}_1\dot{x}_1`$ and some $`p\dot{P}_2`$, i.e. $`p(\dot{x}_1)=\dot{y}_1y_2`$ for some $`y_2`$. If the next input to the system is $`x_2`$, $`p`$ will respond with $`p(\dot{x}_1x_2)=\dot{y}_1y_2y_3`$ for some $`y_3`$ depending on $`x_2`$. We write $`y_3(x_2)`$ in the following<sup>5</sup><sup>5</sup>5Dependency on dotted words like $`\dot{x}_1`$ is not shown as the dotted words are fixed.. The numerator of (1) simplifies to
$$\underset{q\dot{Q}_2}{}\mu (q)C_{23}(p,q)=\underset{q:q(\dot{y}_1)=\dot{x}_1}{}\mu (q)C_{23}(p,q)=\underset{x_2x_3}{}(c(x_2)+c(x_3))\underset{q:q(\dot{y}_1y_2y_3(x_2))=\dot{x}_1x_2x_3}{}\mu (q)=$$
$$=\underset{x_2x_3}{}(c(x_2)+c(x_3))\mu (\dot{y}_1\underset{¯}{\overset{\dot{}}{x}}_1y_2\underset{¯}{x}_2y_3(x_2)\underset{¯}{x}_3)$$
In the first equality we inserted the definition of $`\dot{Q}_2`$. In the second equality we split the sum over $`q`$ by first summing over $`q`$ with fixed $`x_2x_3`$. This allows us to pull $`C_{23}=c(x_2)+c(x_3)`$ out of the inner sum. Then we sum over $`x_2x_3`$. Further, we have inserted $`p`$, i.e. replaced $`p`$ by $`y_2`$ and $`y_3()`$. In the last equality we used (10). The denominator reduces to
$$\underset{q\dot{Q}_2}{}\mu (q)=\underset{q:q(\dot{y}_1)=\dot{x}_1}{}\mu (q)=\mu (\dot{y}_1\underset{¯}{\overset{\dot{}}{x}}_1).$$
For the quotient we get
$$C_{23}(p|\dot{y}_1\dot{x}_1)=\underset{x_2x_3}{}(c(x_2)+c(x_3))\mu (\dot{y}_1\dot{x}_1y_2\underset{¯}{x}_2y_3(x_2)\underset{¯}{x}_3)$$
We have seen that the relevant behaviour of $`p\dot{P}_2`$ in cycle 2 and 3 is completely determined by $`y_2`$ and the function $`y_3()`$
$$\underset{p\dot{P}_2}{\mathrm{max}}C_{23}(p|\dot{y}_1\dot{x}_1)=\underset{y_2}{\mathrm{max}}\underset{y_3()}{\mathrm{max}}\underset{x_2x_3}{}(c(x_2)+c(x_3))\mu (\dot{y}_1\dot{x}_1y_2\underset{¯}{x}_2y_3(x_2)\underset{¯}{c}_3)=$$
$$=\underset{y_2}{\mathrm{max}}\underset{x_2}{}\underset{y_3}{\mathrm{max}}\underset{x_3}{}(c(x_2)+c(x_3))\mu (\dot{y}_1\dot{x}_1y_2\underset{¯}{x}_2y_3\underset{¯}{x}_3)$$
In the last equality we have used the fact that the functional minimization over $`y_3()`$ reduces to a simple minimization over the word $`y_3`$ when interchanging with the sum over its arguments $`(\mathrm{max}_{y_3()}_{x_2}_{x_2}\mathrm{max}_{y_3})`$. In the functional case $`\dot{y}_2`$ is therefore determined by
$$\dot{y}_2=\underset{y_2}{maxarg}\underset{x_2}{}\underset{y_3}{\mathrm{max}}\underset{x_3}{}(c(x_2)+c(x_3))\mu (\dot{y}_1\dot{x}_1y_2\underset{¯}{x}_2y_3\underset{¯}{x}_3)$$
This is identical to the iterative definition (9) with $`k=2`$ and $`m_2=3`$ $``$.
### Factorizable $`\mu `$:
Up to now we have made no restrictions on the form of the prior probability $`\mu `$ apart from being a chronological probability distribution. On the other hand, we will see that, in order to prove rigorous credit bounds, the prior probability must satisfy some separability condition to be defined later. Here we introduce some very strong form of separability, when $`\mu `$ factorizes into products. We start with a factorization into two factors. Let us assume that $`\mu `$ is of the form
$$\mu (y\underset{¯}{x}_{1:n})=\mu _1(y\underset{¯}{x}_{<l})\mu _2(y\underset{¯}{x}_{l:n})$$
(11)
for some fixed $`l`$ and sufficiently large $`nm_k`$. For this $`\mu `$ the output $`\dot{y}_k`$ in cycle $`k`$ of the AI$`\mu `$ system (9) for $`kl`$ depends on $`\dot{y}\dot{x}_{l:k1}`$ and $`\mu _2`$ only and is independent of $`\dot{y}\dot{x}_{<l}`$ and $`\mu _1`$. This is easily seen when inserting
$$\mu (\dot{y}\dot{x}_{<k}y\underset{¯}{x}_{k:m_k})=\underset{1}{\underset{}{\mu _1(\dot{y}\dot{x}_{<l})}}\mu _2(\dot{y}\dot{x}_{l:k1}y\underset{¯}{x}_{k:m_k})$$
(12)
into (9). For $`k<l`$ the output $`\dot{y}_k`$ depends on $`\dot{y}\dot{x}_{<k}`$ (this is trivial) and $`\mu _1`$ only (trivial if $`m_k<l`$) and is independent of $`\mu _2`$. The non-trivial case, where the horizon $`m_kl`$ reaches into the region $`\mu _2`$, can be proved as follows (we abbreviate $`m:=m_k`$ in the following). Inserting (11) into the definition of $`C_{lm}^{}(yx_{<l})`$ the factor $`\mu _1`$ is $`1`$ as in (12). We abbreviate $`C_{lm}^{}:=C_{lm}^{}(yx_{<l})`$ as it is independent of its arguments. One can decompose
$$C_{km}^{}(yx_{<k})=C_{k,l1}^{}(yx_{<k})+C_{lm}^{}$$
(13)
For $`k=l`$ this is true because the first term on the r.h.s. is zero. For $`k<l`$ we prove the decomposition by induction from $`k+1`$ to $`k`$.
$$C_{km}^{}(yx_{<k})=\underset{y_k}{\mathrm{max}}\underset{x_k}{}[c(x_k)+C_{k+1,l1}^{}(yx_{1:k})+C_{lm}^{}]\mu _1(yx_{<k}y\underset{¯}{x}_k)=$$
$$=\underset{y_k}{\mathrm{max}}\left[\underset{x_k}{}(c(x_k)+C_{k+1,l1}^{}(yx_{<k}))\mu _1(yx_{<k}y\underset{¯}{x}_k)+C_{lm}^{}\right]=$$
$$=C_{k,l1}^{}(yx_{<k})+C_{lm}^{}$$
Inserting (13), valid for $`k`$ by induction hypothesis, into (7) gives the first equality. In the second equality we have performed the $`x_k`$ sum for the $`C_{lm}^{}\mu _1`$ term which is now independent of $`y_k`$. It can therefore be pulled out of $`\mathrm{max}_{y_k}`$. In the last equality we used again the definition (7). This completes the induction step and proves (13) for $`k<l`$. $`\dot{y}_k`$ can now be represented as
$$\dot{y}_k=\underset{y_k}{maxarg}C_{km}^{}(\dot{y}\dot{x}_{<k}y_k)=\underset{y_k}{maxarg}C_{k,l1}^{}(\dot{y}\dot{x}_{<k}y_k)$$
(14)
where (8) and (13) and the fact that an additive constant $`C_{lm}^{}`$ does not change $`maxarg_{y_k}`$ has been used. $`C_{k,l1}^{}(\dot{y}\dot{x}_{<k}y_k)`$ and hence $`\dot{y}_k`$ is independent of $`\mu _2`$ for $`k<l`$. Note, that $`\dot{y}_k`$ is also independent of the choice of $`m`$, as long as $`ml`$.
In the general case the cycles are grouped into independent episodes $`r=1,2,3,\mathrm{}`$, where each episode $`r`$ consists of the cycles $`k=n_r+1,\mathrm{},n_{r+1}`$ for some $`0=n_0<n_1<\mathrm{}<n_s=n`$:
$$\mu (y\underset{¯}{x}_{1:n})=\underset{r=0}{\overset{s1}{}}\mu _r(y\underset{¯}{x}_{n_r+1:n_{r+1}})$$
(15)
In the simplest case, when all episodes have the same length $`l`$ then $`n_r=rl`$. $`\dot{y}_k`$ depends on $`\mu _r`$ and $`x`$ and $`y`$ of episode $`r`$ only, with $`r`$ such that $`n_r<kn_{r+1}`$.
$$\dot{y}_k=\underset{y_k}{maxarg}\underset{x_k}{}\mathrm{}\underset{y_t}{\mathrm{max}}\underset{x_t}{}(c(x_k)+\mathrm{}+c(x_t))\mu _r(\dot{y}\dot{x}_{n_r+1:k1}y\underset{¯}{x}_{k:n_{r+1}})$$
(16)
with $`t:=\mathrm{min}\{m_k,n_{r+1}\}`$. The different episodes are completely independent in the following sense. The inputs $`x_k`$ of different episodes are statistically independent and depend only on $`y_k`$ of the same episode. The outputs $`y_k`$ depend on the $`x`$ and $`y`$ of the corresponding episode $`r`$ only, and are independent of the actual I/O of the other episodes.
If all episodes have a length of at most $`l`$, i.e. $`n_{r+1}n_rl`$ and if we choose the horizon $`h_k`$ to be at least $`l`$, then $`m_kk+l1n_r+ln_{r+1}`$ and hence $`t=n_{r+1}`$ independent of $`m_k`$. This means that for factorizable $`\mu `$ there is no problem in taking the limit $`m_k\mathrm{}`$. Maybe this limit can also be performed in the more general case of a separable $`\mu `$. The (problem of the) choice of $`m_k`$ will be discussed in more detail later.
Although factorizable $`\mu `$ are too restrictive to cover all AI problems, it often occurs in practice in the form of repeated problem solving, and hence, is worth being studied. For example, if the system has to play games like chess repeatedly, or has to minimize different functions, the different games/functions might be completely independent, i.e. the environmental probability factorizes, where each factor corresponds to a game/function minimization. For details, see the appropriate sections on strategic games and function minimization.
Further, for factorizable $`\mu `$ it is probably easier to derive suitable credit bounds for the universal AI$`\xi `$ model defined in the next section, than for the general separable case which will be introduced later. This could be a first step toward a definition and proof for the general case of separable problems. One goal of this paragraph was to show, that the notion of a factorizable $`\mu `$ could be the first step toward a definition and analysis of the general case of separable $`\mu `$.
### Constants and Limits:
We have in mind a universal system with complex interactions that is as least as intelligent and complex as a human being. One might think of a system whose input $`y_k`$ comes from a digital video camera, the output $`x_k`$ is some image to a monitor<sup>6</sup><sup>6</sup>6Humans can only simulate a screen as output device by drawing pictures., only for the valuation we might restrict to the most primitive binary one, i.e. $`c_kIB`$. So we think of the following constant sizes:
$$\begin{array}{ccccccccc}1& & l(y_kx_k)& & k& & T& & |Y\times X|\\ 1& & 2^{16}& & 2^{24}& & 2^{32}& & 2^{65536}\end{array}$$
The first two limits say that the actual number $`k`$ of inputs/outputs should be reasonably large, compared to the typical size $`l`$ of the input/output words, which itself should be rather sizeable. The last limit expresses the fact that the total lifetime $`T`$ (number of I/O cycles) of the system is far too small to allow every possible input to occur, or to try every possible output, or to make use of identically repeated inputs or outputs. We do not expect any useful outputs for $`kl`$. More interesting than the lengths of the inputs is the complexity $`K(x_1\mathrm{}x_k)`$ of all inputs until now, to be defined later. The environment is usually not ”perfect”. The system could either interact with a non-perfect human or tackle a non-deterministic world (due to quantum mechanics or chaos) world<sup>7</sup><sup>7</sup>7Whether there exist stochastic processes at all is a difficult question. At least the quantum indeterminacy comes very close to it.. In either case, the sequence contains some noise, leading to $`Klk`$. The complexity of the probability distribution of the input sequence is something different. We assume that this noisy world operates according to some simple computable, though not finite rules. $`K(\mu _k)lk`$, i.e. the rules of the world can be highly compressed. On the other hand, there may appear new aspects of the environment for $`k\mathrm{}`$ causing a non-bounded $`K(\mu _k)`$.
In the following we never use these limits, except when explicitly stated. In some simpler models and examples the size of the constants will even violate these limits (e.g. $`l(x_k)=l(y_k)=1`$), but it is the limits above that the reader should bear in mind. We are only interested in theorems which do not degenerate under the above limits.
### Sequential decision theory:
In the following we clarify the connection of (7) and (8) to sequential decision theory and discuss similarities and differences. With probability $`M_{ij}^a`$, the system under consideration should reach (environmental) state $`iS`$ when taking action $`aA`$ depending on the current state $`jS`$. If the system receives reward $`R(i)`$, the optimal policy $`p^{}`$, maximizing expected utility (defined as sum of future rewards), and the utility $`U(i)`$ of policy $`p^{}`$ are
$$p^{}(i)=\underset{a}{maxarg}\underset{j}{}M_{ij}^aU(j),U(i)=R(i)+\underset{a}{\mathrm{max}}\underset{j}{}M_{ij}^aU(j)$$
(17)
See for details and further references. Let us identify
$$\begin{array}{c}S=(Y\times X)^{},A=Y,a=y_k,M_{ij}^a=\mu (yx_{<k}y\underset{¯}{x}_k),\\ i=yx_{<k},R(i)=c(x_{k1}),U(i)=C_{k1,m}^{}(yx_{<k})=c(x_{k1})+C_{km}^{}(yx_{<k}),\\ j=yx_{1:k},R(j)=c(x_k),U(j)=C_{km}^{}(yx_{1:k})=c(x_k)+C_{k+1,m}^{}(yx_{1:k}),\end{array}$$
where we further set $`M_{ij}^a=0`$ if $`i`$ is not a starting substring of $`j`$ or if $`ay_k`$. This ensures the sum over $`j`$ in (17) to reduce to a sum over $`x_k`$. If we set $`m_k=m`$ and use $`C_{km}^{}(yx_{<k}y_k)=_{x_k}C_{km}^{}(yx_{1:k})`$ in (8), it is easy to see that (17) coincides with (7) and (8).
Note that despite of this formal equivalence, we were forced to use the complete history $`yx_{<k}`$ as environmental state $`i`$. The AI$`\mu `$ model neither assumes stationarity, nor Markov property, nor complete accessibility of the environment, as any assumption would restrict the applicability of AI$`\mu `$. The consequence is that every state occurs at most once in the lifetime of the system. Every moment in the universe is unique! Even if the state space could be identified with the input space $`X`$, inputs would usually not occur twice by assumption $`k|X|`$, made in the last subsection. Further, there is no (obvious) universal similarity relation on $`(X\times Y)^{}`$ allowing an effective reduction of the size of the state space. Although many algorithms (e.g. value and policy iteration) have problems in solving (17) for huge or infinite state spaces in practice, there is no principle problem in determining $`p^{}`$ and $`U`$, as long as $`\mu `$ is known and $`|X|`$, $`|Y|`$ and $`m`$ are finite.
Things dramatically change if $`\mu `$ is unknown. Reinforcement learning algorithms are commonly used in this case to learn the unknown $`\mu `$. They succeed if the state space is either small or has effectively been made small by so called generalization techniques. In any case, the solutions are either ad hoc, or work in restricted domains only, or have serious problems with state space exploration versus exploitation, or have non-optimal learning rate. There is no universal and optimal solution to this problem so far. In the next section we present a new model and argue that it formally solves all these problems in an optimal way. It will not concern with learning of $`\mu `$ directly. All we do is to replace the true prior probability $`\mu `$ by a universal probability $`\xi `$, which is shown to converge to $`\mu `$ in a sense.
## 4 The Universal AI$`\xi `$ Model
### Induction and Algorithmic Information theory:
One very important and highly non-trivial aspect of intelligence is inductive inference. Before formulating the AI$`\xi `$ model, a short introduction to the history of induction is given, culminating into the sequence prediction theory by Solomonoff. We emphasize only those aspects which will be of importance for the development of our universal AI$`\xi `$ model.
Simply speaking, induction is the process of predicting the future from the past or, more precisely, it is the process of finding rules in (past) data and using these rules to guess future data. On the one hand, induction seems to happen in every day life by finding regularities in past observations and using them to predict the future. On the other hand, this procedure seems to add knowledge about the future from past observations. But how can we know something about the future? This dilemma and the induction principle in general have a long philosophical history
* Hume’s negation of Induction (1711-1776) ,
* Epicurus’ principle of multiple explanations (342?-270? BC),
* Occams’ razor (simplicity) princple (1290?-1349?),
* Bayes’ rule for conditional probabilites
and a short but important mathematical history: a clever unification of all these aspects into one formal theory of inductive inference has been done by Solomonoff based on Kolmogorov’s definition of complexity. For an excellent introduction into Kolmogorov complexity and Solomonoff induction one should consult the book of Li and Vitányi . In the rest of this subsection we state all results which are needed or generalized later.
Let us choose some universal prefix Turing machine $`U`$ with unidirectional binary input and output tapes and a bidirectional working tape. We can then define the (prefix) Kolomogorov complexity as the shortest prefix program $`p`$, for which $`U`$ outputs $`x=x_{1:n}`$ with $`x_iIB`$:
$$K(x):=\underset{p}{\mathrm{min}}\{l(p):U(p)=x\}$$
The universal semimeasure $`\xi (\underset{¯}{x})`$ is defined as the probability that the output of the universal Turing machine $`U`$ starts with $`x`$ when provided with fair coin flips on the input tape . It is easy to see that this is equivalent to the formal definition
$$\xi (\underset{¯}{x}):=\underset{p:U(p)=x}{}2^{l(p)}$$
(18)
where the sum is over minimal programs $`p`$ for which $`U`$ outputs a string starting with $`x`$. $`U`$ might be non-terminating. As the shortest programs dominate the sum, $`\xi `$ is closely related to $`K(x)`$ ($`\xi (\underset{¯}{x})=2^{K(x)+O(K(l(x))}`$). $`\xi `$ has the important universality property , that it majorizes every computable probability distribution $`\rho `$ up to a multiplicative factor depending only on $`\rho `$ but not on $`x`$:
$$\xi (\underset{¯}{x})\stackrel{\times }{}\mathrm{\hspace{0.33em}2}^{K(\rho )}\rho (\underset{¯}{x}).$$
(19)
A ’$`\times `$’ above an (in)equality denotes (in)equality within a universal multiplicative constant, a ’$`+`$’ above an (in)equality denotes (in)equality within a universal additive constant, both depending only on the choice of the universal reference machine $`U`$. $`\xi `$ itself is not a probability distribution<sup>8</sup><sup>8</sup>8It is possible to normalize $`\xi `$ to a probability distribution as has been done in by giving up the enumerability of $`\xi `$. Error bounds (20) and (22) hold for both definitions.. We have $`\xi (\underset{¯}{x0})+\xi (\underset{¯}{x1})<\xi (\underset{¯}{x})`$ because there are programs $`p`$, which output just $`x`$, neither followed by $`0`$ nor $`1`$. They just stop after printing $`x`$ or continue forever without any further output. We will call a function $`\rho 0`$ with the properties $`\rho (ϵ)1`$ and $`_{x_n}\rho (\underset{¯}{x}_{1:n})\rho (\underset{¯}{x}_{<n})`$ a semimeasure. $`\xi `$ is a semimeasure and (19) actually holds for all enumerable semimeasures $`\rho `$.
(Binary) sequence prediction algorithms try to predict the continuation $`x_n`$ of a given sequence $`x_1\mathrm{}x_{n1}`$. In the following we will assume that the sequences are drawn according to a probability distribution and that the true prior probability of $`x_{1:n}`$ is $`\mu (\underset{¯}{x_1\mathrm{}x_n})`$. The probability of $`x_n`$ given $`x_{<n}`$ hence is $`\mu (x_{<n}\underset{¯}{x}_n)`$. The best possible system predicts the $`x_n`$ with higher probability. Usually $`\mu `$ is unknown and the system can only have some belief $`\rho `$ about the true prior probability $`\mu `$. Let SP$`\rho `$ be a probabilistic sequence predictor, predicting $`x_n`$ with probability $`\rho (x_{<n}\underset{¯}{x}_n)`$. Further we define a deterministic sequence predictor SP$`\mathrm{\Theta }_\rho `$ predicting the $`x_n`$ with higher $`\rho `$ probability. $`\mathrm{\Theta }_\rho (x_{<n}\underset{¯}{x}_n):=1`$ if $`\rho (x_{<n}\underset{¯}{x}_n)>\frac{1}{2}`$ and $`\mathrm{\Theta }_\rho (x_{<n}\underset{¯}{x}_n):=0`$ otherwise. If $`\rho `$ is only a semimeasure the SP$`\rho `$ and SP$`\mathrm{\Theta }_\rho `$ systems might refuse any output in some cycles $`n`$. The SP$`\mathrm{\Theta }_\mu `$ is the best prediction scheme when $`\mu `$ is known.
If $`\rho (x_{<n}\underset{¯}{x}_n)`$ converges quickly to $`\mu (x_{<n}\underset{¯}{x}_n)`$ the number of additional prediction errors introduced by using $`\mathrm{\Theta }_\rho `$ instead of $`\mathrm{\Theta }_\mu `$ for prediction should be small in some sense. Now the universal probability $`\xi `$ comes into play as it has been proved by Solomonoff that the $`\mu `$ expected Euclidean distance betweewn $`\xi `$ and $`\mu `$ is finite
$$\underset{k=1}{\overset{\mathrm{}}{}}\underset{x_{1:k}}{}\mu (\underset{¯}{x}_{1:k})(\xi (x_{<k}\underset{¯}{x}_k)\mu (x_{<k}\underset{¯}{x}_k))^2\stackrel{+}{<}\frac{1}{2}\mathrm{ln}2K(\mu )$$
(20)
The ’$`+`$’ atop ’$`<`$’ means up to additive terms of order 1. So indeed the difference does tend to zero, i.e. $`\xi (x_{<n}\underset{¯}{x}_n)\stackrel{n\mathrm{}}{}\mu (x_{<n}\underset{¯}{x}_n)`$ with $`\mu `$ probability $`1`$ for any computable probability distribution $`\mu `$. The reason for the astonishing property of a single (universal) function to converge to any computable probability distribution lies in the fact that the set of $`\mu `$ random sequences differ for different $`\mu `$. The universality property (19) is the central ingredient for proving (20).
Let us define the total number of expected erroneous predictions the SP$`\rho `$ system makes for the first $`n`$ bits
$$E_{n\rho }:=\underset{k=1}{\overset{n}{}}\underset{x_{1:k}}{}\mu (\underset{¯}{x}_{1:k})(1\rho (x_{<k}\underset{¯}{x}_k))$$
(21)
The SP$`\mathrm{\Theta }_\mu `$ system is best in the sense that $`E_{n\mathrm{\Theta }_\mu }E_{n\rho }`$ for any $`\rho `$. In it has been shown that SP$`\mathrm{\Theta }_\xi `$ is not much worse
$$E_{n\mathrm{\Theta }_\xi }E_{n\rho }H+\sqrt{4E_{n\rho }H+H^2}=O(\sqrt{E_{n\rho }}),H\stackrel{+}{<}\mathrm{ln}2K(\mu )$$
(22)
with the tightest bound for $`\rho =\mathrm{\Theta }_\mu `$. For finite $`E_{\mathrm{}\mathrm{\Theta }_\mu }`$, $`E_{\mathrm{}\mathrm{\Theta }_\xi }`$ is finite too. For infinite $`E_{\mathrm{}\mathrm{\Theta }_\mu }`$, $`E_{n\mathrm{\Theta }_\xi }/E_{n\mathrm{\Theta }_\mu }\stackrel{n\mathrm{}}{}1`$ with rapid convergence. One can hardly imagine any better prediction algorithm without extra knowledge about the environment. In , (20) and (22) have been generalized from binary to arbitrary alphabet. Apart from computational aspects, which are of course very important, the problem of sequence prediction could be viewed as essentially solved.
### Definition of the AI$`\xi `$ Model:
We have developed enough formalism to suggest our universal AI$`\xi `$ model<sup>9</sup><sup>9</sup>9Speak ’aixi’ and write AIXI without Greek letters.. All we have to do is to suitably generalize the universal semimeasure $`\xi `$ from the last subsection and replace the true but unknown prior probability $`\mu ^{AI}`$ in the AI$`\mu `$ model by this generalized $`\xi ^{AI}`$. In what sense this AI$`\xi `$ model is universal will be discussed later.
In the functional formulation we define the universal probability $`\xi ^{AI}`$ of an environment $`q`$ just as $`2^{l(q)}`$
$$\xi (q):=\mathrm{\hspace{0.33em}2}^{l(q)}$$
The definition could not be easier<sup>10</sup><sup>10</sup>10It is not necessary to use $`2^{K(q)}`$ or something similar as some reader may expect at this point. The reason is that for every program $`q`$ there exists a functionally equivalent program $`q^{}`$ with $`K(q^{})=l(q^{})`$.!<sup>11</sup><sup>11</sup>11Here and later we identify objects with their coding relative to some fixed Turing machine $`U`$. For example, if $`q`$ is a function $`K(q):=K(q)`$ with $`q`$ being a binary coding of $`q`$ such that $`U(q,y):=q(y)`$. On the other hand, if $`q`$ already is a binary string we define $`q(y):=U(q,y)`$. Collecting the formulas of section 2 and replacing $`\mu (q)`$ by $`\xi (q)`$ we get the definition of the AI$`\xi `$ system in functional form. Given the history $`\dot{y}\dot{x}_{<k}`$ the functional AI$`\xi `$ system outputs
$$\dot{y}_k:=\underset{y_k}{maxarg}\underset{p:p(\dot{x}_{<k})=\dot{y}_{<k}y_k}{\mathrm{max}}\underset{q:q(\dot{y}_{<k})=\dot{x}_{<k}}{}2^{l(q)}C_{km_k}(p,q)$$
(23)
in cycle $`k`$, where $`C_{km_k}(p,q)`$ is the total credit of cycles $`k`$ to $`m_k`$ when system $`p`$ interacts with environment $`q`$. We have dropped the denominator $`_q\mu (q)`$ from (1) as it is independent of the $`p\dot{P}_k`$ and a constant multiplicative factor does not change $`maxarg`$.
For the iterative formulation the universal probability $`\xi `$ can be obtained by inserting the functional $`\xi (q)`$ into (10)
$$\xi (y\underset{¯}{x}_{1:k})=\underset{q:q(y_{1:k})=x_{1:k}}{}2^{l(q)}$$
(24)
Replacing $`\mu `$ by $`\xi `$ in (9) the iterative AI$`\xi `$ system outputs
$$\dot{y}_k=\underset{y_k}{maxarg}\underset{x_k}{}\underset{y_{k+1}}{\mathrm{max}}\underset{x_{k+1}}{}\mathrm{}\underset{y_{m_k}}{\mathrm{max}}\underset{x_{m_k}}{}(c(x_k)+\mathrm{}+c(x_{m_k}))\xi (\dot{y}\dot{x}_{<k}y\underset{¯}{x}_{k:m_k})$$
(25)
in cycle $`k`$ given the history $`\dot{y}\dot{x}_{<k}`$.
One subtlety has been passed over. Like in the SP case, $`\xi `$ is not a probability distribution but satisfies only the weaker inequalities
$$\underset{x_n}{}\xi (y\underset{¯}{x}_{1:n})\xi (y\underset{¯}{x}_{<n}),\xi (ϵ)\mathrm{\hspace{0.33em}1}$$
(26)
Note, that the sum on the l.h.s. is not independent of $`y_n`$ unlike for chronological probability distributions. Nevertheless, it is bounded by something (the r.h.s) which is independent of $`y_n`$. The reason is that the sum in (24) runs over (partial recursive) chronological functions only and the functions $`q`$ which satisfy $`q(y_{1:n})=x_{<n}`$ are a subset of the functions satisfying $`q(y_{<n})=x_{<n}`$. Therefore we will in general call functions satisfying (26) chronological semimeasures. The important point is that the conditional probabilities (4) are $`1`$ like for true probability distributions.
The equivalence of the functional and iterative AI model proven in section 3 is true for every chronological semimeasure $`\rho `$, esp. for $`\xi `$, hence we can talk about the AI$`\xi `$ model in this respect. It (slightly) depends on the choice of universal Turing machine. $`l(q)`$ is defined only up to an additive constant. It also depends on the choice of $`X=C\times X^{}`$ and $`Y`$, but we do not expect any bias when the spaces are chosen sufficiently simple, e.g. all strings of length $`2^{16}`$. Choosing $`IN`$ as word space would be optimal, but whether the maxima (suprema) exist in this case, has to be shown beforehand. The only non-trivial dependence is on the horizon function $`m_k`$ which will be discussed later. So apart from $`m_k`$ and unimportant details the AI$`\xi `$ system is uniquely defined by (23) or (25). It doesn’t depend on assumptions about the environment apart from being generated from some computable (but unknown!) probability distribution.
### Universality of $`\xi ^{AI}`$:
In which sense the AI$`\xi `$ model is optimal will be clarified later. In this and the next two subsections we show that $`\xi ^{AI}`$ defined in (24) is universal and converges to $`\mu ^{AI}`$ analog to the SP case (19) and (20). The proofs are generalizations from the SP case. The $`y`$ are pure spectators and cause no difficulties in the generalization. The replacement of the binary alphabet $`IB`$ used in SP by the (possibly infinite) alphabet $`X`$ is possible, but needs to be done with care. In (19) $`U(p)=x`$ produces strings starting with $`x`$, whereas in (24) we can demand $`q`$ to output exactly $`n`$ words $`x_{1:n}`$ as $`q`$ knows $`n`$ from the number of input words $`y_1\mathrm{}y_n`$. For proofs of (19) and (20) see and .
There is an alternative definition of $`\xi `$ which coincides with (24) within a multiplicative constant of $`O(1)`$,
$$\xi (y\underset{¯}{x}_{1:n})\stackrel{\times }{=}\underset{\rho }{}2^{K(\rho )}\rho (y\underset{¯}{x}_{1:n})$$
(27)
where the sum runs over all enumerable chronological semimeasures. The $`2^{K(\rho )}`$ weighted sum over probabilistic environments $`\rho `$, coincides with the sum over $`2^{l(q)}`$ weighted deterministic environments $`q`$, as will be proved below. In the next subsection we show that an enumeration of all enumerable functions can be converted into an enumeration of enumerable chronological semimeasures $`\rho `$. $`K(\rho )`$ is co-enumerable, therefore $`\xi `$ defined in (27) is itself enumerable. The representation (24) is also enumerable. As $`_\rho 2^{K(\rho )}1`$ and the $`\rho ^{}s`$ satisfy (26), $`\xi `$ is a chronological semimeasure as well. If we pick one $`\rho `$ in (27) we get the universality property ”for free”
$$\xi (y\underset{¯}{x}_{1:n})\stackrel{\times }{}\mathrm{\hspace{0.33em}2}^{K(\rho )}\rho (y\underset{¯}{x}_{1:n})$$
(28)
$`\xi `$ is a universal element in the sense of (28) in the set of all enumerable chronological semimeasures.
To prove universality of $`\xi `$ in the form (24) we have to show that for every enumerable chronological semimeasure $`\rho `$ there exists a Turing machine $`T`$ with
$$\rho (y\underset{¯}{x}_{1:n})=\underset{q:T(qy_{1:n})=x_{1:n}}{}2^{l(q)}\text{and}l(T)\stackrel{+}{=}K(\rho ).$$
(29)
This will not be done here. Given $`T`$ the universality of $`\xi `$ follows from
$$\xi (y\underset{¯}{x}_{1:n})=\underset{q:U(qy_{1:n})=x_{1:n}}{}2^{l(q)}\underset{q:U(Tq^{}y_{1:n})=x_{1:n}}{}\mathrm{\hspace{0.33em}2}^{l(Tq^{})}=\mathrm{\hspace{0.33em}2}^{l(T)}\underset{q:T(q^{}y_{1:n})=x_{1:n}}{}2^{l(q^{})}\stackrel{\times }{=}2^{K(\rho )}\rho (y\underset{¯}{x}_{1:n})$$
The first equality and (24) are identical by definition. In the inequality we have restricted the sum over all $`q`$ to $`q`$ of the form $`q=Tq^{}`$. The third relation is true as running $`U`$ on $`Tz`$ is a simulation of $`T`$ on $`z`$. The last equality follows from (29). All enumerable, universal, chronological semimeasures coincide up to a multiplicative constant, as they mutually dominate each other. Hence, definitions (24) and (27) are, indeed, equivalent.
### Converting general functions into chronological semi-measures:
To complete the proof of the universality (28) of $`\xi `$ we need to convert enumerable functions $`\psi :IB^{}IR^+`$ into enumerable chronological semi-measures $`\rho :(Y\times X)^{}IR^+`$ with certain additional properties. Every enumerable function like $`\psi `$ and $`\rho `$ can be approximated from below by definition<sup>12</sup><sup>12</sup>12Defining enumerability as the supremum of total primitive recursive functions is more suitable for our purpose than the equivalent definition as a limit of monotone increasing partial recursive functions. In terms of Turing machines, the recursion parameter is the time after which a computation is terminated. by primitive recursive functions $`\phi :IB^{}\times INIQ^+`$ and $`\varphi :(Y\times X)^{}\times INIQ^+`$ with $`\psi (s)=sup_t\phi (s,t)`$ and $`\rho (s)=sup_t\varphi (s,t)`$ and recursion parameter $`t`$. For arguments of the form $`s=yx_{1:n}`$ we recursively (in $`n`$) construct $`\varphi `$ from $`\phi `$ as follows:
$`\phi ^{}(yx_{1:n},t)`$ $`:=`$ $`\{\begin{array}{cc}\hfill \phi (yx_{1:n},t)\text{for}& x_n<t\hfill \\ \hfill 0\text{for}& x_nt\hfill \end{array},\phi ^{}(ϵ,t):=\phi (ϵ,t)`$ (32)
$`\varphi (ϵ,t)`$ $`:=`$ $`\underset{0it}{\mathrm{max}}\left\{\phi ^{}(ϵ,i):\phi ^{}(ϵ,i)1\right\}`$ (33)
$`\varphi (y\underset{¯}{x}_{1:n},t)`$ $`:=`$ $`\underset{0it}{\mathrm{max}}\left\{\phi ^{}(yx_{1:n},i):_{x_n}\phi ^{}(yx_{1:n},i)\varphi (y\underset{¯}{x}_{<n},t)\right\}`$ (34)
With $`x_n<t`$ we mean that the natural number associated with string $`x_n`$ is smaller than $`t`$. According to (32) with $`\phi `$ also $`\phi ^{}`$ as well as $`_{x_n}\phi ^{}`$ are primitive recursive functions. Further, if we allow $`t=0`$ we have $`\phi ^{}(s,0)=0`$. This ensures that $`\varphi `$ is a total function.
In the following we prove by induction over $`n`$ that $`\varphi `$ is a primitive recursive chronological semimeasure monotone increasing in $`t`$. All necessary properties hold for $`n=0`$ ($`yx_{1:0}=ϵ`$) according to (33). For general $`n`$ assume that the induction hypothesis is true for $`\varphi (y\underset{¯}{x}_{<n},t)`$. We can see from (34) that $`\varphi (y\underset{¯}{x}_{1:n},t)`$ is monotone increasing in $`t`$. $`\varphi `$ is total as $`\phi ^{}(yx_{1:n},i=0)=0`$ satisfies the inequality. By assumption $`\varphi (yx_{<n},t)`$ is primitive recursive, hence with $`_{x_n}\phi ^{}`$ also the order relation $`\phi ^{}\varphi `$ is primitive recursive. This ensures that the non-empty finite set $`\{\phi ^{}:\phi ^{}\varphi \}_i`$ and its maximum $`\varphi (y\underset{¯}{x}_{1:n},t)`$ are primitive recursive. Further, $`\varphi (y\underset{¯}{x}_{1:n},t)=\phi ^{}(yx_{1:n},i)`$ for some $`i`$ with $`it`$ independent of $`x_n`$. Thus, $`_{x_n}\varphi (y\underset{¯}{x}_{1:n},t)`$ $`=`$ $`_{x_n}\phi ^{}(yx_{1:n},i)`$ $``$ $`\varphi (y\underset{¯}{x}_{<n},t)`$ which is the condition for $`\varphi `$ being a chronological semimeasure. Inductively we have proved that $`\varphi `$ is indeed a primitive recursive chronological semimeasure monotone increasing in $`t`$.
In the following we show that every (total)<sup>13</sup><sup>13</sup>13Semimeasures are, by definition, total functions. enumerable chronological semimeasure $`\rho `$ can be enumerated by some $`\varphi `$. By definition of enumerability there exist primitive recursive functions $`\stackrel{~}{\phi }`$ with $`\rho (s)=sup_t\stackrel{~}{\phi }(s,t)`$. The function $`\phi (s,t):=(1^1/_t)\mathrm{max}_{i<t}\stackrel{~}{\phi }(s,i)`$ also enumerates $`\rho `$ but has the additional advantage of being strictly monotone increasing in $`t`$.
$`\phi ^{}(yx_{1:n},\mathrm{})=\phi (yx_{1:n},\mathrm{})=\rho (yx_{1:n})`$ by definition (32). $`\varphi (ϵ,t)=\phi ^{}(ϵ,t)`$ by (33) and the fact that $`\phi ^{}(ϵ,i1)<\phi ^{}(ϵ,i)\phi (ϵ,i)\rho (ϵ)1`$, hence $`\varphi (ϵ,\mathrm{})=\rho (ϵ)`$. $`\varphi (y\underset{¯}{x}_{1:n},t)\phi ^{}(yx_{1:n},t)`$ by (34), hence $`\varphi (y\underset{¯}{x}_{1:n},\mathrm{})\rho (y\underset{¯}{x}_{1:n})`$. We prove the opposite direction $`\varphi (y\underset{¯}{x}_{1:n},\mathrm{})\rho (yx_{1:n})`$ by induction over $`n`$. We have
$$\underset{x_n}{}\phi ^{}(yx_{1:n},i)\underset{x_n}{}\phi (yx_{1:n},i)<\underset{x_n}{}\phi (yx_{1:n},\mathrm{})=\underset{x_n}{}\rho (yx_{1:n})\rho (y\underset{¯}{x}_{<n})$$
(35)
The strict monotony of $`\phi `$ and the semimeasure property of $`\rho `$ have been used. By induction hypothesis $`lim_t\mathrm{}\varphi (y\underset{¯}{x}_{<n},t)\rho (y\underset{¯}{x}_{<n})`$ and (35) for sufficiently large $`t`$ we have $`\varphi (y\underset{¯}{x}_{<n},t)>_{x_n}\phi ^{}(yx_{1:n},i)`$. The condition in (34) is, hence, satisfied and therefore $`\varphi (y\underset{¯}{x}_{1:n},t)\phi ^{}(yx_{1:n},i)`$ for sufficiently large $`t`$, especially $`\varphi (y\underset{¯}{x}_{1:n},\mathrm{})\phi ^{}(yx_{1:n},i)`$ for all $`i`$. Taking the limit $`i\mathrm{}`$ we get $`\varphi (y\underset{¯}{x}_{1:n},\mathrm{})\phi ^{}(yx_{1:n},\mathrm{})=\rho (y\underset{¯}{x}_{1:n})`$.
Combining all results, we have shown that the constructed $`\varphi (,t)`$ are primitive recursive chronological semimeasures monotone increasing in $`t`$, which converge to the enumerable chronological semimeasure $`\rho `$. This finally proves the enumerability of the set of enumerable chronological semimeasures.
### Convergence of $`\xi ^{AI}`$ to $`\mu ^{AI}`$:
In the following inequality is proved
$$2\underset{i=1}{\overset{|X|}{}}y_i(y_iz_i)^2\underset{i=1}{\overset{|X|}{}}y_i\mathrm{ln}\frac{y_i}{z_i}\text{with}\underset{i=1}{\overset{|X|}{}}y_i=1,\underset{i=1}{\overset{|X|}{}}z_i1$$
(36)
If we identify $`i=x_k`$ and $`y_i=\mu (yx_{<k}y\underset{¯}{x}_k)`$ and $`z_i=\xi (yx_{<k}y\underset{¯}{x}_k)`$, multiply both sides with $`\mu (y\underset{¯}{x}_{<k})`$, take the sum over $`x_{<k}`$, then the sum over $`k`$ and use Bayes’ rule $`\mu (y\underset{¯}{x}_{<k})\mu (yx_{<k}y\underset{¯}{x}_k)=\mu (y\underset{¯}{x}_{1:k})`$ we get
$$2\underset{k=1}{\overset{n}{}}\underset{x_{1:k}}{}\mu (y\underset{¯}{x}_{1:k})\left(\mu (yx_{<k}\underset{¯}{x}_k)\xi (yx_{<k}\underset{¯}{x}_k)\right)^2\underset{k=1}{\overset{n}{}}\underset{x_{1:k}}{}\mu (y\underset{¯}{x}_{1:k})\mathrm{ln}\frac{\mu (yx_{<k}\underset{¯}{x}_k)}{\xi (yx_{<k}\underset{¯}{x}_k)}=\mathrm{}$$
(37)
In the r.h.s. we can replace $`_{x_{1:k}}\mu (y\underset{¯}{x}_{1:k})`$ by $`_{x_{1:n}}\mu (y\underset{¯}{x}_{1:n})`$ as the argument of the logarithm is independent of $`x_{k+1:n}`$. The $`k`$ sum can now be brought into the logarithm and converts to a product. Using Bayes’ rule (4) for $`\mu `$ and $`\xi `$ we get
$$\mathrm{}=\underset{x_{1:n}}{}\mu (y\underset{¯}{x}_{1:n})\mathrm{ln}\underset{k=1}{\overset{n}{}}\frac{\mu (yx_{<k}\underset{¯}{x}_k)}{\xi (yx_{<k}\underset{¯}{x}_k)}=\underset{x_{1:n}}{}\mu (y\underset{¯}{x}_{1:n})\mathrm{ln}\frac{\mu (y\underset{¯}{x}_{1:n})}{\xi (y\underset{¯}{x}_{1:n})}\stackrel{+}{<}\mathrm{ln}2K(\mu )$$
(38)
where we have used the universality property (28) of $`\xi `$ in the last step. The main complication for generalizing (20) to (37,38) was the generalization of (36) from $`|X|=2`$ to a general alphabet, the $`y`$ are, again, pure spectators. This will change when we analyze error/credit bounds analog to (22).
(37,38) shows that the $`\mu `$ expected squared difference of $`\mu `$ and $`\xi `$ is finite for computable $`\mu `$. This, in turn, shows that $`\xi (yx_{<k}y\underset{¯}{x}_k)`$ converges to $`\mu (yx_{<k}y\underset{¯}{x}_k)`$ for $`k\mathrm{}`$ with $`\mu `$ probability 1. If we take a finite product of $`\xi ^{}s`$ and use Bayes’ rule, we see that also $`\xi (yx_{<k}y\underset{¯}{x}_{k:k+r})`$ converges to $`\mu (yx_{<k}y\underset{¯}{x}_{k:k+r})`$. More generally, in case of a bounded horizon $`h_k`$, it follows that
$$\xi (yx_{<k}y\underset{¯}{x}_{k:m_k})\stackrel{k\mathrm{}}{}\mu (yx_{<k}y\underset{¯}{x}_{k:m_k})\text{if}h_km_kk+1h_{max}<\mathrm{}$$
(39)
This gives makes us confident that the outputs $`\dot{y}_k`$ of the AI$`\xi `$ model (25) could converge to the outputs $`\dot{y}_k`$ from the AI$`\mu `$ model (9), at least for bounded horizon.
We want to call an AI model universal, if it is $`\mu `$ independent (unbiased, model-free) and is able to solve any solvable problem and learn any learnable task. Further, we call a universal model, universally optimal, if there is no program, which can solve or learn significantly faster (in terms of interaction cycles). As the AI$`\xi `$ model is parameterless, $`\xi `$ converges to $`\mu `$ (39), the AI$`\mu `$ model is itself optimal, and we expect no other model to converge faster to AI$`\mu `$ by analogy to SP (22),
we expect AI$`\xi `$ to be universally optimal.
This is our main claim. In a sense, the intention of the remaining (sub)sections is to define this statement more rigorously and to give further support.
### Intelligence order relation:
We define the $`\xi `$ expected credit in cycles $`k`$ to $`m`$ of a policy $`p`$ similar to (1) and (23). We extend the definition to programs $`p\dot{P}_k`$ which are not consistent with the current history.
$$C_{km}^\xi (p|\dot{y}\dot{x}_{<k}):=\frac{1}{𝒩}\underset{q:q(\dot{y}_{<k})=\dot{x}_{<k}}{}2^{l(q)}C_{km}(\stackrel{~}{p},q)$$
(40)
The normalization $`𝒩`$ is again only necessary for interpreting $`C_{km}`$ as the expected credit but otherwise unneeded. For consistent policies $`p\dot{P}_k`$ we define $`\stackrel{~}{p}:=p`$. For $`p\dot{P}_k`$, $`\stackrel{~}{p}`$ is a modification of $`p`$ in such a way that its output is consistent with the current history $`\dot{y}\dot{x}_{<k}`$, hence $`\stackrel{~}{p}\dot{P}_k`$, but unaltered for the current and future cycles $`k`$. Using this definition of $`C_{km}`$ we could take the maximium over all systems $`p`$ in (23), rather than only the consistent ones.
We call $`p`$ more or equally intelligent than $`p^{}`$ if
$$pp^{}:k\dot{y}\dot{x}_{<k}:C_{km_k}^\xi (p|\dot{y}\dot{x}_{<k})C_{km_k}^\xi (p^{}|\dot{y}\dot{x}_{<k})$$
(41)
i.e. if $`p`$ yields in any circumstance higher $`\xi `$ expected credit than $`p^{}`$. As the algorithm $`p^{}`$ behind the AI$`\xi `$ system maximizes $`C_{km_k}^\xi `$ we have $`p^{}p`$ for all $`p`$. The AI$`\xi `$ model is hence the most intelligent system w.r.t. $``$. $``$ is a universal order relation in the sense that it is free of any parameters (except $`m_k`$) or specific assumptions about the environment. A proof, that $``$ is a reliable intelligence order (what we believe to be true), would prove that AI$`\xi `$ is universally optimal. We could further ask: how useful is $``$ for ordering policies of practical interest with intermediate intelligence, or how can $``$ help to guide toward constructing more intelligent systems with reasonable computation time. An effective intelligence order relation $`^c`$ will be defined in section 10, which is more useful from a practical point of view.
### Credit bounds and separability concepts:
The credits $`C_{km}`$ associated with the AI systems correspond roughly to the negative error measure $`E_{n\rho }`$ of the SP systems. In SP, we were interested in small bounds for the error excess $`E_{n\mathrm{\Theta }_\xi }E_{n\rho }`$. Unfortunately, simple credit bounds for AI$`\xi `$ in terms of $`C_{km}`$ analog to the error bound (22) do not hold. We even have difficulties in specifying what we can expect to hold for AI$`\xi `$ or any AI system which claims to be universally optimal. Consequently, we cannot have a proof if we don’t know what to prove. In SP, the only important property of $`\mu `$ for proving error bounds was its complexity $`K(\mu )`$. We will see that in the AI case, there are no useful bounds in terms of $`K(\mu )`$ only. We either have to study restricted problem classes or consider bounds depending on other properties of $`\mu `$, rather than on its complexity only. In the following, we will exhibit the difficulties by two examples and introduce concepts which may be useful for proving credit bounds. Despite the difficulties in even claiming useful credit bounds, we nevertheless, firmly believe that the order relation (41) correctly formalizes the intuitive meaning of intelligence and, hence, that the AI$`\xi `$ system is universally optimal.
In the following, we choose $`m_k=T`$. We want to compare the true, i.e. $`\mu `$ expected credit $`C_{1T}^\mu `$ of a $`\mu `$ independent universal policy $`p^{best}`$ with any other policy $`p`$. Naively, we might expect the existence of a policy $`p^{best}`$ which maximizes $`C_{1T}^\mu `$, apart from additive corrections of lower order for $`T\mathrm{}`$
$$C_{1T}^\mu (p^{best})C_{1T}^\mu (p)o(\mathrm{})\mu ,p$$
(42)
Note, that the policy $`p^\xi `$ of the AI$`\xi `$ system maximizes $`C_{1T}^\xi `$ by definition ($`p^\xi p`$). As $`C_{1T}^\xi `$ is thought to be a guess of $`C_{1T}^\mu `$, we might expect $`p^{best}=p^\xi `$ to approximately maximize $`C_{1T}^\mu `$, i.e. (42) to hold. Let us consider the problem class (set of environments) $`\{\mu _0,\mu _1\}`$ with $`Y=C=\{0,1\}`$ and $`c_k=\delta _{iy_1}`$ in environment $`\mu _i`$. The first output $`y_1`$ decides whether you go to heaven with all future credits $`c_k`$ being $`1`$ (good) or to hell with all future credits being $`0`$ (bad). It is clear, that if $`\mu _i`$, i.e. $`i`$ is known, the optimal policy $`p^{\mu _i}`$ is to output $`y_1=i`$ in the first cycle with $`C_{1T}^\mu (p^{\mu _i})=T`$. On the other hand, any unbiased policy $`p^{best}`$ independent of the actual $`\mu `$ either outputs $`y_1=1`$ or $`y_1=0`$. Independent of the actual choice $`y_1`$, there is always an environment ($`\mu =\mu _{1y_1}`$) for which this choice is catastrophic ($`C_{1T}^\mu (p^{best})=0`$). No single system can perform well in both environments $`\mu _0`$ and $`\mu _1`$. The r.h.s. of (42) equals $`To(T)`$ for $`p=p^\mu `$. For all $`p^{best}`$ there is a $`\mu `$ for which the l.h.s. is zero. We have shown that no $`p^{best}`$ can satisfy (42) for all $`\mu `$ and $`p`$, so we cannot expect $`p^\xi `$ to do so. Nevertheless, there are problem classes for which (42) holds, for instance SP and CF. For SP, (42) is just a reformulation of (22) with an appropriate choice for $`p^{best}`$ (which differs from $`p^\xi `$, see next section). We expect (42) to hold for all inductive problems in which the environment is not influenced<sup>14</sup><sup>14</sup>14Of course, the credit feedback $`c_k`$ depends on the system’s output. What we have in mind is, like in sequence prediction, that the true sequence is not influenced by the system by the output of the system. We want to call these $`\mu `$, passive or inductive environments. Further, we want to call $`\mu `$ satisfying (42) with $`p^{best}=p^\xi `$ pseudo passive. So we expect inductive $`\mu `$ to be pseudo passive.
Let us give a further example to demonstrate the difficulties in establishing credit bounds. Let $`C=\{0,1\}`$ and $`|Y|`$ be large. We consider all (deterministic) environments in which a single complex output $`y^{}`$ is correct ($`c=1`$) and all others are wrong ($`c=0`$). The problem class $`M`$ is defined by
$$M:=\{\mu :\mu (yx_{<k}y_k\underset{¯}{1})=\delta _{y_ky^{}},y^{}Y,K(y^{})=_{}\mathrm{log}_2|Y|_{}\}$$
There are $`N\stackrel{\times }{=}|Y|`$ such $`y^{}`$. The only way a $`\mu `$ independent policy $`p`$ can find the correct $`y^{}`$ , is by trying one $`y`$ after the other in a certain order. In the first $`N1`$ cycles at most, $`N1`$ different $`y`$ are tested. As there are $`N`$ different possible $`y^{}`$, there is always a $`\mu M`$ for which $`p`$ gives erroneous outputs in the first $`N1`$ cycles. The number of errors are $`E_\mathrm{}pN1\stackrel{\times }{=}|Y|\stackrel{\times }{=}2^{K(y^{})}\stackrel{\times }{=}2^{K(\mu )}`$ for this $`\mu `$. As this is true for any $`p`$, it is also true for the AI$`\xi `$ model, hence $`E_{k\xi }2^{K(\mu )}`$ is the best possible error bound we can expect, which depends on $`K(\mu )`$ only. Actually, we will derive such a bound in section 5 for SP. Unfortunately, as we are mainly interested in the cycle region $`k|Y|\stackrel{\times }{=}2^{K(\mu )}`$ (see section 3) this bound is trivial. There are no interesting bounds depending on $`K(\mu )`$ only, unlike the SP case for deterministic $`\mu `$. Bounds must either depend on additional properties of $`\mu `$ or we have to consider specialized bounds for restricted problem classes. The case of probabilistic $`\mu `$ is similar. Whereas for SP there are useful bounds in terms of $`E_{k\mu }`$ and $`K(\mu )`$, there are no such bounds for AI$`\xi `$. Again, this is not a drawback of AI$`\xi `$ since for no unbiased AI system the errors/credits could be bound in terms of $`K(\mu )`$ and the errors/credits of AI$`\mu `$ only.
There is a way to make use of gross (e.g. $`2^{K(\mu )}`$) bounds. Assume that after a reasonable number of cycles $`k`$, the information $`\dot{x}_{<k}`$ perceived by the AI$`\xi `$ system contains a lot of information about the true environment $`\mu `$. The information in $`\dot{x}_{<k}`$ might be coded in any form. Let us assume that the complexity $`K(\mu |\dot{x}_{<k})`$ of $`\mu `$ under the condition that $`\dot{x}_{<k}`$ is known, is of order 1. Consider a theorem, bounding the sum of credits or of other quantities over cycles $`1\mathrm{}\mathrm{}`$ in terms of $`f(K(\mu ))`$ for a function $`f`$ with $`f(O(1))=O(1)`$, like $`f(n)=2^n`$. Then, there will be a bound for cycles $`k\mathrm{}\mathrm{}`$ in terms of $`f(K(\mu |\dot{x}_{<k}))=O(1)`$. Hence, a bound like $`2^{K(\mu )}`$ can be replaced by small bound $`2^{K(\mu |\dot{x}_{<k})}=O(1)`$ after a reasonable number of cycles. All one has to show/ensure/assume is that enough information about $`\mu `$ is presented (in any form) in the first $`k`$ cycles. In this way, even a gross bound could become useful. In section 8 we use a similar argument to prove that AI$`\xi `$ is able to learn supervised.
In the following, we weaken (42) in the hope of getting a bound applicable to wider problem classes than the passive one. Consider the I/O sequence $`\dot{y}_1\dot{x}_1\mathrm{}\dot{y}_n\dot{x}_n`$ caused by AI$`\xi `$. On history $`\dot{y}\dot{x}_{<k}`$, AI$`\xi `$ will output $`\dot{y}_k\dot{y}_k^\xi `$ in cycle $`k`$. Let us compare this to $`\dot{y}_k^\mu `$ what AI$`\mu `$ would output, still on the same history $`\dot{y}\dot{x}_{<k}`$ produced by AI$`\xi `$. As AI$`\mu `$ maximizes the $`\mu `$ expected credit, AI$`\xi `$ causes lower (or at best equal) $`C_{km_k}^\mu `$, if $`\dot{y}_k^\xi `$ differs from $`\dot{y}_k^\mu `$. Let $`D_{n\mu \xi }:=_{k=1}^n1\delta _{\dot{y}_k^\mu ,\dot{y}_k^\xi }_\mu `$ be the $`\mu `$ expected number of suboptimal choices of AI$`\xi `$, i.e. outputs different from AI$`\mu `$ in the first $`n`$ cycles. One might weigh the deviating cases by their severity. Especially when the $`\mu `$ expected credits $`C_{km_k}^\mu `$ for $`\dot{y}_k^\xi `$ and $`\dot{y}_k^\mu `$ are equal or close to each other, this should be taken into account in the definition of $`D_{n\mu \xi }`$. These details do not matter in the following qualitative discussion. The important difference to (42) is that here we stick on the history produced by AI$`\xi `$ and count a wrong decision as, at most, one error. The wrong decision in the Heaven&Hell example in the first cycle no longer counts as losing $`T`$ credits, but counts as one wrong decision. In a sense, this is fairer. One shouldn’t blame somebody too much who makes a single wrong decision for which he just has too little information available, in order to make a correct decision. The AI$`\xi `$ model would deserve to be called asymptotically optimal, if the probability of making a wrong decision tends to zero, i.e. if
$$D_{n\mu \xi }/n0\text{for}n\mathrm{},\text{i.e.}D_{n\mu \xi }=o(n).$$
(43)
We say that $`\mu `$ can be asymptotically learned (by AI$`\xi `$) if (43) is satisfied. We claim that AI$`\xi `$ (for $`m_k\mathrm{}`$) can asymptotically learn every problem $`\mu `$ of relevance, i.e. AI$`\xi `$ is asymptotically optimal. We included the qualifier of relevance, as we are not sure whether there could be strange $`\mu `$ spoiling (43) but we expect those $`\mu `$ to be irrelevant from the perspective of AI. In the field of Learning, there are many asymptotic learnability theorems, often not too difficult to prove. So a proof of (43) might also be accessible. Unfortunately, asymptotic learnability theorems are often too weak to be useful from a practical point. Nevertheless, they point in the right direction.
From the convergence (39) of $`\mu \xi `$ we might expect $`C_{km_k}^\xi C_{km_k}^\mu `$ and hence, $`\dot{y}_k^\xi `$ defined in (25) to converge to $`\dot{y}_k^\mu `$ defined in (9) with $`\mu `$ probability 1 for $`k\mathrm{}`$. The first problem is, that if the $`C_{km_k}`$ for the different choices of $`y_k`$ are nearly equal, then even if $`C_{km_k}^\xi C_{km_k}^\mu `$, $`\dot{y}_k^\xi \dot{y}_k^\mu `$ is possible due to the non-continuity of $`maxarg_{y_k}`$. This can be cured by a weighted $`D_{n\mu \xi }`$ as described above. More serious is the second problem we explain for $`h_k=1`$ and $`X=C=\{0,1\}`$. For $`\dot{y}_k^\xi maxarg_{y_k}\xi (\dot{y}\dot{c}_{<k}y_k\underset{¯}{1})`$ to converge to $`\dot{y}_k^\mu maxarg_{y_k}\mu (\dot{y}\dot{c}_{<k}y_k\underset{¯}{1})`$, it is not sufficient to know that $`\xi (\dot{y}\dot{c}_{<k}\dot{y}\underset{¯}{\overset{\dot{}}{c}}_k)\mu (\dot{y}\dot{c}_{<k}\dot{y}\underset{¯}{\overset{\dot{}}{c}}_k)`$ as has been proved in (39). We need convergence not only for the true output $`\dot{y}_k`$ and credit $`\dot{c}_k`$, but also for alternate outputs $`y_k`$ and credit 1. $`\dot{y}_k^\xi `$ converges to $`\dot{y}_k^\mu `$ if $`\xi `$ converges uniformly to $`\mu `$, i.e. if in addition to (39)
$$\left|\mu (yx_{<k}y_k^{}\underset{¯}{x}_k^{})\xi (yx_{<k}y_k^{}\underset{¯}{x}_k^{})\right|<c\left|\mu (yx_{<k}y\underset{¯}{x}_k)\xi (yx_{<k}y\underset{¯}{x}_k)\right|y_k^{}x_k^{}$$
(44)
holds for some constant $`c`$ (at least in some $`\mu `$ expected sense). We call $`\mu `$ satisfying (44) uniform. For uniform $`\mu `$ one can show (43) with appropriately weighted $`D_{n\mu \xi }`$ and bounded horizon $`h_k<h_{max}`$. Unfortunately there are relevant $`\mu `$ which are not uniform. Details will be given elsewhere.
In the following, we briefly mention some further concepts. A Markovian $`\mu `$ is defined as depending only on the last output, i.e. $`\mu (yx_{<k}y\underset{¯}{x}_k)=\mu _k(y\underset{¯}{x}_k)`$. We say $`\mu `$ is generalized Markovian, if $`\mu (yx_{<k}y\underset{¯}{x}_k)=\mu _k(yx_{kl:k1}y\underset{¯}{x}_k)`$ for fixed $`l`$. This property has some similarities to factorizable $`\mu `$ defined in (15). If further $`\mu _k\mu _1k`$, $`\mu `$ is called stationary. Further, for all enumerable $`\mu `$, $`\mu (yx_{<k}y\underset{¯}{x}_k)`$ and $`\xi (yx_{<k}y\underset{¯}{x}_k)`$ get independent of $`yx_{<l}`$ for fixed $`l`$ and $`k\mathrm{}`$ with $`\mu `$ probability 1. This property, which we want to call forgetfulness, will be proved elsewhere. Further, we say $`\mu `$ is farsighted, if $`lim_{m_k\mathrm{}}\dot{y}_k^{(m_k)}`$ exists. More details will be given in the next subsection, where we also give an example of a possibly relevant $`\mu `$, which is not farsighted.
We have introduced several concepts, which might be useful for proving credit bounds, including forgetful, relevant, asymptotically learnable, farsighted, uniform, (generalized) Markovian, factorizable and (pseudo) passive $`\mu `$. We have sorted them here, approximately in the order of decreasing generality. We want to call them separability concepts. The more general (like relevant, asymptotically learnable and farsighted) $`\mu `$ will be called weakly separable, the more restrictive (like (pseudo) passive and factorizable) $`\mu `$ will be called strongly separable, but we will use these qualifiers in a more qualitative, rather than rigid sense. Other (non-separability) concepts are deterministic $`\mu `$ and, of course, the class of all chronological $`\mu `$.
### The choice of the horizon:
The only significant arbitrariness in the AI$`\xi `$ model lies in the choice of the horizon function $`h_km_kk+1`$. We discuss some choices which seem to be natural and give preliminary conclusions at the end. We will not discuss ad hoc choices of $`h_k`$ for specific problems (like the discussion in section 6 in the context of finite games). We are interested in universal choices of $`m_k`$.
If the lifetime of the system is known to be $`T`$, which is in practice always large but finite, then the choice $`m_k=T`$ maximizes correctly the expected future credit. $`T`$ is usually not known in advance, as in many cases the time we are willing to run a system depends on the quality of its outputs. For this reason, it is often desirable that good outputs are not delayed too much, if this results in a marginal credit increase only. This can be incorporated by damping the future credits. If, for instance, we assume that the survival of the system in each cycle is proportional to the past credit an exponential damping $`c_k:=c_k^{}e^{\lambda k}`$ is appropriate, where $`c_k^{}`$ are bounded, e.g. $`c_k^{}[0,1]`$. The expression (25) converges for $`m_k\mathrm{}`$ in this case. But this does not solve the problem, as we introduced a new arbitrary time-scale $`{}_{}{}^{1}/_{\lambda }^{}`$. Every damping introduces a time-scale.
Even the time-scale invariant damping factor $`k^\alpha `$ introduces a dynamic time-scale. In cycle $`k`$ the contribution of cycle $`2^{1/\alpha }k`$ is damped by a factor $`\frac{1}{2}`$. The effective horizon $`h_k`$ in this case is $`k`$. The choice $`h_k=\beta k`$ with $`\beta 2^{1/\alpha }`$ qualitatively models the same behaviour. We have not introduced an arbitrary time-scale $`T`$, but limited the farsightedness to some multiple (or fraction) of the length of the current history. This avoids the pre-selection of a global time-scale $`T`$ or $`{}_{}{}^{1}/_{\lambda }^{}`$. This choice has some appeal, as it seems that humans of age $`k`$ years usually do not plan their lives for more than, perhaps, the next $`k`$ years ($`\beta _{human}=1`$). From a practical point of view this model might serve all needs, but from a theoretical point we feel uncomfortable with such a limitation in the horizon from the very beginning. Note, that we have to choose $`\beta =O(1)`$ because otherwise we would again introduce a number $`\beta `$, which has to be justified.
The naive limit $`m_k\mathrm{}`$ in (25) may turn out to be well defined and the previous discussion superfluous. In the following, we define a limit which is always well defined (for finite $`|Y|`$). Let $`\dot{y}_k^{(m)}`$ be defined as in (25) with $`m_k`$ replaced by $`m`$. Further, let $`\dot{Y}_k^{(m)}:=\{\dot{y}_k^{(m)}:m_km\}`$ be the set of outputs in cycle $`k`$ for the choices $`m_k=m,m+1,m+2,\mathrm{}`$. Because $`\dot{Y}_k^{(m)}\dot{Y}_k^{(m+1)}\{\}`$, we have $`\dot{Y}_k^{(\mathrm{})}:=_{m=k}^{\mathrm{}}\dot{Y}_k^{(m)}\{\}`$. We define the $`m_k=\mathrm{}`$ model to output any $`\dot{y}_k^{(\mathrm{})}\dot{Y}_k^{(\mathrm{})}`$. This is the best output consistent with any choice of $`m_k`$, esp. $`m_k\mathrm{}`$. Choosing the lexicographically smallest $`\dot{y}_k^{(\mathrm{})}\dot{Y}_k^{(\mathrm{})}`$ would correspond to the limes inferior $`\underset{¯}{lim}_m\mathrm{}\dot{y}_k^{(m)}`$. $`\dot{y}_k^{(\mathrm{})}`$ is unique, i.e. $`|\dot{Y}_k^{(\mathrm{})}|=1`$ iff the naive limit $`lim_m\mathrm{}\dot{y}_k^{(m)}`$ exists. Note, that the limit $`lim_m\mathrm{}C_{km}^{}(yx_{<k})`$ needs not to exist for this construction.
The construction above leads to a mathematically elegant, no-parameter AI$`\xi `$ model. Unfortunately this is not the end of the story. The limit $`m_k\mathrm{}`$ can cause undesirable results in the AI$`\mu `$ model for special $`\mu `$ which might also happen in the AI$`\xi `$ model whatever we define $`m_k\mathrm{}`$. Consider $`Y=C=\{0,1\}`$ and $`X^{}=\{\}`$. Output $`y_k=0`$ shall give credit $`c_k=0`$, output $`y_k=1`$ shall give $`c_k=1`$ iff $`\dot{y}_{kl\sqrt{l}}\mathrm{}\dot{y}_{kl}=0\mathrm{}0`$ for some $`l`$. I.e. the system can achieve $`l`$ consecutive positive credits if there was a sequence of length at least $`\sqrt{l}`$ with $`y_k=c_k=0`$. If the lifetime of the AI$`\mu `$ system is $`T`$, it outputs $`\dot{y}_k=0`$ in the first $`r`$ cycles and then $`\dot{y}_k=1`$ for the remaining $`r^2`$ cycles with $`r`$ such that $`r+r^2=T`$. This will lead to the highest possible total credit $`C_{1T}=\sqrt{T+^1/_4}^1/_2`$. Any fragmentation of the $`0`$ and $`1`$ sequences would reduce this. For $`T\mathrm{}`$ the AI$`\mu `$ system can and will delay the point $`r`$ of switching to $`\dot{y}_k=1`$ indefinitely and always output $`0`$ with total credit $`0`$, obviously the worst possible behaviour. The AI$`\xi `$ system will explore the above rule after a while of trying $`y_k=0/1`$ and then applies the same behaviour as the AI$`\mu `$ system, since the simplest rules covering past data dominate $`\xi `$. For finite $`T`$ this is exactly what we want, but for infinite $`T`$ the AI$`\xi `$ model fails just as the AI$`\mu `$ model does. The good point is, that this is not a weakness of the AI$`\xi `$ model, as AI$`\mu `$ fails too and no system can be better than AI$`\mu `$. The bad point is that $`m_k\mathrm{}`$ has far reaching consequences, even when starting from an already very large $`m_k=T`$. The reason being that the $`\mu `$ of this example is highly non-local in time, i.e. it may violate one of our weak separability conditions.
In the last paragraph we have considered the consequences of $`m_k\mathrm{}`$ in the AI$`\mu `$ model. We now consider whether the AI$`\xi `$ model is a good approximation of the AI$`\mu `$ model for large $`m_k`$. Another objection against too large choices of $`m_k`$ is that $`\xi (yx_{<k}y\underset{¯}{x}_{k:m_k})`$ has been proved to be a good approximation of $`\mu (yx_{<k}y\underset{¯}{x}_{k:m_k})`$ only for $`kh_k`$, which is never satisfied for $`m_k=T`$ or $`m_k=\mathrm{}`$. We have seen that, for factorizable $`\mu `$, the limit $`h_k\mathrm{}`$ causes no problem, as from a certain $`h_k`$ on the output $`\dot{y}_k`$ is independent of $`h_k`$. As $`\xi \mu `$ for bounded $`h_k`$, $`\xi `$ will develop this separability property too. So, from a certain $`k_0`$ on the limit $`h_k\mathrm{}`$ might also be safe for $`\xi `$. Therefore, taking the limit from the very beginning worsens the behaviour of AI$`\xi `$ maybe only for finitely many cycles $`kk_0`$, which would be acceptable. We suppose that the valuations $`c_k^{}`$ for $`k^{}k`$, where $`\xi `$ can no longer be trusted as a good approximation to $`\mu `$, are in some sense randomly disturbed with decreasing influence on the choice of $`\dot{y}_k`$. This claim is supported by the forgetfulness property of $`\xi `$.
We are not sure whether the choice of $`m_k`$ is of marginal importance, as long as $`m_k`$ is chosen sufficiently large and of low complexity, $`m_k=2^{2^{16}}`$ for instance, or whether the choice of $`m_k`$ will turn out to be a central topic for the AI$`\xi `$ model or for the planning aspect of any AI system in general. We suppose that the limit $`m_k\mathrm{}`$ for the AI$`\xi `$ model results in correct behaviour for weakly separable $`\mu `$, and that even the naive limit exists, but to prove this would probably give interesting insights.
## 5 Sequence Prediction (SP)
We have introduced the AI$`\xi `$ model as a unification of the ideas of decision theory and universal probability distribution. We might expect AI$`\xi `$ to behave identically to SP$`\mathrm{\Theta }_\xi `$, when faced with a sequence prediction problem, but things are not that simple, as we will see.
### Using the AI$`\mu `$ Model for Sequence Prediction:
We have seen in the last section how to predict sequences for known and unknown prior distribution $`\mu ^{SP}`$. Here we consider binary sequences<sup>15</sup><sup>15</sup>15We use $`z_k`$ to avoid notational conflicts with the systems inputs $`x_k`$. $`z_1z_2z_3\mathrm{}IB^{\mathrm{}}`$ with known prior probability $`\mu ^{SP}(\underset{¯}{z_1z_2z_3\mathrm{}})`$.
We want to show how the AI$`\mu `$ model can be used for sequence prediction. We will see that it gives the same prediction as the SP$`\mathrm{\Theta }_\mu `$ system. First, we have to specify how the AI$`\mu `$ model should be used for sequence prediction. The following choice is natural:
The systems output $`y_k`$ is interpreted as a prediction for the $`k^{th}`$ bit $`z_k`$ of the string, which has to be predicted. This means that $`y_k`$ is binary ($`y_kIB=:Y`$). As a reaction of the environment, the system receives credit $`c_k=1`$ if the prediction was correct ($`y_k=z_k`$), or $`c_k=0`$ if the prediction was erroneous ($`y_kz_k`$). The question is what the input $`x_k^{}`$ of the next cycle should be. One choice would be to inform the system about the correct $`k^{th}`$ bit of the last cycle of the string and set $`x_k^{}=z_k`$. But as from the credit $`c_k`$ in conjunction with the prediction $`y_k`$, the true bit $`z_k=\delta _{y_kc_k}`$ can be inferred, this information is redundant. $`\delta `$ is the Kronecker symbol, defined as $`\delta _{ab}=1`$ for $`a=b`$ and $`0`$ otherwise. There is no need for this additional feedback. So we set $`x_k^{}=ϵX=\{ϵ\}`$ thus having $`x_kc_k`$. The system’s performance does not change when we include this redundant information, it merely complicates the notation. The prior probability $`\mu ^{AI}`$ of the AI$`\mu `$ model is
$$\mu ^{AI}(y_1\underset{¯}{x}_1\mathrm{}y_k\underset{¯}{x}_k)=\mu ^{AI}(y_1\underset{¯}{c}_1\mathrm{}y_k\underset{¯}{c}_k)=\mu ^{SP}(\underset{¯}{\delta _{y_1c_1}\mathrm{}\delta _{y_kc_k}})=\mu ^{SP}(\underset{¯}{z_1\mathrm{}z_k})$$
(45)
In the following, we will drop the superscripts of $`\mu `$ because they are clear from the arguments of $`\mu `$ and the $`\mu `$ equal in any case.
The formula (7) for the expected credit reduces to
$$C_{km}^{}(yx_{<k})=\underset{y_k}{\mathrm{max}}\underset{c_k}{}[c_k+C_{k+1,m}^{}(yx_{1:k})]\mu (\delta _{y_1c_1}\mathrm{}\delta _{y_{k1}c_{k1}}\underset{¯}{\delta _{y_kc_k}})$$
(46)
The first observation we can make, is that for this special $`\mu `$, $`C_{km}^{}`$ only depends on $`\delta _{y_ic_i}`$, i.e. replacing $`y_i`$ and $`c_i`$ simultaneously with their complements does not change the value of $`C_{km}^{}`$. We have a symmetry in $`y_ic_i`$. For $`k=m+1`$ this is definitely true as $`C_{m+1,m}^{}=0`$ in this case (see (6)). For $`km`$ we prove it by induction. The r.h.s. of (46) is symmetric in $`y_ic_i`$ for $`i<k`$ because $`\mu `$ possesses this symmetry and $`C_{k+1,m}^{}`$ possesses it by induction hypothesis, so the symmetry holds for the l.h.s., which completes the proof. The prediction $`\dot{y}_k`$ is
$$\dot{y}_k=\underset{y_k}{maxarg}C_{km_k}^{}(\dot{y}\dot{x}_{<k}y_k)=\underset{y_k}{maxarg}\underset{c_k}{}[c_k+C_{k+1,m_k}^{}(yx_{1:k})]\mu (\mathrm{}\underset{¯}{\delta _{y_kc_k}})=$$
(47)
$$=\underset{y_k}{maxarg}\underset{c_k}{}c_k\mu (\delta _{\dot{y}_1\dot{c}_1}\mathrm{}\underset{¯}{\delta _{y_kc_k}})=\underset{y_k}{maxarg}\mu (\dot{z}_1\mathrm{}\dot{z}_{k1}\underset{¯}{y}_k)=\underset{z_k}{maxarg}\mu (\dot{z}_1\mathrm{}\dot{z}_{k1}\underset{¯}{z}_k)$$
The first equation is the definition of the system’s prediction (8). In the second equation, we have inserted (5) which gives the r.h.s. of (46) with $`\mathrm{max}_{y_k}`$ replaced by $`maxarg_{y_k}`$. $`_cf(\mathrm{}\delta _{yc}\mathrm{})`$ is independent of $`y`$ for any function, depending on the combination $`\delta _{yc}`$ only. Therefore, the $`_cC^{}\mu `$ term is independent of $`y_k`$ because $`C_{k+1,m}^{}`$ as well as $`\mu `$ depend on $`\delta _{y_kc_k}`$ only. In the third equation, we can therefore drop this term, as adding a constant to the argument of $`maxarg_{y_k}`$ does not change the location of the maximum. In the second last equation we evaluated the $`_{c_k}`$. Further, if the true credit to $`\dot{y}_i`$ is $`\dot{c}_i`$ the true $`i^{th}`$ bit of the string must be $`\dot{z}_i=\delta _{\dot{y}_i\dot{c}_i}`$. The last equation is just a renaming.
So, the AI$`\mu `$ model predicts that $`z_k`$ that has maximal $`\mu `$ probability, given $`\dot{z}_1\mathrm{}\dot{z}_{k1}`$. This prediction is independent of the choice of $`m_k`$. It is exactly the prediction scheme of the deterministic sequence prediction with known prior SP$`\mathrm{\Theta }_\mu `$ described in the last section. As this model was optimal, AI$`\mu `$ is optimal, too, i.e. has minimal number of expected errors (maximal expected credit) as compared to any other sequence prediction scheme.
From this, it is already clear that the total expected credit $`C_{km}`$ must be related to the expected sequence prediction error $`E_{m\mathrm{\Theta }_\mu }`$ (21). Let us prove directly that $`C_{1m}(ϵ)+E_{m\mathrm{\Theta }_\mu }=m`$. We rewrite $`C_{km}^{}`$ in (46) as a function of $`z_i`$ instead of $`y_ic_i`$ as it is symmetric in $`y_ic_i`$. Further, we can pull $`C_{km}^{}`$ out of the maximization, as it is independent of $`y_k`$ similar to (47). Renaming the bounded variables $`y_k`$ and $`c_k`$ we get
$$C_{km}^{}(z_{<k})=\underset{z_k}{\mathrm{max}}\mu (z_{<k}\underset{¯}{z}_k)+\underset{z_k}{}C_{k+1,m}^{}(z_{1:k})\mu (z_{<k}\underset{¯}{z}_k)$$
(48)
Recursively inserting the l.h.s. into the r.h.s. we get
$$C_{km}^{}(z_{<k})=\underset{i=k}{\overset{m}{}}\underset{z_{k:i1}}{}\underset{z_i}{\mathrm{max}}\mu (z_{<k}\underset{¯}{z_{k:i}})$$
(49)
This is most easily proven by induction. For $`k=m`$ we have $`C_{mm}^{}(z_{<m})=\mathrm{max}_{z_m}\mu (z_{<m}\underset{¯}{z}_m)`$ from (48) and (6), which equals (49). By induction hypothesis, we assume that (49) is true for $`k`$. Inserting this into (48) we get
$$C_{km}^{}(z_{<k})=\underset{z_k}{\mathrm{max}}\mu (z_{<k}\underset{¯}{z}_k)+\underset{z_k}{}\left[\underset{i=k+1}{\overset{m}{}}\underset{z_{k+1:i1}}{}\underset{z_i}{\mathrm{max}}\mu (z_{1:k}\underset{¯}{z}_{k+1:i})\right]\mu (z_{<k}\underset{¯}{z}_k)=$$
$$=\underset{z_k}{\mathrm{max}}\mu (z_{<k}\underset{¯}{z}_k)+\underset{i=k+1}{\overset{m}{}}\underset{z_{k:i1}}{}\underset{z_i}{\mathrm{max}}\mu (z_{<k}\underset{¯}{z}_{k:i})$$
which equals (49). This was the induction step and hence (49) is proven.
By setting $`k=0`$ and slightly reformulating (49), we get the total expected credit in the first $`m`$ cycles
$$C_{1:m}^{}(ϵ)=\underset{i=1}{\overset{m}{}}\underset{z_{<i}}{}\mu (\underset{¯}{z}_{<i})\mathrm{max}\{\mu (z_{<i}\underset{¯}{0}),\mu (z_{<i}\underset{¯}{1})\}=mE_{m\mathrm{\Theta }_\mu }$$
with $`E_{m\mathrm{\Theta }_\mu }`$ defined in (21).
### Using the AI$`\xi `$ Model for Sequence Prediction:
Now we want to use the universal AI$`\xi `$ model instead of AI$`\mu `$ for sequence prediction and try to derive error bounds analog to (22). Like in the AI$`\mu `$ case, the systems output $`y_k`$ in cycle $`k`$ is interpreted as a prediction for the k<sup>th</sup> bit $`z_k`$ of the string, which has to be predicted. The credit is $`c_k=\delta _{y_kz_k}`$ and there are no other inputs $`x_k=ϵ`$. What makes the analysis more difficult is that $`\xi `$ is not symmetric in $`y_ic_i(1y_i)(1c_i)`$ and (45) does not hold for $`\xi `$. On the other hand, $`\xi ^{AI}`$ converges to $`\mu ^{AI}`$ in the limit (39), and (45) should hold asymptotically for $`\xi `$ in some sense. So we expect that everything proven for AI$`\mu `$ holds approximately for AI$`\xi `$. The AI$`\xi `$ model should behave similarly to SP$`\mathrm{\Theta }_\xi `$, the deterministic variant of Solomonoff prediction. Especially we expect error bounds similar to (22). Making this rigorous seems difficult. Some general remarks have been made in the last section.
Here we concentrate on the special case of a deterministic computable environment, i.e. the environment is a sequence $`\dot{z}=\dot{z}_1\dot{z}_2\mathrm{}`$, $`K(\dot{z}_1\mathrm{}\dot{z}_n)K(\dot{z})<\mathrm{}`$. Furthermore, we only consider the simplest horizon model $`m_k=k`$, i.e. maximize only the next credit. This is sufficient for sequence prediction, as the credit of cycle $`k`$ only depends on output $`y_k`$ and not on earlier decisions. This choice is in no way sufficient and satisfactory for the full AI$`\xi `$ model, as one single choice of $`m_k`$ should serve for all AI problem classes. So AI$`\xi `$ should allow good sequence prediction for some universal choice of $`m_k`$ and not only for $`m_k=k`$, which definitely does not suffice for more complicated AI problems. The analysis of this general case is a challenge for the future. For $`m_k=k`$ the AI$`\xi `$ model (25) with $`x_i^{}=ϵ`$ reduces to
$$\dot{y}_k=\underset{y_k}{maxarg}\underset{c_k}{}c_k\xi (\dot{y}\dot{c}_{<k}y\underset{¯}{c}_k)=\underset{y_k}{maxarg}\xi (\dot{y}\dot{c}_{<k}y_k\underset{¯}{1})=\underset{y_k}{maxarg}\xi (\dot{y}\underset{¯}{\overset{\dot{}}{c}}_{<k}y_k\underset{¯}{1})$$
(50)
The environmental response $`\dot{c}_k`$ is given by $`\delta _{\dot{y}_k\dot{z}_k}`$; it is 1 for a correct prediction $`(\dot{y}_k=\dot{z}_k)`$ and 0 otherwise. In the following, we want to bound the number of errors this prediction scheme makes. We need the following inequality
$$\xi (y\underset{¯}{c}_1\mathrm{}y\underset{¯}{c}_k)>\mathrm{\hspace{0.33em}2}^{K(\delta _{y_1c_1}\mathrm{}\delta _{y_kc_k})O(1)}$$
(51)
We have to find a short program in the sum (24) calculating $`c_1\mathrm{}c_k`$ from $`y_1\mathrm{}y_k`$. If we knew $`z_i:=\delta _{y_ic_i}`$ for $`1ik`$ a program of size $`O(1)`$ could calculate $`c_1\mathrm{}c_k=\delta _{y_1z_1}\mathrm{}\delta _{y_kz_k}`$. So combining this program with a shortest coding of $`z_1\mathrm{}z_k`$ leads to a program of size $`K(z_1\mathrm{}z_k)+O(1)`$, which proves (51).
Let us now assume that we make a wrong prediction in cycle $`k`$, i.e. $`\dot{c}_k=0`$, $`\dot{y}_k\dot{z}_k`$. The goal is to show that $`\dot{\xi }`$ defined by
$$\dot{\xi }_k:=\xi (\dot{y}\underset{¯}{\overset{\dot{}}{c}}_{1:k})=\xi (\dot{y}\underset{¯}{\overset{\dot{}}{c}}_{<k}\dot{y}_k\underset{¯}{0})\xi (\dot{y}\underset{¯}{\overset{\dot{}}{c}}_{<k})\xi (\dot{y}\underset{¯}{\overset{\dot{}}{c}}_{<k}\dot{y}_k\underset{¯}{1})<\dot{\xi }_{k1}\alpha $$
decreases for every wrong prediction, at least by some $`\alpha `$. The $``$ arises from the fact that $`\xi `$ is only a semimeasure.
$$\xi (\dot{y}\underset{¯}{\overset{\dot{}}{c}}_1\mathrm{}\dot{y}\underset{¯}{1})>\xi (\dot{y}_1\underset{¯}{\overset{\dot{}}{c}}_1\mathrm{}(1\dot{y}_k)\underset{¯}{1})\stackrel{\times }{>}\mathrm{\hspace{0.33em}2}^{K(\delta _{\dot{y}_1\dot{c}_1}\mathrm{}\delta _{(1\dot{y}_k)1})}=$$
$$=\mathrm{\hspace{0.33em}2}^{K(\dot{z}_1\mathrm{}\dot{z}_k)}>\mathrm{\hspace{0.33em}2}^{K(\dot{z})O(1)}=:\alpha $$
In the first inequality we have used the fact that $`\dot{y}_k`$ maximizes by definition (50) the argument, i.e. $`1\dot{y}_k`$ has lower probability than $`\dot{y}_k`$. (51) has been applied in the second inequality. The equality holds, because $`\dot{z}_i=\delta _{\dot{y}_i\dot{c}_i}`$ and $`\delta _{(1\dot{y}_k)1}=\delta _{\dot{y}_k0}=\delta _{\dot{y}_k\dot{c}_k}=\dot{z}_k`$. The last inequality follows from the definition of $`\dot{z}`$.
We have shown that each erroneous prediction reduces $`\dot{\xi }`$ by at least the $`\alpha `$ defined above. Together with $`\dot{\xi }_0=1`$ and $`\dot{\xi }_k>0`$ for all $`k`$ this shows that the system can make at most $`1/\alpha `$ errors, since otherwise $`\dot{\xi }_k`$ would become negative. So the number of wrong predictions $`E_{n\xi }^{AI}`$ of system (50) is bounded by
$$E_{n\xi }^{AI}<\frac{1}{\alpha }=\mathrm{\hspace{0.33em}2}^{K(\dot{z})+O(1)}<\mathrm{}$$
(52)
for a computable deterministic environment string $`\dot{z}_1\dot{z}_2\mathrm{}`$. The intuitive interpretation is that each wrong prediction eliminates at least one program $`p`$ of size $`l(p)\stackrel{+}{<}K(\dot{z})`$. The size is smaller than $`K(\dot{z})`$, as larger policies could not mislead the system to a wrong prediction, since there is a program of size $`K(\dot{z})`$ making a correct prediction. There are at most $`2^{K(\dot{z})+O(1)}`$ such policies, which bounds the total number of errors.
We have derived a finite bound for $`E_{n\xi }^{AI}`$, but unfortunately, a rather weak one as compared to (22). The reason for the strong bound in the SP case was that every error at least halves $`\dot{\xi }`$ because the sum of the $`maxarg_{x_k}`$ arguments was 1. Here we have
$$\begin{array}{c}\xi (\dot{y}_1\dot{c}_1\mathrm{}\dot{y}_{k1}\dot{c}_{k1}0\underset{¯}{0})+\xi (\dot{y}_1\dot{c}_1\mathrm{}\dot{y}_{k1}\dot{c}_{k1}0\underset{¯}{1})=1\\ \xi (\dot{y}_1\dot{c}_1\mathrm{}\dot{y}_{k1}\dot{c}_{k1}1\underset{¯}{0})+\xi (\dot{y}_1\dot{c}_1\mathrm{}\dot{y}_{k1}\dot{c}_{k1}1\underset{¯}{1})=1\end{array}$$
but $`maxarg_{y_k}`$ runs over the right top and right bottom $`\xi `$, for which no sum criterion holds.
The AI$`\xi `$ model would not be sufficient for realistic applications if the bound (52) were sharp, but we have the strong feeling (but only weak arguments) that better bounds proportional to $`K(\dot{z})`$ analog to (22) exist. The technique used above may not be appropriate for achieving this. One argument for a better bound is the formal similarity between $`maxarg_{z_k}(\dot{z}_{<k}z_k)`$ and (50), the other is that we were unable to construct an example sequence for which (50) makes more than $`O(K(\dot{z}))`$ errors.
## 6 Strategic Games (SG)
### Introduction:
A very important class of problems are strategic games, like chess. In fact, what is subsumed under game theory nowadays, is so general, that it includes not only a huge variety of games, from simple games of chance like roulette, combined with strategy like Backgammon, up to purely strategic games like chess or checkers or go. Game theory can also describe political and economic competitions and coalitions, even Darwinism and many more have been modeled within game theory. It seems that nearly every AI problem could be brought into the form of a game. Nevertheless, the intention of a game is that several players perform some actions with (partial) observable consequences. The goal of each player is to maximize some utility function (e.g. to win the game). The players are assumed to be rational, taking into account all information they posses. The different goals of the players are usually in conflict. For an introduction into game theory, see .
If we interpret the AI system as one player and the environment models the other rational player and the environment provides the reinforcement feedback $`c_k`$, we see that the system-environment configuration satisfies all criteria of a game. On the other hand, we know that the AI system can handle more general situations, since it interacts optimally with an environment, even if the environment is not a rational player with conflicting goals.
### Strictly competitive strategic games:
In the following, we restrict ourselves to deterministic, strictly competitive strategic<sup>16</sup><sup>16</sup>16In game theory, games like chess are often called ’extensive’, whereas ’strategic’ is reserved for a different kind of game. games with alternating moves. Player 1 makes move $`y_k^{}`$ in round $`k`$, followed by the move $`x_k^{}`$ of player 2. So a game with $`n`$ rounds consists of a sequence of alternating moves $`y_1^{}x_1^{}y_2^{}x_2^{}\mathrm{}y_n^{}x_n^{}`$. At the end of the game in cycle $`n`$ the game or final board state is evaluated with $`C(y_1^{}x_1^{}\mathrm{}y_n^{}x_n^{})`$. Player 1 tries to maximize $`C`$, whereas player 2 tries to minimize $`C`$. In the simplest case, $`C`$ is $`1`$ if player 1 won the game, $`C=1`$ if player 2 won and $`C=0`$ for a draw. We assume a fixed game length $`n`$ independent of the actual move sequence. For games with variable length but maximal possible number of moves $`n`$, we could add dummy moves and pad the length to $`n`$. The optimal strategy (Nash equilibrium) of both players is a minimax strategy
$$\dot{x}_k^{}=\underset{x_k^{}}{minarg}\underset{y_{k+1}^{}}{\mathrm{max}}\underset{x_{k+1}^{}}{\mathrm{min}}\mathrm{}\underset{y_n^{}}{\mathrm{max}}\underset{x_n^{}}{\mathrm{min}}C(\dot{y}_1^{}\dot{x}_1^{}\mathrm{}\dot{y}_k^{}x_k^{}\mathrm{}y_n^{}x_n^{})$$
(53)
$$\dot{y}_k^{}=\underset{y_k^{}}{maxarg}\underset{x_k^{}}{\mathrm{min}}\mathrm{}\underset{y_n^{}}{\mathrm{max}}\underset{x_n^{}}{\mathrm{min}}C(\dot{y}_1^{}\dot{x}_1^{}\mathrm{}\dot{y}_{k1}^{}\dot{x}_{k1}^{}y_k^{}x_k^{}\mathrm{}y_n^{}x_n^{})$$
(54)
But note, that the minimax strategy is only optimal if both players behave rationally. If, for instance, player 2 has limited capabilites or makes errors and player 1 is able to discover these (through past moves) he could exploit these and improve his performance by deviating from the minimax strategy. At least, the classical game theory of Nash equilibria does not take into account limited rationality, whereas the AI$`\xi `$ system should.
### Using the AI$`\mu `$ model for game playing:
In the following, we demonstrate the applicability of the AI model to games. The AI system takes the position of player 1. The environment provides the evaluation $`C`$. For a symmetric situation we could take a second AI system as player 2, but for simplicity we take the environment as the second player and assume that this environmental player behaves according to the minimax strategy (53). The environment serves as a perfect player and as a teacher, albeit a very crude one as it tells the system at the end of the game, only whether it won or lost.
The minimax behaviour of player 2 can be expressed by a (deterministic) probability distribution $`\mu ^{SG}`$ as the following
$$\mu ^{SG}(y_1^{}\underset{¯}{x}_1^{}\mathrm{}y_n^{}\underset{¯}{x}_n^{}):=\{\begin{array}{c}1\text{if}x_k^{}=\underset{x_k^{\prime \prime }}{minarg}\mathrm{}\underset{y_n^{\prime \prime }}{\mathrm{max}}\underset{x_n^{\prime \prime }}{\mathrm{min}}C(y_1^{}\mathrm{}x_{k1}^{}y_k^{\prime \prime }\mathrm{}x_n^{\prime \prime })\mathrm{\hspace{0.33em}1}kn\hfill \\ 0\text{otherwise}\hfill \end{array}$$
(55)
The probability that player 2 makes move $`x_k^{}`$ is $`\mu ^{SG}(\dot{y}_1^{}\dot{x}_1^{}\mathrm{}\dot{y}_k^{}\underset{¯}{x}_k^{})`$ which is 1 for $`x_k^{}=\dot{x}_k^{}`$ as defined in (53) and 0 otherwise.
Clearly, the AI system receives no feedback, i.e. $`c_1=\mathrm{}=c_{n1}=0`$, until the end of the game, where it should receive positive/negative/neutral feedback on a win/loss/draw, i.e. $`c_n=C(\mathrm{})`$. The environmental prior probability is therefore
$$\mu ^{AI}(y_1\underset{¯}{x}_1\mathrm{}y_n\underset{¯}{x}_n)=\{\begin{array}{cc}\mu ^{SG}(y_1^{}\underset{¯}{x}_1^{}\mathrm{}y_n^{}\underset{¯}{x}_n^{})& \text{if}c_1=\mathrm{}=c_{n1}=0\text{and}c_n=C(y_1^{}x_1^{}\mathrm{}y_n^{}x_n^{})\hfill \\ 0& \text{otherwise}\hfill \end{array}$$
(56)
where $`y_i=y_i^{}`$ and $`x_i=c_ix_i^{}`$. If the environment is a minimax player (53) plus a crude teacher $`C`$, i.e. if $`\mu ^{AI}`$ is the true prior probability, the question now is, what is the behaviour $`\dot{y}_k^{AI}`$ of the AI$`\mu `$ system. It turns out that if we set $`m_k=n`$ the AI$`\mu `$ system is also a minimax player (54) and hence optimal
$$\dot{y}_k^{AI}=\underset{y_k}{maxarg}\underset{x_k^{}}{}\mathrm{}\underset{y_n}{\mathrm{max}}\underset{x_n^{}}{}C(\dot{y}\dot{x}_{<k}^{}yx_{k:n}^{})\mu ^{SG}(\dot{y}\dot{x}_{<k}^{}y\underset{¯}{x}_{k:n}^{})=$$
$$=\underset{y_k}{maxarg}\underset{x_k^{}}{}\mathrm{}\underset{y_{n1}}{\mathrm{max}}\underset{x_{n1}^{}}{}\underset{y_n}{\mathrm{max}}\underset{x_n^{}}{\mathrm{min}}C(\dot{y}\dot{x}_{<k}^{}yx_{k:n}^{})\mu ^{SG}(\dot{y}\dot{x}_{<k}^{}y\underset{¯}{x}_{k:n1}^{})=$$
(57)
$$=\mathrm{}=\underset{y_k}{maxarg}\underset{x_{k+1}^{}}{\mathrm{min}}\mathrm{}\underset{y_n}{\mathrm{max}}\underset{x_n^{}}{\mathrm{min}}C(\dot{y}\dot{x}_{<k}^{}yx_{k:n}^{})=\dot{y}_k^{SG}$$
In the first line we inserted $`m_k=n`$ and (56) into the definition (9) of $`\dot{y}_k^{AI}`$. This removes all sums over the $`c_k`$. Further, the sum over $`x_n^{}`$ gives only a contribution for $`x_n^{}=minarg_{x_n^{}}C(\dot{x}_1^{}\dot{y}_1^{}\mathrm{}x_n^{}y_n^{})`$ by definition (55) of $`\mu ^{SG}`$. Inserting this $`x_n^{}`$ gives the second line. $`\mu ^{SG}`$ is effectively reduced to a lower number of arguments and the sum over $`x_n^{}`$ replaced by $`\mathrm{min}_{x_n^{}}`$. Repeating this procedure for $`x_{n1}^{},\mathrm{},x_{k+1}^{}`$ leads to the last line, which is just the minimax strategy of player 1 defined in (54).
Let us now assume that the game under consideration is played $`s`$ times. The prior probability then is
$$\mu ^{AI}(y\underset{¯}{x}_1\mathrm{}y\underset{¯}{x}_{sn})=\underset{r=0}{\overset{s1}{}}\mu _1^{AI}(y\underset{¯}{x}_{rn+1}\mathrm{}y\underset{¯}{x}_{(r+1)n})$$
(58)
where we have renamed the prior probability (56) for one game to $`\mu _1^{AI}`$. (58) is a special case of a factorizable $`\mu `$ (15) with identical factors $`\mu _r=\mu _1^{AI}`$ for all $`r`$ and equal episode lengths $`n_{r+1}n_r=n`$. The AI$`\mu `$ system (58) for repeated game playing also implements the minimax strategy,
$$\dot{y}_k^{AI}=\underset{y_k}{maxarg}\underset{x_k^{}}{\mathrm{min}}\mathrm{}\underset{y_{(r+1)n}}{\mathrm{max}}\underset{x_{(r+1)n}^{}}{\mathrm{min}}C(\dot{y}\dot{x}_{rn+1:k1}^{}\mathrm{}yx_{k:(r+1)n}^{})$$
(59)
with $`r`$ such that $`rn<k(r+1)n`$ and for any choice of $`m_k`$ as long as the horizon $`h_kn`$. This can be proved by using (16) and (57). See section (4) for a discussion on separable and factorizable $`\mu `$.
### Games of variable length:
In the unrepeated case we have argued that games of variable but bounded length can be padded to a fixed length without effect. We now analyze in a sequence of games the effect of replacing the games with fixed length by games of variable length. The sequence $`y_1^{}x_1^{}\mathrm{}y_n^{}x_n^{}`$ can still be grouped into episodes corresponding to the moves of separated consecutive games, but now the length and total number of games that fit into the $`n`$ moves depend on the actual moves taken<sup>17</sup><sup>17</sup>17If the sum of game lengths do not fit exactly into $`n`$ moves, we pad the last game appropriately.. $`C(y_1^{}x_1^{}\mathrm{}y_n^{}x_n^{})`$ equals the number of games where the system wins, minus the number of games where the environment wins. Whenever a loss, win or draw has been achieved by the system or the environment, a new game starts. The player whose turn it would next be, begins the next game. The games are still separated in the sense that the behaviour and credit of the current game does not influence the next game. On the other hand, they are slightly entangled, because the length of the current game determines the time of start of the next. As the rules of the game are time invariant, this does not influence the next game directly. If we play a fixed number of games, the games are completely independent, but if we play a fixed number of total moves $`n`$, the number of games depends on their lengths. This has the following consequences: the better player tries to keep the games short, to win more games in the given time $`n`$. The poorer player tries to draw the games out, in order loose less games. The better player might further prefer a quick draw, rather than to win a long game. Formally, this entanglement is represented by the fact that the prior probability $`\mu `$ does no longer factorize. The reduced form (59) of $`\dot{y}_k^{AI}`$ to one episode is no longer valid. Also, the behaviour $`\dot{y}_k^{AI}`$ of the system depends on $`m_k`$, even if the horizon $`h_k`$ is chosen larger than the longest possible game (unless $`m_kn`$). The important point is that the system realizes that keeping games short/long can lead to increased credit. In practice, a horizon much larger than the average game length should be sufficient to incorporate this effect. The details of games in the distant future do not affect the current game and can, therefore, be ignored. A more quantitative analysis could be interesting, but would lead us too far astray.
### Using the AI$`\xi `$ model for game playing:
When going from the specific AI$`\mu `$ model, where the rules of the game have been explicitly modeled into the prior probability $`\mu ^{AI}`$, to the universal model AI$`\xi `$ we have to ask whether these rules can be learned from the assigned credits $`c_k`$. Here, another (actually the main) reason for studying the case of repeated games, rather than just one game arises. For a single game there is only one cycle of non-trivial feedback namely the end of the game - too late to be useful except when there are further games following.
Even in the case of repeated games, there is only very limited feedback, at most $`\mathrm{log}_23`$ bits of information per game if the 3 outcomes win/loss/draw have the same frequency. So there are at least $`O(K(game))`$ number of games necessary to learn a game of complexity $`K(game)`$. Apart from extremely simple games, even this estimate is far too optimistic. As the AI$`\xi `$ system has no information about the game to begin with, its moves will be more or less random and it can win the first few games merely by pure luck. So the probability that the system looses is near to one and hence the information content $`I`$ in the feedback $`c_k`$ at the end of the game is much less than $`\mathrm{log}_23`$. This situation remains for a very large number of games. On the other hand, in principle, every game should be learnable after a very long sequence of games even with this minimal feedback only, as long as $`I0`$.
The important point is that no other learning scheme with no extra information can learn the game more quickly. We expect this to be true as $`\mu ^{AI}`$ factorizes in the case of games of fixed length, i.e. $`\mu ^{AI}`$ satisfies a strong separability condition. In the case of variable game length the entanglement is also low. $`\mu ^{AI}`$ should still be sufficiently separable allowing to formulate and prove good credit bounds for AI$`\xi `$.
To learn realistic games like tic-tac-toe (noughts and crosses) in realistic time one has to provide more feedback. This could be achieved by intermediate help during the game. The environment could give positive(negative) feedback for every good(bad) move the system makes. The demand on whether a move is to be valued as good should be adopted to the gained experience of the system in such a way that approximately half of the moves are valuated as good and the other half as bad, in order to maximize the information content of the feedback.
For more complicated games like chess, even more feedback is necessary from a practical point of view. One way to increase the feedback far beyond a few bits per cycle is to train the system by teaching it good moves. This is called supervised learning. Despite the fact that the AI model has only a credit feedback $`c_k`$, it is able to learn by teaching, as will be shown in section 8. Another way would be to start with more simple games containing certain aspects of the true game and to switch to the true game when the system has learned the simple game.
No other difficulties are expected when going from $`\mu `$ to $`\xi `$. Eventually $`\xi ^{AI}`$ will converge to the minimax strategy $`\mu ^{AI}`$. In the more realistic case, where the environment is not a perfect minimax player, AI$`\xi `$ can detect and exploit the weakness of the opponent.
Finally, we want to comment on the input/output space $`X`$/$`Y`$ of the AI system. In practical applications, $`Y`$ will possibly include also illegal moves. If $`Y`$ is the set of moves of e.g. a robotic arm, the system could move a wrong figure or even knock over the figures. A simple way to handle illegal moves $`y_k`$ is by interpreting them as losing moves, which terminate the game. Further, if e.g. the input $`x_k`$ is the image of a video camera which makes one shot per move, $`X`$ is not the set of moves by the environment but includes the set of states of the game board. The discussion in this section handles this case as well. There is no need to explicitly design the systems I/O space $`X/Y`$ for a specific game.
The discussion above on the AI$`\xi `$ system was rather informal for the following reason: game playing (the SG$`\xi `$ system) has (nearly) the same complexity as fully general AI, and quantitative results for the AI$`\xi `$ system are difficult (but not impossible) to obtain.
## 7 Function Minimization (FM)
### Applications/Examples:
There are many problems that can be reduced to a minimization problem (FM). The minimum of a (real valued) function $`f:YIR`$ over some domain $`Y`$ or a good approximate of it has to be found, usually with some limited resources.
One popular example is the traveling salesman problem (TSP). $`Y`$ is the set of different routes between towns and $`f(y)`$ the length of route $`yY`$. The task is to find a route of minimal length visiting all cities. This problem is NP hard. Getting good approximations in limited time is of great importance in various applications. Another example is the minimization of production costs (MPC), e.g. of a car, under several constraints. $`Y`$ is the set of all alternative car designs and production methods compatible with the specifications and $`f(y)`$ the overall cost of alternative $`yY`$. A related example is finding materials or (bio)molecules with certain properties (MAT). E.g. solids with minimal electrical resistance or maximally efficient chlorophyll modifications or aromatic molecules that taste as close as possible to strawberry.We can also ask for nice paintings (NPT). $`Y`$ is the set of all existing or imaginable paintings and $`f(y)`$ characterizes how much person $`A`$ likes painting $`y`$. The system should present paintings, which $`A`$ likes.
For now, these are enough examples. The TSP is very rigorous from a mathematical point of view, as $`f`$, i.e. an algorithm of $`f`$, is usually known. In principle, the minimum could be found by extensive search, were it not for computational resource limitations. For MPC, $`f`$ can often be modeled in a reliable and sufficiently accurate way. For MAT you need very accurate physical models, which might be unavailable or too difficult to solve or implement. For NPT the most we have is the judgement of person $`A`$ on every presented painting. The evaluation function $`f`$ cannot be implemented without scanning $`A^{}s`$ brain, which is not possible with todays technology.
So there are different limitations, some depending on the application we have in mind. An implementation of $`f`$ might not be available, $`f`$ can only be tested at some arguments $`y`$ and $`f(y)`$ is determined by the environment. We want to (approximately) minimize $`f`$ with as few function calls as possible or, conversely, find an as close as possible approximation for the minimum within a fixed number of function evaluations. If $`f`$ is available or can quickly be inferred by the system and evaluation is quick, it is more important to minimize the total time needed to imagine new trial minimum candidates plus the evaluation time for $`f`$. As we do not consider computational aspects of AI$`\xi `$ till section 10 we concentrate on the first case, where $`f`$ is not available or dominates the computational requirements.
### The Greedy Model FMG$`\mu `$ :
The FM model consists of a sequence $`\dot{y}_1\dot{z}_1\dot{y}_2\dot{z}_2\mathrm{}`$ where $`\dot{y}_k`$ is a trial of the FM system for a minimum of $`f`$ and $`\dot{z}_k=f(\dot{y}_k)`$ is the true function value returned by the environment. We randomize the model by assuming a probability distribution $`\mu (f)`$ over the functions. There are several reasons for doing this. We might really not know the exact function $`f`$, as in the NPT example, and model our uncertainty by the probability distribution $`\mu `$. More importantly, we want to parallel the other AI classes, like in the SP$`\mu `$ model, where we always started with a probability distribution $`\mu `$ that was finally replaced by $`\xi `$ to get the universal Solomonoff prediction SP$`\xi `$. We want to do the same thing here. Further, the probabilistic case includes the deterministic case by choosing $`\mu (f)=\delta _{ff_0}`$, where $`f_0`$ is the true function. A final reason is that the deterministic case is trivial when $`\mu `$ and hence $`f_0`$ is known, as the system can internally (virtually) check all function arguments and output the correct minimum from the very beginning.
We will assume that $`Y`$ is countable or finite and that $`\mu `$ is a discrete measure, e.g. by taking only computable functions. The probability that the function values of $`y_1,\mathrm{},y_n`$ are $`z_1,\mathrm{},z_n`$ is then given by
$$\mu ^{FM}(y_1\underset{¯}{z}_1\mathrm{}y_n\underset{¯}{z}_n):=\underset{f:f(y_i)=z_i1in}{}\mu (f)$$
(60)
We start with a model that minimizes the expectation $`z_k`$ of the function value $`f`$ for the next output $`y_k`$, taking into account previous information:
$$\dot{y}_k:=\underset{y_k}{minarg}\underset{z_k}{}z_k\mu (\dot{y}_1\dot{z}_1\mathrm{}\dot{y}_{k1}\dot{z}_{k1}y_k\underset{¯}{z}_k)$$
This type of greedy algorithm, just minimizing the next feedback, was sufficient for sequence prediction (SP) and is also sufficient for classification (CF). It is, however, not sufficient for function minimization as the following example demonstrates.
Take $`f:\{0,1\}\{1,2,3,4\}`$. There are 16 different functions which shall be equiprobable, $`\mu (f)=\frac{1}{16}`$. The function expectation in the first cycle
$$z_1:=\underset{z_1}{}z_1\mu (y_1\underset{¯}{z}_1)=\frac{1}{4}\underset{z_1}{}z_1=\frac{1}{4}(1+2+3+4)=\mathrm{\hspace{0.33em}2.5}$$
is just the arithmetic average of the possible function values and is independent of $`y_1`$. Therefore, $`\dot{y}_1=0`$, as $`minarg`$ is defined to take the lexicographically first minimum in an ambiguous case. Let us assume that $`f_0(0)=2`$, where $`f_0`$ is the true environment function, i.e. $`\dot{z}_1=2`$. The expectation of $`z_2`$ is then
$$z_2:=\underset{z_2}{}z_2\mu (02y_2\underset{¯}{z}_2)=\{\begin{array}{cc}\hfill 2\text{for}& y_2=0\hfill \\ \hfill 2.5\text{for}& y_2=1\hfill \end{array}$$
For $`y_2=0`$ the system already knows $`f(0)=2`$, for $`y_2=1`$ the expectation is, again, the arithmetic average. The system will again output $`\dot{y}_2=0`$ with feedback $`\dot{z}_2=2`$. This will continue forever. The system is not motivated to explore other $`y^{}s`$ as $`f(0)`$ is already smaller than the expectation of $`f(1)`$. This is obviously not what we want. The greedy model fails. The system ought to be inventive and try other outputs when given enough time.
The general reason for the failure of the greedy approach is that the information contained in the feedback $`z_k`$ depends on the output $`y_k`$. A FM system can actively influence the knowledge it receives from the environment by the choice in $`y_k`$. It may be more advantageous to first collect certain knowledge about $`f`$ by an (in greedy sense) non-optimal choice for $`y_k`$, rather than to minimize the $`z_k`$ expectation immediately. The non-minimality of $`z_k`$ might be over-compensated in the long run by exploiting this knowledge. In SP, the received information is always the current bit of the sequence, independent of what SP predicts for this bit. This is the reason why a greedy strategy in the SP case is already optimal.
### The general FM$`\mu /\xi `$ Model:
To get a useful model we have to think more carefully about what we really want. Should the FM system output a good minimum in the last output in a limited number of cycles $`T`$, or should the average of the $`z_1,\mathrm{},z_T`$ values be minimal, or does it suffice that just one of the $`z`$ is as small as possible? Let us define the FM$`\mu `$ model as to minimize the $`\mu `$ averaged weighted sum $`\alpha _1z_1+\mathrm{}+\alpha _Tz_T`$ for some given $`\alpha _k0`$. Building the $`\mu `$ average by summation over the $`z_i`$ and minimizing w.r.t. the $`y_i`$ has to be performed in the correct chronological order. With a similar reasoning as in (5) to (9) we get
$$\dot{y}_k^{FM}=\underset{y_k}{minarg}\underset{z_k}{}\mathrm{}\underset{y_T}{\mathrm{min}}\underset{z_T}{}(\alpha _1z_1+\mathrm{}+\alpha _Tz_T)\mu (\dot{y}_1\dot{z}_1\mathrm{}\dot{y}_{k1}\dot{z}_{k1}y_k\underset{¯}{z}_k\mathrm{}y_T\underset{¯}{z}_T)$$
(61)
If we want the final output $`\dot{y}_T`$ to be optimal we should choose $`\alpha _k=0`$ for $`k<T`$ and $`\alpha _T=1`$ (final model FMF$`\mu `$). If we want to already have a good approximation during intermediate cycles, we should demand that the output of all cycles together are optimal in some average sense, so we should choose $`\alpha _k=1`$ for all $`k`$ (sum model FMS$`\mu `$). If we want to have something in between, for instance, increase the pressure to produce good outputs, we could choose the $`\alpha _k=e^{\gamma (kT)}`$ exponentially increasing for some $`\gamma >0`$ (exponential model FME$`\mu `$). For $`\gamma \mathrm{}`$ we get the FMF$`\mu `$, for $`\gamma 0`$ the FMS$`\mu `$ model. If we want to demand that the best of the outputs $`y_1\mathrm{}y_k`$ is optimal, we must replace the $`\alpha `$ weighted $`z`$-sum by $`\mathrm{min}\{z_1,\mathrm{},z_T\}`$ (minimum Model FMM$`\mu `$). We expect the behaviour to be very similar to the FMF$`\mu `$ model, and do not consider it further.
By construction, the FM$`\mu `$ models guarantee optimal results in the usual sense that no other model knowing only $`\mu `$ can be expected to produce better results. The variety of FM variants is not a fault of the theory. They just reflect the fact that there is some interpretational freedom of what is meant by minimization within $`T`$ function calls. In most applications, probably FMF is appropriate. In the NPT application one might prefer the FMS model.
The interesting case (in AI) is when $`\mu `$ is unknown. We define for this case, the FM$`\xi `$ model by replacing $`\mu (f)`$ with some $`\xi (f)`$, which should assign high probability to functions $`f`$ of low complexity. So we might define<sup>18</sup><sup>18</sup>18$`\xi ^{FM}(f)`$ is a true probability distribution if we include partial functions in the domain. So normalization is not necessary. $`\xi (f)=_{q:x[U(qx)=f(x)]}2^{l(q)}`$. The problem with this definition is that it is, in general, undecidable whether a TM $`q`$ is an implementation of a function $`f`$. $`\xi (f)`$ defined in this way is uncomputable, not even approximable. As we only need a $`\xi `$ analog to the l.h.s. of (60), the following definition is natural
$$\xi ^{FM}(y_1\underset{¯}{z}_1\mathrm{}y_n\underset{¯}{z}_n):=\underset{q:q(y_i)=z_i1in}{}2^{l(q)}$$
(62)
$`\xi ^{FM}`$ is actually equivalent to inserting the incomputable $`\xi (f)`$ into (60). $`\xi ^{FM}`$ is an enumerable semi-measure and universal, relative to all probability distributions of the form (60). We will not prove this here.
Alternatively, we could have constrained the sum in (62) by $`q(y_1\mathrm{}y_n)=z_1\mathrm{}z_n`$ analog to (24), but these two definitions are not equivalent. Definition (62) ensures the symmetry<sup>19</sup><sup>19</sup>19See for a discussion on symmetric universal distributions on unordered data. in its arguments and $`\xi ^{FM}(\mathrm{}y\underset{¯}{z}\mathrm{}y\underset{¯}{z}^{}\mathrm{})=0`$ for $`zz^{}`$. It incorporates all general knowledge we have about function minimization, whereas (24) does not. But this extra knowledge has only low information content (complexity of $`O(1)`$), so we do not expect FM$`\xi `$ to perform much worse when using (24) instead of (62). But there is no reason to deviate from (62) at this point.
We can now define an ”error” measure $`E_{T\mu }^{FM}`$ as (61) with $`k=1`$ and $`minarg_{y_1}`$ replaced by $`\mathrm{min}_{y_1}`$ and, additionally, $`\mu `$ replaced by $`\xi `$ for $`E_{T\xi }^{FM}`$. We expect $`|E_{T\xi }^{FM}E_{T\mu }^{FM}|`$ to be bounded in a way that justifies the use of $`\xi `$ instead of $`\mu `$ for computable $`\mu `$, i.e. computable $`f_0`$ in the deterministic case. The arguments are the same as for the AI$`\xi `$ model.
### Is the general model inventive?
In the following we will show that FM$`\xi `$ will never cease searching for minima, but will test an infinite set of different $`y^{}s`$ for $`T\mathrm{}`$.
Let us assume that the system tests only a finite number of $`y_iAY`$, $`|A|<\mathrm{}`$. Let $`t1`$ be the cycle in which the last new $`yA`$ is selected (or some later cycle). Selecting $`y^{}s`$ in cycles $`kt`$ a second time, the feedback $`z`$ does not provide any new information, i.e. does not modify the probability $`\xi ^{FM}`$. The system can minimize $`E_{T\xi }^{FM}`$ by outputting in cycles $`kt`$ the best $`yA`$ found so far (in the case $`\alpha _k=0`$, the output does not matter). Let us fix $`f`$ for a moment. Then we have
$$E^a:=\alpha _1z_1+\mathrm{}+\alpha _Tz_T=\underset{k=1}{\overset{t1}{}}\alpha _kf(y_k)+f_1\underset{k=t}{\overset{T}{}}\alpha _k,f_1:=\underset{1k<t}{\mathrm{min}}f(y_k)$$
Let us now assume that the system tests one additional $`y_tA`$ in cycle $`t`$, but no other $`yA`$. Again, it will keep to the best output for $`k>t`$, which is either the one of the previous system or $`y_t`$.
$$E^b=\underset{k=1}{\overset{t}{}}\alpha _kf(y_k)+\mathrm{min}\{f_1,f(y_t)\}\underset{k=t+1}{\overset{T}{}}\alpha _k$$
The difference can be represented in the form
$$E^aE^b=\left(\underset{k=t}{\overset{T}{}}\alpha _k\right)f^+\alpha _tf^{},f^\pm :=\mathrm{max}\{0,\pm (f_1f(y_t))\}\mathrm{\hspace{0.33em}0}$$
As the true FM$`\xi `$ strategy is the one which minimizes $`E`$, assumption $`a`$ is ruled out if $`E^a>E^b`$. We will say that $`b`$ is favored over $`a`$, which does not mean that $`b`$ is the correct strategy, only that $`a`$ is not the true one. For probability distributed $`f`$, $`b`$ is favored over $`a`$ when
$$E^aE^b=\left(\underset{k=t}{\overset{T}{}}\alpha _k\right)f^+\alpha _tf^{}>\mathrm{\hspace{0.33em}0}\underset{k=t}{\overset{T}{}}\alpha _k>\alpha _t\frac{f^{}}{f^+}$$
where $`f^\pm `$ is the $`\xi `$ expectation of $`\pm f_1f(y_t)`$ under the condition that $`\pm f_1\pm f(y_t)`$ and under the constrains imposed in cycles $`1\mathrm{}t1`$. As $`\xi `$ assigns a strictly positive probability to every non-empty event, $`f^+0`$. Inserting $`\alpha _k=e^{\gamma (kT)}`$, assumption $`a`$ is ruled out in model FME$`\xi `$ if
$$Tt>\frac{1}{\gamma }\mathrm{ln}\left[1+\frac{f^{}}{f^+}(e^\gamma 1)\right]1\{\begin{array}{cc}\hfill 0\text{for}& \gamma \mathrm{}\text{ (FMF}\xi \text{ model)}\hfill \\ \hfill f^{}/f^+1\text{for}& \gamma 0\text{ (FMS}\xi \text{ model)}\hfill \end{array}$$
We see that if the condition is not satisfied for some $`t`$, it will remain wrong for all $`t^{}>t`$. So the FMF$`\xi `$ system will test each $`y`$ only once up to a point from which on it always outputs the best found $`y`$. Further, for $`T\mathrm{}`$ the condition always gets satisfied. As this is true for any finite $`A`$, the assumption of a finite $`A`$ is wrong. For $`T\mathrm{}`$ the system tests an increasing number of different $`y^{}s`$, provided $`Y`$ is infinite. The FMF$`\xi `$ model will never repeat any $`y`$ except in the last cycle $`T`$ where it chooses the best found $`y`$. The FMS$`\xi `$ model will test a new $`y_t`$ for fixed $`T`$, only if the expected value of $`f(y_t)`$ is not too large.
The above does not necessarily hold for different choices of $`\alpha _k`$. The above also holds for the FMF$`\mu `$ system if $`f^+0`$. $`f^+=0`$ if the system can already exclude that $`y_t`$ is a better guess, so there is no reason to test it explicitly.
Nothing has been said about the quality of the guesses, but for the FM$`\mu `$ system they are optimal by definition. If $`K(\mu )`$ for the true distribution $`\mu `$ is finite, we expect the FM$`\xi `$ system to solve the ”exploration versus exploitation” problem in a universally optimal way, as $`\xi `$ converges to $`\mu `$.
### Using the AI models for Function Mininimization:
The AI model can be used for function minimization in the following way. The output $`y_k`$ of cycle $`k`$ is a guess for a minimum of $`f`$, like in the FM model. The credit $`c_k`$ should be high for small function values $`z_k=f(y_k)`$. The credit should also be weighted with $`\alpha _k`$ to reflect the same strategy as in the FM case. The choice of $`c_k=\alpha _kz_k`$ is natural. Here, the feedback is not binary but $`c_kCIR`$, with $`C`$ being a countable subset of $`IR`$, e.g. the computable reals or all rational numbers. The feedback $`x_k^{}`$ should be the function value $`f(y_k)`$. So we set $`x_k^{}=z_k`$. Note, that there is a redundancy if $`\alpha _{()}`$ is a computable function with no zeros, as $`c_k=\alpha _kx_k^{}`$. So, for small $`K(\alpha _{()})`$ like in the FMS model, one might set $`x_kϵ`$. If we keep $`x_k^{}`$ the AI prior probability is
$$\mu ^{AI}(y_1\underset{¯}{x}_1\mathrm{}y_n\underset{¯}{x}_n)=\{\begin{array}{cc}\mu ^{FM}(y_1\underset{¯}{z}_1\mathrm{}y_n\underset{¯}{z}_n)& \text{for }c_k=\alpha _kz_k,x_k^{}=z_k,x_k=c_kx_k^{}\hfill \\ 0& \text{else}.\hfill \end{array}$$
(63)
Inserting this into (9) with $`m_k=T`$ we get
$$\dot{y}_k^{AI}=\underset{y_k}{maxarg}\underset{x_k}{}\mathrm{}\underset{y_T}{\mathrm{max}}\underset{x_T}{}(c_k+\mathrm{}+c_T)\mu ^{AI}(\dot{y}_1\dot{x}_1\mathrm{}y_k\underset{¯}{x}_k\mathrm{}y_T\underset{¯}{x}_T)=$$
$$=\underset{y_k}{minarg}\underset{z_k}{}\mathrm{}\underset{y_T}{\mathrm{min}}\underset{z_T}{}(\alpha _kz_k+\mathrm{}+\alpha _Tz_T)\mu ^{FM}(\dot{y}_1\dot{z}_1\mathrm{}y_k\underset{¯}{z}_k\mathrm{}y_T\underset{¯}{z}_T)=\dot{y}_k^{FM}$$
where $`\dot{y}_k^{FM}`$ has been defined in (61). The proof of equivalence was so simple because the FM model has already a rather general structure, which is similar to the full AI model.
One might expect no problems when going from the already very general FM$`\xi `$ model to the universal AI$`\xi `$ model (with $`m_k=T`$), but there is a pitfall in the case of the FMF model. All credits $`c_k`$ are zero in this case, except for the last one being $`c_T`$. Although there is a feedback $`z_k`$ in every cycle, the AI$`\xi `$ system cannot learn from this feedback as it is not told that in the final cycle $`c_T`$ will equal to $`z_T`$. There is no problem in the FM$`\xi `$ model because in this case this knowledge is hardcoded into $`\xi ^{FM}`$. The AI$`\xi `$ model must first learn that it has to minimize a function but it can only learn if there is a non-trivial credit assignment $`c_k`$. FMF works for repeated minimization of (different) functions, such as minimizing $`N`$ functions in $`NT`$ cycles. In this case there are $`N`$ non-trivial feedbacks and AI$`\xi `$ has time to learn that there is a relation between $`c_{kT}`$ and $`x_{kT}^{}`$ every T<sup>th</sup> cycle. This situation is similar to the case of strategic games discussed in section 6.
There is no problem in applying AI$`\xi `$ to FMS because the $`c`$ feedback provides enough information in this case. The only thing the AI$`\xi `$ model has to learn, is to ignore the $`x`$ feedbacks as all information is already contained in $`c`$. Interestingly the same argument holds for the FME model if $`K(\gamma )`$ and $`K(T)`$ are small<sup>20</sup><sup>20</sup>20If we set $`\alpha _k=e^{\gamma k}`$ the condition on $`K(T)`$ can be dropped.. The AI$`\xi `$ model has additionally only to learn the relation $`c_k=e^{\gamma (kT)}x_k^{}`$. This task is simple as every cycle provides one data point for a simple function to learn. This argument is no longer valid for $`\gamma \mathrm{}`$ as $`K(\gamma )\mathrm{}`$ in this case.
### Remark:
TSP seems to be trivial in the AI$`\mu `$ model but non-trivial in the AI$`\xi `$ model. The reason being that (61) just implements an internal complete search as $`\mu (f)=\delta _{ff^{TSP}}`$ contains all necessary information. AI$`\mu `$ outputs from the very beginning, the exact minimum of $`f^{TSP}`$. This ”solution” is, of course, unacceptable from performance perspective. As long as we give no efficient approximation $`\xi ^c`$ of $`\xi `$, we have not contributed anything to a solution of the TSP by using AI$`\xi ^c`$. The same is true for any other problem where $`f`$ is computable and easily accessible. Therefore, TSP is not (yet) a good example because all we have done is to replace a NP complete problem with the uncomputable AI$`\xi `$ model or by a computable AI$`\xi ^c`$ model, for which we have said nothing about computation time yet. It is simply an overkill to reduce simple problems to AI$`\xi `$. TSP is a simple problem in this respect, until we consider the AI$`\xi ^c`$ model seriously. For the other examples, where $`f`$ is inaccessible or complicated, AI$`\xi ^c`$ provides a true solution to the minimization problem as an explicit definition of $`f`$ is not needed for AI$`\xi `$ and AI$`\xi ^c`$.
## 8 Supervised Learning by Examples (EX)
The AI models provide a frame for reinforcement learning. The environment provides a feedback $`c`$, informing the system about the quality of its last output $`y`$; it assigns credit $`c`$ to output $`y`$. In this sense, reinforcement learning is explicitly integrated into the AI$`\rho `$ model. For $`\rho =\mu `$ it maximizes the true expected credit, whereas the AI$`\xi `$ model is a universal, environment independent, reinforcement learning algorithm.
There is another type of learning method: Supervised learning by presentation of examples (EX). Many problems learned by this method are association problems of the following type. Given some examples $`xRX`$, the system should reconstruct, from a partially given $`x^{}`$, the missing or corrupted parts, i.e. complete $`x^{}`$ to $`x`$ such that relation $`R`$ contains $`x`$. In many cases, $`X`$ consists of pairs $`(z,v)`$, where $`v`$ is the possibly missing part.
### Applications/Examples:
Learning functions by presenting $`(z,f(z))`$ pairs and asking for the function value of $`z`$ by presenting $`(z,\mathrm{?})`$ also falls into this category.
A basic example is learning properties of geometrical objects coded in some way. E.g. if there are 18 different objects characterized by their size (small or big), their colors (red, green or blue) and their shapes (square, triange, circle), then $`(object,property)R`$ if the $`object`$ possesses the $`property`$. Here, $`R`$ is a relation which is not the graph of a single valued function.
When teaching a child, by pointing to objects and saying ”this is a tree” or ”look how green” or ”how beautiful”, one establishes a relation of $`(object,property)`$ pairs in $`R`$. Pointing to a (possibly different) tree later and asking ”what is this ?” corresponds to a partially given pair $`(object,\mathrm{?})`$, where the missing part ”?” should be completed by the child saying ”tree”.
A final example we want to give is chess. We have seen that, in principle, chess can be learned by reinforcement learning. In the extreme case the environment only provides credit $`c=1`$ when the system wins. The learning rate is completely inacceptable from a practical point of view. The reason is the very low amount of information feedback. A more practical method of teaching chess is to present example games in the form of sensible $`(board\text{-}state,move)`$ sequences. They contain information about legal and good moves (but without any explanation). After several games have been presented, the teacher could ask the system to make its own move by presenting $`(board\text{-}state,\mathrm{?})`$ and then evaluate the answer of the system.
### Supervised learning with the AI$`\mu /\xi `$ model:
Let us define the EX model as follows: The environment presents inputs $`x_k^{}=z_kv_k(z_k,v_k)R(Z\times \{\mathrm{?}\})Z\times (Y\{\mathrm{?}\})=X^{}`$ to the system in cycle $`k`$. The system is expected to output $`y_{k+1}`$ in the next cycle, which is evaluated with $`c_{k+1}=1`$ if $`(z_k,y_{k+1})R`$ and 0 otherwise. To simplify the discussion, an output $`y_k`$ is expected and evaluated even when $`v_k(\mathrm{?})`$ is given. To complete the description of the environment, the probability distribution $`\mu _R(\underset{¯}{x_1^{}\mathrm{}x_n^{}})`$ of the examples $`x_i^{}`$ (depending on $`R`$) has to be given. Wrong examples should not occur, i.e. $`\mu _R`$ should be 0 if $`x_i^{}R`$ for some $`1in`$. The relations $`R`$ might also be probability distributed with $`\sigma (\underset{¯}{R})`$. The example prior probability in this case is
$$\mu (\underset{¯}{x_1^{}\mathrm{}x_n^{}})=\underset{R}{}\mu _R(\underset{¯}{x_1^{}\mathrm{}x_n^{}})\sigma (\underset{¯}{R})$$
(64)
The knowledge of the valuation $`c_k`$ on output $`y_k`$ restricts the possible relations $`R`$, consistent with $`R(z_k,y_{k+1})=c_{k+1}`$, where $`R(z,y):=1`$ if $`(z,y)R`$ and 0 otherwise. The prior probability for the input sequence $`x_1\mathrm{}x_n`$ if the output sequence is $`y_1\mathrm{}y_n`$, is therefore
$$\mu ^{AI}(y_1\underset{¯}{x}_1\mathrm{}y_n\underset{¯}{x}_n)=\underset{R:1i<n[R(z_i,y_{i+1})=c_{i+1}]}{}\mu _R(\underset{¯}{x_1^{}\mathrm{}x_n^{}})\sigma (\underset{¯}{R})$$
where $`x_i=c_ix_i^{}`$ and $`x_{i1}^{}=z_iv_i`$ with $`v_iY\{\mathrm{?}\}`$. In the I/O sequence $`y_1x_1y_2x_2\mathrm{}=y_1c_1z_2v_2y_2c_2z_3v_3\mathrm{}`$ the $`c_1y_1`$ are dummies, after which regular behaviour starts.
The AI$`\mu `$ model is optimal by construction of $`\mu ^{AI}`$. For computable prior $`\mu _R`$ and $`\sigma `$, we expect a near optimal behavior of the universal AI$`\xi `$ model if $`\mu _R`$ additionally satisfies some separability property. In the following, we give some motivation why the AI$`\xi `$ model takes into account the supervisor information contained in the examples and why it learns faster than by reinforcement.
We keep $`R`$ fixed and assume $`\mu _R(x_1^{}\mathrm{}x_n^{})=\mu _R(x_1^{})\mathrm{}\mu _R(x_n^{})0x_i^{}R(Z\times \{\mathrm{?}\})i`$ to simplify the discussion. Short codes $`q`$ contribute mostly to $`\xi ^{AI}(y_1\underset{¯}{x}_1\mathrm{}y_n\underset{¯}{x}_n)`$. As $`x_1^{}\mathrm{}x_n^{}`$ is distributed according to the computable probability distribution $`\mu _R`$, a short code of $`x_1^{}\mathrm{}x_n^{}`$ for large enough $`n`$ is a Huffman coding w.r.t. the distribution $`\mu _R`$. So we expect $`\mu _R`$ and hence $`R`$ coded in the dominant contributions to $`\xi ^{AI}`$ in some way, where the plausible assumption was made that the $`y`$ on the input tape do not matter. Much more than one bit per cycle will usually be learned, hence, relation $`R`$ can be learned in $`nK(R)`$ cycles by appropriate examples. This coding of $`R`$ in $`q`$ evolves independently of the feedbacks $`c`$. To maximize the feedback $`c_k`$, the system has to learn to output a $`y_{k+1}`$ with $`(z_k,y_{k+1})R`$. The system has to invent a program extension $`q^{}`$ to $`q`$, which extracts $`z_k`$ from $`x_k=z_kv_k`$ and searches for and outputs a $`y_{k+1}`$ with $`(z_k,y_{k+1})R`$. As $`R`$ is already coded in $`q`$, $`q^{}`$ can re-use this coding of $`R`$ in $`q`$. The size of the extension $`q^{}`$ is, therefore, of $`O(1)`$. To learn this $`q^{}`$, the system requires feedback $`c`$ with information content of $`O(1)=K(q^{})`$ only.
Let us compare this with reinforcement learning, where only $`x_k^{}=(z_k,\mathrm{?})`$ pairs are presented. A coding of $`R`$ in a short code $`q`$ for $`x_1^{}\mathrm{}x_n^{}`$ is of no use and will therefore be absent. Only the credits $`c`$ force the system to learn $`R`$. $`q^{}`$ is therefore expected to be of size $`K(R)`$. The information content in the $`c^{}s`$ must be of the order $`K(R)`$. In practice, there are often only very few $`c_k=1`$ at the beginning of the learning phase and the information content in $`c_1\mathrm{}c_n`$ is much less than $`n`$ bits. The required number of cycles to learn $`R`$ by reinforcement is, therefore, at least but in many cases much larger than $`K(R)`$.
Although AI$`\xi `$ was never designed or told to learn supervised, it learns how to take advantage of the examples from the supervisor. $`\mu _R`$ and $`R`$ are learned from the examples, the credits $`c`$ are not necessary for this process. The remaining task of learning how to learn supervised is then a simple task of complexity $`O(1)`$, for which the credits $`c`$ are necessary.
## 9 Other AI Classes
### Other aspects of intelligence:
In AI, a variety of general ideas and methods have been developed. In the last sections, we have seen how several problem classes can be formulated within AI$`\xi `$. As we claim universality of the AI$`\xi `$ model, we want to enlight which of, and how the other AI methods are incorporated in the AI$`\xi `$ model, by looking its structure. Some methods are directly included, others are or should be emergent. We do not claim the following list to be complete.
Probability theory and utility theory are the heart of the AI$`\mu /\xi `$ models. The probabilities are the true/universal behaviours of the environment. The utility function is what we called total credit, which should be maximized. Maximization of an expected utility function in a probabilistic environment is usually called sequential decision theory, and is explicitly integrated in full generality in our model. This includes probabilistic (a generalization of deterministic) reasoning, where the object of reasoning are not true or false statements, but the prediction of the environmental behaviour. Reinforcement Learning is explicitly built in, due to the credits. Supervised learning is an emergent phenomenon (section 8). Algorithmic information theory leads us to use $`\xi `$ as a universal estimate for the prior probability $`\mu `$.
For horizon $`>1`$, the alternative series of expectimax series in (16) and the process of selecting maximal values can be interpreted as abstract planning. This expectimax series also includes informed search, in the case of AI$`\mu `$, and heuristic search, for AI$`\xi `$, where $`\xi `$ could be interpreted as a heuristic for $`\mu `$. The minimax strategy of game playing in case of AI$`\mu `$ is also subsumed. The AI$`\xi `$ model converges to the minimax strategy if the environment is a minimax player but it can also take advantage of environmental players with limited rationality. Problem solving occurs (only) in the form of how to maximize the expected future credit.
Knowledge is accumulated by AI$`\xi `$ and is stored in some form not specified further on the working tape. Any kind of information in any representation on the inputs $`y`$ is exploited. The problem of knowledge engineering and representation appears in the form of how to train the AI$`\xi `$ model. More practical aspects, like language or image processing have to be learned by AI$`\xi `$ from scratch.
Other theories, like fuzzy logic, possibility theory, Dempster-Shafer theory, … are partly outdated and partly reducible to Bayesian probability theory . The interpretation and effects of the evidence gap $`g:=1_{x_k}\xi (yx_{<k}y\underset{¯}{x}_k)>0`$ in $`\xi `$ may be similar to those in Dempster-Shafer theory. Boolean logical reasoning about the external world plays, at best, an emergent role in the AI$`\xi `$ model.
Other methods, which don’t seem to be contained in the AI$`\xi `$ model might also be emergent phenomena. The AI$`\xi `$ model has to construct short codes of the environmental behaviour, the AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ (see next section) has to construct short action programs. If we would analyze and interpret these programs for realistic environments, we might find some of the unmentioned or unused or new AI methods at work in these algorithms. This is, however, pure speculation at this point. More important: when trying to make AI$`\xi `$ practically usable, some other AI methods, like genetic algorithms or neural nets, may be useful.
The main thing we wanted to point out is that the AI$`\xi `$ model does not lack any important known property of intelligence or known AI methodology. What is missing, however, are computational aspects, which are addressed, in the next section.
## 10 Time Bounds and Effectiveness
### Introduction:
Until now, we have not bothered with the non-computability of the universal probability distribution $`\xi `$. As all universal models in this paper are based on $`\xi `$, they are not effective in this form. In this section, We will outline how the previous models and results can be modified/generalized to the time-bounded case. Indeed, the situation is not as bad as it could be. $`\xi `$ and $`C`$ are enumerable and $`\dot{y}_k`$ is still approximable or computable in the limit. There exists an algorithm, that will produce a sequence of outputs eventually converging to the exact output $`\dot{y}_k`$, but we can never be sure whether we have already reached it. Besides this, the convergence is extremely slow, so this type of asymptotic computability is of no direct (practical) use, but will nevertheless, be important later.
Let $`\stackrel{~}{p}`$ be a program which calculates within a reasonable time $`\stackrel{~}{t}`$ per cycle, a reasonable intelligent output, i.e. $`\stackrel{~}{p}(\dot{x}_{<k})=\dot{y}_{1:k}`$. This sort of computability assumption, that a general purpose computer of sufficient power is able to behave in an intelligent way, is the very basis of AI, justifying the hope to be able to construct systems which eventually reach and outperform human intelligence. For a contrary viewpoint see . It is not necessary to discuss here, what is meant by ’reasonable time/intelligence’ and ’sufficient power’. What we are interested in, in this section, is whether there is a computable version AI$`\xi ^{\stackrel{~}{t}}`$ of the AI$`\xi `$ system which is superior or equal to any $`p`$ with computation time per cycle of at most $`\stackrel{~}{t}`$. With ’superior’, we mean ’more intelligent’, so what we need is an order relation (like) (41) for intelligence.
The best result we could think of would be an AI$`\xi ^{\stackrel{~}{t}}`$ with computation time $`\stackrel{~}{t}`$ at least as intelligent as any $`p`$ with computation time $`\stackrel{~}{t}`$. If AI is possible at all, we would have reached the final goal, the construction of the most intelligent algorithm with computation $`\stackrel{~}{t}`$. Just as there is no universal measure in the set of computable measures (within time $`\stackrel{~}{t}`$), such an AI$`\xi ^t`$ may neither exist.
What we can realistically hope to construct, is an AI$`\xi ^{\stackrel{~}{t}}`$ system of computation time $`c\stackrel{~}{t}`$ per cycle for some constant $`c`$. The idea is to run all programs $`p`$ of length $`\stackrel{~}{l}:=l(\stackrel{~}{p})`$ and time $`\stackrel{~}{t}`$ per cycle and pick the best output. The total computation time is $`2^{\stackrel{~}{l}}\stackrel{~}{t}`$, hence $`c=2^{\stackrel{~}{l}}`$. This sort of idea of ’typing monkeys’ with one of them eventually writing Shakespeare, has been applied in various forms and contexts in theoretical computer science. The realization of this best vote idea, in our case, is not straightforward and will be outlined in this section. An idea related to this, is that of basing the decision on the majority of algorithms. This ’democratic vote’ idea has been used in for sequence prediction, and is referred to as ’weighted majority’ there.
### Time limited probability distributions:
In the literature one can find time limited versions of Kolmogorov complexity and the time limited universal semimeasure . In the following, we utilize and adapt the latter and see how far we get. One way to define a time-limited universal chronological semimeasure is as a sum over all enumerable chronological semimeasures computable within time $`\stackrel{~}{t}`$ and of size at most $`\stackrel{~}{l}`$ similar to the unbounded case (27).
$$\xi ^{\stackrel{~}{t}\stackrel{~}{l}}(y\underset{¯}{x}_{1:n}):=\underset{\rho :l(\rho )\stackrel{~}{l}t(\rho )\stackrel{~}{t}}{}2^{l(\rho )}\rho (y\underset{¯}{x}_{1:n})$$
(65)
Let us assume that the true environmental prior probability $`\mu ^{AI}`$ is equal to or sufficiently accurately approximated by a $`\rho `$ with $`l(\rho )\stackrel{~}{l}`$ and $`t(\rho )\stackrel{~}{t}`$ with $`\stackrel{~}{t}`$ and $`\stackrel{~}{l}`$ of reasonable size. There are several AI problems that fall into this class. In function minimization of section 7, the computation of $`f`$ and $`\mu ^{FM}`$ are usually feasible. In many cases, the sequences of section 5 which should be predicted, can be easily calculated when $`\mu ^{SP}`$ is known. In a classifier problem, the probability distribution $`\mu ^{CF}`$, according to which examples are presented, is, in many cases, also elementary. But not all AI problems are of this ’easy’ type. For the strategic games of section 6, the environment is usually, itself, a highly complex strategic player with a difficult to calculate $`\mu ^{SG}`$ that is difficult to calculate, although one might argue that the environmental player may have limited capabilities too. But it is easy to think of a difficult to calculate physical (probabilistic) environment like the chemistry of biomolecules.
The number of interesting applications makes this restricted class of AI problems, with time and space bounded environment $`\mu ^{\stackrel{~}{t}\stackrel{~}{l}}`$, worth being studied. Superscripts to a probability distribution except for $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ indicate their length and maximal computation time. $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ defined in (65), with a yet to be determined computation time, multiplicatively dominates all $`\mu ^{\stackrel{~}{t}\stackrel{~}{l}}`$ of this type. Hence, an AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ model, where we use $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ as prior probability, is universal, relative to all AI$`\mu ^{\stackrel{~}{t}\stackrel{~}{l}}`$ models in the same way as AI$`\xi `$ is universal to AI$`\mu `$ for all enumerable chronological semimeasures $`\mu `$. The $`maxarg_{y_k}`$ in (25) selects a $`y_k`$ for which $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ has the highest expected utility $`C_{km_k}`$, where $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ is the weighted average over the $`\rho ^{\stackrel{~}{t}\stackrel{~}{l}}`$. $`\dot{y}_k^{AI\xi ^{\stackrel{~}{t}\stackrel{~}{l}}}`$ is determined by a weighted majority. We expect $`AI\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ to outperform all (bounded) $`AI\rho ^{\stackrel{~}{t}\stackrel{~}{l}}`$, analog to the unrestricted case.
In the following we analyze the computability properties of $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ and AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$, i.e. of $`\dot{y}_k^{AI\xi ^{\stackrel{~}{t}\stackrel{~}{l}}}`$. To compute $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ according to the definition (65) we have to enumerate all chronological enumerable semimeasures $`\rho ^{\stackrel{~}{t}\stackrel{~}{l}}`$ of length $`\stackrel{~}{l}`$ and computation time $`\stackrel{~}{t}`$. This can be done similarly to the unbounded case (32-34). All $`2^{\stackrel{~}{l}}`$ enumerable functions of length $`\stackrel{~}{l}`$, computable within time $`\stackrel{~}{t}`$ have to be converted to chronological probability distributions. For this, one has to evaluate each function for $`|X|k`$ different arguments. Hence, $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ is computable within time<sup>21</sup><sup>21</sup>21We assume that a TM can be simulated by another in linear time. $`t(\xi ^{\stackrel{~}{t}\stackrel{~}{l}}(y\underset{¯}{x}_{1:k}))=O(|X|k2^{\stackrel{~}{l}}\stackrel{~}{t})`$. The computation time of $`\dot{y}_k^{AI\xi ^{\stackrel{~}{t}\stackrel{~}{l}}}`$ depends on the size of $`X`$, $`Y`$ and $`m_k`$. $`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ has to be evaluated $`|Y|^{h_k}|X|^{h_k}`$ times in (25). It is possible to optimize the algorithm and perform the computation within time
$$t(\dot{y}_k^{AI\xi ^{\stackrel{~}{t}\stackrel{~}{l}}})=O(|Y|^{h_k}|X|^{h_k}2^{\stackrel{~}{l}}\stackrel{~}{t})$$
(66)
per cycle. If we assume that the computation time of $`\mu ^{\stackrel{~}{t}\stackrel{~}{l}}`$ is exactly $`\stackrel{~}{t}`$ for all arguments, the brute force time $`\overline{t}`$ for calculating the sums and maxs in (9) is $`\overline{t}(\dot{y}_k^{AI\mu ^{\stackrel{~}{t}\stackrel{~}{l}}})|Y|^{h_k}|X|^{h_k}\stackrel{~}{t}`$. Combining this with (66), we get
$$t(\dot{y}_k^{AI\xi ^{\stackrel{~}{t}\stackrel{~}{l}}})=O(2^{\stackrel{~}{l}}\overline{t}(\dot{y}_k^{AI\mu ^{\stackrel{~}{t}\stackrel{~}{l}}}))$$
This result has the proposed structure, that there is a universal AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ system with computation time $`2^{\stackrel{~}{l}}`$ times the computation time of a special AI$`\mu ^{\stackrel{~}{t}\stackrel{~}{l}}`$ system.
Unfortunately, the class of AI$`\mu ^{\stackrel{~}{t}\stackrel{~}{l}}`$ systems with brute force evaluation of $`\dot{y}_k`$, according to (9) is completely uninteresting from a practical point of view. E.g. in the context of chess, the above result says that the AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ is superior within time $`2^{\stackrel{~}{l}}\stackrel{~}{t}`$ to any brute force minimax strategy of computation time $`\stackrel{~}{t}`$. Even if the factor of $`2^{\stackrel{~}{l}}`$ in computation time would not matter, the AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ system is, nevertheless practically useless, as a brute force minimax chess player with reasonable time $`\stackrel{~}{t}`$ is a very poor player.
Note, that in the case of sequence prediction ($`h_k=1`$, $`|Y|=|X|=2`$) the computation time of $`\rho `$ coincides with that of $`\dot{y}_k^{AI\rho }`$ within a factor of 2. The class AI$`\rho ^{\stackrel{~}{t}\stackrel{~}{l}}`$ includes all non-incremental sequence prediction algorithms of size $`\stackrel{~}{l}`$ and computation time $`\stackrel{~}{t}/2`$. With non-incremental, we mean that no information of previous cycles is taken into account for the computation of $`\dot{y}_k`$ of the current cycle.
The shortcomings (mentioned and unmentioned ones) of this approach are cured in the next subsection, by deviating from the standard way of defining a timebounded $`\xi `$ as a sum over functions or programs.
### The idea of the best vote algorithm:
A general cybernetic or AI system is a chronological program $`p(x_{<k})=y_{1:k}`$. This form, introduced in section 2, is general enough to include any AI system (and also less intelligent systems). In the following, we are interested in programs $`p`$ of length $`\stackrel{~}{l}`$ and computation time $`\stackrel{~}{t}`$ per cycle. One important point in the time-limited setting is that $`p`$ should be incremental, i.e. when computing $`y_k`$ in cycle $`k`$, the information of the previous cycles stored on the working tape can be re-used. Indeed, there is probably no practically interesting, non-incremental AI system at all.
In the following, we construct a policy $`p^{}`$, or more precisely, policies $`p_k^{}`$ for every cycle $`k`$ that outperform all time and length limited AI systems $`p`$. In cycle k, $`p_k^{}`$ runs all $`2^{\stackrel{~}{l}}`$ programs $`p`$ and selects the one with the best output $`y_k`$. This is a ’best vote’ type of algorithm, as compared to the ’weighted majority’ like algorithm of the last subsection. The ideal measure for the quality of the output would be the $`\xi `$ expected credit
$$C_{km}(p|\dot{y}\dot{x}_{<k}):=\underset{q\dot{Q}_k}{}2^{l(q)}C_{km}(p,q),C_{km}(p,q):=c(x_k^{pq})+\mathrm{}+c(x_m^{pq})$$
(67)
The program $`p`$ which maximizes $`C_{km_k}`$ should be selected. We have dropped the normalization $`𝒩`$ unlike in (40), as it is independent of $`p`$ and does not change the order relation which we are solely interested in here. Furthermore, without normalization, $`C_{km}`$ is enumerable, which will be important later.
### Extended chronological programs:
In the (functional form of the) AI$`\xi `$ model it was convenient to maximize $`C_{km_k}`$ over all $`p\dot{P}_k`$, i.e. all $`p`$ consistent with the current history $`\dot{y}\dot{x}_{<k}`$. This was no restriction, because for every possibly inconsistent program $`p`$ there exists a program $`p^{}\dot{P}_k`$ consistent with the current history and identical to $`p`$ for all future cycles $`k`$. For the time limited best vote algorithm $`p^{}`$ it would be too restrictive to demand $`p\dot{P}_k`$. To prove universality, one has to compare all $`2^{\stackrel{~}{l}}`$ algorithms in every cycle, not just the consistent ones. An inconsistent algorithm may become the best one in later cycles. For inconsistent programs we have to include the $`\dot{y}_k`$ into the input, i.e. $`p(\dot{y}\dot{x}_{<k})=y_{1:k}^p`$ with $`\dot{y}_iy_i^p`$ possible. For $`p\dot{P}_k`$ this was not necessary, as $`p`$ knows the output $`\dot{y}_ky_k^p`$ in this case. The $`c_i^{pq}`$ in the definition of $`C_{km}`$ are the valuations emerging in the I/O sequence, starting with $`\dot{y}\dot{x}_{<k}`$ (emerging from $`p^{}`$) and then continued by applying $`p`$ and $`q`$ with $`\dot{y}_i:=y_i^p`$ for $`ik`$.
Another problem is that we need $`C_{km_k}`$ to select the best policy, but unfortunately $`C_{km_k}`$ is uncomputable. Indeed, the structure of the definition of $`C_{km_k}`$ is very similar to that of $`\dot{y}_k`$, hence a brute force approach to approximate $`C_{km_k}`$ requires too much computation time as for $`\dot{y}_k`$. We solve this problem in a similar way, by supplementing each $`p`$ with a program that estimates $`C_{km_k}`$ by $`w_k^p`$ within time $`\stackrel{~}{t}`$. We combine the calculation of $`y_k^p`$ and $`w_k^p`$ and extend the notion of a chronological program once again to
$$p(\dot{y}\dot{x}_{<k})=w_1^py_1^p\mathrm{}w_k^py_k^p$$
(68)
with chronological order $`w_1^py_1^p\dot{y}_1\dot{x}_1w_2^py_2^p\dot{y}_2\dot{x}_2\mathrm{}`$.
### Valid approximations:
$`p`$ might suggest any output $`y_k^p`$ but it is not allowed to rate it with an arbitrarily high $`w_k^p`$ if we want $`w_k^p`$ to be a reliable criterion for selecting the best $`p`$. We demand that no policy is allowed to claim that it is better than it actually is. We define a (logical) predicate VA($`p`$) called valid approximation, which is true if, and only if, $`p`$ always satisfies $`w_k^pC_{km_k}(p)`$, i.e. never overrates itself.
$$\text{VA}(p)kw_1^py_1^p\dot{y}_1\dot{x}_1\mathrm{}w_k^py_k^p:p(\dot{y}\dot{x}_{<k})=w_1^py_1^p\mathrm{}w_k^py_k^pw_k^pC_{km_k}(p|\dot{y}\dot{x}_{<k})$$
(69)
In the following, we restrict our attention to programs $`p`$, for which VA($`p`$) can be proved in some formal axiomatic system. A very important point is that $`C_{km_k}`$ is enumerable. This ensures the existence of sequences of program $`p_1,p_2,p_3,\mathrm{}`$ for which VA($`p_i`$) can be proved and $`lim_i\mathrm{}w_k^{p_i}=C_{km_k}(p)`$ for all $`k`$ and all I/O sequences. The approximation is not uniform in $`k`$, but this does not matter as the selected $`p`$ is allowed to change from cycle to cycle.
Another possibility would be to consider only those $`p`$ which check $`w_k^pC_{km_k}(p)`$ online in every cycle, instead of the pre-check VA($`p`$), either by constructing a proof (on the working tape) for this special case, or it is already evident by the construction of $`w_k^p`$. In cases where $`p`$ cannot guarantee $`w_k^pC_{km_k}(p)`$ it sets $`w_k=0`$ and, hence, trivially satisfies $`w_k^pC_{km_k}(p)`$. On the other hand, for these $`p`$ it is also no problem to prove VA($`p`$) as one has simply to analyze the internal structure of $`p`$ and recognize that $`p`$ shows the validity internally itself, cycle by cycle, which is easy by assumption on $`p`$. The cycle by cycle check is, therefore, a special case of the pre-proof of VA($`p`$).
### Effective intelligence order relation:
In section 4 we have introduced an intelligence order relation $``$ on AI systems, based on the expected credit $`C_{km_k}(p)`$. In the following we need an order relation $`^c`$ based on the claimed credit $`w_k^p`$ which might be interpreted as an approximation to $``$. We call $`p`$ effectively more or equally intelligent than $`p^{}`$ if
$$\begin{array}{c}p^cp^{}:k\dot{y}\dot{x}_{<k}w_{1:n}w_{1:n}^{}:\\ p(\dot{y}\dot{x}_{<k})=w_1\mathrm{}w_kp^{}(\dot{y}\dot{x}_{<k})=w_1^{}\mathrm{}w_k^{}w_kw_k^{}\end{array}$$
(70)
i.e. if $`p`$ always claims higher credit estimate $`w`$ than $`p^{}`$. $`^c`$ is a co-enumerable partial order relation on extended chronological programs. Restricted to valid approximations it orders the policies w.r.t. the quality of their outputs and their ability to justify their outputs with high $`w_k`$.
### The universal time bounded AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ system:
In the following we, describe the algorithm $`p^{}`$ underlying the universal time bounded AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ system. It is essentially based on the selection of the best algorithms $`p_k^{}`$ out of the time $`\stackrel{~}{t}`$ and length $`\stackrel{~}{l}`$ bounded $`p`$, for which there exists a proof of VA($`p`$) with length $`l_P`$.
1. Create all binary strings of length $`l_P`$ and interpret each as a coding of a mathematical proof in the same formal logic system in which VA($``$) has been formulated. Take those strings which are proofs of VA($`p`$) for some $`p`$ and keep the corresponding programs $`p`$.
2. Eliminate all $`p`$ of length $`>\stackrel{~}{l}`$.
3. Modify all $`p`$ in the following way: all output $`w_k^py_k^p`$ is temporarily written on an auxiliary tape. If $`p`$ stops in $`\stackrel{~}{t}`$ steps the internal ’output’ is copied to the output tape. If $`p`$ does not stop after $`\stackrel{~}{t}`$ steps a stop is forced and $`w_k=0`$ and some arbitrary $`y_k`$ is written on the output tape. Let $`P`$ be the set of all those modified programs.
4. Start first cycle: $`k:=1`$.
5. Run every $`pP`$ on extended input $`\dot{y}\dot{x}_{<k}`$, where all outputs are redirected to some auxiliary tape: $`p(\dot{y}\dot{x}_{<k})=w_1^py_1^p\mathrm{}w_k^py_k^p`$.
6. Select the program $`p`$ with highest claimed credit $`w_k^p`$: $`p_k^{}:=maxarg_pw_k^p`$.
7. Write $`\dot{y}_k:=y_k^{p_k^{}}`$ to the output tape.
8. Receive input $`\dot{x}_k`$ from the environment.
9. Begin next cycle: $`k:=k+1`$, goto step 5.
It is easy to see that the following theorem holds.
### Main theorem:
Let $`p`$ be any extended chronological (incremental) program like (68) of length $`l(p)\stackrel{~}{l}`$ and computation time per cycle $`t(p)\stackrel{~}{t}`$, for which there exists a proof of VA($`p`$) defined in (69) of length $`l_P`$. The algorithm $`p^{}`$ constructed in the last subsection, depending on $`\stackrel{~}{l}`$, $`\stackrel{~}{t}`$ and $`l_P`$ but not on $`p`$, is effectively more or equally intelligent, according to $`^c`$ defined in (70) than any such $`p`$. The size of $`p^{}`$ is $`l(p^{})=O(\mathrm{ln}(\stackrel{~}{l}\stackrel{~}{t}l_P))`$, the setup-time is $`t_{setup}(p^{})=O(l_P2^{l_P})`$, the computation time per cycle is $`t_{cycle}(p^{})=O(2^{\stackrel{~}{l}}\stackrel{~}{t})`$.
Roughly speaking, the theorem says, that if there exists a computable solution to some AI problem at all, the explicitly constructed algorithm $`p^{}`$ is such a solution. Although this theorem is quite general, there are some limitations and open questions which we discuss in the following.
### Limitations and open questions:
* Formally, the total computation time of $`p^{}`$ for cycles $`1\mathrm{}k`$ increases linearly with $`k`$, i.e. is of order $`O(k)`$ with a coefficient $`2^{\stackrel{~}{l}}\stackrel{~}{t}`$. The unreasonably large factor $`2^{\stackrel{~}{l}}`$ is a well known drawback in best/democratic vote models and will be taken without further comments, whereas the factor $`\stackrel{~}{t}`$ can be assumed to be of reasonable size. If we don’t take the limit $`k\mathrm{}`$ but consider reasonable $`k`$, the practical usefulness of the timebound on $`p^{}`$ is somewhat limited, due to the additional additive constant $`O(l_P2^{l_P})`$. It is much larger than $`k2^{\stackrel{~}{l}}\stackrel{~}{t}`$ as typically $`l_Pl(`$VA$`(p))l(p)\stackrel{~}{l}`$.
* $`p^{}`$ is superior only to those $`p`$ which justify their outputs (by large $`w_k^p`$). It might be possible that there are $`p`$ which produce good outputs $`y_k^p`$ within reasonable time, but it takes an unreasonably long time to justify their outputs by sufficiently high $`w_k^p`$. We do not think that (from a certain complexity level onwards) there are policies where the process of constructing a good output is completely separated from some sort of justification process. But this justification might not be translatable (at least within reasonable time) into a reasonable estimate of $`C_{km_k}(p)`$.
* The (inconsistent) programs $`p`$ must be able to continue strategies started by other policies. It might happen that a policy $`p`$ steers the environment to a direction for which it is specialized. A ’foreign’ policy might be able to displace $`p`$ only between loosely bounded episodes. There is probably no problem for factorizable $`\mu `$. Think of a chess game, where it is usually very difficult to continue the game/strategy of a different player. When the game is over, it is usually advantageous to replace a player by a better one for the next game. There might also be no problem for sufficiently separable $`\mu `$.
* There might be (efficient) valid approximations $`p`$ for which VA($`p`$) is true but not provable, or for which only a very long ($`>l_P`$) proof exists.
### Remarks:
* The idea of suggesting outputs and justifying them by proving credit bounds implements one aspect of human thinking. There are several possible reactions to an input. Each reaction possibly has far reaching consequences. Within a limited time one tries to estimate the consequences as well as possible. Finally, each reaction is valued and the best one is selected. What is inferior to human thinking is, that the estimates $`w_k^p`$ must be rigorously proved and the proofs are constructed by blind extensive search, further, that all behaviours $`p`$ of length $`\stackrel{~}{l}`$ are checked. It is inferior ’only’ in the sense of necessary computation time but not in the sense of the quality of the outputs.
* In practical applications there are often cases with short and slow programs $`p_s`$ performing some task $`T`$, e.g. the computation of the digits of $`\pi `$, for which there also exist long and quick programs $`p_l`$ too. If it is not too difficult to prove that this long program is equivalent to the short one, then it is possible to prove $`K(T)l(p_s)`$ within time $`t(p_l)`$. Similarly, the method of proving bounds $`w_k`$ for $`C_{km_k}`$ can give high lower bounds without explicitly executing these short and slow programs, which mainly contribute to $`C_{km_k}`$.
* Dovetailing all length and time-limited programs is a well known elementary idea (typing monkeys). The crucial part which has been developed here, is the selection criterion for the most intelligent system.
* By construction of AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ and due to the enumerability of $`C_{km_k}`$, ensuring arbitrary close approximations of $`C_{km_k}`$ we expect that the behaviour of AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ converges to the behaviour of AI$`\xi `$ in the limit $`\stackrel{~}{t},\stackrel{~}{l}\mathrm{}`$ in a sense.
* Depending on what you know/assume that a program $`p`$ of size $`\stackrel{~}{l}`$ and computation time per cycle $`\stackrel{~}{t}`$ is able to achieve, the computable AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ model will have the same capabilities. For the strongest assumption of the existence of a Turing machine, which outperforms human intelligence, the AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ will do too, within the same time frame up to a (unfortunately very large) constant factor.
## 11 Outlook & Discussion
This section contains some discussion of otherwise unmentioned topics and some (more personal) remarks. It also serves as an outlook to further research.
### Miscellaneous:
* In game theory one often wants to model the situation of simultaneous actions, whereas the AI$`\xi `$ models has serial I/O. Simultaneity can be simulated by withholding the environment from the current system’s output $`y_k`$, until $`x_k`$ has been received by the system. Formally, this means that $`\xi (yx_{<k}y\underset{¯}{x}_k)`$ is independent of $`y_k`$. The AI$`\xi `$ system is already of simultaneous type in an abstract view if the behaviour $`p`$ is interpreted as the action. In this sense, AI$`\xi `$ is the action $`p^{}`$ which maximizes the utility function (credit), under the assumption that the environment acts according to $`\xi `$. The situation is different from game theory as the environment is not modeled to be a second ’player’ that tries to optimize his own utility although it might actually be a rational player (see section 6).
* In various examples we have chosen differently specialized input and output spaces $`X`$ and $`Y`$. It should be clear that, in principle, this is unnecessary, as large enough spaces $`X`$ and $`Y`$, e.g. $`2^{32}`$ bit, serve every need and can always be Turing reduced to the specific presentation needed internally by the AI$`\xi `$ system itself. But it is clear that using a generic interface, such as camera and monitor for, learning tic-tac-toe for example, adds the task of learning vision and drawing.
### Outlook:
* Rigorous proofs for credit bounds are the major theoretical challenge are – general ones as well as tighter bounds for special environments $`\mu `$. Of special importance are suitable (and acceptable) conditions to $`\mu `$, under which $`\dot{y}_k`$ and finite credit bounds exist for infinite $`Y`$, $`X`$ and $`m_k`$.
* A direct implementation of the AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ model is ,at best, possible for toy environments due to the large factor $`2^{\stackrel{~}{l}}`$ in computation time. But there are other applications of the AI$`\xi `$ theory. We have seen in several examples how to integrate problem classes into the AI$`\xi `$ model. Conversely, one can downscale the AI$`\xi `$ model by using more restricted forms of $`\xi `$. This could be done in the same way as the theory of universal induction has been downscaled with many insights to the Minimum Description Length principle or to the domain of finite automata . The AI$`\xi `$ model might similarly serve as a super model or as the very definition of (universal unbiased) intelligence, from which specialized models could be derived.
* With a reasonable computation time, the AI$`\xi `$ model would be a solution of AI (see next point if you disagree). The AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ model was the first step, but the elimination of the factor $`2^{\stackrel{~}{l}}`$ without giving up universality will (almost certainly) be a very difficult task. One could try to select programs $`p`$ and prove VA($`p`$) in a more clever way than by mere enumeration, to improve performance without destroying universality. All kinds of ideas like, genetic algorithms, advanced theorem provers and many more could be incorporated. But now we are in trouble. We seem to have transferred the AI problem just to a different level. This shift has some advantages (and also some disadvantages) but presents, in no way, a solution. Nevertheless, we want to stress that we have reduced the AI problem to (mere) computational questions. Even the most general other systems the author is aware of, depend on some (more than computational) assumptions about the environment or it is far from clear whether they are, indeed, universal and optimal. Although computational questions are themselves highly complicated, this reduction is a non-trivial result. A formal theory of something, even if not computable, is often a great step toward solving a problem and has also merits of its own, and AI should not be different (see previous item).
* Many researchers in AI believe that intelligence is something complicated and cannot be condensed into a few formulas. It is more a combining of enough methods and much explicit knowledge in the right way. From a theoretical point of view, we disagree as the AI$`\xi `$ model is simple and seems to serve all needs. From a practical point of view we agree to the following extent. To reduce the computational burden one should provide special purpose algorithms (methods) from the very beginning, probably many of them related to reduce the complexity of the input and output spaces $`X`$ and $`Y`$ by appropriate preprocessing methods.
* There is no need to incorporate extra knowledge from the very beginning. It can be presented in the first few cycles in any format. As long as the algorithm to interpret the data is of size $`O(1)`$, the AI$`\xi `$ system will ’understand’ the data after a few cycles (see section 8). If the environment $`\mu `$ is complicated but extra knowledge $`z`$ makes $`K(\mu |z)`$ small, one can show that the bound (20) reduces to $`\frac{1}{2}\mathrm{ln}2K(\mu |z)`$ when $`x_1z`$, i.e. when $`z`$ is presented in the first cycle. The special purpose algorithms could be presented in $`x_1`$, too, but it would be cheating to say that no special purpose algorithms had been implemented in AI$`\xi `$. The boundary between implementation and training is unsharp in the AI$`\xi `$ model.
* We have not said much about the training process itself, as it is not specific to the AI$`\xi `$ model and has been discussed in literature in various forms and disciplines. A serious discussion would be out of place. To repeat a truism, it is, of course, important to present enough knowledge $`x_k^{}`$ and evaluate the system output $`y_k`$ with $`c_k`$ in a reasonable way. To maximize the information content in the credit, one should start with simple tasks and give positive reward $`c_k=1`$ to approximately half of the outputs $`y_k`$.
### The big questions:
This subsection is devoted to the big questions of AI in general and the AI$`\xi `$ model in particular with a personal touch.
* There are two possible objections to AI in general and, therefore, also against AI$`\xi `$ in particular we want to comment on briefly. Non-computable physics (which is not too weird) could make Turing computable AI impossible. As at least the world that is relevant for humans seems mainly to be computable we do not believe that it is necessary to integrate non-computable devices into an AI system. The (clever and nearly convincing) ’Gödel’ argument by Penrose that non-computational physics must exist and is relevant to the brain, has (in our opinion convincing) loopholes.
* A more serious problem is the evolutionary information gathering process. It has been shown that the ’number of wisdom’ $`\mathrm{\Omega }`$ contains a very compact tabulation of $`2^n`$ undecidable problems in its very first $`n`$ binary digits . $`\mathrm{\Omega }`$ is only enumerable with computation time increasing more rapidly with $`n`$, than any recursive function. The enormous computational power of evolution could have developed and coded something like $`\mathrm{\Omega }`$ into our genes, which significantly guides human reasoning. In short: Intelligence could be something complicated and evolution toward it from an even cleverly designed algorithm of size $`O(1)`$ could be too slow. As evolution has already taken place, we could add the information from our genes or brain structure to any/our AI system, but this means that the important part is still missing and a simple formal definition of AI is principally impossible.
* For the probably biggest question about consciousness we want to give a physical analogy. Quantum (field) theory is the most accurate and universal physical theory ever invented. Although already developed in the 1930ies the big question regarding the interpretation of the wave function collapse is still open. Although extremely interesting from a philosophical point of view, it is completely irrelevant from a practical point of view<sup>22</sup><sup>22</sup>22In the theory of everything, the collapse might become of ’practical’ importance and must or will be solved.. We believe the same to be true for consciousness in the field of Artificial Intelligence. Philosophically highly interesting but practically unimportant. Whether consciousness will be explained some day is another question.
## 12 Conclusions
All tasks which require intelligence to be solved can naturally be formulated as a maximization of some expected utility in the framework of agents. We gave a functional (2) and an iterative (9) formulation of such a decision theoretic agent, which is general enough to cover all AI problem classes, as has been demonstrated by several examples. The main remaining problem is the unknown prior probability distribution $`\mu ^{AI}`$ of the environment(s). Conventional learning algorithms are unsuitable, because they can neither handle large (unstructured) state spaces, nor do they converge in the theoretically minimal number of cycles, nor can they handle non-stationary environments appropriately. On the other hand, the universal semimeasure $`\xi `$ (18), based on ideas from algorithmic information theory, solves the problem of the unknown prior distribution for induction problems. No explicit learning procedure is necessary, as $`\xi `$ automatically converges to $`\mu `$. We unified the theory of universal sequence prediction with the decision theoretic agent by replacing the unknown true prior $`\mu ^{AI}`$ by an appropriately generalized universal semimeasure $`\xi ^{AI}`$. We gave strong arguments that the resulting AI$`\xi `$ model is the most intelligent, parameterless and environmental/application independent model possible. We defined an intelligence order relation (41) to give a rigorous meaning to this claim. Furthermore, possible solutions to the horizon problem have been discussed. We outlined for a number of problem classes in sections 58, how the AI$`\xi `$ model can solve them. They include sequence prediction, strategic games, function minimization and, especially, how AI$`\xi `$ learns to learn supervised. The list could easily be extended to other problem classes like classification, function inversion and many others. The major drawback of the AI$`\xi `$ model is that it is uncomputable, or more precisely, only asymptotically computable, which makes an implementation impossible. To overcome this problem, we constructed a modified model AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$, which is still effectively more intelligent than any other time $`\stackrel{~}{t}`$ and space $`\stackrel{~}{l}`$ bounded algorithm. The computation time of AI$`\xi ^{\stackrel{~}{t}\stackrel{~}{l}}`$ is of the order $`\stackrel{~}{t}2^{\stackrel{~}{l}}`$. Possible further research has been discussed. The main directions could be to prove general and special credit bounds, use AI$`\xi `$ as a super model and explore its relation to other specialized models and finally improve performance with or without giving up universality.
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# 1 Interaction of Q-balls with matter
## 1 Interaction of Q-balls with matter
### 1.1 Interaction with matter of Q-balls type SECS
SECS (Supersymmetric Electrically Charged Solitons) are Q-balls with a net positive electric charge which tends to be mainly in the outer layer. The charge of SECS originates from the unequal rates of absorption in the condensate<sup>1</sup>. The positive electric charge could be of one unit up to several tens. This positive electric charge may be neutralized by a surrounding cloud of electrons.
For small size Q-balls the positive charge interacts with matter (electrons and nuclei) via elastic or quasi elastic collisions<sup>2</sup>. The cross section is similar to the Bohr cross section of hydrogenoid atoms<sup>3-4</sup>
$$\sigma =\pi a_0^210^{16}cm^2$$
(1)
where $`a_0`$ is the Bohr radius.
The main energy losses<sup>5</sup> of SECS passing through matter with velocities in the range $`10^4<\beta <10^2`$ are due to two contributions: $`\left(i\right)`$ the interaction of the SECS positive charge with the nuclei (nuclear contribution), and, $`\left(ii\right)`$ with the electrons of the traversed medium (electronic contribution). The total energy loss is the sum of the two contributions<sup>5-6</sup>.
Electronic energy losses of SECS: The electronic<sup>5-7</sup> contribution to the energy loss of SECS may be computed with the following formula
$$\frac{dE}{dx}=\frac{8\pi a_0e^2\beta }{\alpha }\frac{Z_Q^{7/6}N_e}{\left(Z_Q^{2/3}+Z^{2/3}\right)^{3/2}}$$
(2)
where $`\alpha `$ is the fine structure constant, $`\beta =v/c`$, $`Z_Q`$ is the positive charge of SECS, $`Z`$ is the atomic number of the medium and $`N_e`$ is the density of electrons in the medium. Electronic losses dominate for $`\beta >10^4`$ for the case of the dE/dx in the earth<sup>5</sup>.
Nuclear energy losses of SECS: The nuclear<sup>5</sup> contribution to the energy loss of SECS is due to the interaction of the SECS positive core with the nuclei of the medium and it is given by
$$\frac{dE}{dx}=\frac{\pi a^2\gamma NE}{ϵ}S_n\left(ϵ\right)$$
(3)
where
$$S_n\left(ϵ\right)\frac{0.56ln\left(1.2ϵ\right)}{1.2ϵ\left(1.2ϵ\right)^{0.63}},ϵ=\frac{aME}{Z_QZe^2M_Q}$$
(4)
and
$$a=\frac{0.885a_0}{\left(\sqrt{Z_Q}+\sqrt{Z}\right)^{2/3}},\gamma =\frac{4M}{M_Q}$$
(5)
$`M_Q`$ is the mass of the incident Q-ball; $`M`$ is the mass of the target nucleus; $`Z_Qe`$ and $`Ze`$ are their electric charges; we assume that $`M_Q>>M`$. Nuclear energy losses dominates for $`\beta 10^4`$.
Using the energy losses of SECS discussed above, we can computed, for a specific velocity $`v=250km/s`$, the angular acceptance<sup>6</sup> of a detector located at the underground Gran Sasso Laboratory in Italy (MACRO experiment).
### 1.2 Interaction with matter of Q-balls type SENS
The Q-ball interior of SENS (Supersymmetric Electrically Neutral Solitons) is characterized by a large Vacuum Expectation Value (VEV) of squarks, and may be of sleptons and Higgs fields. The $`SU\left(3\right)_c`$ symmetry is broken and deconfinement takes place inside the Q-ball. If a nucleon enters this region of deconfinement, it dissociates into three quarks, some of which may then become absorbed in the supersymmetric condensate. The reaction looks like<sup>3</sup>
$$\left(Q\right)+Nucleon\left(Q+1\right)+pions$$
(6)
or less probably,
$$\left(Q\right)+Nucleon\left(Q+1\right)+kaons$$
(7)
The nucleon enter the $`\stackrel{~}{q}`$ condensate and it gives rise to the process
$$qq\stackrel{~}{q}\stackrel{~}{q}$$
(8)
If it is assumed that the energy released in (6) and (7) is the same as in typical hadronic processes (about 1 GeV per nucleon), this energy is carried by 2 or 3 pions (or kaons). The cross section is determined by the Q-ball radius
$$\sigma =\pi R_Q^2=\frac{16\pi ^2}{9}M_Q^2Q^26\times 10^{34}Q^{1/2}\left(\frac{1TeV}{M_S}\right)^2cm^2$$
(9)
The corresponding mean free path $`\lambda `$ is
$$\lambda =\frac{1}{\sigma N}$$
(10)
The energy loss of SENS moving with velocities in the range $`10^4<\beta <10^2`$ is constant and is given by
$$\frac{dE}{dx}\frac{\zeta }{\lambda }=\sigma N\zeta $$
(11)
where $`N`$ is the number of atoms per $`cm^3`$ in the traversed material; $`\zeta =1\left(\frac{\rho }{1g.cm^3}\right)GeV`$ is the energy released in one reaction. Large mass SENS lose a small fraction of their kinetic energy and are able to traverse the earth without attenuation for all masses of our interest.
## 2 Energy losses of Q-balls in the earth
In general the energy losses in the earth are calculated using the density profile of the earth interior. One may observe three layers: the nucleus, the mantle and the crust. For our purposes it is sufficient to use a simpler model, in which the density and composition of each layer is uniform.
In the earth interior model the nucleus is made of iron, with a density of 11.5 $`\text{g/cm}^3`$ and a conductivity of $`1.6\times 10^{16}s^1`$; the mantle is made of Si, with a density of 4.3 $`\text{g/cm}^3`$. The radius of the nucleus is 0.54 earth radii. The crust may be neglected as long as we consider Q-balls arriving at underground detectors from below. The rock above is assumed to have the same composition of the mantle.
The energy losses of SECS in the earth mantle and earth core can be easily computed for different $`\beta `$-ranges and for different positive electric charges of the Q-ball core.
## 3 Energy losses of SECS in detectors
### 3.1 Light Yield of SECS in Scintillators
For SECS we distinguish two contributions to the light yield in scintillators: the primary light yield and the secondary light yield.
The primary light yield is due to the direct excitation (and ionization that occurs only for $`\beta >10^3`$) produced by the SECS in the medium. The energy loss in the MACRO liquid scintillator is computed from the energy loss of protons in hydrogen and carbon<sup>5-8</sup>
* For electric charge $`q=1e`$ the energy loss of SECS is calculated for two cases
i) For $`10^5<\beta <5\times 10^3`$ we have the following formula
$$\left(\frac{dE}{dx}\right)_{SECS}=1.3\times 10^5\beta [1\mathrm{exp}\left(\frac{\beta }{7\times 10^4}\right)^2]\frac{MeV}{cm}$$
(12)
ii) For $`5\times 10^3<\beta <10^2`$ we used the following formula
$$SP=SP_H+SP_C=\left(\frac{dE}{dx}\right)_{SECS}$$
(13)
where
$$SP_H=\frac{SL_H\times SH_H}{SL_H+SH_H}$$
(14)
$$SP_C=\frac{SL_C\times SH_C}{SL_C+SH_C}$$
(15)
and
$$SL=A_1E^{0.45},SH=A_2Ln\left(1+\frac{A_3}{E}+A_4E\right)$$
(16)
where ($`A_{i=1,4}`$) are constants obtained from experimental data, and $`E`$ is the energy of a proton with velocity $`\beta `$.
* For SECS with electric charge $`q=Z_1e`$ the energy losses for $`10^5<\beta <10^2`$ are given by<sup>9-11</sup>
$$\left(\frac{dE}{dx}\right)_{SECS}=\frac{8\pi e^2a_0\beta }{\alpha }\frac{Z_1^{7/6}N_e}{\left(Z_1^{2/3}+Z_2^{2/3}\right)^{3/2}}[1\mathrm{exp}(\frac{\beta }{7\times 10^4})^2]$$
(17)
where $`Z_2`$ is the atomic number of the target atom, $`N_e`$ the density of electrons and $`\alpha `$ is the fine structure constant.
The primary light yield of SECS is given by
$$\left(\frac{dL}{dx}\right)_{SECS}=A\left[\frac{1}{1+AB\frac{dE}{dx}}\right]\frac{dE}{dx}$$
(18)
where $`dE/dx`$ is the total energy loss of SECS; $`A`$ is a constant of conversion of the energy losses in photons (light yield) and $`B`$ is the parameter describing the saturation of the light yield; both parameters depend only on the velocity of SECS.
The secondary light yield arises from recoiling particles: we consider the elastic or quasi-elastic recoil of hydrogen and carbon nuclei. The light yield $`L_p`$ from a hydrogen or carbon nucleus of given initial energy $`E`$ is computed as
$$L_p\left(E\right)=_0^E\frac{dL}{dx}\left(ϵ\right)S_{tot}^1𝑑ϵ$$
(19)
where $`S_{tot}`$ is the sum of electronic and nuclear energy losses. The secondary light yield is then
$$\left(\frac{dL}{dx}\right)_{\text{secondary}}=N_0^{T_m}L_p\left(T\right)\frac{d\sigma }{dT}𝑑T$$
(20)
where $`N`$ is the number density of atoms in the medium $`T_m`$ is the maximum energy transferred and $`\frac{d\sigma }{dT}`$ is the differential scattering cross section.
### 3.2 Energy losses of SECS in streamer tubes
The composition of the gas in the MACRO limited streamer tubes is 73% helium and 27% n-pentane, in volume. The pressure is about one atmosphere and the resulting density is low (in comparison with the density of the other detectors): $`\rho _{gas}=0.856\text{mg/cm}^3`$.
The ionization energy losses of SECS with $`10^3<\beta <10^2`$ in the MACRO streamer tubes are computed with the same general procedure used for scintillators, using the density and the chemical composition of streamer tubes.
The threshold for ionizing n-pentane occurs for $`\beta 2\times 10^3`$.
### 3.3 Restricted Energy Losses of SECS in the Nuclear Track Detectors
The relevant parameter for nuclear track detectors is the Restricted Energy Loss (REL), that is, the energy deposited within $``$ 100 Å from the track.
The chemical composition of CR39 nuclear track detector is $`\left(\text{C}_{12}\text{H}_{18}\text{O}_7\right)_n`$, and the density is 1.31 $`\text{g/cm}^3`$. For the computation of the REL, only energy transfers to atoms larger than $`12`$ eV are considered, because it is estimated that $`12`$ eV are necessary to break the molecular bonds.
At low velocities ($`3\times 10^5<\beta <10^2`$) there are two contributions to REL: the ionization and the atomic recoil contributions.
The ionization contribution, which become important only for $`\beta >2\times 10^3`$, was computed with Ziegler’s fit to the experimental data.
The atomic recoil contribution, important for low $`\beta `$ values, and was calculated using the interaction potential between an atom and a SECS which is
$$V\left(r\right)=\frac{Z_1Z_2e^2}{r}\varphi \left(r\right)$$
(21)
where $`r`$ is the distance between the core of SECS and the target atom, $`Z_1e`$ is the electric charge of the SECS core, $`Z_2`$ is the atomic number of the target atom.
The Restricted Energy Losses are finally obtained by integrating the transferred energies as
$$\frac{dE}{dx}=N\sigma \left(K\right)𝑑K$$
(22)
where $`N`$ is the number density of atoms in the medium, $`\sigma \left(K\right)`$ is the differential cross section as function of the transferred kinetic energy K.
## 4 Conclusions
We computed for a large range of velocity the energy losses of Q-balls of type SENS and SECS in matter. Using these energy losses and a rough model of the earth’s composition and density profiles, we have computed the energy losses in the earth interior.
We also calculated the energy deposited in scintillators, streamer tubes and nuclear track detectors by SECS, in forms useful for their detection. In particular we computed the light yield in scintillators, the ionization in streamer tubes and the REL in nuclear track detectors.
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# Analysis of the Metallic Phase of Two-Dimensional Holes in SiGe in Terms of Temperature Dependent Screening
## Abstract
We find that temperature dependent screening can quantitatively explain the metallic behaviour of the resistivity on the metallic side of the so-called metal-insulator transition in p-SiGe. Interference and interaction effects exhibit the usual insulating behaviour which is expected to overpower the metallic background at sufficiently low temperatures. We find empirically that the concept of a Fermi-liquid describes our data in spite of the large $`r_s8`$.
Conductivity measurements give experimental access to the nature of the electronic ground state of disordered conductors. At zero temperature ($`T=0`$) extended wave functions at the Fermi-level result in metallic behaviour (zero or finite resistivity $`\rho `$), while localised wave-functions lead to insulating properties (infinite resistivity) . Since in experiments the absolute zero is not accessible, the criterion $`d\rho /dT>0`$ ($`d\rho /dT<0`$) measured at the lowest accessible temperatures is thought to indicate a metallic (insulating) ground state.
In two-dimensional (2D) systems it was suggested theoretically and supported experimentally that the ground state at zero temperature is of insulating nature. This belief has been challenged by the interpretation of experiments on high mobility Si-MOSFETs . Later on, experiments on GaAs hole and electron gases , AlAs electron gases and SiGe hole gases were shown to exhibit similar features as observed in the Si-MOSFET systems . As a consequence, the existence of a metal-insulator transition (MIT) in 2D-systems has been controversially discussed .
In this paper we discuss the metallic behaviour observed in our p-SiGe quantum wells . Interference corrections to the conductivity are extracted from weak localisation studies at low magnetic fields and interaction corrections to the conductivity are obtained from the temperature dependence of the Hall resistivity, similar to Ref. . In addition, following the suggestion in Ref. , we compare the remaining (Drude) part of the conductivity with the theory for temperature dependent screening suitable for our system . We arrive at the following conclusions: (I) In the metallic regime the 2D-hole gas in our SiGe samples behaves like an ordinary Fermi liquid and exhibits localising interference and interaction corrections to the conductivity which can be described by conventional theory. (II) The metallic temperature dependence of the resistivity can be described by the theory of temperature dependent screening entering the Drude part of the resistivity.
Our samples were grown by MBE (molecular beam epitaxy). They comprise a 20 nm Si<sub>0.85</sub>Ge<sub>0.15</sub> quantum well where the two-dimensional hole gas is formed, sandwiched between undoped Si layers. The structures are remotely doped with boron at a distance of 15 nm above the quantum well and gated with a Ti/Al Schottky gate. For more details see Ref. . Measurements were carried out at temperatures between 180 mK and 15 K using four-terminal AC and DC techniques.
The SiGe quantum well is compressively strained normal to the growth direction, which leads, together with the confinement, to a splitting between heavy-hole (HH) and light-hole (LH) bands of the order of 25 meV . The 2D hole gas resides in the lowest HH subband with effective mass $`m^{}=0.25m_0`$ as determined from Shubnikov-de Haas oscillations. The hole density in ungated areas of the device was $`4.3\times 10^{11}`$ cm<sup>-2</sup>. With the gate the hole density could be tuned between $`1.12.6\times 10^{11}`$cm<sup>-2</sup>, i.e. the Fermi energy was below 2 meV. The hole mobility in these structures increases with carrier concentration from 1000 cm<sup>2</sup>/Vs at the lowest to 7800 cm<sup>2</sup>/Vs at the highest density similar to other p-SiGe structures . The ratio of the Drude scattering time $`\tau _e`$ and the quantum lifetime $`\tau _q`$ extracted from Shubnikov-de Haas oscillations is close to one indicating that transport is dominated by short range scattering potentials. It was confirmed in other studies that large angle scattering at interface charges dominates the mobility in such structures .
Recently Coleridge et al. and we have shown that p-SiGe exhibits the characteristic features of the MIT. On increasing the hole density $`d\rho /dT`$ changes sign . At the critical density $`p_c`$, we found $`r_s=1/a_B^{}\sqrt{\pi p}8`$, where $`a_B^{}`$ is the effective Bohr radius. The critical resistance was $`\rho _ch/e^2`$. The $`r_s`$ value is a measure of the importance of interaction effects and $`\rho `$ is a measure of the degree of disorder in the system. Like in other systems temperature and electric field scaling could be performed allowing the interpretation in terms of a quantum phase transition. We have found in Ref. that in the metallic regime weak localisation (WL) reduces the metallic behaviour without destroying it, whereas localising interference corrections dominate the zero field resistivity in the insulating regime.
In the following we analyse magnetotransport data in the metallic regime by first extracting the contributions of WL and interaction corrections to the temperature dependence of the conductivity; then the effect of temperature dependent screening is quantitatively discussed.
In Fig. 1a we show magnetoconductivities measured at the highest densities $`p=2.6\times 10^{11}`$ cm<sup>-2</sup> (achieved under the gated part of the sample) and $`p=4.3\times 10^{11}`$ cm<sup>-2</sup> (ungated part of the sample) for various temperatures. Fig. 1b presents the magnetoconductivity curves of the higher density for the lowest and highest temperature. Fig. 1c shows magnetoconductivity curves of the lower density together with theoretical fits (solid lines) to be discussed later in this paper. It can be seen in Fig. 1c that the temperature dependence of the conductivity minimum around zero magnetic field, which we attribute to the WL corrections, is overpowered by the much stronger temperature dependence of the background conductivity. We therefore analyse the zero-field conductivity $`\sigma (T)`$ in the spirit of Ref.
$$\sigma (T)=\sigma _D(T)+\delta \sigma _{WL}(T)+\delta \sigma _I(T),$$
(1)
where $`\sigma _D(T)`$ is the Drude conductivity, $`\delta \sigma _{WL}(T)`$ and $`\delta \sigma _I(T)`$ are the WL and the interaction contributions, respectively. Strictly speaking, eq. (1) is based on the validity of the Fermi-liquid description and $`r_s1`$ is required. Nevertheless we empirically apply this concept to our system, since phenomenologically the magnetoconductivity exhibits all the features also found in samples with low $`r_s`$.
In the well established theory of WL spin-orbit relaxation mechanisms like the Elliott-Yafet mechanism, the Dyakanov-Perel mechanism or the Rashba-effect have been taken into account perturbatively , which is appropriate for conduction band electrons. WL in 2D p-type systems such as p-GaAs has only recently been studied in detail . In hole gases the valence band is strongly influenced by spin-orbit interaction and strain. The spin relaxation time is of the same order as the momentum relaxation time and can therefore no longer be treated perturbatively . It is therefore not a priori clear that the theories in Refs. can be applied to our system.
Here we apply the theory developed in Ref. for p-type 2D hole gases which takes the valence band structure into account and includes the effect of HH-LH mixing. Fig. 1c shows theoretical curves fitted to the measured conductivity according to
$`\delta `$ $`\sigma _{WL}(B,T)\delta \sigma _{WL}(0,T)={\displaystyle \frac{e^2}{2\pi ^2\mathrm{}}}\times `$ (3)
$`\left[f_2\left({\displaystyle \frac{B}{B_\phi +B_{||}}}\right)+{\displaystyle \frac{1}{2}}f_2\left({\displaystyle \frac{B}{B_\phi +B_{}}}\right){\displaystyle \frac{1}{2}}f_2\left({\displaystyle \frac{B}{B_\phi }}\right)\right]`$
Here $`f_2(x)=\mathrm{ln}x+\psi (1/2+1/x)`$, $`\psi `$ is the digamma function and $`B_i=\mathrm{}/(4De\tau _i)`$ ($`i=||,,\phi `$), with the diffusion constant $`D`$ and the phase-coherence time $`\tau _\phi `$. The other relaxation times are given in Ref. to be $`1/\tau _{||}=1/\tau _e\left(k_Fa/\pi \right)^4I_{||}`$ and $`1/\tau _{}=1/\tau _e\left(k_Fa/\pi \right)^6I_{}`$, where $`\tau _e`$ is the transport relaxation time, $`a`$ is the width of the quantum well and $`k_F`$ is the Fermi wavevector. The quantities $`I_{||/}`$ depend on the ratio of the LH and HH mass and have to be computed numerically. In our samples the holes are effectively confined by a triangular well due to the asymmetric doping and we estimate $`k_Fa/\pi 0.2`$. We use $`B_\phi ,B_{}`$ and $`B_{||}`$ as fitting parameters and find that $`B_{||},B_{}B_\phi `$, i.e. HH-LH mixing seems to be insignificant in our sample. This result is in contrast to the measurements on p-GaAs, where stronger HH-LH mixing leads to appreciable values for $`B_{}`$ and $`B_{||}`$ . In our case eq. (3) reduces to the result for negligible spin-orbit scattering already used in Ref. . As it was evaluated there the phase breaking rate is $`1/\tau _\phi T`$ and it increases with decreasing hole concentration as expected from theory.
Interaction corrections to the conductivity can be extracted from the temperature dependence of the Hall-resistance at small magnetic fields . It is predicted that the WL correction does not affect the Hall resistance, while interaction corrections obey $`\delta \sigma _I=\sigma \delta R_H/2R_H=e^2/(2\pi ^2\mathrm{})(13F^{}/4)\mathrm{ln}(kT\tau _e/\mathrm{})`$, where $`R_H=d\rho _{xy}/dB`$ is the Hall constant. Figure 2 shows $`\delta \sigma _I(T)`$ determined in this way for two different densities in the metallic regime. We find that $`\delta \sigma _I`$ is negative and determine $`F^{}=0.91`$ in agreement with Ref. . This means that interaction corrections to the conductivity give another contribution of insulating behaviour to the total conductivity in the metallic phase. This result which agrees with Ref. for p-GaAs is of special importance, since interaction corrections have been suggested to lead to metallic behaviour.
Now that we have experimentally determined $`\delta \sigma _{WL}(T)`$ and $`\delta \sigma _I(T)`$ we are in the position to extract the bare $`\sigma _D(T)`$ according to eq. (1) by subtracting the two corrections from the measured zero-field conductivity (Fig. 3). The obtained Drude conductivity shows a linear metallic temperature dependence which we will address in the following.
The temperature dependent Drude conductivity is
$`\sigma _D(T)={\displaystyle \frac{pe^2\tau _e(T)}{m^{}}}`$
where $`p`$ is the sheet density of the holes and the average Drude scattering time $`\tau _e(T)`$ has to be calculated according to
$`\tau _e(T)={\displaystyle \frac{𝑑EE\tau _e(E,T)(df/dE)}{𝑑EE(df/dE)}},`$
with the Fermi distribution function $`f(E,T)`$ and the energy dependent scattering rate
$`{\displaystyle \frac{\mathrm{}}{\tau _e(E,T)}}=2\pi N_i{\displaystyle \frac{d^2k^{}}{(2\pi )^2}\left|\frac{V(q)}{\epsilon (q,T)}\right|^2(1\mathrm{cos}\theta )\delta \left(E_kE_k^{}\right)}`$
Here $`N_i`$ is the density of ionized impurities, $`V(q)`$ is the matrix element for scattering by a wavevector $`q=k_F\sqrt{2(1cos\theta )}`$ and $`\epsilon (q,T)`$ is Lindhard’s dielectric function. For dominant large angle scattering (i.e. scattering for $`q2k_F`$) it has been shown that
$$\sigma _D(T)=\sigma _D(0)\left[1C(p)\frac{T}{T_F}\right]+𝒪\left[\left(\frac{T}{T_F}\right)^{3/2}\right].$$
(4)
The linear term stems from the temperature dependence of the dielectric function $`\epsilon (q,T)`$, which becomes particularly important for large angle scattering. It is a direct consequence of electron-electron interactions causing a temperature dependent polarizability of the 2D hole gas. Theory predicts values for the constant $`C(p)`$ which depend on the scattering mechanism and on the hole density $`p`$. If we apply eq. (4) to $`\sigma _D(T)`$ in Fig. 3 we determine $`C=3.1`$ in reasonable agreement with a predicted value of $`C2.8`$ for charged interface impurity scattering at low $`p`$. Similar agreement has been found for all densities in the metallic regime. This demonstrates that temperature dependent screening can indeed explain the metallic temperature dependence of the resistivity in our p-SiGe system without invoking a novel metallic phase. It also implies that at sufficiently low temperatures insulating behaviour with $`d\rho /dT<0`$ is expected to be recovered. Similar results in p-SiGe were obtained for much higher electron densities in Ref. .
For higher carrier densities $`p`$ we expect $`\sigma _D(T)`$ to exhibit a weaker temperature dependence, since $`C(p)`$ decreases . In Fig. 1a we therefore show the magnetic field and the temperature dependence measured at two densities. The lower surface shows the curves measured at $`p=2.6\times 10^{11}`$ cm<sup>-2</sup>. At large $`\left|B\right|`$ the metallic temperature dependence can be seen. At $`B=0`$ the WL peak tends to counteract this metallic behaviour without overpowering it in the range of temperatures shown. The upper surface was measured at a density of $`4.3\times 10^{11}`$ cm<sup>-2</sup> on an ungated part of the same sample. The temperature dependence at large $`\left|B\right|`$ is weaker than for the other surface ($`C=0.8`$), in agreement with temperature dependent screening. At $`B=0`$ however, the WL correction is strong enough to restore the insulating temperature dependence of the conductivity. Such a reentrant insulating behaviour is consistent with the observations made in Ref. .
In the following, we address the range of validity of our analysis. The theories of linear screening, WL and interaction corrections are expected to hold as long as the disorder in the system is low enough, i.e. as long as $`k_Fl>1`$. Therefore, our analysis is limited to the metallic regime. Non-linear screening models have to be invoked at the MIT, where our system undergoes the transition from WL ($`k_Fl>1`$) to strong localisation ($`k_Fl<1`$). We have empirically applied concepts developed for weakly interacting systems with $`r_s1`$. This attempt was successful and a Fermi-liquid description of the 2D hole gas appears to be consistent with the experiment, in spite of $`r_s8`$. The theoretical understanding and interpretation of this finding remains an open issue.
How relevant is this analysis for other systems in which the MIT was observed? Our analysis was simplified by the fact that large angle scattering is dominant in p-SiGe. For this special case, the analytical results of Ref. apply. Candidates for a similar analysis are therefore low-mobility Si-MOSFETS (for which this theory was originally developed) or the n-GaAs system with InAs quantum dots near the 2D electron gas . However, Kravchenko et al. stated in Ref. explicitly that temperature dependent screening can not account for the metallic behaviour in their high-mobility Si-MOSFET samples. Performing a similar analysis on the n-type GaAs samples of Ref. we find that the temperature dependence in the metallic range is much too large compared to the theory used above. For p-type GaAs samples exhibiting the MIT temperature dependent screening was suggested in Refs. as the relevant mechanism. Under the condition of small angle scattering this effect is much weaker than in our case.
In conclusion, we have analysed the magnetoresistance and the Hall resistance in p-SiGe samples in terms of interference and interaction corrections to the conductivity and found the measurements to be consistent with ordinary Fermi-liquid behaviour in spite of $`r_s8`$. Both corrections tend to localise the system as the temperature is lowered. The temperature dependence of the background Drude-conductivity has been found to depend linearly on temperature. Its behaviour is in good agreement with the theory of temperature dependent screening for systems in which large angle scattering dominates. The analysis applies to densities where the system is in the metallic regime and can not easily be extended to the transition region, where $`k_Fl1`$. Although these observations can not exclude the possibility of a novel metallic ground state in our or in other systems unambiguously, we are inclined to discard this exciting possible interpretation for p-SiGe hole gases on the basis of our experiments and analysis. The theoretical understanding of experimental consistency with Fermi-liquid behaviour at these large $`r_s`$-values remains a challenging topic for further research.
We have enjoyed fruitful discussions with V. Dolgopolov and J.L. Pichard. Financial support from ETH Zurich and the Schweizerischer Nationalfonds is gratefully acknowledged.
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# BCP3: Summary of Theory
## 1 CKM Fits: Where Do We Stand?
During BCP3 there was much discussion of how well we can determine the parameters in the CKM matrix. In practice, this means finding the allowed region, at a certain confidence level, in the $`\rho \eta `$ plane - with $`\rho `$ and $`\eta `$ being the usual Wolfenstein $`^\mathrm{?}`$ parameters. The ingredients of the global fits discussed here particularly by Stocchi $`^\mathrm{?}`$, as well as by Kim $`^\mathrm{?}`$ and Eigen $`^\mathrm{?}`$, are crucially dependent on the errors one assigns to the values one infers experimentally for $`|V_{ub}|`$ and $`|V_{cb}|`$. $`^\mathrm{?}`$ $`^\mathrm{?}`$ Although the data on $`\mathrm{sin}2\beta `$ coming from the Tevatron $`^\mathrm{?}`$, and now also from ALEPH, $`^\mathrm{?}`$ $`^\mathrm{?}`$ is not accurate enough to have much impact on the allowed region in the $`\rho \eta `$ plane, the sharpened LEP bound on $`\mathrm{\Delta }m_s`$, presented here by Willocq, $`^\mathrm{?}`$ provides a significant restriction.
When considering the result of the CKM fits, it is important to distinguish three different ingredients that enter into these fits
i) There are a number of experimental inputs which have quite negligible experimental errors. These include: the CP-violating parameter $`|ϵ|`$; the $`B_d\overline{B}_d`$ mass difference, $`\mathrm{\Delta }m_d`$; the $`\mathrm{\Delta }m_s`$ bound mentioned above; $`^\mathrm{?}`$ and the value of the sine of the Cabibbo angle, $`\mathrm{sin}\theta _C\lambda `$.
ii) There are a number of parameters which have, what one may call controllable errors. These include $`|V_{cb}|`$ ; the ratio $`|V_{ub}|/|V_{cb}|`$ and the $`SU(3)`$-breaking parameter $`\xi =f_{B_s}\sqrt{B_{B_s}}/f_{B_d}\sqrt{B_{B_d}}`$.
iii) There are also associated theoretical parameters, whose errors are due to theoretical uncertaintes in determining appropriate hadronic matrix elements. These include $`B_K`$, the parameter which measures how different the matrix element $`<K|(\overline{d}s)_{VA}^2|\overline{K}>`$ is from the vacuum saturation result $`(B_K=1)`$; and the analogous parameter for the $`B_d`$ system, $`B_df_{B_d}^2`$, which is connected to the matrix element $`<B_d|(\overline{d}b)_{VA}^2|\overline{B}_d>`$.
The parameters with controllable errors, in principle, are the ones which can be improved using more experimental information. <sup>a</sup><sup>a</sup>aI will discuss explicitly below how one can effect this error reduction in the case of $`|V_{cb}|`$. Similar ideas can be brought to bear on $`|V_{ub}|/|V_{cb}|`$ and also on $`\xi `$. In contrast, it is difficult to reduce, or even correctly estimate, the theoretical errors associated with parameters like $`B_K`$, because they are connected with the way one attempts to calculate the hadronic matrix elements. Although Soni $`^\mathrm{?}`$ in this meeting has given 15% errors on $`B_K`$ and $`f_{B_d}\sqrt{B_{B_d}}`$ $`[B_K=0.85\pm 0.13`$$`f_{B_d}\sqrt{B_{B_d}}=(230\pm 30)MeV]`$, the systematic error on each of these quantities is much harder to pin down, since it depends on the uncertainties associated with going from a quenched to a fully unquenched calculation of the relevant matrix elements.
The fit of Stocchi shown in Fig. 1 uses the parameters detailed in Table 1. Although I find this set of parameters (and the fit!) perfectly reasonable, I believe that one should not take the 95% confidence limit region shown in Fig. 1 strictly as such. There is enough ”theoretical error” in this whole business that the most prudent approach is to transform a ”formal” 95% confidence limit region into an ”effective” 68% limit region. Even after doing such an unorthodox thing, one cannot but be impressed by how well the CKM model fits the data!
As I mentioned earlier, the new result from LEP reported here by Willocq $`^\mathrm{?}`$ on the $`B_s\overline{B}_s`$ mass difference is of considerable importance for the final $`\rho \eta `$ plane fit. Last year this parameter was bounded a bit more weakly $`(\mathrm{\Delta }m_s>12.4\mathrm{ps}^1)`$ so that the final $`\rho \eta `$ plane fit permitted values of $`\rho <0`$$`^\mathrm{?}`$ With the present bound, however, one cannot contemplate any longer negative values for $`\rho `$— a point emphasized by Deshpande $`^\mathrm{?}`$ at this meeting. If $`\rho >0`$, it follows that the CKM CP-violating phase $`\gamma `$ cannot be as large as $`90^0`$. In turn, this implies that one cannot any longer contemplate the superweak solution for the unitarity triangle angles $`\alpha `$ and $`\beta `$, in which $`\mathrm{sin}2\beta =\mathrm{sin}2\alpha `$. Indeed, the fit of the CKM parameters presented by Stocchi here leads to a best value for $`\mathrm{sin}2\beta `$ and $`\mathrm{sin}2\alpha `$ which are quite different from each other: $`^\mathrm{?}`$
$$\mathrm{sin}2\beta =0.75\pm 0.06:\mathrm{sin}2\alpha =0.22\pm 0.24.$$
(1)
I remark, that the 95% confidence limit on $`\mathrm{sin}2\beta `$ which one infers from the CKM analysis,
$$0.62\mathrm{sin}2\beta 0.88,$$
(2)
is precisely in the range determined by present experiments. The value reported here for $`\mathrm{sin}2\beta `$, obtained by ALEPH $`^\mathrm{?}`$ $`^\mathrm{?}`$
$$\mathrm{sin}2\beta =0.93_{0.880.24}^{+0.64+0.36},$$
(3)
when combined with the CDF result $`^\mathrm{?}`$ \[$`\mathrm{sin}2\beta =0.79_{0.44}^{+0.41}`$\], leads to an average value
$$<\mathrm{sin}2\beta >=0.91\pm 0.35,$$
(4)
which is perfectly compatible with the range for $`\mathrm{sin}2\beta `$ obtained indirectly from the CKM analysis. Furthermore, this experimental result, by itself, is already almost significant statistically.
It is possible to imagine improving considerably the allowed value for $`\mathrm{sin}2\beta `$ and other unitarity triangle parameters by improving the controllable errors on $`|V_{ub}|/|V_{cb}|`$ and on the other parameters which enter in the CKM analysis. Let me illustrate how these improvements can come about by focusing specifically on the case of $`|V_{cb}|`$. If one could trust the parton model, and one knew the mass of the b-quark precisely, one could directly extract $`|V_{cb}|`$ from the semileptonic width for b-quarks to decay into final states containing c-quarks, given by the standard formula:
$$\mathrm{\Gamma }(bc\mathrm{}\overline{\nu }_e)=|V_{cb}|^2\frac{G_F^2}{192\pi ^3}m_b^5.$$
(5)
This formula, however, is of no use since $`m_b`$ itself is not known precisly enough and, furthermore, there are corrections to the parton model!
What one does, in practice, is to use the Heavy Quark Effective Theory (HQET) to replace $`m_b`$ by the mass of the B-meson. In so doing the rate (5) above is corrected by the addition of terms involving the matrix elements of the operators $`0_i`$ which enter in the operator produced expansion underlying the HQET. $`^\mathrm{?}`$ To $`0(1/m_b^2)`$ one has two new operators contributing to the total width, which are characterized by two parameters $`\lambda _1`$, and $`\lambda _2`$ given by
$`\lambda _1`$ $`=`$ $`{\displaystyle \frac{1}{2M_B}}<B|\overline{b}(iD)^2b|B>`$ (6)
$`\lambda _2`$ $`=`$ $`{\displaystyle \frac{g_3}{2M_B}}<B|\overline{b}\sigma _{\mu \nu }G^{\mu \nu }b|B>,`$
where $`G^{\mu \nu }`$ is the gluon field strength tensor. In addition, to this order in the heavy quark expansion, $`^\mathrm{?}`$ the Fermi momentum $`\overline{\mathrm{\Lambda }}`$ enters in the formulas relating the mass of the b-quark to that of the $`B`$ and $`B^{}`$ mesons:
$`M_B`$ $`=`$ $`m_b+\overline{\mathrm{\Lambda }}{\displaystyle \frac{(\lambda _1+3\lambda _2)}{m_b}}`$ (7)
$`M_B^{}`$ $`=`$ $`m_b+\overline{\mathrm{\Lambda }}{\displaystyle \frac{(\lambda _1\lambda _2)}{m_b}}.`$
Using the physical masses for $`M_B`$ and $`M_B^{}`$ yields directly a value for $`\lambda _2[\lambda _20.12\mathrm{GeV}^2]`$.
Using the heavy quark expansion one can obtain a formula for the semileptonic rate $`\mathrm{\Gamma }(BX_c\mathrm{}\overline{\nu }_{\mathrm{}})`$ which involves $`M_B`$ rather than $`m_b`$, at the price of some corrections involving $`\lambda _1,\lambda _2`$, and $`\overline{\mathrm{\Lambda }}`$. One finds, including leading order QCD effects, $`^\mathrm{?}`$ the formula:
$`\mathrm{\Gamma }(BX_c\mathrm{}\overline{\nu }_{\mathrm{}})`$ $`=`$ $`|V_{cb}|^2{\displaystyle \frac{G_F^2}{192\pi ^3}}M_B^5\mathrm{\Phi }({\displaystyle \frac{M_D}{M_B}})[11.54{\displaystyle \frac{\alpha _s}{\pi }}+\mathrm{}]`$
$``$ $`\{11.65{\displaystyle \frac{\overline{\mathrm{\Lambda }}}{M_B}}3.18{\displaystyle \frac{\lambda _1}{M_B^2}}+0.02{\displaystyle \frac{\lambda _2}{M_B^2}}+\mathrm{}\}.`$
Here $`\mathrm{\Phi }(M_D/M_B)`$ is a phase space factor, while the first square bracket contains the $`0(\alpha _s)`$ perturbative QCD corrections. The terms in the curly brackets are the dominant corrections coming from the heavy quark expansion. If, besides $`\lambda _2`$, one knew $`\lambda _1`$ and $`\overline{\mathrm{\Lambda }}`$ with good accuracy, one could determine $`|V_{cb}|`$ from the semileptonic width with a small $`0(1/m_b^3)`$ theory error. The parameters $`\lambda _1,\lambda _2`$, and $`\overline{\mathrm{\Lambda }}`$ are, in principle, controllable since they can be inferred from properties of the lepton spectrum and of the hadronic mass spectrum in non-leptonic $`B`$-decays . $`^\mathrm{?}`$ Specifically, moments of the lepton spectrum and/or the hadronic mass spectrum can be used to determine $`\lambda _1`$ and $`\overline{\mathrm{\Lambda }}`$.
Unfortunately, as shown in Fig. 2, a preliminary analysis by the CLEO collaboration — presented here by Thaler $`^\mathrm{?}`$ — give inconsistent results. The values for $`\lambda _1`$ and $`\overline{\mathrm{\Lambda }}`$ determined from the leptonic spectrum do not seem to agree with those computed from an analysis of the moments of the hadronic mass spectrum. In my view, the latter values $`[\lambda _1=(0.13\pm 0.01\pm 0.06)\mathrm{GeV}^2`$$`\overline{\mathrm{\Lambda }}=(0.33\pm 0.02\pm 0.08)`$ GeV\] are probably more trustworthy, since they are less sensitive to some of the necessary experimental cuts. Indeed, these values are quite close to the values reported by W. M. Zhang $`^\mathrm{?}`$ at this meeting, coming from a B-bound state analysis using an explicit phenomenological wavefunction for the B-meson. At any rate, using these techniques, it is likely that one will be able to eventually reduce the theoretical error on $`|V_{cb}|`$ to about $`1\times 10^3(2.5\%)`$ and on $`|V_{ub}|`$ to about $`4\times 10^4(10\%)`$, with a corresponding reduction of the allowed region in the $`\rho \eta `$ plane. As Misiak $`^\mathrm{?}`$ discussed in BCP3, this same approach can be used in other contexts — in particular, to better constrain the dependence on the photon energy of the branching ratio for $`BX_s\gamma `$ above some minimum photon energy cut.
## 2 Lifetimes of Heavy Hadrons
There is reasonable theoretical understanding of the specific pattern of lifetimes of heavy hadrons containing $`c`$ and $`b`$ quarks. This topic was reviewed here by Bigi $`^\mathrm{?}`$ and specific aspects were considered by Melic $`^\mathrm{?}`$ and by Yang $`^\mathrm{?}`$. However, as the data on these lifetimes is now rather precise, $`^\mathrm{?}`$ a puzzle has emerged. I want to discuss this briefly here.
One can understand the total decay width – and hence the lifetime – of a heavy hadron theoretically via the HQET. The total width is given by the discontinuity of the correlator of the weak Hamiltonian responsible for the decay with itself. This discontinuity is given by a sum of matrix elements of operators of increasing dimensonality, arising in an operator product expansion of the correlator. Schematically, one finds in this way the formula $`^\mathrm{?}`$
$`\mathrm{\Gamma }_{\mathrm{tot}}G_F^2m_Q^5`$ $`[<H_Q|\overline{Q}Q|H_Q>+{\displaystyle \frac{C_g}{m_Q^2}}<H_Q|\overline{Q}\sigma _{\mu \nu }G^{\mu \nu }Q|H_Q>`$
$`+{\displaystyle \frac{C_{4f}}{m_Q^3}}<H_Q|\overline{Q}\mathrm{\Gamma }q\overline{q}\mathrm{\Gamma }Q|H_Q>+\mathrm{}],`$
where $`m_Q`$ is the mass of the heavy quark. As can be gleaned from the above, the expansion parameter for the total width is of order $`\mathrm{\Lambda }_{QCD}/m_Q`$. Hence one expects that the effects of the non-leading terms in Eq. (9) should be bigger in charm decays than in bottom decays. Indeed, experimentally $`\tau (D^+)/\tau (D^o)2.5`$ while the $`B^+`$ and $`B_d^o`$ lifetimes are nearly equal $`[\tau (B^+)/\tau (B_d^o)=1.066\pm 0.024`$ $`^\mathrm{?}`$\]
More specifically, the situation regarding the decays of charmed hadrons seems rather satisfactory. $`^\mathrm{?}`$ The dominant effect that causes the large difference between the lifetimes of the $`D^+`$ and $`D^0`$ mesons can be traced to the Pauli interference $`^\mathrm{?}`$ originating from the 4-fermion term in Eq. (9). These 4-fermion terms are also crucial to explain the pattern of charmed baryon lifetimes, as detailed by Melic $`^\mathrm{?}`$ at this meeting. The fact that $`\tau (D_s)/\tau (D^o)`$ is near to 1.2 $`[\tau (D^o)=408\pm 4.1\mathrm{fs};\tau (D_s)=486\pm 15\mathrm{fs}`$ $`^\mathrm{?}`$\] is probably evidence for some weak annihilation contribution $`(c\overline{s}W^+)`$ in the $`D_s`$ width. $`^\mathrm{?}`$ If so, perhaps one should be able to see some evidence for glueball or $`\eta ^{}`$ decays of the $`D_s`$, processes like $`D_s`$ glueball $`\mathrm{}\overline{\nu }_{\mathrm{}}`$ ; $`D_s\eta ^{}\mathrm{}\overline{\nu }_{\mathrm{}}`$.
Although the prediction that the lifetime differences in the $`B`$-sector should be small is well borne out by the data, there is a nagging problem with the $`\mathrm{\Lambda }_b`$ lifetime. $`^\mathrm{?}`$ The Pauli interference effects suppress the lifetime of the $`\mathrm{\Lambda }_b`$ relative to that of the $`B_d`$-meson and one can write
$$\frac{\tau (\mathrm{\Lambda }_b)}{\tau (B_d)}=1\mathrm{\Delta }$$
(10)
Bigi $`^\mathrm{?}`$ estimates $`\mathrm{\Delta }0.05`$, while Melic $`^\mathrm{?}`$ can push $`\mathrm{\Delta }`$ to perhaps $`\mathrm{\Delta }0.1`$ by taking a very light $`b`$-quark mass. However, experimentally this ratio is closer to 0.8, with the latest results giving $`\tau (\mathrm{\Lambda }_b)/\tau (B_d)=0.79\pm 0.05`$$`^\mathrm{?}`$ It is possible that this signals a real discrepancy between theory and experiment that needs explaining. However, in my view, given the fine agreement of all the other heavy hadron lifetimes with the predictions coming from HQET, perhaps one should not worry overmuch at this stage about a (2-3) $`\sigma `$ discrepancy. Hopefully, future data from the Tevatron should help resolve this issue.
## 3 Charmless $`B`$-decays: Rates and CP-Asymmetries
At BCP3 there was lots of discussion of 2-body charmless $`B`$-decay. Data on these decays was presented by J. Smith $`^\mathrm{?}`$ and J. Alexander $`^\mathrm{?}`$, while various theoretical aspects were touched upon by H-Y Cheng $`^\mathrm{?}`$; C. D. Lu $`^\mathrm{?}`$, H.-n. Li $`^\mathrm{?}`$, M. Suzuki $`^\mathrm{?}`$and C. Bhattacharya. $`^\mathrm{?}`$
H.- Y. Cheng $`^\mathrm{?}`$ in particular, summarized the recent considerable progress achieved in calculating within QCD the relative rates for the decays of $`B`$-meson into 2 pseudoscalar states $`[BPP]`$, 2 vector states $`[BVV]`$ and a vector and a pseudoscalar state $`[BVP]`$ . The present treatment formalizes (and, in a sense, justifies and explains) the old diagrammatic approach to this problem of Ling Lie Chau. $`^\mathrm{?}`$ The basic idea of these calculations is well known. One starts with an effective weak Hamiltonian given by a sum of operators
$$H_{\mathrm{eff}}=\underset{i}{}C_i(\mu )0_i(\mu ),$$
(11)
and one tries to estimate the matrix elements of these operators by using factorization. However, then one correct this procedure by including (some) non-factorizable pieces.
Technically, when one uses factorization, in the process of splitting up the operators as a product of currents the scale dependence of the matrix elements of $`0_i(\mu )`$ is not preserved:
$$<M_1M_2|0_i(\mu )|B>=<M_1|J_i|B><M_2|J_i|B>.$$
(12)
However, one can reorganize the Wilson coefficient expansion $`^\mathrm{?}`$ to effectively recover $`\mu `$-independent coefficients $`C_i^{\mathrm{eff}}`$ \- at least to $`0(\alpha _s)`$. For Penguin operators these effects are rather large, amounting to about a $`50\%`$ increase for $`C_6^{\mathrm{eff}}`$$`^\mathrm{?}`$ Non-factorizable contributions are incorporated as $`1/N_c`$ corrections. As Cheng $`^\mathrm{?}`$ discussed, the recent advance comes from understanding that, depending on which operator $`0_i`$ one is considering, the correction factors have different $`N_c`$-factors associated with them . This has been justified in the heavy quark limit recently by Beneke, Buchalla, Neubert and Sachrajda. $`^\mathrm{?}`$ In particular, for $`(VA)(VA)`$ operators $`N_c^{\mathrm{eff}}2`$, while for $`(V+A)(VA)`$ operators $`N_c^{\mathrm{eff}}6`$.
H.- Y. Cheng $`^\mathrm{?}`$ compared the theoretical predictions resulting from the above considerations with experimental results. The overall comparison is quite good, but there are two extant problems. These are:
i) It is difficult to get the branching ratio $`BR(BKK)0.25BR(BK\pi )`$ as is observed experimentally. Typically, one finds much larger $`B\pi \pi `$ branching fractions than are observed.
ii) It is very hard to push the $`BR(B\eta ^{}K)`$ to the very large level observed $`(BR80\times 10^6)`$ Typically, theoretically one can reach at most $`BR(B\eta ^{}K)50\times 10^6`$. In fact, to get to values as large as this, it is necessary that $`N_c^{\mathrm{eff}}=2`$ for $`(VA)(VA)`$ operators, as is suggested theoretically.
J. Smith $`^\mathrm{?}`$ discussed a, theoretically less sophisticated, fit of the measured charmless $`B`$-decays. In this fit, which was carried out by Hou, Smith and W$`\ddot{\mathrm{u}}`$rtheim, $`^\mathrm{?}`$ one simply considers the amplitude for the B-decays in question to be given by a Tree plus a Penguin contribution, neglecting throughout any strong rescattering phases. Further, these authors use (naive) factorization to calculate the relevant matrix elements. Since the Tree amplitudes depended on the CKM phase $`\gamma `$ $`[T=|T|e^{i\gamma }]`$, this fit determines this phase. Remarkably the fit discussed by Smith $`^\mathrm{?}`$ determines $`\gamma `$ quite accurately:
$$\gamma =(114_{21}^{+25})^o.$$
(13)
However, it is difficult to judge the reliability of the approach. Furthermore, also in this case the branching fraction for the decay $`B\eta ^{}K`$ is too low $`[BR(B\eta ^{}K)30\times 10^6]`$.
The situation regarding direct CP-asymmetries is much less clear, both experimentally and theoretically, Experimentally, as Alexander $`^\mathrm{?}`$ indicated, one is statistics limited so that $`\delta A_{\mathrm{CPviol}.}`$ is typically of $`0(0.2)`$. For example, one has $`^\mathrm{?}`$
$`A_{\mathrm{CPviol}.}(K^\pm \pi ^{})`$ $`=`$ $`0.04\pm 0.16`$ (14)
$`A_{\mathrm{CPviol}.}(K^0\pi ^\pm )`$ $`=`$ $`+0.18\pm 0.24`$
Improvements will scale as $`1/\sqrt{N}`$ and one will need an integrated luminosity of $`0(100\mathrm{fb}^1)`$ to get to $`\delta A_{\mathrm{CPviol}.}0.04!`$
Theoretically to be able to predict these direct CP-asymmetries one needs a good estimate of the strong rescattering phases. If one writes for the amplitudes of two charge -conjugate processes
$$T=T_1+T_2e^{i\delta _w}e^{i\delta _s};\overline{T}=T_1+T_2e^{i\delta _w}e^{i\delta _s},$$
(15)
one sees that $`A_{\mathrm{CPviol}.}`$ vanishes if there is no strong rescattering phase $`\delta _s`$ between the two amplitudes:
$$A_{\mathrm{CPviol}.}=\frac{2r\mathrm{sin}\delta _w\mathrm{sin}\delta _s}{1+r^2+2r\mathrm{cos}\delta _w\mathrm{cos}\delta _s},$$
(16)
where $`r=T_1/T_2`$. For sizeable effects, in addition to having a non-negligible rescattering phase, one needs a relatively large weak CP-violating phase $`\delta _w`$ (which is probably OK in the CKM model) and a ratio $`r1`$. Since, typically, the two amplitudes involved are Penguin and Tree amplitudes, one needs these amplitudes to be comparable in size and to have a large rescattering phase between them — not too likely a possibility! Numerically, for example, if $`r=0.25,\delta _w\gamma =60^o`$ and $`\delta _s=30^o`$, one obtains $`A_{\mathrm{CPviol}.}0.1`$.
At BCP3 three different approaches were discussed to try to estimate the rescattering phases $`\delta _s`$ which one might expect. All three approaches have some inherent difficulties, demonstrating how challenging really it is to have a reliable estimate for $`\delta _s`$. C. D. Lu $`^\mathrm{?}`$ used the old idea of Bander, Silverman and Soni $`^\mathrm{?}`$ to extract a rescattering phase simply from the phase associated with the discontinuity of Penguin graphs. This phase depends on the momentum transfer carried by the gluon, $`\delta _s(k^2)`$. However, since the relevant $`k^2`$ are rather low, it is natural to question whether one can trust the discontinuity calculated in a pure quark picture to such values of $`k^2`$. In contrast, H.-n. Li $`^\mathrm{?}`$ in his talk at BCP3 estimated $`\delta _s`$ as the rescattering phase which emerges in the Brodsky-Lepage bound state formalism $`^\mathrm{?}`$ by means of factorization. In this case, the question is whether one can really use these techniques given the large energy release in the decay process $`BM_1M_2`$.
Mahiko Suzuki $`^\mathrm{?}`$ presented a more general discussion of the problem based on hadronic methods. He argued, I believe correctly, that it is really difficult to apply perturbative QCD ideas — even if one includes some resummation — for the process at hand, since the effective scale is only of order $`k_{\mathrm{eff}}\sqrt{\mathrm{\Lambda }_{QCD}m_b}1.5`$ GeV. However, Suzuki also pointed out that the rescattering phase $`\delta _s`$ is also difficult to estimate by hadronic methods. This is because the phase associated with, say, the $`|K\pi >`$ final state is not simply the phase associated with elastic $`K\pi `$ scattering. For the energies in question
$$|K\pi >=\underset{f}{}e^{i\delta _f}|\alpha _f>,$$
(17)
and most of the states that contribute in the above sum are inelastic (roughtly 80%). Given this fact, Suzuki tried to estimate the effective phase $`\delta _s`$ that emerges by using a random phase approximation. Using this approximation, he was able to correlate the resulting phase $`\delta _s`$ with the elasticity of the process, with the result depending on how favored or disfavored is the factorization of the amplitude. When factorization is not favoreda $`\delta _s`$ is bigger. However, if one really has a large rescattering phase $`\delta _s`$ there are also sizable distorsions of the amplitude — a result which Suzuki $`^\mathrm{?}`$ points out was first understood by Fermi! Basically, not only is there a change in the imaginary part associated with the amplitude, but also the real part is affected:
$$TT\mathrm{exp}\left[\frac{P}{\pi }\frac{ds^{}\delta _s(s^{})}{s^{}s}\right].$$
(18)
These considerations make it unlikely that one will be able to really ever get a reliable theoretical prediction for a direct CP-asymmetry. However, this fact should not discourage experimentalists from looking for such CP-asymmetries.
## 4 Imputs from the Strange Quark Sector
In the last year, the new results on $`ϵ^{}/ϵ`$ announced by KTeV $`^\mathrm{?}`$ and by the NA48 Collaboration $`^\mathrm{?}`$ have generated an enormous amount of interest. Not surprisingly, this subject was also a topic of considerable discussion at BCP3.
### 4.1 $`ϵ^{}/ϵ`$ Results and their Interpretation
In BCP3 the experimental situation regarding $`ϵ^{}/ϵ`$ was reviewed by A. Roodman $`^\mathrm{?}`$ and by A. Nappi $`^\mathrm{?}`$, while the theoretical aspets of these recent results were discussed by Bertolini, $`^\mathrm{?}`$ Soni, $`^\mathrm{?}`$ Masiero, $`^\mathrm{?}`$ and Soldan. $`^\mathrm{?}`$ In my view, the nicest aspect of the new KTEV/NA48 results is that they resolve the discrepancy that existed betwen the old Fermilab result, coming from the E731 experiment, $`^\mathrm{?}`$ and the results of the old CERN experiment NA31. $`^\mathrm{?}`$ The present results for $`ϵ^{}/ϵ`$:
$$\mathrm{Re}ϵ^{}/ϵ=\{\begin{array}{c}(28.0\pm 4.1)\times 10^4\mathrm{Ref}[35]\\ (18.5\pm 7.3)\times 10^4\mathrm{Ref}[36]\end{array}$$
(19)
when combined with the E731 and NA31 results lead to a world average value for this quantity:
$$<\mathrm{Re}ϵ^{}/ϵ>=(21.2\pm 4.6)\times 10^4,$$
(20)
which clearly establishes the existence of $`\mathrm{\Delta }S=1`$ CP-violation. This result provides the first direct confirmation that $`\mathrm{sin}\gamma 0`$ or, equivalently, that $`\eta 0`$— something which could only be inferred from the CKM analysis we discussed earlier.
This said, however, it is difficult to extract from the present measurement of $`ϵ^{}/ϵ`$ a precise value for $`\eta `$. This is because the value of $`ϵ^{}/ϵ`$, although proportional to $`\eta `$, is rather uncertain due to uncertainties in the calculation of the relevant hadronic matrix element. This point was emphasized by Bertolini $`^\mathrm{?}`$ at this meeting. His argument can be illustrated by making use of an approximate formula for $`ϵ^{}/ϵ`$ due to Buras and his collaborators. $`^\mathrm{?}`$ As is well known, both matrix elements of gluonic Penguin operators and of electroweak Penguin operators contribute to the process $`sd\overline{q}q`$. Although the gluonic Penguin contributions are of $`0(\alpha _s)`$, and hence enhanced with respect to the electroweak Penguin contributions which are of $`O(\alpha )`$, they are suppressed by the $`\mathrm{\Delta }I=1/2`$ rule since they are proportional to $`A_2/A_0`$$`^\mathrm{?}`$ The electroweak Penguin contributions, on the other hand, although naturally small do not suffer from the $`\mathrm{\Delta }I=1/2`$ suppresion and, furthermore, are enhanced by a factor of $`m_t^2/M_W^2`$$`^\mathrm{?}`$ As a result, for $`ϵ^{}/ϵ`$ both Penguin operator contributions are of a similar magnitude and, because they enter with opposite sign, they render the theoretical value for this quantity rather uncertain.
This can be appreciated quite nicely from the approximate formula for $`ϵ^{}/ϵ`$ derived by Buras and his collaborators. $`^\mathrm{?}`$ One has
$$ϵ^{}/ϵ34\eta \left[B_6^{1/2}0.53B_8^{1/2}\right]\left(\frac{110MeV}{m_s(2GeV)}\right)^2\times 10^4.$$
(21)
Here $`B_6`$ and $`B_8`$ are the matrix elements of the gluonic Penguin operator and of the electroweak Penguin operator, respectively — normalized so that in the vacuum insertion approximation $`B_6=B_8=1`$. Because, as we have seen, $`\eta 0.3`$, the above formula in this approximation gives $`ϵ^{}/ϵ5\times 10^4`$, which is much below the value measured by KTeV and NA48. Indeed, to get agreement with the present world average \[cf Eq. (20)\], one needs to stretch all the parameters in Buras’s formula. Namely, one needs to maximize $`\eta `$ \[ $`\eta 0.4]`$; decrease the value one assumes for $`m_s(2GeV)`$ \[perhaps to $`m_s`$ as low as $`m_s(2GeV)90`$ MeV\]; increase $`B_6^{1/2}`$ from unity — something which is not clear one can obtain in lattice calculations, $`^\mathrm{?}`$ but which appears to be true in the chiral quark model $`^\mathrm{?}`$; and decrease $`B_8^{1/2}`$ — something which emerges naturally both in lattice and $`1/N`$ calculations. $`^\mathrm{?}`$
In this meeting, Bertolini $`^\mathrm{?}`$ emphasized that what improves the post-dictions of $`ϵ^{}/ϵ`$ is to incorporate in the calculation the trend observed in the chiral quark model and in the $`1/N`$ approximation that $`B_6^{1/2}/B_8^{1/2}2`$, and not unity as expected in the vacuum insertion approximation. However, this is just a phenomenological observation and one will not really trust a theoretical result for $`ϵ^{}/ϵ`$ until the lattice results quantitatively arrive at a value for $`ϵ^{}/ϵ`$. Unfortunately, at the moment, there is considerable controversy on what this value might be. This was clear from Soni’s talk $`^\mathrm{?}`$ at BCP3.
Soni $`^\mathrm{?}`$ argued forcefully that the extant lattice calculations for $`ϵ^{}/ϵ`$ which use staggered fermions are not to be trusted, since they have trouble correctly incorporating the needed counterterms and cannot account well for chiral mixing. However, when one uses domain wall fermions, which have the correct chiral behaviour, one arrives at a result for $`ϵ^{}/ϵ`$ which is difficult to believe: $`^\mathrm{?}`$
$$[ϵ^{}/ϵ]_{DWF}=(120\pm 60)\times 10^4.$$
(22)
Not only is the sign of $`ϵ^{}/ϵ`$ reversed from what one measures experimentally — due to a large negative contribution from the, so called, ”eye diagrams” $`^\mathrm{?}`$ — but the overall magnitude is also rather large.
Given the theoretical disarray concerning $`ϵ^{}/ϵ`$, it appears to me premature to try to invoke the presence of some new physics to $`\mathrm{`}\mathrm{`}`$explain” the experimental value of $`ϵ^{}/ϵ`$. However, considering new physics with correlated predictions is interesting. For example, as discussed by Masiero $`^\mathrm{?}`$ here, in supersymmetric models with a large $`\delta _{RL}`$ phase, not only does one get $`ϵ^{}/ϵ`$ to be large, but one also predicts largish values for the electron dipole moment and the lepton flavor violating process $`\mu e\gamma `$. As Pakvasa $`^\mathrm{?}`$ observed, such models also tend to give a rather large CP-violating contribution in $`\mathrm{\Lambda }`$-decays.
### 4.2 Rare K-decays
In contrast to $`ϵ^{}/ϵ`$, as D’Ambrosio $`^\mathrm{?}`$ discussed in BCP3, there are rare K-decays where the theoretical prediction are on much firmer ground. In particular, the ”golden modes” $`K_L\pi ^0\nu \overline{\nu }`$ (whose branching ratio is proportional to $`\eta ^2`$) and $`K^+\pi ^+\nu \overline{\nu }`$ (whose branching ratio is proportional to $`|V_{td}|^2`$) have theoretical errors of order 1% and 5%, respectively. The Brookhaven experiment BNL E787 has observed one event corresponding to this latter process and one infers a branching ratio $`^\mathrm{?}`$
$$BR(K^+\pi ^+\nu \overline{\nu })=(1.5_{1.3}^{+3.5})\times 10^{10},$$
(23)
which is totally consistent with the expectations of the standard model $`[(BR)_{SM}=(0.82\pm 0.32)\times 10^{10}]`$. There are strong hopes $`^\mathrm{?}`$ that the proposed BNL experiment E949 will be able to actually pin down a value for $`|V_{td}|`$, since they expect of order 10 events at the SM level of sensitivity.
Clearly, since the branching ratio for $`K_L\pi ^0\nu \overline{\nu }`$ is directly proportional to $`\eta ^2`$ and it has a negligible theoretical error, a measurement of this process would have a significant impact on our knowledge of the CKM matrix. Unfortunately, the decay $`K_L\pi ^0\nu \overline{\nu }`$ is extremely challenging experimentally, since one expects a very small branching ratio in the SM $`[(BR)_{SM}3\times 10^{11}]`$ and, further, one has an all neutral final state. Nevertheless, as discussed here by Hsiung $`^\mathrm{?}`$ there are a number of proposed experiments in a planning stage, both at Brookhaven \[KOPIO\], Fermilab \[KAMI\] and KEK \[E391\]. However, to determine $`\eta `$ at a significant level (with error of order $`\delta \eta /\eta 0.1`$) one needs to measure the $`K_L\pi ^0\nu \overline{\nu }`$ branching ratio to an accuracy of order 20%. This is very hard indeed!
As Pakvasa $`^\mathrm{?}`$ emphasized in his talk at BCP3, the theoretical predictions for CP-violation in Hyperon decays are likely to have less uncertainty than those for $`ϵ^{}/ϵ`$. However, in Hyperon decays, at least in the standard model, the expectations for CP-violating phenomena are very small and appear to be well below the present experimental reach. $`^\mathrm{?}`$ Typically, one expects theoretically for $`\mathrm{\Lambda }`$-decays a CP-violating contribution of order $`A_{\mathrm{CPviol}.}(\mathrm{\Lambda })10^5`$, while experimentally one can reach only a level of order $`A_{\mathrm{CPviol}.}(\mathrm{\Lambda })|_{exp}10^310^4`$.
## 5 Looking at the Future: Determing the Phases in the Unitarity Triangle
The most interesting challenge of the B-factories now entering into operation, and of future collider B experiments, is to try to pin down the angles in the unitarity triangle. As Trischuk $`^\mathrm{?}`$ discussed at this meeting for the most favorable mode, involving the decay $`B_d\psi K_s`$, one can expect to reach the level $`\delta \mathrm{sin}2\beta 0.1`$ at the B-factories, with an integrated luminosity of $`50\mathrm{fb}^1`$. To reach this same level of accuracy at the Tevatron, one probably will need about $`2\mathrm{fb}^1`$ of integrated luminosity. These luminosity numbers are likely to be achieved in the next five years. Although the ”uncertainty” $`\delta \mathrm{sin}2\beta 0.1`$ is precisely that which obtains from present day CKM fits, $`^\mathrm{?}`$ the results on $`\mathrm{sin}2\beta `$ both at the B-factories and at the Tevatron will involve an actual measurement of a CP-violating asymmetry. Thus they are extremely important— even if they probably will not significantly improve the knowledge of $`\mathrm{sin}2\beta `$ obtained indirectly through CKM fits.
In the same vein, it is also important to check that the value of $`\mathrm{sin}2\beta `$ measured in a variety of different physical processes is in fact the same. Indeed, as remarked by a number of people at this meeting , $`^\mathrm{?}`$ $`^\mathrm{?}`$ this may well be the best way to look for physics beyond the standard model. For instance, if there were to be an extra $`bs`$ Penguin phase $`\mathrm{\Phi }_P`$, the process $`B_d\psi K_s`$ would still approximately measure $`\mathrm{sin}2\beta `$, while the process $`B_d\mathrm{\Phi }K_s`$ (which is Penguin dominated) would measure $`\mathrm{sin}(2\beta +\mathrm{\Phi }_P)`$.
Extracting the two other angles in the unitarity triangle, $`\alpha `$ and $`\gamma `$, from experiment is likely to be much more challenging. This will require measuring many processes, as Sheldon Stone $`^\mathrm{?}`$ emphasized in his nice overview at BCP3. Furthermore, one will need to use, at the same time, theory input rather judiciously. $`^\mathrm{?}`$ $`^\mathrm{?}`$ That this is so can be appreciated in a number of ways. For example, as is well known, $`^\mathrm{?}`$ the process $`B_d\pi ^+\pi ^{}`$ does not measure purely $`\mathrm{sin}2\alpha `$ since there is likely significant Penguin pollution. In principle, one could imagine estimating these effects by studying $`B_d\pi ^o\pi ^o`$$`^\mathrm{?}`$ However, this decay is estimated to have a $`BR10^7`$, which is too small to make this an effective means to control Penguin pollution. Deshpande $`^\mathrm{?}`$ at this meeting suggested what may be a useful alternative. Namely, using a theoretical calculation to extract $`\alpha `$ from the value $`\alpha _{\mathrm{meaus}}`$ gotten experimentally from the $`B_d\pi ^+\pi ^{}`$ asymmetry. The model calculations he presented appeared rather encouraging.
Gronau $`^\mathrm{?}`$ discussed, analogously, how theoretical input — in this case the $`SU(3)`$ transformation properties of the weak Hamiltonian — can be used to constrain $`\gamma `$ for the $`\mathrm{\Delta }S=1`$ $`BK\pi `$ process. The piece of the decay amplitude which transforms as the $`\overline{15}`$ dimensional representation in the Tree amplitude essentially fixes the electroweak-Penguin amplitude
$$T_{EWP}(\overline{15})=\delta _{EW}T_{\mathrm{tree}}(\overline{15}),$$
(24)
with $`\delta _{EW}=0.65\pm 0.15`$ a calculable number in the standard model. Then using that $`|K^o\pi ^+>+\sqrt{2}|K^+\pi ^o>`$ transforms as the 27 representation one can obtain useful interelations among $`BK\pi `$ process, from which one can extract $`\gamma `$. These techniques can be used also in other contexts and Gronau $`^\mathrm{?}`$ estimated that they may lead to a determination of $`\gamma `$ with an error $`\delta \gamma =\pm 20^o`$.
Obviously, with new data coming in from the B-factories, and soon also from the Tevatron, we should expect interesting news for BCP4!
## Acknowledgements
I would like to thank both George Hou and Hai-Yang Cheng for their splendid hospitality in Taipei. This work was supported in part by the Department of Energy under Contract DE-FE03-91ER40662, Task C.
## References
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# 1 Introduction
## 1 Introduction
Since most of the path integrals cannot be computed exactly, different methods of approximation have been developed. The perturbation expansion is the most familiar but it diverges for all couplings strength. An another useful procedure found by Kleinert is the convergent variational perturbation expansion. In field theory the ’exact’ Renormalization Group (RG) equations like the Wegner and Houghton one have led to numerous non-perturbative results. In the path integral approach a one particle system is similar to a one dimensional spin lattice , then the RG method could be also useful to extract some non-perturbative results for this system. In this paper we derive the RG equation for the potential of a quantum particle which at finite temperature is not a closed expression contrary the field theory case. We show also that we can improve the Feynman-Kleinert’s variational method, but cannot compete with the efficiency of Kleinert’s systematic variational perturbation theory.
## 2 Renormalization Group Equations in Quantum Mechanics
Consider the euclidean action of a quantum particle at a finite temperature
$$S(x)=_0^T\left(\frac{1}{2}M\left(\frac{d}{dt}x(t)\right)^2+V(x(t))\right),$$
(1)
with $`M`$ the mass, $`V`$ the potential and $`T=\overline{h}\beta `$.
The effective potential is defined as a constrained path integration over periodic paths with period $`T=\overline{h}\beta `$
$$\mathrm{exp}\left(\beta V_0(x_0)\right)=𝒟x\delta (\overline{x}x_0)\mathrm{exp}\left(\frac{1}{\overline{h}}S(x)\right).$$
(2)
where $`\overline{x}=\frac{1}{T}𝑑tx(t)`$.
The partition function is
$$Z=\frac{dx_0}{\sqrt{\frac{2\pi \overline{h}T}{M}}}\mathrm{exp}\left(\beta V_0(x_0)\right),$$
(3)
We consider the Feynman path integral with a discretized time $`t_n=\frac{nT}{N+1}=nϵ`$ with $`N`$ an arbitrary large number, and $`n=0,\mathrm{},N+1`$. The Fourier decomposition of a periodic path $`x(t_n)`$ contains only a finite number of Fourier modes
$$x(t_n)=x_0+\frac{1}{\sqrt{N+1}}\stackrel{}{}\mathrm{exp}(i\omega _mt_n)x_m+H.C,$$
where $`^{}`$ is from $`1`$ to $`\frac{N}{2}`$ if $`N`$ is even and from $`1`$ to $`\frac{N1}{2}`$ if $`N`$ is odd. The $`x_m`$ are the Fourier modes and $`\omega _m^2=\frac{22cos\frac{2\pi m}{N+1}}{ϵ^2}`$.
The discretized partition function is then
$$Z=\frac{dx_0}{\sqrt{\frac{2\pi \overline{h}ϵ}{M}}}\mathrm{\Pi }_1^{\frac{N}{2}}\frac{dx_md\overline{x}_m}{\frac{2\pi ϵ\overline{h}}{M}}\mathrm{exp}\left(\frac{1}{\overline{h}}S_N\right),$$
where the discretized action is given by
$$S_N(x)=ϵ\underset{0}{\overset{\frac{N}{2}}{}}M\omega _m^2|x_m|^2ϵ\underset{n=1}{\overset{N+1}{}}V(x(t_n)).$$
Using the fact that $`\mathrm{\Pi }_1^{\frac{N}{2}}ϵ^2\omega _m^2=\sqrt{N+1}`$ and $`T=(N+1)ϵ`$, we define the effective potential as (see ):
$$\mathrm{exp}\left(\beta V_0(x_0)\right)=\mathrm{\Pi }_1^{\frac{N}{2}}\frac{dx_md\overline{x}_m}{\frac{2\pi ϵ\overline{h}}{ϵ^2\omega _m^2M}}\mathrm{exp}\left(\frac{1}{\overline{h}}S_N\right).$$
(4)
Instead of evaluating (4) in one step we integrate mode after mode. Let us denote $`V_m`$ the ’running potential’ obtained after $`\frac{N}{2}m`$ integrations, and $`V_{\frac{N}{2}}`$ the initial potential $`V`$. To find the potential $`V_{m1}`$ with respect to $`V_m`$ we consider paths with only one Fourier mode:
$$x(t_n)=x_0+\frac{1}{\sqrt{N+1}}exp(i\omega _mt_n)x_m+H.C.$$
We thus obtain the relation:
$`\mathrm{exp}\left(\beta V_{m1}(x_0)\right)`$
$`={\displaystyle }{\displaystyle \frac{dx_md\overline{x}_m}{\frac{2\pi ϵ\overline{h}}{ϵ^2\omega _m^2M}}}\mathrm{exp}({\displaystyle \frac{ϵ}{\mathrm{}}}(M\omega _m^2|x_m|^2+{\displaystyle \underset{n=0}{\overset{N+1}{}}}V_m(x_0+{\displaystyle \frac{exp(i\omega _mt_n)x_m}{\sqrt{N+1}}}+H.C.)))`$
(5)
Expanding the potential $`V_m`$ around the point $`x_0`$ and summing over n, yields the following RG equation for the potential
$`V_{m1}(x_0)`$ $`=`$ $`V_m(x_0)+{\displaystyle \frac{1}{\beta }}\mathrm{log}(1+{\displaystyle \frac{V_m^{(2)}(x_0)}{\omega _m^2M}})+{\displaystyle \frac{1}{\beta }}\mathrm{log}(1+{\displaystyle \underset{n2}{}}{\displaystyle \frac{n!}{\beta ^{n1}}}`$
$`{\displaystyle \underset{k_1p_1+\mathrm{}+k_np_n=n}{}}(1)^{k_1+\mathrm{}+k_n}{\displaystyle \frac{\left(V_m^{(2p_1)}(x_0)\right)^{k_1}}{k_1!(p_1!)^{2k_1}}}\mathrm{}{\displaystyle \frac{\left(V_m^{(2p_n)}(x_0)\right)^{k_n}}{k_n!(p_n!)^{2k_n}}}P_m^n),`$
where all the $`p_i2`$ are different and the propagator $`P_m=\frac{1}{M\omega _m^2+V_m^{(2)}(x_0)}`$.
Note that in this derivation we have ignored the evolution of higher derivative interactions. The equation (LABEL:poteff1) has not the closed form of the ’exact’ WH differential equation in field theory which takes into account only the one loop contributions resumed in a logarithm term. In field theory the higher order loop terms are negligible even at finite temperature but their influence is taken into account through the running of the parameters of the theory. In (LABEL:poteff1) the one loop contributions are resumed in the first logarithm term but the second one contains the higher loop terms. To compare with the usual perturbation expansion we can rewrite the RG equation in the familiar cumulant expansion. Let us define $`𝒜_m=ϵ_{n=0}^{N+1}\left(V_m(x(t_n))V_m(x_0)\frac{1}{2}V_m^{(2)}(x_0)x^2(t_n)\right)`$, and denote the expectation value with respect to the gaussian weight $`M\omega _m^2+V_m^{(2)}(x_0)`$ by $`<>_m`$. It is straightforward to derive:
$$\begin{array}{c}\text{ }\frac{1}{l!}(\frac{1}{\overline{h}})^l<𝒜_m^l>_m=\underset{n2}{}\frac{n!}{\beta ^{n1}}\text{ }\hfill \\ \text{ }\underset{k_1p_1+\mathrm{}+(lk_1\mathrm{}k_{n1})p_n=n}{}(1)^{k_1+\mathrm{}+k_n}\frac{\left(V_m^{(2p_1)}(x_0)\right)^{k_1}}{k_1!(p_1!)^{2k_1}}\mathrm{}\frac{\left(V_m^{(2p_n)}(x_0)\right)^{lk_1\mathrm{}k_{n1}}}{k_n!(p_n!)^{2k_n}}P_m^n.\text{ }\hfill \end{array}$$
Then defining the cumulants by
$$<𝒜_m^n>_{m,c}=<\left(𝒜_m<𝒜_m>_m\right)^n>_m,$$
the Renormalization Group equation becomes:
$`V_{m1}(x_0)`$ $`=`$ $`V_m(x_0)+{\displaystyle \frac{1}{\beta }}\mathrm{log}\left(1+{\displaystyle \frac{V_m^{(2)}(x_0)}{\omega _m^2M}}\right)+{\displaystyle \frac{1}{\overline{h}\beta }}<𝒜_m>_m`$ (7)
$``$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{l2}{}}{\displaystyle \frac{1}{l!}}({\displaystyle \frac{1}{\overline{h}}})^l<𝒜_m^l>_{m,c}`$
In this formulation the equation has the disadvantage to mix the powers of $`\frac{1}{\beta }`$.
Clearly the RG equation seems not to be useful to extract nonperturbative results. Actually it contains an infinite number of terms, negligible only if $`(\frac{V_m}{\beta })^p`$ are small (for finite temperature the coupling constants must be small). Nevertheless it is interesting to discuss the various limiting cases.
### 2.1 Infinite temperature
For a large temperature the quantum fluctuations are small and the system is close to the classical one. The flow of the potential stays near the classical one (the potential energy of the action) and will be obtained after a relatively small number of iterations of the RG equation. In particular in the limit $`\beta 0`$, $`\frac{V_m^{(2)}(x_0)}{\omega _m^2M}`$ and $`P`$ are of order $`\beta ^2`$ and we obtain
$$V_{m1}(x_0)=V_m(x_0)$$
a running potential constant along the flow. As a consequence the quantum partition function reduces to the classical one.
### 2.2 Zero temperature
In this opposite case due to quantum fluctuations the effective potential is expected to be different from the classical one and will be obtained after a huge number of iterations of the RG equation (typically $`10^8`$ for $`\beta =10^5`$). The higher loop contributions are negligible, then we can extract informations for large coupling constants. In particular, in the limit $`\beta \mathrm{}`$ the propagator $`P_m`$ is of order $`\frac{1}{V_m^{(2)}(x_0)}`$. Then keeping only the one loop contributions in (LABEL:poteff1) we obtain the closed form:
$$V_{m1}(x_0)=V_m(x_0)+\frac{1}{\beta }\mathrm{log}\left(1+\frac{V_m^{(2)}(x_0)}{\omega _m^2M}\right).$$
(8)
Recall that for large $`N`$, $`\omega _m^2`$ is equivalent to $`\left(\frac{2\pi m}{\overline{h}\beta }\right)^2`$ and then runs from $`0`$ to $`\left(\frac{\pi }{ϵ}\right)^2`$. Rewriting (8) as
$$V_{m1}(x_0)=V_m(x_0)+\frac{\overline{h}}{2\pi }\frac{\pi }{ϵ}\frac{2}{N+1}\mathrm{log}\left(1+\frac{V_m^{(2)}(x_0)}{\left(\frac{\pi }{ϵ}\frac{2m}{N+1}\right)^2M}\right),$$
the limit $`N\mathrm{}`$ gives a continuous equation
$$V_{k\mathrm{\Delta }k}(x_0)=V_k(x_0)+\frac{\overline{h}\mathrm{\Delta }k}{2\pi }\mathrm{log}(1+\frac{V_k^{(2)}(x_0)}{Mk^2}),$$
(9)
where we have introduced the notations $`k^2=\omega _m^2`$, $`\mathrm{\Delta }k=\frac{2\pi m}{ϵ(N+1)}`$, and neglected the higher loop contributions which are of order $`(\mathrm{\Delta }k)^2`$. This equation is the one dimensional Wegner Houghton equation for the effective potential . It is useful to extract informations even for very large coupling constants as shown below.
## 3 Examples
We solve exactly the Renormalization Group equation for some easy examples and recover known results.
### 3.1 The free particle
This case is trivial because $`V_{\frac{N}{2}}(x_0)=0`$. It’s easy to show recursively that at each step $`V_{m1}(x_0)=V_m(x_0)=0`$. Then, the effective potential $`V_0(x_0)=0`$.
### 3.2 The harmonic oscillator
As a second example, consider the case: $`V_{\frac{N}{2}}(x_0)=\frac{M}{2}\mathrm{\Omega }^2x_0^2`$. Once again, we obtain recursively: $`V_{m1}(x_0)=\frac{M}{2}\mathrm{\Omega }^2x_0^2+C_m`$ with $`C_m`$ a constant. Inserting this result in (LABEL:poteff1), we obtain the solution
$$V_0(x_0)=\frac{M}{2}\mathrm{\Omega }^2x_0^2+\frac{1}{\beta }\mathrm{log}\left(\mathrm{\Pi }_1^{\frac{N}{2}}\left(\frac{\omega _m^2+\mathrm{\Omega }^2}{\omega _m^2}\right)\right),$$
which gives the partition function (see ):
$$Z=\frac{1}{\overline{h}\beta \mathrm{\Omega }}\mathrm{\Pi }_1^{\frac{N}{2}}\left(\frac{\omega _m^2}{\omega _m^2+\mathrm{\Omega }^2}\right)$$
### 3.3 Perturbation expansion
Replacing $`V_m`$ by $`V_{\frac{N}{2}}`$ in (7) one recovers the usual perturbative expansion
$`V_0(x_0)`$ $`=`$ $`V_{\frac{N}{2}}(x_0)+{\displaystyle \underset{m=1}{\overset{\frac{N}{2}}{}}}{\displaystyle \frac{1}{\beta }}\mathrm{log}\left(1+{\displaystyle \frac{V_{\frac{N}{2}}^{(2)}(x_0)}{\omega _m^2M}}\right)+{\displaystyle \frac{1}{\overline{h}\beta }}{\displaystyle \underset{m=1}{\overset{\frac{N}{2}}{}}}<𝒜_{\frac{N}{2}}>_m`$
$`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{l2}{}}{\displaystyle \frac{1}{l!}}({\displaystyle \frac{1}{\overline{h}}})^l{\displaystyle \underset{m=1}{\overset{\frac{N}{2}}{}}}<𝒜_{\frac{N}{2}}^l>_{m,c}`$
## 4 The anharmonic oscillator
In this part we compute the ground state energy of the anharmonic oscillator at zero temperature for various values of the coupling constant by using the RG equation, and compare with Kleinert’s variational results. The ground state energy is the value of the effective potential $`V_0(x_0)`$ at its minimum $`x_0=0`$. Unfortunately it is a hard task to use the RG equation (8) because at each step we must find the second derivative of the potential by fitting it numerically. Due to the very slow convergence we must repeat the procedure a huge number of time for each point of the potential. We leave this method for an another paper and try instead to compute the flow equations for the coupling constants.
We define the n-th coupling constant at the scale $`m`$ by $`g_m^{(n)}=\frac{d^n}{dx_0^n}_{|x_0=0}V_m(x_0)`$.
We give a general formulation for the flow of the couplings where we neglect terms of higher order in $`\frac{1}{\beta }`$:
$$g_{m1}^{(0)}=g_m^{(0)}+\frac{1}{\beta }\mathrm{log}\left(1+\frac{g_m^{(2)}}{\omega _m^2M}\right).$$
(10)
For $`k0`$ we have
$$g_{m1}^{(k)}=g_m^{(k)}+\frac{1}{\beta }\left(\underset{p}{}(1)^{p1}P_m(0)^p\left(\underset{\alpha _1+\mathrm{}\alpha _p=k,\alpha _i>0}{}\frac{k!A_p}{p}\frac{g_m^{(\alpha _1+2)}}{\alpha _1!}\mathrm{}\frac{g_m^{(\alpha _p+2)}}{\alpha _p!}\right)\right),$$
(11)
where $`A_p`$ is the combinatorial factor of $`\frac{g_m^{(\alpha _1+2)}}{\alpha _1!}\mathrm{}\frac{g_m^{(\alpha _p+2)}}{\alpha _p!}`$ in the series expansion of $`\left(V^{(2)}(x_0)\right)^p`$ in powers of $`x_0`$ and $`P_m(0)=\frac{1}{\omega _m^2M+g_m^{(2)}}`$.
For the initial potential we choose
$$V_{\frac{N}{2}}(x)=\frac{1}{2}M\mathrm{\Omega }^2x^2+\frac{\lambda }{4!}x^4$$
(12)
The numerical solution of the flow equations (10) (11) shows a good convergence for $`N=10^8`$ and $`\beta =10^5`$. There are no significant improvement for larger values of $`N`$ or $`\beta `$. At first, we have tried different truncations of the polynomial potential and found the best values by keeping only the couplings $`g^{(4)}`$ and $`g^{(6)}`$ as is shown in table 1 where we can also read the different values of the couplings at the end of the flow when $`m=0`$. Adding more coupling does not improve the result, on the contrary. In fact in order to get better results we can not truncate the effective potential because in quantum mechanics each coupling constant is relevant. Even when the couplings $`g^{(8)}`$ and $`g^{(10)}`$ take large negative values they do not influence very much the flows of the other couplings. Then, it seems that around its minimum the potential is well fitted by a polynomial interaction of order six.
A detailed comparison of the truncation at $`x^6`$ with Kleinert’s variational method is given in table 2. All the couplings are relevant because they grow as we iterate the RG equation. Surprisingly, the truncation of the potential at the six order gives always better results than those of the variational method. The advantage of the latter lies on the fact that one can improve it systematically with a perturbation expansion which is convergent (see ). In our case the only way to improve it is to consider the potential as a whole. This requires a numerical study which is difficult to achieve due to the very slow convergence of the iterations in the zero temperature case.
### 4.1 The double well potential
In this part we consider the case of the double potential with a frequency $`\mathrm{\Omega }=1`$. The results for some values of the coupling $`\lambda `$ (see table 3) show that $`g_{m=0}^{(2)}>0`$. It means that the quantum fluctuations smear out completely the double well as it must be, because the effective potential is a convex quantity at $`T=0`$. For $`\lambda =2.4`$ our truncation is not accurate enough to get a positive term in the logarithm of equation (8). For all larger couplings our approximation is enough to show the spreading of the double potential by the quantum fluctuations.
## 5 The Feynman Kleinert method
In this section, we give some comments concerning the FK method improved by the renormalization group. In the FK method, one tries to find a quadratic potential at each point $`x_0`$ fitting at best the effective potential. One improves this procedure by looking for a quadratic potential $`M\mathrm{\Omega }_mx^2`$ fitting the potential $`V_m(x_0)`$ at each step of the renormalization group flow. Then by improving in such manner the FK method one will take into account some contributions of the Kleinert’s variational perturbation expansion . To achieve this, we follow Kleinert’s method and insert in the right hand side of (5) the trial frequency $`\mathrm{\Omega }_m^2(x_0)`$ (see ).
$`\mathrm{exp}\left({\displaystyle \frac{T}{\overline{h}}}V_{m1}(x_0)\right)={\displaystyle \frac{dx_md\overline{x_m}}{\frac{2\pi ϵ\overline{h}}{ϵ^2\omega _m^2M}}\mathrm{exp}\left(\frac{ϵ}{\mathrm{}}M\left(\omega _m^2+\mathrm{\Omega }_m^2(x_0)\right)|x_m|^2\right)}`$
$`\times \mathrm{exp}({\displaystyle \frac{ϵ}{\mathrm{}}}({\displaystyle \underset{n=0}{\overset{N+1}{}}}V_m(x_0+{\displaystyle \frac{exp(i\omega _mt_n)x_m}{\sqrt{N+1}}}+H.C.)M\mathrm{\Omega }_m^2(x_0)|x_m|^2))`$
Then, using the Jensen-Peierls inequality
$$𝑑\mu (x)\mathrm{exp}\left(O(x)\right)\mathrm{exp}\left(𝑑\mu (x)O(x)\right)$$
for any positive measure $`\mu `$ normalized to one, and following the same steps as in Kleinert’s book , we obtain the following RG inequality for the potential
$$\begin{array}{c}\text{ }V_{m1}(x_0)\text{ }\hfill \\ \text{ }\frac{1}{\beta }\mathrm{log}\left(1+\frac{\mathrm{\Omega }_m^2(x_0)}{\omega _m^2}\right)\frac{1}{\beta }\frac{\mathrm{\Omega }_m^2}{\mathrm{\Omega }_m^2+\omega _m^2}+\frac{dx}{\sqrt{2\pi a_m^2(x_0)}}\mathrm{exp}\left((xx_0)^2/2a_m^2(x_0)\right)V_m(x).\text{ }\hfill \end{array}$$
where
$$a_m^2(x_0)=\frac{2}{\beta M}\frac{1}{\omega _m^2+\mathrm{\Omega }_m^2}.$$
Let us call $`V_{a_m^2}`$ the integral on the right hand side of the inequality. Minimizing the right hand side in the variable $`\mathrm{\Omega }_m^2`$ we obtain the best approximation for $`V_{m1}`$. The frequency
$$\mathrm{\Omega }_m^2(x_0)=\frac{2}{M}\frac{}{a_m^2}V_{a_m^2}(x_0),$$
(13)
gives the following equation, which can be solved recursively, starting from $`\mathrm{\Omega }_m^2(x_0)=0`$
$$\mathrm{\Omega }_m^2(x_0)=\frac{dx}{\sqrt{2\pi }}\mathrm{exp}\left(x^2/2\right)\frac{1}{2\sqrt{a_m^2}}V_m^{(1)}(x_0+\sqrt{a_m^2(x_0)}x).$$
(14)
Keeping this value for $`\mathrm{\Omega }_m^2`$, enables us to rewrite the approximation
$$V_{m1}(x_0)=\frac{1}{\beta }\mathrm{log}\left(1+\frac{\mathrm{\Omega }_m^2(x_0)}{\omega _m^2}\right)\frac{1}{\beta }\frac{\mathrm{\Omega }_m^2(x_0)}{\mathrm{\Omega }_m^2(x_0)+\omega _m^2}+V_{a_m^2}(x),$$
(15)
which is the variational RG equation. For finite $`\beta `$, it is possible to solve (15), firstly by computing recursively $`\mathrm{\Omega }_m^2(x_0)`$, secondly by calculating $`V_{a_m^2}`$ and inserting this result in the Renormalization group equation. We plan to use this equation in an another paper. In the following section we check its validity in the limiting cases.
### 5.1 Infinite temperature
Recall that for large $`N`$, $`\omega _m^2`$ is equivalent to $`\left(\frac{2\pi m}{\beta }\right)^2`$. Then for $`\beta 0`$, $`\frac{\mathrm{\Omega }_m^2(x_0)}{\omega _m^2}`$ and $`\frac{\mathrm{\Omega }_m^2}{\mathrm{\Omega }_m^2+\omega _m^2}`$ are of order $`\beta ^2`$ and $`a_m^2(x_0)`$ is of order $`\beta `$. Inserting these results in the equation for $`V_{a_m^2}`$ gives $`V_{a_m^2}=V_m(x_0)`$. For $`\beta 0`$ we recover the usual invariance for the flow of the effective potential:
$$V_{m1}(x_0)=V_m(x_0).$$
### 5.2 Zero temperature
In the limit $`\beta \mathrm{}`$, $`a_m^2(x_0)`$ is now of order $`\frac{1}{\beta }`$. To obtain $`\mathrm{\Omega }_m^2`$ we expand the potential $`V_m`$ in a series of $`a_m^2`$. To lowest order we obtain: $`\mathrm{\Omega }_m^2(x_0)=\frac{V_m^{(2)}(x_0)}{M}`$. In the same way, we get the lowest order: $`V_{a_m^2}(x_0)=V_m(x_0)`$.
Then keeping only the relevant terms, the limit for $`\beta 0`$ gives again the Wegner Houghton equation:
$$V_{m1}(x_0)=V_m(x_0)+\frac{1}{\beta }\mathrm{log}\left(1+\frac{V_m^{(2)}(x_0)}{\omega _m^2M}\right)$$
To understand why we recover the WH equation, one compares with the Kleinert’s variational perturbation expansion. In the RG equation for the potential at $`T=0`$, the contributions of all the higher loops terms are included in the flow of the couplings. By choosing a variational frequency which flows in the RG equation we automatically take into account all the higher loops of the variational perturbation expansion and the result must be independent of the variational frequency. In other words, by integrating mode after mode we automatically resume the perturbation expansion and we cannot improve it. A possibility to resume partially the higher loops would be to integrate a lot of modes at each step, and to deduce the variational frequency for the resulting potential.
## 6 Conclusion
The RG equation for the potential of a quantum particle at finite temperature was derived. This equation does not allow us to compute non-perturbative quantities. Then we suggest to use the Feynman Kleinert’s variational method improved by the RG. At zero temperature we recovered the Wegner Houghton equation which was used to compute the ground state energy of the anharmonic oscillator. It would be also interesting to compute the RG equation for the effective action. Our preliminary work show the generation of non local interactions. We also plan to apply the RG procedure for systems with many quantum particles in the discretized path integral representation.
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# On the formation of Wigner molecules in small quantum dots
## I Introduction
If the number of electrons artificially confined on a quasi two-dimensional electron island (made for example in a semiconductor heterostructure) is very large, many properties of such a so-called “quantum dot” or “artificial atom” can be described from what is known about the limit of the infinite (two-dimensional) electron gas (2DEG). Until now, most experiments were performed at electron densities which are slightly below the equilibrium density of the 2DEG. The liquid-like properties then still dominate. For only a few trapped particles (as experimentally realized in vertical quantum dots ), pronounced addition energy maxima as a consequence of shell structure and aligned spins in the mid-shell regions due to Hund’s rules were observed in close analogy to atomic physics . Even the simplest picture of $`N`$ non-interacting particles in a two-dimensional harmonic trap could explain many features of the conductance spectra. For larger systems, mean field approaches like Hartree-Fock or density functional methods were applied. In the small-$`N`$ limit, much theoretical work focused on exact diagonalization techniques . This approach was mostly used for dots in magnetic fields, where correlations become increasingly important with field strengths. It was particularly successful in the (integer and fractional) quantum Hall regime where one can restrict the basis set to the lowest spin-polarized Landau level . Quantum Monte Carlo (QMC) methods provide alternative approaches yielding energies whose accuracy reaches that of exact diagonalizations.
When the electron density is lowered and the Coulomb energy increases relative to the kinetic energy, correlations begin to strongly dominate the electronic structure also in the absence of magnetic fields. For densities smaller than a certain critical value a Wigner crystal will be formed, in which the Coulomb interaction distributes the single electrons classically on a lattice. For the homogeneous two-dimensional electron gas, such crystallization is expected at very low densities. Monte-Carlo calculations indicate that in the 2D bulk a transition to a Wigner crystal-like state, preceeded by a transition to a polarized phase , occurs only at densities corresponding to Wigner-Seitz radii $`r_{s,2D}>37a_B^{}`$ (the density $`n_0`$ and $`r_{s,2D}`$ are related by $`n_0=1/(\pi r_{s,2D}^2)`$), whereas in 3D the classical limit lies as high as $`r_{s,3D}=100a_B^{}`$ . (In the following, for simplicity we write $`r_{s,2D}=r_s`$.) Chui and Tanatar found that in 2D systems without translational invariance the critical density parameter for a fluid-solid transition is shifted to a considerably smaller value ($`r_s7.5a_B^{}`$). Could this be important for finite systems such as the above mentioned lateral or vertical semiconductor quantum dot structures? This question was recently posed and it was argued that in finite systems confining only a few particles localization would indeed occur at significantly higher densities than in the 2D bulk. In the Wigner limit the few electrons in the trap would distribute such that their electrostatic repulsion is minimized. The internal structure of the wave function of the many body system should then have the symmetry of the corresponding classical charge distribution. Wigner crystallization was found to be particularly pronounced in quantum dots with steep walls and polygonal geometry . Egger et al. have performed quantum Monte Carlo studies using a multilevel blocking algorithm . For parabolic quantum dots with azimuthal symmetry they reported that at a critical density of $`r_s=4a_B^{}`$ the formation of Wigner molecule-like ground states should become energetically favorable. Hartree-Fock (HF) methods by definition fail to accurately describe the correlated regime. When performed in an unrestricted scheme, however, spontaneous symmetry breaking and localization in the spatial distribution of the electronic densities of quantum dots and lateral quantum dot molecules at $`r_s3.5a_B^{}`$ was attributed to the onset of Wigner crystallization .
In the present article we report numerically exact configuration interaction calculations. This method has a long history in quantum chemistry, and was applied to quantum dots by many different groups . Much of the previous work, however, concentrated on the electronic structure in large magnetic fields where the electron gas is polarized. Our purpose here is a comparison of the exact diagonalization results to the above mentioned recent predictions of localized states in the low-density limit and zero magnetic field. We first give a brief outline of the configuration interaction method and then turn to a discussion of the many-body spectra of a six-electron quantum dot at zero angular momentum as a function of the average electron density in the dot. Calculations for different $`r_s`$-values indicate that the ground state remains unpolarized. At values of $`r_s`$ which are accessible to exact diagonalization techniques, for a dot confining six electrons clear signals of formation of a Wigner molecule could not be observed. We conclude with a brief comparison to results of density functional theory (DFT) in the local density approximation (LSDA) describing the electronic ground state structures.
## II Method and Convergence
Consider $`N`$ interacting electrons trapped in a circularly symmetric harmonic well $`V(r)=m^{}\omega _0^2r^2/2`$, where $`r^2=x^2+y^2`$. (In the quasi two-dimensional limit one assumes that the confinement in $`z`$-direction is much stronger than in the $`x`$-$`y`$–plane. Then, only the lowest subband in $`z`$-direction is populated.) We write for the Hamiltonian
$$H=\underset{i=1}{\overset{N}{}}\left[\frac{\mathrm{}^2}{2m^{}}_i^2+V(r_i)\right]+\underset{i<j}{\overset{N}{}}\frac{e^2}{4\pi \epsilon _0\epsilon }\frac{1}{𝐫_i𝐫_j}.$$
(1)
Here, $`m^{}`$ and $`\epsilon `$ are the effective mass and the dielectric constant. The calculations are done for different values of the density parameter $`r_s`$ which determines the average particle density in the dot, $`n_0=1/(\pi r_s^2)`$. The latter is approximated by setting the oscillator parameter $`\omega _0^2=e^2/(4\pi \epsilon _0\epsilon mr_s^3\sqrt{N})`$ . Throughout this paper we use effective atomic units in which the length unit $`a_B^{}`$ is a factor $`\epsilon /m^{}`$ times the Bohr radius $`a_B`$, and the energy is given in effective Hartree, Ha$`{}_{}{}^{}=`$Ha$`(m^{}/\epsilon ^2)`$. (For GaAs for example, $`m^{}=0.067m_e`$ and $`\epsilon =12.4`$, for which the length and energy units then scale to $`a_B^{}=97.9\text{ Å}`$ and Ha$`{}_{}{}^{}=11.9`$ meV.) To diagonalize the Hamiltonian, Eq. (1), the spatial single-particle states of the Fock space are chosen to be eigenstates of the two-dimensional harmonic oscillator with optimized oscillator parameter $`\omega `$. In general, the electron-electron interaction tends to expand the system, and thus $`\omega \omega _0`$. This effect becomes stronger with increasing $`r_s`$ and we used the (empirical) relation $`\omega =\omega _0/\sqrt{r_s}`$ . The two-body matrix elements of the electron-electron interaction are calculated using the addition theorem for $`1/𝐫_1𝐫_2`$ . To set up the Fock states for diagonalization, we use eight lowest oscillator shells containing 36 states and sample over the full space with a fixed number of spin down and spin up electrons, $`N_{}+N^{}=N`$. From this sampling, only those Fock states with a given total orbital angular momentum and a configuration energy (corresponding to the sum of occupied single-particle energies) less than or equal to a specified cutoff energy $`E_\text{c}`$ are included (see Figure 1 below for an example). The purpose was to select only the most important Fock states from the full basis, hereby reducing the matrix dimension to a size $`d\stackrel{<}{}10^5`$. To obtain all the eigenstates we have to set $`N_{}=N_{}=N/2`$ for even particle numbers ($`S_z=0`$, all states with different total spin have this component), and analogously we would have $`N_{}=N_{}\pm 1`$ for odd numbers. Once the active Fock states are at disposal, the Hamiltonian matrix is calculated. For Lanczos diagonalization we use the Arpack library . Finally, the total spin of each eigenvector is determined by calculating the expectation value of the $`\widehat{S}^2`$ operator.
As mentioned above, setting up the Hamiltonian matrix from chosen Fock states and subsequent diagonalization only in principle yields an exact solution of the many-body problem. For reasons of numerical feasibility it is necessary to truncate the set of basis functions to be used in the diagonalization. One then has to make sure that convergence of the spectra is reached with respect to the cut-off. As the required matrix size increases rapidly with $`N`$, computational expenses severely restrict the calculations to only the smallest systems at not too large values of $`r_s`$. Thus, with increasing electron number or $`r_s`$, the results become less accurate due to the restricted number of basis states that can be included in the calculations. (The fact that the ground and excited states are the closer in energy the larger the particle number $`N`$ imposes an additional difficulty.)
For a quantum dot confining $`N=6`$ electrons at a density corresponding to $`r_s=4a_B^{}`$ (the largest value of $`r_s`$ we found accessible within the calculational scheme used here), Figure 1 shows the convergence of the many-body spectra as a function of the cut-off energy $`E_c`$. The lowest possible Fock state for six electrons has two particles in the state $`|n,l=|0,0`$ and four particles in $`|0,\pm 1`$. Thus, the configuration energy equals $`2\mathrm{}\omega +42\mathrm{}\omega `$. This means that for the spectra with different cut-off energies displayed in Figure 1 all excitations up to an energy $`E_c10\mathrm{}\omega `$ are included. The many-body spectra for $`14\mathrm{}\omega `$ and $`16\mathrm{}\omega `$ differ drastically from the results obtained for $`E_c20\mathrm{}\omega `$. Looking at the relative ordering of the levels and the spin sequence, it becomes clear that convergence is only reached for $`E_c>20\mathrm{}\omega `$. The ground-state energy for zero angular momentum is $`E_0=3.049`$ Ha for $`22\mathrm{}\omega `$ and $`E_0=3.045`$ Ha for $`24\mathrm{}\omega `$. An extrapolation to infinite cut-off energy can be made by plotting the total energy as a function of $`(E_c10\mathrm{}\omega )^{3/2}`$ . This gives the estimate 3.043 Ha for the fully converged results at $`r_s=4a_B^{}`$. For $`E_c=14\mathrm{}\omega `$ a too small number of Slater determinants was included to build up the required correlations, such that the polarized $`S=3`$ state appeared as the ground state. (We identify a similar effect in unrestricted HF results mentioned above , where the single Slater determinant that is available incorrectly favors a spin-polarized ground state .) While for $`E_c=22\mathrm{}\omega `$ the matrix dimension $`44181`$ with 21448811 non-zero matrix elements is reasonably small, the value $`E_c=24\mathrm{}\omega `$ already yields a matrix dimension 108375 with 67521121 non-zero matrix elements. As for larger densities the states are less correlated, a smaller number of Slater determinants is needed for an accurate description. Density parameters larger than $`r_s=4a_B^{}`$ or a higher number of particles than $`N=6`$ would go beyond the limits of numerical feasibility and accurate results could not be obtained.
## III Many-body spectra of a six-electron quantum dot
We now analyze the many-body spectra and the sequence of spins for the low-lying states as the two-dimensional density parameter $`r_s`$ is varied. We choose the particle number $`N=6`$ as it corresponds to the smallest dot size for which classically two stable crystalline structures co-exist: A pentagonal ring with one electron at the center, and a slightly distorted six-fold ring . We fix the angular momentum to $`L=0`$ and show in Figure 2 the 50 lowest states for a quantum dot confining $`N=6`$ particles at density parameters between $`r_s=1a_B^{}`$ and $`r_s=4a_B^{}`$ (in steps of $`0.5a_B^{}`$). To obtain a better resolution of the spectra the energies of the eigenstates $`ϵ_i`$ are scaled such that the energy difference between the ground state and 50<sup>th</sup> excited state equals one, i.e. plotted are the dimensionless quantities $`\stackrel{~}{ϵ}_i=(ϵ_iϵ_1)/(ϵ_{50}ϵ_1)`$ . (The total energies of the ground state with $`S=0`$ and the excited state with $`S=3`$ are given in Table 1 below.) At a very large density corresponding to $`r_s=1a_B^{}`$, the ground state has spin $`S=0`$ and is separated from the lowest excited state with spin $`S=2`$ by a gap of 0.49 Ha. This state is followed by a state with $`S=1`$ and again another (excited) spin singlet. The lowest fully polarized state is only found at a fairly high energy, the energy difference to the $`S=0`$ ground state being about 0.95 Ha. (We note that for $`N=6`$ at $`r_s=1.73a_B^{}`$ we obtained excellent agreement of the $`S=0`$ ground state energy with the result of Pederiva et al. .) As $`r_s`$ is increased, the fully polarized $`S=3`$ state moves down in energy. At $`r_s=2.5a_B^{}`$ it has passed the excited singlet, but is still far from competing with the nonpolarized $`S=0`$ ground state. From the evolution of the energy difference of the $`S=3`$ state to the ground state, we do not expect any crossing of the polarized state and the ground state unless $`r_s`$ becomes much larger than $`4a_B^{}`$ . Estimating the decrease in energy of the polarized state with respect to the ground state as $`r_s`$ is increased, our data seem to support the result of Egger et al. that the ground state of the 6 electron dot is not polarized for $`r_s`$-values smaller than about $`8a_B^{}`$.
Table 1 compares the total energies of the ground state with spin zero and the lowest polarized state with spin $`S=3`$ with the corresponding result obtained from DFT,
| | paramagnetic | | ferromagnetic | |
| --- | --- | --- | --- | --- |
| $`r_s[a_B^{}]`$ | exact | LSDA | exact | LSDA |
| 1.0 | 14.27 | 14.30 | 15.22 | 15.30 |
| 1.5 | 8.983 | 8.988 | 9.363 | 9.409 |
| 2.0 | 6.508 | 6.503 | 6.695 | 6.724 |
| 2.5 | 5.084 | 5.073 | 5.188 | 5.204 |
| 3.0 | 4.162 | 4.148 | 4.225 | 4.233 |
| 3.5 | 3.519 | 3.502 | 3.559 | 3.560 |
| 4.0 | 3.045 | 3.027 | 3.071 | 3.068 |
TABLE 1: Table of energies \[Ha\] for the paramagnetic $`(S=0)`$ and ferromagnetic $`(S=3)`$ states in a 6-electron quantum dot for different densities. In the diagonalization we used $`E_c=22\mathrm{}\omega `$ for $`r_s3.5a_B^{}`$ and $`E_c=24\mathrm{}\omega `$ for $`r_s=4a_B^{}`$. For comparison the energies obtained with the local spin density approximation (LSDA) are also shown.
where the exchange-correlation part of the electron-electron interactions is treated in LSDA. For the DFT results we used an interpolation formula for the Tanatar-Ceperley exchange-correlation energy. We refer to for the further details concerning the numerical method. The CI energies of the $`S=0`$ ground states and $`S=3`$ isomer compare well with the LSDA results. For the paramagnetic case, the LSDA gives lower energies than the exact results, when $`r_s2a_B^{}`$. This might be mainly due to the fact that the Tanatar-Ceperley interpolation formula slightly overestimates the correlation energy. Nevertheless, the LSDA gives surprisingly accurately the energy difference between the fully polarized ($`S=3`$) and the paramagnetic ($`S=0`$) state, as seen in Fig. 3. (For comparison, we also show the results for the infinite electron gas.)
Figure 2 shows that for $`r_s=4a_B^{}`$ there are only two $`L=0`$ states between the ferromagnetic $`(S=3)`$ state and the paramagnetic $`(S=0)`$ ground state. However, if we consider all $`L`$-values, there are several states within this energy range. This is shown in Table 2 where energies, spins and angular momenta of all levels up to the first ferromagnetic state are given. Indeed, the lowest excited state has $`L=1`$ and $`S=1`$. It is possible that at very large $`r_s`$ this partially polarized state might become the ground state instead of the fully polarized $`S=3`$ state .
| $`r_s=3a_B^{}`$ | | | $`r_s=4a_B^{}`$ | | |
| --- | --- | --- | --- | --- | --- |
| E \[Ha\] | $`S`$ | $`L`$ | E \[Ha\] | $`S`$ | $`L`$ |
| 4.162 | 0 | 0 | 3.046 | 0 | 0 |
| 4.183 | 1 | 1 | 3.054 | 1 | 1 |
| 4.194 | 1 | 3 | 3.060 | 2 | 0 |
| 4.196 | 2 | 0 | 3.062 | 1 | 3 |
| 4.201 | 1 | 1 | 3.063 | 1 | 1 |
| 4.205 | 0 | 1 | 3.065 | 0 | 1 |
| 4.209 | 1 | 0 | 3.066 | 1 | 0 |
| 4.209 | 0 | 2 | 3.068 | 2 | 1 |
| 4.209 | 0 | 3 | 3.068 | 0 | 2 |
| 4.213 | 1 | 2 | 3.070 | 1 | 2 |
| 4.216 | 2 | 1 | 3.070 | 0 | 3 |
| 4.216 | 2 | 2 | 3.071 | 3 | 0 |
| 4.225 | 3 | 0 | | | |
TABLE 2: Table of energies \[Ha\], spins $`S`$ and angular momenta $`L`$ of all levels up to the lowest ferromagnetic state for $`r_s=3a_B^{}`$ (left) and $`r_s=4a_B^{}`$ (right).
## IV Charge densities and pair correlation
For $`r_s=4a_B^{}`$ the radial densities of the $`S=0`$ ground state and the 3rd excited state of $`L=0`$, which is the lowest fully polarized state with spin $`S=3`$, are shown in Figure 4. The azimuthally symmetric charge density for the polarized case shows a clear maximum at the center surrounded by an outer ring of lower density. In the paramagnetic case, the density profile is more smooth, the maximum density being at about $`r6a_B^{}`$. The LSDA result shows a clear minimum at the origin, while the exact result has a larger density at the center. For comparison with the results of Egger et al. we also show the density $`\rho (r)`$ multiplied by a factor $`2\pi r`$ (cf. lower panel of Fig. 4). For the polarized state, the maximum at the center is now seen as a clear shoulder in the density profile. Note that this is missing in the paramagnetic ground state density. (This is in disagreement with the results of Egger et al. who found for the ground state a density profile with a clear shoulder as in the polarized case.) The azimuthal averages of the density profiles qualitatively have similarities with the broken symmetry solutions of the unrestricted HF which for the paramagnetic case results in a ring of six electrons, and for the ferromagnetic case (ground state in HF) a ring of five electrons with one electron in the center. However, the localization of the electrons is largely exaggerated in HF. Opposite to HF, the LSDA correctly gives the paramagnetic state as the ground state, and its density profile resembles the exact result. LSDA does not break the azimuthal symmetry until $`r_s>8a_B^{}`$ when spin- or charge density wave-like states can occur . Purely classical Monte Carlo computations have shown that for $`N<6`$ the charges are distributed on the perimeter of the dot, and none of the particles occupies the dot center. This changes for $`N=6`$, where the charge distribution with lowest energy consists of five electrons sitting on a ring, with the remaining electron occupying the center of the dot. This configuration is labeled by $`(5,1)`$. If all 6 particles are arranged on the dot perimeter (labeled by $`(6,0)`$), the classical state is stable but has a higher energy than the $`(5,1)`$-configuration.
The classical charge distribution can be arbitrarily oriented. The density from the CI solution, however, must be circularly symmetric. For an azimuthal average of the $`(5,1)`$-pentagon structure, one would expect a pronounced maximum of the electron density in the center, and a less pronounced maximum at the dot radius. Correspondingly, the $`(6,0)`$ configuration should correspond to a minimum of charge density in the center and a maximum at finite radius. A first comparison of exact diagonalization calculations with the results of the mean-field approximation was given by Pfannkuche et al. for ”quantum dot helium”, i.e., quantum dots containing only two electrons. They found from a comparison of exact diagonalizations with Hartree- and Hartree-Fock results that the exchange and correlation contributions are crucial. While the triplet state showed a reasonable agreement between the exact result and HF, the singlet could not be well reproduced. As mentioned above, Yannouleas and Landman reported that in geometrically unrestricted HF at a density corresponding to $`r_s3.5a_B^{}`$ the $`N=6`$ ground state is polarized and shows enhanced localization in the charge density. This $`S=3`$ state exhibits the same geometry than the classical distribution of 6 electrons in a harmonic well: 5 particles are equidistanly localized on the perimeter of the dot, and the 6th particle is trapped in the center of the harmonic well. The non-polarized $`S=0`$ state corresponding to the $`(6,0)`$-configuration is about 0.034 Ha higher in energy. The exact diagonalization results described above do not support these HF results. Although being limited to small $`r_s`$ values due to the necessary restrictions of the basis set, the systematic evolution and energy sequence of CI energies and densities shown in Figs. 2 and 4 seems to indicate that polarization as well as formation of Wigner molecules in circularly symmetric, parabolic wells would be impossible at densities as large as predicted by HF . The geometrically unrestricted solution of the Kohn-Sham equations of tend to overestimate the $`r_s`$-value at which spontaneously broken spin- or charge-symmetries can occur in the internal structure of the wave function . Although calculated in a geometrically unrestricted DFT scheme, the fully converged LSDA densities for $`r_s4a_B^{}`$ shown in Fig. 4 are azimuthally symmetric. Although LSDA suffers from the self-interaction problem, at the densities in question the results are in better agreement with CI studies than the unrestricted HF results.
It is finally of interest to study how the spin and spatial symmetry in the internal structure of the wave function can be recognized in the pair correlation function
$$g_\sigma (\phi )=\widehat{n}_{}(r,0)\widehat{n}_\sigma (r,\phi ),$$
(2)
which describes the probablity to find another (spin up or down) particle if a particle with spin up is placed at $`(r,0)`$. Here $`r`$ is the radius of maximum density and $`\phi `$ the angle between the electrons. Figure 5 shows $`g_{}(\phi )`$ and $`g_{}(\phi )`$ for the ground state (lower panel) and $`g_{}(\phi )`$ for the excited polarized state with $`S=3`$ (upper panel). From the $`\phi `$-values of the maxima in $`g_\sigma (\phi )`$ one clearly concludes that the $`S=0`$ state has 6-fold symmetry with antiferromagetic spin ordering, whereas the fully polarized case shows four maxima, corresponding to a five-fold symmetry. These intrinsic symmetries are in qualitative agreement with the unrestricted HF results, although the crystallization predicted by HF does not yet occur.
## V Conclusions
We commented on the recent conjecture that Wigner molecules would form in quantum dots at rather large electron densities $`r_s\stackrel{>}{}3.5a_B^{}`$ . Our results are essentially exact up to $`r_s4a_B^{}`$ for six confined particles. The many-body spectra, densities and pair correlations obtained for $`N=6`$ clearly illustrate that the onset of formation of Wigner molecules and in particular, polarization of the ground state should be expected at much higher $`r_s`$-values than anticipated from unrestricted HF . The critical density at which such a transition occurs does, in fact, strongly depend also on geometry and on the number of confined particles at fixed average electron density in the dot. For the six-electron dot at $`r_s4a_B^{}`$ in question , the “exact” ground-state clearly prefers $`S=0`$ and shows antiferromagnetic order in the pair correlation. The polarized state with $`C_{5v}`$-symmetry is clearly higher in energy even than the $`S=1`$ state, which would also be a candidate for the classical $`(5,1)`$ ground state configuration . In addition, for values below $`r_s=4a_B^{}`$ we did not find clear signals of rotational structure in the spectra for non-zero angular momenta that would indicate a crystallized ground state. We note that the situation is different for $`N<5`$, where for densities as large as $`r_s=2a_B^{}`$ the low-lying states could be well understood by assuming a square-shaped $`(4,0)`$ Wigner molecule for the internal structure of the wave function and analyzing its rotational structure. This also became clear when comparing the low-energy spectrum of a Heisenberg model with four electrons on a square . For $`N>5`$, however, this simple picture does seem to not hold, as our results for $`N=6`$ clearly point out. The ground-state energies and densities obtained by density functional calculations in the local spin density approximation agree rather well with the results of exact diagonalizations, even though this comparison was restricted to a small particle number where the accuracy of the local density approximation is questionable. This gives some confidence that the method is well suited for describing the ground state electronic structures for larger sizes.
This work was supported by the Academy of Finland, the “Bayerische Staatsministerium für Wissenschaft, Forschung und Kunst” and the TMR programme of the European Community under contract ERBFMBICT972405.
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# Insulator-metal transition in high-Tc superconductors as result of percolation over -U centers
## I Introduction
In 13 years elapsed after the discovery of HTS, numerous models have been suggested (see review in Ref. ) to explain the nature of the ground state and anomalous properties of HTS compounds. However, the lack of any crucial experiment gave no way of choosing between these models.
In this paper, we intend to demonstrate that the mechanism responsible for various anomalous properties of HTS compounds (including the high-temperature superconductivity itself) is apparently based on interaction of electrons with so-called -U centers. To this end, we consider the way the insulator-metal transition occurs in HTS under doping. On the basis of a simple ionic model, we will show that such a transition should pass through a certain range of dopant concentrations; this range corresponds to the situation when local transitions of singlet electron pairs from oxygen ions to a pair of neighboring cations (-U center) become possible in individual microclusters including several unit cells, while the single-electron transitions are still forbidden. In our opinion, it is this range of concentrations that corresponds to the HTS range where interelectron attraction results from interaction of electron pairs with -U centers. Conduction in such a system arises when the concentration of -U centers exceeds the percolation threshold for orbitals of singlet hole pairs. We are going to consider which specific fragments of crystal structure are involved in the formation of -U centers and what is the range of existence of infinite percolation cluster interconnecting singlet hole pair orbitals. On this basis we will construct the phase diagram for Ln-214 (Ln=La, Nd) compounds. In our opinion, a comparison of this phase diagram with the carefully investigated diagram for Ln-214 compound should be taken as the aforementioned crucial experiment for choosing the mechanism responsible for the HTS properties. We discuss the special features of phase diagrams for Ln-214 with $`n`$\- and $`p`$-type of doping, and also we give a somewhat different interpretation of some experimental results obtained for underdoped HTS. Furthermore, we consider the nature of mobile hole carriers and mechanism of their relaxation in HTS. As will be seen the distinctive feature of HTS normal state is the nondegenerate distribution of hole carriers that results in hole-hole scattering to dominate in kinetic processes. It will be shown that this feature may account for the unusual transport and optical properties of HTS observed in experiment.
## II -U center formation in HTS
There are good reasons to believe that electronic spectrum of insulator phase for various HTS compounds in the vicinity of Fermi energy E<sub>F</sub> can be best approximated by the model of charge-transfer insulator. In this model, the upper empty band formed by unfilled orbitals of cations is separated with a gap from the O2p band formed largely by oxygen states (Fig. 1(a)). The gap $`\mathrm{\Delta }_{ct}`$ existing in the spectrum is related to the transfer of electron from oxygen to the neighboring cation and lies in the range of $`1.52`$ eV for all HTS.
The question arises: What is the mechanism of insulator-metal transition in the doped HTS? As an example of HTS, we consider Ln-214. For these compounds, the quantity $`\mathrm{\Delta }_{ct}`$, in terms of simple ionic model, is defined by the following relationship:
$$\mathrm{\Delta }_{ct}\left|\mathrm{\Delta }E_M\right|+A_pI_d.$$
(1)
Here, $`I_d`$ is the second ionization potential of Cu, $`A_p`$ is the electronegativity of oxygen with respect to formation of O<sup>2-</sup>, and $`\left|\mathrm{\Delta }E_M\right|`$ is the difference of the Madelung energy $`E_M`$ between the configuration in which the copper and oxygen atoms are in the state of Cu<sup>2+</sup> and O<sup>2-</sup> and that with these atoms in the states of Cu<sup>1+</sup> and O<sup>1-</sup>. Taking into account that $`I_d20`$ eV, $`A_p0`$ eV and $`\mathrm{\Delta }_{ct}1.52`$ eV, there is a subtle balance between these three quantities.
This balance can be varied by heterovalent doping; for example by doping the La<sub>2</sub>CuO<sub>4</sub> with divalent Sr or by doping Nd<sub>2</sub>CuO<sub>4</sub> with tetravalent Ce. It is a matter of great importance that the charge carriers introduced by doping (so-called doped carriers) are localized (at least for low concentrations) in CuO<sub>2</sub> plane in the vicinity of the dopant ion: either at the O2p<sub>x,y</sub> orbitals (the holes in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>) or at the Cu3d$`_{x^2y^2}`$ orbitals (electrons in Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub>). Taking into account that the interaction of O<sup>2-</sup> and Cu<sup>2+</sup> gives the main contribution to $`E_M`$, an addition both of electrons (to the Cu orbitals) and of holes (to the oxygen orbitals) will results in the same thing: a decrease in $`\left|\mathrm{\Delta }E_M\right|`$ and, correspondingly, a decrease in $`\mathrm{\Delta }_{ct}`$ for other pairs of copper and oxygen ions sited in the vicinity of the doped carrier. For a certain critical concentration $`x_c`$, the gap $`\mathrm{\Delta }_{ct}`$ vanishes throughout the entire crystal. As a result, electron transitions from oxygen to copper become possible, and the material passes into a normal metal.
In generally by this means, we can conceive the transition of a charge-transfer insulator into the metallic state with doping in terms of ionic model. However, we argue that, in HTS compounds, the transition from insulator to metal with x passes through a special range of concentrations $`x_0<x<x_c`$ for which the two-electron transitions from oxygen ions to certain pairs of neighboring cations become possible, while the single-electron transitions are still forbidden. In other words, -U centers are formed on certain pairs of cations at $`x_0<x<x_c`$.
Let us consider a Cu<sub>2</sub>M<sub>2</sub>O<sub>n</sub> cluster where two neighboring Cu ions belong to the CuO<sub>2</sub> plane, and M=Cu in the selfsame CuO<sub>2</sub> plane for Ln-214, M=Cu in the CuO<sub>3</sub> chain for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>, and M=Bi for Bi<sub>2</sub>Sr<sub>2</sub>Ca<sub>2</sub>Cu<sub>2</sub>O<sub>8+δ</sub>. As will be seen from the following consideration, the condition for the formation of -U centers at the neighboring Cu ions in CuO<sub>2</sub> plane consists in the presence of one doped carrier in the vicinity of each M ion (in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> and Bi<sub>2</sub>Sr<sub>2</sub>Ca<sub>2</sub>Cu<sub>2</sub>O<sub>8+δ</sub>, these carriers are in the CuO<sub>3</sub> chains and BiO planes, respectively). In Ln-214, two types of such clusters are possible (Figures 2(a, b)): the projections of the dopant ions onto the CuO<sub>2</sub> plane are separated by either 3$`a`$ or $`a\sqrt{5}`$, where $`a`$ is the lattice constant in the CuO<sub>2</sub> plane. In both cases, the presence of doped carrier in the vicinity of each M ion reduces $`\mathrm{\Delta }_{ct}`$ for the neighboring Cu ions and, as it will be shown, provides the conditions (i.e., forms a local minimum of potential energy) for simultaneous transition of two electrons to the pair of interior Cu ions from O ions surrounding this pair. It should be noted here that, in the intermediate case where the M ion projections are spaced $`a\sqrt{8}`$ apart, such pair of Cu ions does not emerge (Fig. 2(c)).
It is possible to estimate a decrease in $`\mathrm{\Delta }_{ct}`$ for a given Cu ion in La-214 due to the presence of a single hole around the neighboring Cu ion if we assume that this hole is ”distributed” (Fig. 3) over 12 nearest-neighbor oxygen ions (the first and second coordination spheres). This assumption is consistent with experimentally determined limit of substitutability (see below). We take into account only the interaction between the nearest neighbors. Therefore we consider 3 oxygen ions from the total number of 12, where the hole ”feels” the unscreened Cu ion separated by a distance of $`r=a/22`$ Å from this hole. In this case for a given Cu ion, a lowering of the Cu3d<sup>10</sup> state energy amounts to $`\mathrm{\Delta }E(1/4)e^2/r1.8`$ eV; here, $`e`$ is the electron charge. This means that, due to doping, the Cu3d<sup>10</sup> state energy for a given Cu ion is decreased, so that it is now by $`\mathrm{\Delta }E1.8`$ eV below the bottom of the conduction band for undoped material. This value is still smaller by $`0.10.2`$ eV than $`\mathrm{\Delta }_{ct}1.92.0`$ eV for La<sub>2</sub>CuO<sub>4</sub>. In Nd-214 doped electrons increase O2p state energy around M ion that results in a decrease in $`\mathrm{\Delta }_{ct}`$, too.
An additional lowering of the Cu3d<sup>10</sup> state energy is attained owing to formation of a bound state of two electrons at neighboring Cu ions in the presence of two holes in the immediate vicinity of this pair. Such a decrease is possible for a bonding orbital of a singlet hole pair, as is the case for H<sub>2</sub> molecule. Here, this analogy is even more appropriate because the distance between electrons at Cu ions ($`3.8`$ Å) is close to $`R_0ϵ_{\mathrm{}}3.6`$ Å, where $`ϵ_{\mathrm{}}4.5`$ is the high-frequency dielectric constant and $`R_00.8`$ Å is the distance between nuclei in H<sub>2</sub> molecule. Therefore, an additional lowering of the energy $`\delta E_U`$ due to transition of two electrons to neighboring copper ions can be, in the case under consideration, estimated from the relationship $`\delta E_U\delta E_H/ϵ_{\mathrm{}}^2`$ 0.23 eV, where $`\delta E_H=4.75`$ eV is the binding energy in H<sub>2</sub> molecule. However, this value $`\delta E_U`$ is apparently underestimated because oxygen ions in CuO<sub>2</sub> plane efficiently screen the repulsive interaction of electrons at Cu ions and weakly screen the electron-hole attraction.
Thus, we may assume that $`\mathrm{\Delta }_{ct}`$, which amounts to $`1.52.0`$ eV for doped cuprates, vanishes for two-electron transitions to neighboring Cu ions. In this case, holes apparently occupy predominantly $`\pi p_{x,y}`$ orbitals, thus providing naturally the bonding character of the hole-pair orbital owing to configuration of bonds in the CuO<sub>2</sub> plane, which allows the holes to reside in the space between Cu ions (Fig. 4). Here we assume that the states at the top of the oxygen valence band are formed predominantly by $`\pi p_{x,y}`$ orbitals.
It follows from the above consideration that, if an appropriate local minimum in $`\mathrm{\Delta }_{ct}`$ is formed, the bound electron state can arise at the pair of neighboring Cu ions while the single-electron transitions are forbidden (i.e., -U center is formed). In this case, a singlet hole pair is localized in the vicinity of -U center at a distance of $`a/2`$. The region of the hole-pair localization is restricted by the condition that the pair level lies in line with the top of the valence band (the energy of the pair level becomes higher with increasing of the hole-pair localization area). Thus, the energy spectrum of CT-insulator under doping is modified by the addition of electron pair level lying in line with the top of the valence band (Fig. 1(b)). In this case the states at the top of valence band should be consider as formed by the $`\pi p_{x,y}`$ orbitals of oxygen ions surrounding -U centers.
Conduction in such a system occurs if the areas of the localization of singlet hole pairs form a percolation cluster. The localization areas of doped holes may enter into this percolation cluster, too (e.g. in La-214). The percolation threshold for the hole-pair orbitals of the -U centers in Ln-214 is coincident with that for a ensemble of segments of length $`L=3a`$ or $`a\sqrt{5}`$ in square lattice. The region of delocalization of hole carriers in the percolation cluster is also restricted by the condition that the position of the pair level should coincide with the valence-band top. Such a mechanism keeps the pair level exactly at the top of the valence band (Fig. 1(b)).
On the other hand, if projections of two M ions are separated by a distance 2$`a`$ (Fig. 2(d)), $`\mathrm{\Delta }_{ct}`$ vanishes for single-electron transitions to interior Cu ion as well. Such a fragment is a nucleus of ordinary metallic phase. For the relevant concentrations $`xx_c`$, the entire crystal transforms into ordinary metal. This state corresponds to a single-band electron spectrum. In La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> in the ordinary metal phase, the charge carriers are electrons because, due to doping with divalent Sr, the filling of the band $`\rho <1/2`$; under the same conditions, the charge carriers are holes in Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub> because, due to doping with tetravalent Ce, we have $`\rho >1/2`$. The intermediate range of $`x`$ where -U centers and ordinary metal phase are coexisted in the percolation cluster is so-called ”overdoped” region. In this case, the additional carriers from metal phase decrease $`\left|\mathrm{\Delta }E_M\right|`$ and lower the pair level below the top of the valence band.
## III Construction of phase diagram of Ln-214 compounds
We will construct the phase diagram for Ln-214 system using the following assumptions:
(i) the -U centers are formed at the pairs of neighboring Cu ions belonging only to the clusters with $`L=3a`$ and $`L=a\sqrt{5}`$ ;
(ii) the orbitals of hole pairs are in the immediate vicinity of these pairs at a distance of $`a/2`$;
(iii) conduction in the system arises in the case of percolation over segments containing -U centers; and
(iv) localized doping carriers cannot be separated by a distance smaller than 2a.
Assumption (iv) follows from the existence of the substitutability limits in Ln-214 $`x_{lim}0.20.25`$. If these limits are exceeded, a decomposition of single-phase state and/or a change of the oxygen content occur. In our opinion, the existence of the substitutability limit is related to repulsive interaction between localized doped carriers. In turn, such a repulsion would affect the dopant ion distribution, if the mobility of these ions at heat treatment temperature was rather high. Therefore, we believe that the dopant ions (more precisely, their projections onto the CuO<sub>2</sub> plane) cannot be separated by a distance $`L<2a`$ as well.
Taking into account the above assumption, we may define the $`2D`$-percolation threshold as follows. Let us assume that we have a square lattice with the cell parameter $`a`$=1 and let a fraction $`x`$ of the sites be occupied by atoms. In order to determine the percolation threshold for segments with length $`L`$ (i.e., for the pairs of atoms separated by a distance $`L`$) we locate each occupied site at the center of circle with the radius $`L/2`$ (Fig. 5(a)). The sum of the areas of the circles constructed around these atoms is given by $`S=x\pi L^2/4`$. For a square lattice, the percolation sets in when $`S`$ 0.466. Consequently, the concentration corresponding to the percolation threshold is given by $`x_p(L)=0.593/L^2`$. Here, we assume that the distribution of atoms over the sites is random, and L is the smallest distance between atoms for the given concentration. Otherwise, we would have an overlapping of circles (Fig. 5(b)), and the value of $`x_p`$ would be larger than that obtained from the above relationship. In this case a percolation cluster would include not only the segments with a length of L but also those with smaller length (linking the sites separated by smaller distances). The maximal number of segments with length $`L`$ \- $`x_M(L)`$ \- can be attained in the case of ordered arrangement of atoms in a square lattice with the parameter $`L`$; i.e., $`x_M(L)=1/L^2.`$ The values of $`x_p`$ and $`x_M`$ for various values of $`L`$ are listed in Table 1. Therein, in the righthand column, it is indicated which state (insulator, metal, or HTS) would correspond to the case of percolation over the segments with a given $`L`$.
Figure 6(a) shows the percolation ranges for segments with different length $`L`$ (therein, the corresponding value of $`L^2`$ is indicated to left of each rectangle). The left side of each rectangle corresponds to the $`2D`$-percolation threshold for the segments with the length $`L`$ under the condition that the atoms are randomly distributed and there are no segments with the length smaller than $`L`$. The right side of each rectangle corresponds to the point $`x_M`$. Heavy lines indicate the percolation ranges for the segments with $`L=3`$ and $`L=\sqrt{5}`$. It is in these concentration regions we suggest, the high-temperature superconductivity occurs.
As will be evident from what follows, the ratio between order and disorder in the distribution of dopants over the sites is very important and defines all special features of phase diagrams for HTS of Ln-214 type. We believe that a tendency towards ordering is related to the difference in ionic radii of Ln and the dopant. This tendency should be most pronounced for the La/Ba pair and least pronounced for the Nd/Ce pair. The degree of ordering should also increase with increasing $`x`$ (with decreasing $`L`$). In what follows, while on subject of the dopant ordering, we imply a formation of a large number of ordered clusters (with sizes smaller than 100 Å) separated by thin walls with random distribution of the dopant. In these walls, the constraint $`L2`$ is also preserved. Thus, only the short-range order exists in the system. It is assumed that, if the value of $`x_M`$ is exceeded, a transition to random distribution of atoms over the sites occurs with formation of an infinite percolation cluster with a smaller value of $`L`$.
We now consider various ranges of concentrations in Fig. 6(a) assuming that, at each point $`x>0.1`$, no more than two types of segments exist.
1. $`0.20<x<0.25`$. In this case, the $`2D`$-percolation is effected over the segments with $`L=2`$, i.e., over the clusters of ordinary metal.
2. $`0.148<x<0.20`$. In this range, the $`2D`$-percolation threshold for the segments with $`L=2`$ depends on the degree of ordering of the dopant atoms with $`L=\sqrt{5}`$ . In the case of ordering of the atoms with $`L=\sqrt{5}`$, percolation over the segments with $`L=2`$ sets in at $`x=0.2`$, whereas, for random distribution, percolation over the segments with $`L=2`$ is attained at $`x=0.148`$. Therefore, this range of $`x`$ corresponds either to HTS, or to a mixed state of HTS and ordinary metal.
3. $`0.125<x<0.148`$. Here, we have a ”pure” $`2D`$-percolation over the segments with $`L=\sqrt{5}`$ . This range corresponds to HTS.
4. $`0.118<x<0.125`$. If the atoms with $`L=\sqrt{8}`$ are ordered, percolation over the segments with $`L=\sqrt{5}`$ sets in at $`x=0.125`$, whereas, if the atoms with $`L=\sqrt{8}`$ and $`L=\sqrt{5}`$ are randomly distributed, percolation over the segments with $`L=\sqrt{5}`$ sets in at $`x=0.118`$. In this range, we have either insulator (in the former case) or HTS (in the latter case).
5. $`0.111<x<0.118`$. This is the domain of ”pure” $`2D`$-percolation over the segments with $`L=\sqrt{8}`$. This phase corresponds to an insulator.
6. $`0.10<x<0.111`$. In the case of ordering of atoms with $`L=3`$, the $`2D`$-percolation over the segments with $`L=3`$ sets in at $`x=0.10`$, whereas, if the atom pairs with $`L=3`$ and $`L=\sqrt{8}`$ are randomly distributed, there is no $`2D`$\- percolation neither over the two types of segments. In this case, a percolation cluster would include the areas with L=3 (HTS) and with $`L=\sqrt{8}`$ (insulator) and the conduction is possible only through the tunneling between clusters with $`L=3`$ (3D-percolation).
7. $`0.066<x<0.10`$. In this range, there is no $`2D`$-percolation over the segments with $`L=3`$. $`3D`$-percolation (and superconductivity) is still possible for $`x>0.077`$, whereas, for $`x<0.077`$, there is no bulk superconductivity.
For comparison, Figures 6(b - d) show the experimental phase diagrams $`T_c(x)`$ for La<sub>2-x</sub>Ba<sub>x</sub>CuO<sub>4</sub>, La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, and Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub>. Comparing Figures 6(a - d), we can readily see that all particular points in experimental phase diagrams practically coincide with the boundaries of percolation domains corresponding to the segments with different values of $`L`$. Difference between the phase diagrams of La<sub>2-x</sub>Ba<sub>x</sub>CuO<sub>4</sub> and La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> consists in the fact that, in the former case, a dip in $`T_c`$ occurs at $`x_d=0.125`$, whereas, in the latter case, it takes place at $`x_d=0.115`$. It is reasonable to relate the difference in the values of $`x_d`$ for these compounds to a larger degree of ordering in the La/Ba sublattice as compared to that for La/Sr, as a result of which the percolation threshold for the segments with $`L=\sqrt{5}`$ becomes shifted to $`x=0.125`$ (the point of the highest ordering for $`L=\sqrt{8}`$).
The maximum of $`T_c`$ at $`x0.15`$ is related to the fact that the ordinary metal clusters appear within the superconducting phase for $`x>0.148`$. Because of this, superconductivity is ”depressed” to a degree corresponding to the ratio of the volumes occupied by superconducting and metal (nonsuperconducting) phases. This is confirmed by measurements of the Meissner-phase volume as a function of a magnetic field. These measurements show that, for $`x=0.15`$, this volume is virtually independent of the field, whereas, for $`x>0.15`$, a magnetic field reduces significantly the volume of superconducting phase. At the same time, in low fields, the magnitude of the Meissner effect increases with $`x`$ up to $`x0.2`$, which is indicative of ordering of the dopants with $`L=\sqrt{5}`$.
In Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub> (Fig. 6(d)), there is virtually no ordering of Ce owing to a small differences in sizes of Nd and Ce ions. Therefore, the $`2D`$-percolation is possible only for $`x>0.118`$ (for $`L=\sqrt{5}`$). However, due to the absence of ordering, a percolation cluster would also include the regions with $`L=2`$; thus, the percolation threshold for $`L=\sqrt{5}`$ shifts to that for $`L=2`$. This is consistent with experimental phase diagram.
The region $`x<0.12`$ for La<sub>2-x</sub>Ba<sub>x</sub>CuO<sub>4</sub> and La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> deserves special consideration (the region of ”underdoping”). As it follows from Fig. 6a, even the $`3D`$-percolation is not possible for $`x<0.077`$, and only the ”traces” of superconductivity can be observed. This conclusion is consistent with the results reported in where the bulk superconductivity was not observed in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> for $`x=0.08`$. As mentioned above, it is not to be expected that the $`2D`$-percolation can originate (at least, in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>) in the range of $`0.08<x<0.12`$ as well because the percolation regions for three different types of segments with $`L^2=8,9,`$ and $`10`$ are close to each other. Most likely, we can have here a tunneling between clusters with $`L=3`$. This inference accounts for the results reported in Ref. where, in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> at $`T0`$, a logarithmic divergence of resistivity was observed for $`x<0.15`$, with superconductivity being depressed by a magnetic field.
The proposed model allows to give an alternate interpretation of the experiments on the pseudogap observation in under- and optimally- doped HTS. As it follows from experiment the pseudogap has the same symmetry and about the same value as the superconducting gap but it vanishes at $`T^{}>T_c`$ ($`T^{}`$ decrease with $`x`$ down to $`T_c`$).
We believe that the observed pseudogap is nothing but just the same superconducting gap opening at $`T>T_c`$ because of the large fluctuations of the number of particles between valence band and pair level in the short chains of -U centers. The point is that in HTS the mechanism of the superconducting gap suppression is the occupation of pair level with electrons. Therefore the decrease of pair level occupation by fluctuation may result in turning on superconductivity in such chains at $`T^{}>T_c`$ (first-order phase transition.) At small $`x`$ the great bulk of -U centers are grouped on short chains where large relative fluctuations of the number of particles are possible. With $`x`$ the increasing part of -U centers belong to the infinite percolation cluster. Therefore $`T^{}`$ decrease with $`x`$ down to $`T_c`$, unless all -U centers are integrated in infinite percjlation cluster.
Thus, we may conclude that all special features observed in phase diagrams of Ln-214 HTS represent no more than geometric relations in a square lattice and a competition between ordering and disorder in distribution of the dopant ions. These items determine the range of existence of infinite percolation cluster integrating the localized orbitals of singlet hole pairs of -U centers. The agreement between experimental and calculated phase diagrams supports the conclusion that it is the considered fragments involving the pairs of neighboring Cu ions in the CuO<sub>2</sub> plane which are responsible for superconductivity in Ln-214. It is interesting that the percolation cluster in Ln-214 resembles the Little’s polymer while that in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> and Bi<sub>2</sub>Sr<sub>2</sub>Ca<sub>2</sub>Cu<sub>2</sub>O<sub>8+δ</sub> (where doped holes are in the plane parallel to the CuO<sub>2</sub> plane) resembles Ginsburg’s sandwich .
## IV Hole carrier generation and relaxation processes in HTS
Now we consider the process of the generation of hole carriers as well as the features of the transport and optical properties of HTS in the framework of the proposed model of electronic spectrum (Fig. 1(b)). Let there is an infinite cluster including a quantity of -U centers together with the nearest oxygen ions. The two-particle hybridization of the pair states with the O2$`p_{x,y}`$ states of oxygen ions surrounding -U centers results in the broadening both of the pair states and of the O2$`p_{x,y}`$ ones. (In this case the electrons at the top of the valence band resemble the ”marginal” Fermi liquid). The pair level broadening can be expressed as $`\mathrm{\Gamma }\pi (DV)^2kT`$, where $`V`$ is the constant of hybridization, D is the density of states at the top of O2p band. The broadening of band states is $`\gamma \mathrm{\Gamma }/DE_0`$, where $`E_0`$ is the energy width of statistic distribution of pair states over sites. This broadening smoothes the features of the band state density and results in its independence on energy over the interaction range. The occupation of pair level results from transitions of electrons from $`\pi p_{x,y}`$ oxygen orbitals to the -U centers and accompanied by the generation of hole carriers in the oxygen band. We call the phase where the additional hole carriers appear through this mechanism as -U phase. The electron occupancy of -U centers $`\eta `$ as well as the hole concentration $`n`$ in -U phase are determined by the balance of rates of electron pair transitions from oxygen band to pair level and back. If $`N`$ is -U center concentration so $`n=2N\eta `$. The rate of electron transitions from pair level to singlet orbitals is proportional to $`N\eta \mathrm{\Gamma }\eta T`$. The rate of reverse process is determined by the frequency of electron-electron scattering and proportional to $`\gamma ^2(1\eta )T^2(1\eta )`$. Therefore
$$n=2NT/(T_0+T),$$
(2)
where $`T_0`$ is temperature-independent value. Thus $`nT`$ at low temperature and tends to $`2N`$ at high temperature. This is in agreement with the data of Hall measurements in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>, where doped carriers in chains do not contribute to Hall conductivity. As it follows from the above consideration the hole carrier distribution turns out to be nondegenerated owing to the interaction with -U centers. Taking into account the nondegeneracy of carrier distribution (absence of the Pauly blocking) and their high concentration $`(10^{21}10^{22}`$ cm$`{}_{}{}^{3})`$ the electron-electron scattering (more precisely to say hole-hole scattering in this case) is likely to provide the dominant contribution to the relaxation processes in HTS. As far as the interaction of two holes in the system with -U centers corresponds to the effective attraction, it will not be conventional Coulomb scattering. The main mechanism of carrier relaxation in HTS is likely to be similar to that assumed for the metals and alloys with strong electron-phonon coupling. In these materials the electrons in a layer $`k\mathrm{\Theta }_D`$ at the Fermi surface ($`\mathrm{\Theta }_D`$ is Debye temperature) are attracted due to virtual phonon exchange. This electron-electron interaction enhanced by phonon-mediation effects far exceeds the screened Coulomb repulsion. Therefore in the metals with strong electron-phonon interaction the dominant channel of the electron-electron scattering will be determined by the virtual phonon exchange. The contribution of these processes becomes essential for $`T<\mathrm{\Theta }_D`$ with the amplitude of the electron-electron scattering being independent of scattering particle energy $`E`$ for $`Ek\mathrm{\Theta }_D`$ and dropping sharply at $`Ek\mathrm{\Theta }_D`$. For $`E>k\mathrm{\Theta }_D`$ only the Coulomb interaction gives the contribution to the amplitude of electron-electron scattering. The electron-electron contribution to the electrical resistivity $`\rho `$ ($`\rho =AT^2`$) exceeding the electron-phonon one, have been observed experimentally for Al at $`T<4`$ K and for the A15 superconductors at $`T<50`$ K, with the amplitude $`A`$ more than one order exceeding the result of calculations based on the Coulomb mechanism of scattering.
The dominant relaxation process in HTS is the electron-electron scattering as well and the main channel of this scattering is the interaction of hole pairs on -U centers. This process may be considered as an exchange by virtual bosons (exitons) with energy $`W`$. As far as $`W0.11`$ eV (in contrast to the exchange by virtual phonon with energy $`Ek\mathrm{\Theta }_D<0.03`$ eV) the temperature range, where the scattering processes with exchange of virtual bosons are dominant, expands to $`T10^3`$ K. The temperature dependence $`\rho (T)`$ in such a model may be obtained from the Drude formula $`\rho =m^{}\nu /ne^2`$ (here $`m^{}`$ is the effective mass of holes, $`\nu `$ \- the rate of hole-hole scattering). For $`WE`$ the amplitude of scattering doesn’t depend on the particle energy. Assuming that density of states is energy-independent $`\nu `$ will be determined only by the hole concentration and statistical factor in scattering cross-section (i.e. the phase volume available for occupation with scattered particles). The latter is proportional to $`E_1+E_2`$ (here $`E_1`$ and $`E_2`$ are the energies of scattering holes measured from the top of valence band). Thus
$$\nu n(E_1+E_2).$$
(3)
For dc-conductivity $`E_1E_2\gamma T`$ and $`\nu nT`$, so that $`\rho (T)T`$. This kind of dependence has been observed experimentally for the optimally doped samples of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>, La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, Bi<sub>2</sub>Sr<sub>2</sub>Ca<sub>2</sub>Cu<sub>2</sub>O<sub>8+δ</sub> and others. The dependence $`\rho (T)=\rho _0+bT`$ observed frequently in experiments may be explained by the contribution of insulating tunnel barriers interlayering the areas of the -U phase that results in temperature-independent term $`r_0`$. For the overdoped samples where both -U phase and metal phase coexist the additional carriers from metal phase decrease $`\left|\mathrm{\Delta }E_M\right|`$ and lower the pair level $`\delta E`$ below the top of the valence band. In this case the distribution of hole carriers degenerates and $`n`$ is temperature-independent for $`\gamma \delta E`$. The temperature-dependent contribution in resistivity is $`\rho (T)\gamma ^2T^2`$. For $`\gamma \delta E`$ the transition to the linear dependence $`\rho (T)`$ takes place.
The dominant contribution of electron-electron scattering into the scattering process also has an effect upon both of frequency and temperature dependencies of optical conductivity $`\sigma _{opt}`$:
$$\sigma _{opt}=\frac{e^2n}{m^{}}\frac{\nu _{opt}}{\omega ^2+\nu _{opt}^{}{}_{}{}^{2}}$$
(4)
here $`\omega `$ is the photon frequency, $`\nu _{opt}`$ \- the rate of optical relaxation. For the electron-electron scattering (at $`n=10^{21}10^{22}`$ cm<sup>-3</sup>) the collision frequency is $`10^{14}10^{15}`$ s<sup>-1</sup>. Thus $`\nu _{opt}\omega `$ for IR region and the formula for optical conductivity becomes more simple:
$$\sigma _{opt}(\omega ,T)=e^2n/m^{}\nu _{opt}.$$
(5)
For optical relaxation we have $`E_1\omega ,E_2\gamma T`$. From where $`\sigma _{opt}\omega ^1`$ (for $`\omega \gamma )`$ and $`\sigma _{opt}T^1`$ (for $`\omega <\gamma `$). These results are in a good agreement with the data of various experiments.
## V Conclusions
The mechanism of -U center formation in high-$`T_c`$ superconductors (HTS) with doping is considered. It is shown that the transition of HTS from insulator to metal passes through the particular dopant concentration range where the local transfer of singlet electron pairs from oxygen ions to pairs of neighboring cations (-U centers) are allowed while the single-electron transitions are still forbidden. The additional hole carriers are generated as the result of singlet electron pair transitions from the oxygen ions to the -U-centers. The orbitals of the arising singlet hole pairs are localized in the nearest vicinity of -U center. In such a system the hole conductivity starts up at the dopant concentration exceeding the classical 2d-percolation threshold for singlet hole pair orbitals of -U centers. The main features of hole carriers in HTS normal state are found to be the nondegenerate distribution and the dominant contribution of electron-electron scattering to the hole carrier relaxation processes. These features account for the unusual transport and optical properties of HTS.
In the framework of the proposed model the phase diagram Ln-214 HTS compounds is constructed. The remarkable accord between calculated and experimental phase diagrams may be considered as the confirmation of the supposed model. Thus the HTS s may be considered as the special class of solids in-between insulator and metal where the unusual mechanism of superconductivity resulting from interaction of electron pair with -U centers is realized.
## VI Acknowledgments
We are grateful for valuable discussions with L. V. Keldysh, B. A. Volkov, E. G. Maksimov and P. A. Arseev.
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# Untitled Document
ELLIPTICITY CONDITIONS FOR THE LAX OPERATOR
OF THE KP EQUATIONS
Giampiero Esposito<sup>1,2</sup> and Boris G. Konopelchenko<sup>3,4</sup>
<sup>1</sup>Istituto Nazionale di Fisica Nucleare, Sezione di Napoli, Complesso Universitario di Monte S. Angelo, Via Cintia, Edificio N’, 80126 Napoli, Italy <sup>2</sup>Università di Napoli Federico II, Dipartimento di Scienze Fisiche, Complesso Universitario di Monte S. Angelo, Via Cintia, Edificio N’, 80126 Napoli, Italy <sup>3</sup>Dipartimento di Fisica, Università di Lecce, Via Arnesano, 73100 Lecce, Italy <sup>4</sup>Istituto Nazionale di Fisica Nucleare, Sezione di Lecce, 73100, Italy
Abstract. The Lax pseudo-differential operator plays a key role in studying the general set of KP equations, although it is normally treated in a formal way, without worrying about a complete characterization of its mathematical properties. The aim of the present paper is therefore to investigate the ellipticity condition. For this purpose, after a careful evaluation of the kernel with the associated symbol, the majorization ensuring ellipticity is studied in detail. This leads to non-trivial restrictions on the admissible set of potentials in the Lax operator. When their time evolution is also considered, the ellipticity conditions turn out to involve derivatives of the logarithm of the $`\tau `$-function.
PACS numbers: 05.45.Y
1. Introduction
Several important developments in modern mathematical physics are due to the investigation of pseudo-differential operators on $`𝐑^m`$ and on general Riemannian manifolds . For our purposes, it is sufficient to recall the following basic properties. (i) A linear partial differential operator $`P`$ of order $`d`$ can be written in the form
$$P\underset{|\alpha |d}{}a_\alpha (x)D_x^\alpha $$
$`(1.1)`$
where (here $`i\sqrt{1}`$)
$$|\alpha |\underset{k=1}{\overset{m}{}}\alpha _k$$
$`(1.2)`$
$$D_x^\alpha (i)^{|\alpha |}\left(\frac{}{x_1}\right)^{\alpha _1}\mathrm{}\left(\frac{}{x_m}\right)^{\alpha _m}$$
$`(1.3)`$
and $`a_\alpha `$ is a $`C^{\mathrm{}}`$ function on $`𝐑^m`$ for all $`\alpha `$. The associated symbol is, by definition,
$$p(x,\xi )\underset{|\alpha |d}{}a_\alpha (x)\xi ^\alpha $$
$`(1.4)`$
i.e. it is obtained by replacing the differential operator $`D_x^\alpha `$ by the monomial $`\xi ^\alpha `$. The pair $`(x,\xi )`$ may be viewed as defining a point of the cotangent bundle of $`𝐑^m`$, and the action of $`P`$ on the elements of the Schwarz space $`𝒮`$ of smooth complex-valued functions on $`𝐑^m`$ of rapid decrease is given by
$$Pf(x)\mathrm{e}^{i(xy)\xi }p(x,\xi )f(y)𝑑y𝑑\xi $$
$`(1.5)`$
where the $`dy=dy_1\mathrm{}dy_m`$ and $`d\xi =d\xi _1\mathrm{}d\xi _m`$ orders of integration cannot be interchanged, since the integral is not absolutely convergent. (ii) Pseudo-differential operators are instead a more general class of operators whose symbol need not be a polynomial but has suitable regularity properties. More precisely, let $`S^d`$ be the set of all symbols $`p(x,\xi )`$ such that (1) $`p`$ is smooth in $`(x,\xi )`$, with compact $`x`$ support. (2) For all $`(\alpha ,\beta )`$, there exist constants $`C_{\alpha ,\beta }`$ for which
$$\left|D_x^\alpha D_\xi ^\beta p(x,\xi )\right|C_{\alpha ,\beta }(1+|\xi |)^{d|\beta |}$$
$`(1.6)`$
for some real (not necessarily positive) value of $`d`$, where $`|\beta |_{k=1}^m\beta _k`$ (see (1.2)). The associated pseudo-differential operator, defined on the Schwarz space and taking values in the set of smooth functions on $`𝐑^m`$ with compact support:
$$P:𝒮C_c^{\mathrm{}}(𝐑^m)$$
is defined in a way formally analogous to Eq. (1.5). (iii) Let now $`U`$ be an open subset with compact closure in $`𝐑^m`$, and consider an open subset $`U_1`$ whose closure $`\overline{U}_1`$ is properly included into $`U`$: $`\overline{U}_1U`$. If $`p`$ is a symbol of order $`d`$ on $`U`$, it is said to be elliptic on $`U_1`$ if there exists an open set $`U_2`$ which contains $`\overline{U}_1`$ and positive constants $`C_i`$ so that
$$|p(x,\xi )|^1C_1(1+|\xi |)^d$$
$`(1.7)`$
for $`|\xi |C_0`$ and $`xU_2`$, where
$$|\xi |\sqrt{g^{ab}(x)\xi _a\xi _b}=\sqrt{\underset{k=1}{\overset{m}{}}\xi _k^2}.$$
$`(1.8)`$
The corresponding operator $`P`$ is then elliptic. From a mathematical point of view, pseudo-differential operators occur in many problems in global analysis , and recent developments deal with the functional calculus of pseudo-differential boundary-value problems . From a physical point of view, such a formalism is important in quantum gravity and quantum field theory \[4–6\]. In particular, we are here interested in an interdisciplinary field, i.e. a rigorous approach to the Kadomtsev–Petviashvili (hereafter KP) equations. Recall that the KP equation can be written in the form
$$\frac{}{x}\left(\frac{u}{t}+6u\frac{u}{x}+\frac{^3u}{x^3}\right)=3\alpha ^2\frac{^2u}{y^2}.$$
$`(1.9)`$
If $`\alpha ^2=1`$, it describes an Hamiltonian wave system which is exactly solvable but not Liouville integrable. It exhibits a degenerative dispersion, the asymptotic states for $`t\pm \mathrm{}`$ do not coincide and an infinite number of invariants of motion exist. If $`\alpha ^2=1`$, it describes an Hamiltonian wave system exactly solvable and completely integrable. It exhibits non-degenerative dispersion and lack of decay, and the asymptotic states for $`t\pm \mathrm{}`$ coincide upon imposing rapid-decrease boundary conditions. Moreover, the number of invariants of motion remains infinite .
The general set of KP equations may be described by using the first-order operator
$$T\frac{}{x}=_x$$
$`(1.10)`$
with $`x𝐑`$, and the associated Lax pseudo-differential operator
$$LT+\underset{k=1}{\overset{\mathrm{}}{}}u_k(x,t_1,t_2,t_3,\mathrm{})T^k$$
$`(1.11)`$
where the functions $`u_k`$ are here called the ‘potentials’. By doing so, one allows in general for their dependence on an infinite number of time variables $`t(t_1,t_2,\mathrm{},t_p,\mathrm{})`$. On assuming that $`T^1`$ is a well defined inverse operator (see section 2), so that
$$TT^1=T^1T=1\mathrm{I}$$
$`(1.12)`$
one can compose the Lax operator with itself, giving rise to its ‘powers’, i.e. $`L^nLL^{n1}`$, for all $`n=2,3,\mathrm{},\mathrm{}`$. Each such power has a differential part, denoted by $`B_n`$. To begin one sets
$$B_1T$$
$`(1.13)`$
and, by virtue of (1.12), one finds
$$B_2T^2+2u_1$$
$`(1.14)`$
$$B_3T^3+3u_1T+3(u_2+u_{1,x})$$
$`(1.15)`$
and so on. The KP hierarchy of integrable equations is then defined by the generalized Lax equation
$$\frac{L}{t_n}=[B_n,L]=B_nLLB_n$$
$`(1.16)`$
and by the Zakharov–Shabat equation
$$\frac{B_m}{t_n}\frac{B_n}{t_m}=[B_n,B_m]$$
$`(1.17)`$
which may be seen as the compatibility conditions of the linear equations
$$L\psi =\lambda \psi $$
$`(1.18)`$
and
$$\frac{\psi }{t_n}=B_n\psi $$
$`(1.19)`$
for all $`n`$, under the assumption that
$$\frac{\lambda }{t_n}=0.$$
$`(1.20)`$
At this stage, since $`t_1`$ plays the same role as $`x`$, $`t_1`$ or $`x`$ are used without distinction in the literature . Once the equations (1.16) are written for all values of $`n`$, the coefficients of $`T^k`$ are equated, and this leads to an infinite set of equations
$$\frac{u_k}{t_n}=\phi _{kn}$$
$`(1.21)`$
where $`\phi _{kn}`$ are certain differential polynomials in the potentials and their derivatives. For example, from the equations
$$\frac{u_1}{t_2}=u_{1,xx}+2u_{2,x}$$
$`(1.22)`$
$$\frac{u_2}{t_2}=u_{2,xx}+2u_{3,x}+2u_1u_{1,x}$$
$`(1.23)`$
$$\frac{u_1}{t_3}=u_{1,xxx}+3u_{2,xx}+3u_{3,x}+6u_1u_{1,x}$$
$`(1.24)`$
one gets the KP equation for $`u_1`$:
$$\frac{}{x}\left(4\frac{u_1}{t_3}12u_1\frac{u_1}{x}\frac{^3u_1}{x^3}\right)3\frac{^2u_1}{t_2^2}=0.$$
$`(1.25)`$
In section 2 we derive the kernel of the Lax pseudo-differential operator by using a Green-function method and the theory of distributions when all potentials only depend on $`x`$. In section 3 we obtain the symbol from the kernel of section 2, and derive suitable majorizations which ensure ellipticity of the Lax operator. Strong ellipticity is studied in section 4. Behaviour of the ellipticity conditions under KP flows is investigated in section 5, and concluding remarks are presented in section 6, while the appendix describes relevant details.
2. The Lax operator and its kernel
Following , we first consider a ‘restricted’ form of the Lax operator, for which the potentials $`u_k`$ only depend on $`x`$. The general form (1.11) will be restored in section 5, where time evolution is studied (cf section 3 of ).
Once the operator (1.10) is given, the inverse operator $`T^1`$ is an integral operator with kernel given by the Green function $`G_1(x,y)`$ of $`T`$. Its action on any function $`f`$ in its domain reads
$$(T^1f)(x)=_{\mathrm{}}^{\mathrm{}}G_1(x,y)f(y)𝑑y$$
$`(2.1)`$
where the Green function $`G_1`$ obeys the equation
$$T_xG_1(x,y)=\delta (x,y).$$
$`(2.2)`$
More precisely, the Green function $`G_1`$ is a kernel which solves the equation
$$\frac{}{x}G_1(x,y)=0xy$$
$`(2.3)`$
and the jump condition
$$\underset{xy^+}{lim}G_1(x,y)\underset{xy^{}}{lim}G_1(x,y)=1.$$
$`(2.4)`$
The problem described by Eqs. (2.3) and (2.4) is solved by
$$G_1(x,y)=A_{1,1}(y)\mathrm{if}x>y$$
$`(2.5a)`$
$$G_1(x,y)=A_{1,1}(y)1\mathrm{if}x<y$$
$`(2.5b)`$
where $`A_{1,1}(y)`$ is an arbitrary smooth function of $`y`$ unless a suitable boundary condition is specified (see below).
Similarly, the operators $`T^2,T^3`$ and so on are integral operators with kernel given by the Green function of $`T^2,T^3,\mathrm{},`$ respectively. For example, the operator $`T^2=TT`$ has a Green function $`G_2(x,y)`$ satisfying the differential equation
$$\frac{^2}{x^2}G_2(x,y)=0xy$$
$`(2.6)`$
the continuity condition
$$\underset{xy^+}{lim}G_2(x,y)=\underset{xy^{}}{lim}G_2(x,y)$$
$`(2.7)`$
and the jump condition
$$\underset{xy^+}{lim}\frac{G_2}{x}\underset{xy^{}}{lim}\frac{G_2}{x}=1.$$
$`(2.8)`$
Equations (2.6)–(2.8) are solved by
$$G_2(x,y)=A_{1,2}(y)+A_{2,2}(y)x\mathrm{if}x>y$$
$`(2.9a)`$
$$G_2(x,y)=y+A_{1,2}(y)+(A_{2,2}(y)1)x\mathrm{if}x<y$$
$`(2.9b)`$
where now two arbitrary functions $`A_{1,2}`$ and $`A_{2,2}`$ are involved because $`G_2`$ is the Green function of a second-order differential operator.
It is therefore clear that, assuming for the time being that $`u_k`$ only depends on $`x`$, the Lax operator (1.11) can be viewed as an integral operator whose action is given by
$$(L\psi )(x)=_{\mathrm{}}^{\mathrm{}}K(x,y)\psi (y)𝑑y$$
$`(2.10)`$
with kernel
$$K(x,y)=\delta ^{}(x,y)+\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)G_p(x,y)$$
$`(2.11)`$
where we have used the well known distributional action of the first derivative of the Dirac delta functional , and the Green function $`G_p(x,y)`$ can be expressed in the form
$$G_p(x,y)=\underset{r=1}{\overset{p}{}}C_{r,p}(y)x^{r1}.$$
$`(2.12)`$
The ‘coefficients’ $`C_{r,p}`$ are actually functions of $`y`$ obeying a law of the type (see (2.5) and (2.9))
$$C_{r,p}(y)=A_{r,p}(y)\mathrm{if}x>y$$
$`(2.13a)`$
$$C_{r,p}(y)=B_{r,p}(y)\mathrm{if}x<y$$
$`(2.13b)`$
where the coefficients $`B_{r,p}(y)`$ in (2.13b) can be expressed in terms of the coefficients $`A_{r,p}(y)`$ after imposing the continuity conditions
$$\underset{xy^+}{lim}\frac{^qG_p}{x^q}\underset{xy^{}}{lim}\frac{^qG_p}{x^q}=0q=0,1,\mathrm{},p2$$
$`(2.14)`$
and the jump condition
$$\underset{xy^+}{lim}\frac{^{p1}G_p}{x^{p1}}\underset{xy^{}}{lim}\frac{^{p1}G_p}{x^{p1}}=1.$$
$`(2.15)`$
In , the Green function $`G_1(x,y)`$ given by (2.5a) and (2.5b) has been written in the form
$$G_1(x,y)=\frac{1}{2}[\theta (xy)\theta (yx)]$$
$`(2.16)`$
where $`\theta `$ is the step function such that $`\theta (x)=1`$ if $`x>0`$, $`\theta (0)=\frac{1}{2}`$, $`\theta (x)=0`$ if $`x<0`$. Equation (2.16) corresponds to choosing
$$A_{1,1}(y)=\frac{1}{2}$$
$`(2.17)`$
in Eqs. (2.5a) and (2.5b). The operator $`T^1`$ is then the integral operator
$$T^1:f\frac{1}{2}_{\mathrm{}}^xf(y)𝑑y\frac{1}{2}_x^{\mathrm{}}f(y)𝑑y.$$
$`(2.18)`$
Similarly, the operator $`T^2`$ turns out to be the integral operator
$$\begin{array}{ccc}& T^2:f\frac{1}{4}_{\mathrm{}}^x𝑑y_{\mathrm{}}^yf(z)𝑑z\frac{1}{4}_{\mathrm{}}^x𝑑y_y^{\mathrm{}}f(z)𝑑z\hfill & \\ & \frac{1}{4}_x^{\mathrm{}}𝑑y_{\mathrm{}}^yf(z)𝑑z+\frac{1}{4}_x^{\mathrm{}}𝑑y_y^{\mathrm{}}f(z)𝑑z.\hfill & (2.19)\hfill \end{array}$$
By comparison with (2.9a) and (2.9b) this leads to the evaluation of $`A_{1,2}(y)`$ and $`A_{2,2}(y)`$, and the procedure can be iterated (in principle) to obtain all $`A_{r,p}(y)`$ coefficients in (2.13a), while the $`B_{r,p}(y)`$ are obtained after imposing (2.14) and (2.15) as we said before. The details of the construction are indeed a bit involved, and hence it is worth showing what can be done with the integral operator $`T^2`$ given in (2.19). On the one hand, Eqs. (2.9a) and (2.9b) lead to
$$\begin{array}{ccc}\hfill (T^2f)(x)& =_{\mathrm{}}^xdy[A_{1,2}(y)+xA_{2,2}(y)]f(y)\hfill & \\ & +_x^{\mathrm{}}dy[y+A_{1,2}(y)+x(A_{2,2}(y)1)]f(y).\hfill & (2.20)\hfill \end{array}$$
On the other hand, by virtue of (2.19), $`(T^2f)(x)`$ is also given by
$$(T^2f)(x)=_{\mathrm{}}^x𝑑y\frac{1}{4}h(y)+_x^{\mathrm{}}𝑑y\left(\frac{1}{4}h(y)\right)$$
$`(2.21)`$
where
$$h(y)_{\mathrm{}}^yf(z)𝑑z_y^{\mathrm{}}f(z)𝑑z.$$
$`(2.22)`$
Direct comparison of the representations (2.20) and (2.21) of $`(T^2f)(x)`$ yields therefore the equations
$$[A_{1,2}(y)+xA_{2,2}(y)]f(y)=\frac{1}{4}h(y)$$
$`(2.23)`$
$$[y+A_{1,2}(y)+x(A_{2,2}(y)1)]f(y)=\frac{1}{4}h(y).$$
$`(2.24)`$
The addition of (2.23) and (2.24) leads to
$$[2A_{1,2}(y)+2xA_{2,2}(y)+yx]f(y)=0$$
$`(2.25)`$
which is satisfied for all $`f(y)`$ if and only if
$$2A_{1,2}(y)+y=0$$
$`(2.26)`$
$$(2A_{2,2}(y)1)x=0.$$
$`(2.27)`$
Such a system is solved by
$$A_{1,2}(y)=\frac{y}{2}$$
$`(2.28)`$
$$A_{2,2}(y)=\frac{1}{2}$$
$`(2.29)`$
which provides the desired explicit formula for the Green function $`G_2(x,y)`$, upon insertion into (2.9a) and (2.9b).
3. Symbol and ellipticity
Recall now that, if $`L`$ is a pseudo-differential operator defined by a kernel $`K`$, this is related to the symbol $`p(x,\xi )`$ by the equation
$$K(x,y)=(2\pi )^n_{𝐑^n}\mathrm{e}^{i(xy)\xi }p(x,\xi )𝑑\xi .$$
$`(3.1)`$
This equation can be inverted to give a very useful formula for the symbol, i.e. (cf Eq. (2.1.36) in )
$$p(x,\xi )=_{𝐑^n}\mathrm{e}^{iz\xi }K(x,xz)𝑑z.$$
$`(3.2)`$
Equation (3.2) is a key formula for our investigation, because the ellipticity of $`L`$ is defined in terms of its symbol, as we know from the introduction, following .
In our problem, which involves $`x𝐑`$, the integral (3.2) reduces to
$$p(x,\xi )=_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }K(x,xz)𝑑z$$
$`(3.3)`$
where the kernel $`K(x,y)`$ is expressed by (2.11)–(2.13), and we have to check that the inequality (1.7) is satisfied for $`|\xi |C_0`$ to obtain ellipticity. Indeed, the symbol (3.3) turns out to be
$$\begin{array}{ccc}\hfill p(x,\xi )& =_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\delta ^{}(z)𝑑z\hfill & \\ & +_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\hfill & (3.4)\hfill \end{array}$$
and hence obeys the inequality
$$\begin{array}{ccc}\hfill |p(x,\xi )|& \left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\delta ^{}(z)𝑑z\right|\hfill & \\ & \left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\right|.\hfill & (3.5)\hfill \end{array}$$
Of course, the first integral on the right-hand side of (3.5) becomes meaningful within the framework of Fourier transform of distributions . In the simplest possible terms, one has actually to consider the parameter-dependent integral (here $`a>0`$)
$$\begin{array}{ccc}\hfill I_{1,a}(\xi )& _{\mathrm{}}^{\mathrm{}}\mathrm{e}^{az^2}\mathrm{e}^{iz\xi }\delta ^{}(z)𝑑z\hfill & \\ & =_{\mathrm{}}^{\mathrm{}}\delta (z)(2azi\xi )\mathrm{e}^{az^2iz\xi }𝑑z.\hfill & (3.6)\hfill \end{array}$$
By virtue of the property defining the Dirac delta functional, according to which
$$(\delta ,f)=f(0)$$
$`(3.7)`$
the integral (3.6) equals $`i\xi `$, and hence the first term on the right-hand side of (3.5) equals $`|\xi |`$. Now we distinguish two cases, depending on whether
$$f(x,\xi )|\xi |\left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\right|$$
$`(3.8)`$
is positive or negative. If
$$f(x,\xi )>0$$
$`(3.9)`$
holds, the majorization (1.7) for the ellipticity of the restricted Lax operator is satisfied provided that, for $`|\xi |C_0`$,
$$\left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\right||\xi |C_1^1(1+|\xi |)^d.$$
$`(3.10)`$
In contrast, if
$$f(x,\xi )<0$$
$`(3.11)`$
holds, the restricted Lax operator is elliptic provided that
$$\begin{array}{ccc}& \left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\right||\xi |C_1^1(1+|\xi |)^d\hfill & \\ & C_0C_1^1(1+|\xi |)^d\hfill & (3.12)\hfill \end{array}$$
for $`|\xi |C_0`$. If the order $`d`$ of the Lax operator is positive, we can further write that, for $`|\xi |C_0`$, the majorization (3.10) becomes
$$\left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\right||\xi |C_1^1C_0^d.$$
$`(3.13)`$
If we are interested in sufficient conditions we can point out that, since the inequality (3.5) is always satisfied, whereas (1.7) only holds when $`L`$ is elliptic, the sufficient condition for ellipticity of the restricted Lax operator is expressed by
$$C_1^1(1+|\xi |)^d|\xi |\left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\right|$$
$`(3.14)`$
for $`|\xi |C_0`$. This leads in turn to the inequality (3.10).
To sum up, if the function $`f:T^{}(𝐑)𝐑`$ defined by (3.8) has no zeros (which is already a non-trivial requirement on the potentials $`u_p`$), the restricted Lax operator is elliptic provided that either (3.10) or (3.12) is satisfied. The majorization (3.10) is further simplified in the form (3.13) in case of positive order of the Lax operator. A sufficient condition for ellipticity is given instead by (3.14), which coincides with (3.10). In particular, when $`|\xi |=C_0`$ and the equality sign is chosen in (3.14), the order $`d`$ of the restricted Lax operator can be evaluated by the formula
$$d=\frac{\mathrm{log}(C_1(C_0I_\xi ))}{\mathrm{log}(1+C_0)}$$
$`(3.15)`$
where
$$I_\xi \mathrm{sup}_{x𝐑}\left|_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}dz\right|$$
$`(3.16)`$
bearing in mind that the modulus of $`\xi `$ equals the constant $`C_0`$.
4. Strong ellipticity
In a thorough analysis of the ellipticity properties, strong ellipticity should also be studied. For this purpose, following , we assume that the symbol of the restricted Lax operator is polyhomogeneous, in that it admits an asymptotic expansion of the form
$$p(x,\xi )\underset{l=0}{\overset{\mathrm{}}{}}p_{dl}(x,\xi )$$
$`(4.1)`$
where each term $`p_{dl}`$ has the homogeneity property
$$p_{dl}(x,\gamma \xi )=\gamma ^{dl}p_{dl}(x,\xi )$$
$`(4.2)`$
for $`t1`$ and $`|\xi |1`$. The principal symbol $`p^0`$ of the Lax operator is then, by definition,
$$p^0(x,\xi )p_d(x,\xi ).$$
$`(4.3)`$
Strong ellipticity is formulated in terms of the principal symbol, because it requires that
$$\mathrm{Re}p^0(x,\xi )=\frac{1}{2}[p^0(x,\xi )+p^0(x,\xi )^{}]c(x)|\xi |^d$$
$`(4.4)`$
where $`x𝐑,c(x)>0`$ and $`|\xi |1`$. In other words, given a positive function $`c`$, the product $`c(x)|\xi |^d`$ should be always majorized by the real part of the principal symbol of the restricted Lax operator. Indeed, the symbol (3.4) is such that
$$\begin{array}{ccc}& p(x,\gamma \xi )=i\gamma \xi \hfill & \\ & +\gamma ^1_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}\left(x\frac{z}{\gamma }\right)x^{r1}dz.\hfill & (4.5)\hfill \end{array}$$
By virtue of (4.1), (4.2) and (4.5) we find that
$$\begin{array}{ccc}& i\gamma \xi +\gamma ^1_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }\underset{p=1}{\overset{\mathrm{}}{}}u_p(x)\underset{r=1}{\overset{p}{}}C_{r,p}\left(x\frac{z}{\gamma }\right)x^{r1}dz\hfill & \\ & \underset{l=0}{\overset{\mathrm{}}{}}\gamma ^{dl}p_{dl}(x,\xi ).\hfill & (4.6)\hfill \end{array}$$
Moreover, the term on the right-hand side of (4.6) with $`l=0`$ should be the one occurring in the condition (4.4) for strong ellipticity. A mathematical advantage of strong ellipticity lies in the possibility of having a well defined functional trace of the heat semigroup associated to the Lax operator .
5. Behaviour of the ellipticity conditions under KP flows
To study the preservation (or violation) of the ellipticity conditions under KP flows one has to analyze the following problem: suppose that the conditions (3.10) or (3.12) are satisfied for $`t=0`$. Are they still valid for all or some $`t>0\mathrm{?}`$
This means that we consider again the potentials $`u_k`$ as in Eq. (1.11), i.e. as functions of $`x`$ and $`t`$, where $`t`$ is a concise notation for infinitely many time parameters $`(t_1,t_2,\mathrm{})`$. It should be stressed that we consider only one spatial variable and infinitely many time parameters, since otherwise it would be problematic, at least for the authors, to generalize formulae like (3.2) aimed at obtaining the symbol from the kernel of the operator. For our purposes it is convenient to use formulae generating such potentials by means of a single function. This is made possible by the $`\tau `$-function (see appendix), because one finds
$$u_1(x;t)=\frac{^2}{x^2}\mathrm{log}\tau $$
$`(5.1)`$
$$u_2(x;t)=\frac{1}{2}\left(\frac{^3}{x^3}+\frac{^2}{xt_2}\right)\mathrm{log}\tau $$
$`(5.2)`$
$$u_3(x;t)=\frac{1}{6}\left(\frac{^4}{x^4}3\frac{^3}{x^2t_2}+2\frac{^2}{xt_3}\right)\mathrm{log}\tau u_1^2$$
$`(5.3)`$
and infinitely many other equations of the general form
$$u_k(x;t)=F_k(\mathrm{log}\tau )$$
$`(5.4)`$
where $`F_k`$ is, in general, a non-linear function of $`\mathrm{log}\tau `$. Equations (5.4), with $`k`$ ranging from $`1`$ through $`\mathrm{}`$, should be inserted into the ellipticity conditions (3.10) and (3.14), substituting therein $`u_k`$ with $`F_k(\mathrm{log}\tau )`$ for all $`k`$. The resulting majorizations involve non-linear functions of the logarithm of the $`\tau `$-function.
Further progress can be made by considering a ‘truncated’ Lax operator, e.g.
$$\stackrel{~}{L}T+u_1(x;t)T^1+u_2(x;t)T^2.$$
$`(5.5)`$
This should not seem an arbitrary simplification, because the Lax operator is obtained from the $`W`$ operator of the appendix as
$$LWTW^1.$$
$`(5.6)`$
Now both $`W`$ and its ‘truncated version’
$$W_m1+\underset{k=1}{\overset{m}{}}w_kT^k$$
$`(5.7)`$
satisfy the Sato equation (A8), from which the generalized Lax equation (1.16) is eventually obtained. Thus, operators like $`\stackrel{~}{L}`$ in (5.5) can be obtained from (5.6) if $`W`$ is replaced by $`W_m`$ therein. In this case, on defining
$$\begin{array}{ccc}& F(x,z;t)\underset{p=1}{\overset{2}{}}u_p(x;t)\underset{r=1}{\overset{p}{}}C_{r,p}(xz)x^{r1}\hfill & \\ & =u_1(x;t)C_{1,1}(xz)+u_2(x;t)[C_{1,2}(xz)+xC_{2,2}(xz)]\hfill & (5.8)\hfill \end{array}$$
the integral on the right-hand side of the ellipticity condition (3.14) reduces to $`(J_1+J_2+J_3)(x,\xi ;t)`$ where, bearing in mind that (see (2.17), (2.28) and (2.29))
$$C_{1,1}(xz)=C_{2,2}(xz)=\frac{1}{2}\mathrm{if}z>0,\frac{1}{2}\mathrm{if}z<0$$
$`(5.9)`$
$$C_{1,2}(xz)=\frac{x}{2}+\frac{z}{2}\mathrm{if}z>0,\frac{x}{2}\frac{z}{2}\mathrm{if}z<0$$
$`(5.10)`$
one finds
$$\begin{array}{ccc}& J_1(x,\xi ;t)_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }u_1(x;t)C_{1,1}(xz)𝑑z\hfill & \\ & =u_1(x;t)\underset{b\mathrm{}}{lim}\left(\frac{2}{\xi }\mathrm{e}^{ib\frac{\xi }{2}}\mathrm{sin}\frac{b\xi }{2}\right)\hfill & (5.11)\hfill \end{array}$$
$$\begin{array}{ccc}& J_2(x,\xi ;t)_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }u_2(x;t)C_{1,2}(xz)𝑑z\hfill & \\ & =xu_2(x;t)\frac{J_1(x,\xi ;t)}{u_1(x;t)}+u_2(x;t)i\frac{d}{d\xi }\frac{J_1(x,\xi ;t)}{u_1(x;t)}\hfill & (5.12)\hfill \end{array}$$
$$\begin{array}{ccc}& J_3(x,\xi ;t)_{\mathrm{}}^{\mathrm{}}\mathrm{e}^{iz\xi }u_2(x;t)xC_{2,2}(xz)𝑑z\hfill & \\ & =xu_2(x;t)\frac{J_1(x,\xi ;t)}{u_1(x;t)}.\hfill & (5.13)\hfill \end{array}$$
In these equations, the infinite upper limit of integration can be recovered by taking the limit as $`b\mathrm{}`$ of integrals from $`0`$ to $`b`$ (the lower limit being amenable to $`0`$ by virtue of (5.9) and (5.10). However, divergences remain, and hence we are only able to obtain well defined formulae by integrating up to finite values of $`b`$. The cancellation of terms involving $`xu_2(x;t)`$ is thus found to occur on performing the sum, and our integral reads eventually
$$(J_1+J_2+J_3)_b(x,\xi ;t)=u_1(x;t)F_b(\xi )+u_2(x;t)i\frac{d}{d\xi }F_b(\xi )$$
$`(5.14)`$
having defined
$$F_b(\xi )\frac{2}{\xi }\mathrm{e}^{ib\frac{\xi }{2}}\mathrm{sin}\frac{b\xi }{2}.$$
$`(5.15)`$
Now we point out that
$$\begin{array}{ccc}& \left|u_1(x;t)F_b(\xi )+u_2(x;t)iF_b^{}(\xi )\right|\left|u_1(x;t)F_b(\xi )\right|+\left|u_2(x;t)iF_b^{}(\xi )\right|\hfill & \\ & \frac{2}{\xi }\left|u_1(x;t)\right|+\left(\frac{b}{\xi }+\frac{2}{\xi ^2}\right)\left|u_2(x;t)\right|.\hfill & (5.16)\hfill \end{array}$$
Thus, a sufficient condition for the validity of the majorization (3.14) is expressed by
$$|\xi |C_1^1(1+|\xi |)^d+\frac{2}{\xi }|u_1(x;t)|+\left(\frac{b}{\xi }+\frac{2}{\xi ^2}\right)|u_2(x;t)|.$$
$`(5.17)`$
It should be stressed that divergent integrals occur already in the ‘time-independent’ ellipticity condition (3.14). When the time parameters are introduced, it may be easier or harder to fulfill ellipticity, depending on the behaviour of $`u_1,u_2,\mathrm{}`$ (which, in turn, all depend on the $`\tau `$-function).
6. Concluding remarks
Our paper has been motivated by the need to obtain a deeper understanding of the basic structures of modern non-linear physics. It is indeed well known that the Sato equation (A8) generates the generalized Lax equation (1.16), the Zakharov–Shabat equation (1.17) and the inverse spectral transform scheme . Moreover, an infinite number of nonlinear evolution equations (i.e. the KP hierarchy), of which the KP equation is the simplest nontrivial one, share solutions, and the $`\tau `$-function makes it possible to express all such solutions.
The work in the present paper is the first step towards a rigorous investigation of the ellipticity properties of the Lax pseudo-differential operator. We have found that, to achieve ellipticity, including its strong form, the various potentials $`u_k(x)`$ are no longer arbitrary, but should be chosen in such a way that the following conditions hold: (i) The function $`f:T^{}𝐑𝐑`$ defined in (3.8) has no zeros. (ii) The majorization (3.14) holds (some care is actually necessary to deal with the integrals on the right-hand side of (3.14), as is clear from the analysis performed in section 5). (iii) The asymptotic expansion (4.6) can be obtained. Note, however, that violation of (4.2) (i.e. lack of homogeneity) for $`|\xi |<1`$ can cause logarithmic terms in the asymptotic expansion of the kernel defined by (2.11)–(2.13) (cf ).
Moreover, on allowing for the time evolution of the potentials in the Lax operator, now viewed as functions of $`x`$ and of infinitely many time variables, we have obtained an explicit ellipticity condition in terms of the $`\tau `$-function, when attention is restricted to the ‘truncated’ Lax operator (5.5). It now remains to be seen how to deal with the infinite sum over $`p`$ in the ellipticity condition (3.14) when the potentials $`u_p(x;t)`$ appropriate for the ‘full’ Lax operator (1.11) are instead considered. The form of $`u_1`$ and $`u_2`$ remains the one given in (5.1) and (5.2), but the occurrence of an infinite number of such potentials makes it hard to re-express the time-dependent form of the majorization (3.14). The work in has indeed obtained a very useful formula for the $`\tau `$-function, but this remains of little help when the infinite sum over all potentials is performed.
A further interesting issue is the investigation of ellipticity for the formulation of KP hierarchy considered in , where the main new technique, when compared to the traditional approach to the generalized Lax equation, consists of replacing the Lax operator by an $`n^{\mathrm{th}}`$-order formal pseudo-differential operator
$$L_nT^n+\underset{j=\mathrm{}}{\overset{n2}{}}q_jT^jn2.$$
$`(6.1)`$
The authors of have been able to factorize $`L_n`$ into $`n1`$ first-order formal differential operators $`A_k,1kn1`$, and one first-order formal pseudo-differential operator $`\stackrel{~}{A}_n`$, i.e.
$$L_n=\stackrel{~}{A}_nA_{n1}\mathrm{}A_2A_1$$
$`(6.2)`$
where
$$A_kT+\eta _{k,x}\mathrm{\hspace{0.33em}\hspace{0.33em}1}kn$$
$`(6.3)`$
$$\underset{k=1}{\overset{n}{}}\eta _{k,x}=0$$
$`(6.4)`$
$$\stackrel{~}{A}_nA_n+\underset{j=\mathrm{}}{\overset{1}{}}b_{n,j}T^j.$$
$`(6.5)`$
The results and unsolved problems described so far seem to show that new exciting developments might be obtained from the effort of combining some key techniques of non-linear physics and the tools of linear and pseudo-differential operator theory. In particular, the mathematical requirement of ellipticity in the various forms considered in sections 3–5 restricts the potentials $`u_k`$ in a form not previously considered in the literature to our knowledge, which might be used to select the realizations of the Lax operator one is interested in.
Acknowledgment
This work has been partially supported by PRIN97 ‘Sintesi’. Correspondence with Gerd Grubb has been helpful for the authors, as well as constant encouragement from Giuseppe Marmo.
Appendix
Since the general reader is not necessarily familiar with the theory of $`\tau `$-functions, we summarize the main properties hereafter. Following , we study for the operator $`W_m`$ defined in (5.7) the ordinary differential equation
$$W_m^mf(x)=(^m+w_1(x)^{m1}+\mathrm{}+w_m(x))f(x)=0$$
$`(A1)`$
which has $`m`$ linearly independent solutions $`f^{(1)}(x),\mathrm{},f^{(m)}(x)`$. On writing equation (A1) $`m`$ times with $`f=f^{(1)}(x),\mathrm{},f=f^{(m)}(x)`$, one finds a linear system which can be solved for $`w_j(x)`$, for all $`j=1,\mathrm{},m`$, and hence $`W_m`$ is found from the definition (5.7).
When the $`w_j`$ are taken to depend also on infinitely many time variables $`(t_1,\mathrm{},t_p,\mathrm{})`$, the operator $`W_m(x;t)`$ is found to be
$$W_m(x;t)=\frac{\mathrm{det}A}{\tau (x;t)}$$
$`(A2)`$
where, given the functions $`h_0^{(j)}(x;t)`$ and $`h_n^{(j)}(x;t)`$ such that
$$\left(\frac{}{t_n}\frac{^n}{x^n}\right)h_0^{(j)}(x;t)=0n=1,2,\mathrm{}$$
$`(A3)`$
$$h_0^{(j)}(x;0)=f^{(j)}(x)$$
$`(A4)`$
$$h_n^{(j)}(x;t)=\frac{}{t_n}h_0^{(j)}(x;t)=\frac{^n}{x^n}h_0^{(j)}(x;t)$$
$`(A5)`$
one has (see the definition (1.10))
$$A\left(\begin{array}{cccc}h_0^{(1)}& \mathrm{}& h_0^{(m)}& T^m\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ h_{m1}^{(1)}& \mathrm{}& h_{m1}^{(m)}& T^1\\ h_m^{(1)}& \mathrm{}& h_m^{(m)}& 1\end{array}\right)$$
$`(A6)`$
and
$$\tau (x;t)\mathrm{det}\left(\begin{array}{ccc}h_0^{(1)}& \mathrm{}& h_0^{(m)}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ h_{m1}^{(1)}& \mathrm{}& h_{m1}^{(m)}\end{array}\right).$$
$`(A7)`$
The function $`\tau `$ is the $`\tau `$-function used in Eqs. (5.1)–(5.4), and the time evolution of the operator $`W_m(x;t)`$ is determined by the Sato equation
$$\frac{W_m}{t_n}=B_nW_mW_mT^n$$
$`(A8)`$
where the operators $`B_n`$ are the same occurring in the generalized Lax equation (1.16).
References
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Seeley R T 1969 Topics in pseudo-differential operators, C.I.M.E., in Conf. on Pseudo-Differential Operators, ed. L. Nirenberg (Roma: Edizioni Cremonese)
Grubb G 1996 Functional Calculus of Pseudodifferential Boundary Problems (Progress of Mathematics 65) (Boston: Birkhäuser)
Esposito G 1999 Class. Quantum Grav. 16 1113
Esposito G 1999 Class. Quantum Grav. 16 3999
Esposito G and Stornaiolo C 2000 Int. J. Mod. Phys. A15 449
Ablowitz M J and Clarkson P A 1991 Solitons, Nonlinear Evolution Equations and Inverse Scattering (Cambridge: Cambridge University Press)
Zakharov V E 1991 What is Integrability? (Berlin: Springer-Berlag)
Date E, Kashiwara M, Jimbo M and Miwa T 1983 Transformation groups for soliton equations, in Proc. of RIMS Symposium on Nonlinear Integrable Systems, Classical Theory and Quantum Theory, eds M Jimbo and T Miwa (Singapore: World Scientific)
Ohta Y, Satsuma J, Takahashi D and Tokihiro T 1988 Prog. Theor. Phys. Suppl. 94 210
Lanczos C 1961 Linear Differential Operators (London: Van Nostrand)
Roos B W 1969 Analytic Functions and Distributions in Physics and Engineering (New York: John Wiley)
Kolmogorov A N and Fomin S V 1980 Elements of the Theory of Functions and Functional Analysis (Moscow: Mir)
Grubb G and Seeley R T 1995 Inv. Math. 121 481
Chau L L, Shaw J C and Yen H C 1992 Commun. Math. Phys. 149 263
Gesztesy F and Unterkofler K 1995 Diff. Int. Eqs. 8 797
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# Limits on radio emission from pulsar wind nebulae
## 1 Introduction
Almost all radio pulsars have rotational periods which are steadily increasing with time. This spin down corresponds to a loss of rotational kinetic energy $`\dot{E}4\pi ^2I\dot{P}/P^3`$, where $`I`$ is the moment of inertia of the neutron star (assumed to be $`10^{45}`$ g cm<sup>2</sup>) and $`P`$ is its period; for the known pulsar population $`\dot{E}`$ falls in the range $`10^{28}10^{39}`$ erg s<sup>-1</sup>. The bolometric luminosity of the radio pulses themselves is in all cases a tiny fraction of the spin-down luminosity $`\dot{E}`$, and it is thought that most of the pulsar’s spin-down luminosity is dissipated via a magnetised wind populated by relativistic electrons and positrons \[Rees & Gunn 1974, Michel 1982, Kennel & Coroniti 1984\]. However it is still not well understood how this wind is produced, how it evolves as it flows away from the pulsar, what it is composed of, how its properties depend on those of the pulsar itself, or how it changes as the pulsar ages.
Particles in the wind move along magnetic field lines as they stream away from the pulsar magnetosphere, and produce no observable emission. At some distance from the pulsar, the pressure of the wind is eventually balanced by an external pressure, and the resulting shock randomises the pitch angles of the relativistic particles. These particles consequently gyrate in the local magnetic field and produce synchrotron emission. The properties of the resulting pulsar wind nebula (PWN) can then be used to determine various parameters of the pulsar wind which produces it. At radio wavelengths, PWN are characterised by an amorphous or filled-centre morphology, a moderate degree of linear polarization ($``$20%), and relatively flat spectra ($`\alpha 0.3`$, $`S_\nu \nu ^\alpha `$) \[Weiler & Panagia 1978\].
Various types of PWN are produced, depending on the source of confinement of the wind. Young pulsars are often still located inside their associated supernova remnants (SNRs), and the hot gas produced by the SNR blast-wave provides the confining pressure. These PWN, also known as “plerions”, are typified by the Crab Nebula. If the SNR has dissipated, the confining pressure is then that of the ambient interstellar medium (ISM). This results in much larger “ghost remnants”, which have been proposed but not observed \[Blandford et al. 1973, Cohen et al. 1983\]. In cases where a pulsar has a high space velocity, the ram pressure resulting from its motion can dominate the ambient gas pressure, resulting in a bow-shock PWN (e.g. Frail & Kulkarni 1991).
While PWN can tell us much about pulsar winds, the number of sources we have to study is small – at radio wavelengths, fewer than 10 pulsars have observable PWN. All of these pulsars are very young and have high values of $`\dot{E}`$. While $``$20% of SNRs have a “plerionic” component, in most cases the associated pulsar has not been detected, and without knowing the pulsar parameters it is difficult to constrain the properties of the PWN and the corresponding wind. Thus it is of considerable interest to target candidate pulsars with the intention of either finding new PWN, or determining upper limits on such emission.
Many searches for radio PWN around energetic or fast-moving pulsars have been carried out, at varying resolutions and surface-brightness sensitivities, but usually with no success (e.g. Schönhardt 1974; Weiler, Goss & Schwarz 1974; Cohen et al. 1983). The most recent and comprehensive of these searches, and the only one to specifically target young, energetic or high velocity pulsars, was the recent survey of Frail & Scharringhausen \[Frail & Scharringhausen 1997\], hereafter FS97. FS97 imaged regions around 35 pulsars with the Very Large Array (VLA) at 8.4 GHz, and found no nebular emission associated with any of their targets. Their stringent upper limits allowed them to conclude that most pulsars put less than $`10^6`$ of $`\dot{E}`$ into radio emission from a PWN, $`100`$ times less than observed for the young, high $`\dot{E}`$, pulsars which power detected radio nebulae. Based on this result, FS97 concluded that pulsar winds change in some way as pulsars age and slow down, such that they are no longer efficient at producing radio emission.
However, despite these apparently constraining limits, in hindsight the observing parameters for this search were probably not ideal for looking for PWN. First, FS97 argued that for ambient densities of 1 cm<sup>-3</sup> and pulsar velocities of 150 km s<sup>-1</sup>, PWN around almost all their sources were likely to be unresolved, even at their high spatial resolution of $`0\stackrel{}{.}8`$. However, other choices of ambient density and pulsar velocity can produce PWN with much larger angular extents, resulting in a flux density limit much poorer than estimated by FS97. Furthermore, the maximum scale to which FS97 were sensitive was only $`20^{\prime \prime }`$; they could not detect PWN larger than this at any flux density. Secondly, in 40% of their sample FS97 detected a point source at the position of the pulsar, but had no way of distinguishing between the compact PWN they were looking for and the pulsars themselves. By extrapolating pulsar flux densities from much lower frequencies, FS97 argued that the flux densities they were detecting at 8.4 GHz were consistent with the expected pulsed fluxes, and hence concluded that they were not detecting any PWN in these data. Finally, a typical PWN (spectral index $`\alpha =0.3`$) has a significantly lower flux density at their observing frequency of 8.4 GHz than at lower frequencies.
Motivated by these points, we have undertaken an extensive survey for PWN at both northern and southern declinations, with observing parameters chosen to give much greater sensitivity to PWN. First, we have observed at 1.4 GHz, at which frequency PWN can be expected to be $``$70% brighter than at 8.4 GHz. Secondly, we have observed at a reduced spatial resolution of $``$12<sup>′′</sup>, to give better sensitivity to extended structure. Thirdly, our observations are in telescope configurations whose shortest spacings correspond to a spatial scale of many arcmin. Finally, and most importantly, all our observations have employed pulsar-gating, in which images are made from data taken only when the pulsar is off; we are thus sensitive to compact and unresolved PWN, which might otherwise be masked by the pulsars themselves.
Initial results from this survey, using the Australia Telescope Compact Array (ATCA), have been presented in two previous papers. Gaensler et al. \[Gaensler et al. 1998\], hereafter GSFJ98, reported the discovery of a faint PWN associated with PSR B0906–49, while Stappers, Gaensler & Johnston \[Stappers, Gaensler & Johnston 1999\], hereafter SGJ99, presented non-detections of PWN towards four pulsars. We here report on the remainder of this survey, consisting of 1.4 GHz pulsar-gated observations of 27 more pulsars, using the VLA and ATCA. In §2, we describe our observations and analysis, while in §3 we present non-detections of PWN towards these sources. In §4 we quantify the improvement in sensitivity of the current survey, and discuss the constraints we can put on the radio luminosities of pulsar winds from our data.
## 2 Observations and Reduction
The 27 pulsars observed are listed in Table 1; all were chosen for their high $`\dot{E}`$ and/or space velocity. From this sample 22 pulsars were observed with the VLA \[Napier, Thompson & Ekers 1983\], while the remaining 5 pulsars were observed with the ATCA \[Frater, Brooks & Whiteoak 1992\]. All observations were carried out at frequencies near 1.4 GHz.
The VLA observations were made in the C configuration, using a bandwidth of 25 MHz for $`00^\mathrm{h}<\mathrm{RA}<12^\mathrm{h}`$ and 12.5 MHz otherwise, and with a phase centre corresponding to the catalogued pulsar position. The observing time for each pulsar was typically 15 min. Amplitudes were calibrated using observations of 3C 286 and 3C 48, assuming 1.4 GHz flux densities of 14.9 Jy and 16.3 Jy respectively (where 1 Jy $`=10^{26}`$ W m<sup>-2</sup> Hz<sup>-1</sup>). ATCA observations were made in the 6C configuration, using a bandwidth of 128 MHz (further subdivided into 32 spectral channels); amplitude calibration was carried out using PKS B1934–638 and assuming a 1.4 GHz flux density of 14.9 Jy. For ATCA observations, the phase centre was offset from the pulsar’s catalogued position by $`1^{}`$; each pulsar was observed for approximately 12 hr. Antenna gains and instrumental polarization were calibrated using observations of strong unresolved sources in the vicinity of each pulsar; all four Stokes parameters were recorded.
All observations were gated at the pulsar period in order to look for off-pulse emission at the pulsar position, using ephemerides supplied by A.G. Lyne. For the VLA, gating was carried out by phasing up the array on a nearby calibrator, then integrating on the pulsar for a few minutes. The analogue sum of the signals from all antennas was formed from these data, and then folded at the apparent pulse period to give an un-dedispersed pulse profile. A gate was then set on-line, such that one IF recorded on-pulse data while the other recorded off-pulse data, effectively giving two bins of possibly uneven size. The smearing due to dispersion across the band was sufficiently small for all pulsars that it was always possible to completely separate on- and off-pulse emission when choosing the gate. For the ATCA, visibilities were recorded at high time-resolution (typically 32 bins per period), and then folded at the apparent pulse period before being written to disk. Dedispersion (of 32 channels across the 128 MHz bandwidth) was carried out during data reduction, and appropriate phase bins were then chosen to generate on- and off-pulse images.
Data were edited and calibrated using the MIRIAD and AIPS packages according to standard procedures \[Greisen 1996, Sault & Killeen 1998\]. On- and off-pulse images of a field containing each pulsar were formed using uniform weighting. Each image was deconvolved using either the CLEAN algorithm (for fields containing primarily point sources) or a maximum entropy algorithm (for fields containing significant extended emission), and then smoothed with a Gaussian restoring beam. For some sources the region was significantly confused by extended structure; in these cases, we constrained the deconvolution process by generating CLEAN boxes from lower resolution Molonglo Galactic Plane Survey (MGPS) or Northern VLA Sky Survey (NVSS) data \[Green et al. 1999, Condon et al. 1998\].
For each source, the position of the pulsar was determined by fitting in the $`uv`$ plane to the difference of the on- and off-pulse data. For VLA data, a flux density for each pulsar was determined by measuring a flux in the on-pulse image, then scaling using the width of the gate used. For ATCA data, flux densities were measured directly from ungated data. The sensitivity of each image was determined by measuring the off-pulse RMS at the pulsar position in each case.
## 3 Results
The positions, fluxes and off-pulse sensitivities for the 27 pulsars observed are listed in Table 1, along with the range of spatial scales to which each image was sensitive. In most cases, the pulsar was clearly detected in the on-pulse image, and was completely gated out in the off-pulse image.
As a test of our sensitivity we included PSR B1706–44 in our VLA observations, a pulsar which is known to be embedded in a candidate PWN \[Frail, Goss & Whiteoak 1994\]. This nebula was easily detected in our image, with a surface brightness and spatial extent similar to that obtained in previous data.
All but three of the pulsars in our sample were detected in these observations. For detections the measured interferometric position was compared with the timing position given in the pulsar catalogue \[Taylor, Manchester & Lyne 1993\]. Our gated position for PSR B1706–44 differs significantly from the original timing position \[Johnston et al. 1992\], but is in reasonable agreement with the interferometric measurements of Frail, Goss & Whiteoak \[Frail, Goss & Whiteoak 1994\] and the new timing position of Wang et al. \[Wang et al. 2000\]. Our position for PSR B1754–24 has significantly smaller uncertainties than the catalogued position. The error in the latter is certainly much greater than the beamwidth in the observations of Kijak et al. \[Kijak et al. 1998\], and can account for their 4.9 GHz non-detection of this flat-spectrum pulsar. The majority of the remaining positions are consistent, to within $`3\sigma `$, with positions determined from pulsar timing.
For the majority of sources in our sample, the off-pulse image showed no emission, extended or point-like, at or near the position of the pulsar. Sources of note are discussed individually below:
B0114+58: On-pulse data were corrupted for this source by hardware problems during the observations, and so no detection of the pulsar was made. The off-pulse image was uncorrupted, and shows no emission at the catalogued pulsar position.
B0136+57: An unresolved source of flux density $`2.6\pm 0.3`$ mJy, approximately half of the observed pulsar flux density, is seen at the pulsar’s position in an off-pulse image. This source is $`90`$% linearly polarized, which is similar to the degree of polarisation seen for the pulsar \[Gould & Lyne 1998\]. Thus the source probably corresponds to a component of the pulse-profile which was not properly gated out when the gate was set during observations.
J0631+1036: Images of the region were badly corrupted by sidelobes from the source 4C+10.20 (flux density 2.5 Jy), $`19^{}`$ distant. The pulsar was not detected, its tabulated 1.4 GHz flux density of 0.8 mJy being below the sensitivity of the data. No other emission at or near the pulsar’s position was detected, down to the sensitivity limit.
B0656+14: No off-pulse emission is apparent at the pulsar’s position, but an extended, polarized source is seen $`2^{}`$ south of the pulsar, which Cordova et al. \[Cordova et al. 1989\] argue is possibly associated with the pulsar. We note that we see no connecting structure between the pulsar and this source, despite being more sensitive to extended emission than Cordova et al. \[Cordova et al. 1989\]. It thus seems likely that this source is unrelated to the pulsar. We note that an X-ray PWN associated with this source was claimed by Kawai & Tamura \[Kawai & Tamura 1996\], but has been discredited by higher resolution data \[Becker et al. 1999\].
B1356–60: The region around this pulsar is shown in Fig 1 – the pulsar lies on the western rim of a shell of radio emission. This shell, which we designate G311.28+1.09, is approximately circular, with a diameter $``$9 arcmin and a flux density $`0.04\pm 0.01`$ Jy at 1.4 GHz. In linear polarization there is significant confusion from other sources in the region, and so it is not possible to determine whether the shell is polarized. Based on morphology alone, we thus consider G311.28+1.09 to be a possible new SNR, although further observations will be required to confirm this. However, given the 300-kyr characteristic age of the pulsar, it is unlikely that there is any physical association between PSR B1356–60 and G311.28+1.09. In an ungated image, no PWN is apparent around the pulsar, although the sensitivity to such a source is poor due to the presence of G311.28+1.09.
B1508–57: This pulsar is in a confused region of the Galactic Plane. We were forced to discard short $`uv`$ spacings in order to image the region, limiting the largest spatial scale to which these observations were sensitive to only $`1\stackrel{}{.}8`$. No off-pulse emission was observed, subject to this constraint.
B1634–45: The pulse profile for this source reveals an interpulse separated in phase by 180 degrees from the main pulse, a result which has been confirmed by recent timing observations (F. Crawford, private communication). Gating out both the pulse and interpulse reveals an unresolved off-pulse source at the pulsar’s position of flux density $`0.8\pm 0.4`$ mJy. This source is $``$80% linearly polarized, similar to that measured for the main pulsed component. Since no PWN has ever been observed to be so highly polarized, we think it unlikely that this off-pulse source corresponds to an extended nebula; it is more likely that this emission comes from the pulsar itself. While we cannot rule out an error in the gating hardware or software, we note that no other gated ATCA observations have shown such an effect (see SGJ99). Alternatives are that there is a low-level bridge of pulsed emission connecting the two main components of the pulse profile (cf. PSR B1259–63; Manchester & Johnston 1995), or that the pulse profile contains an underlying unpulsed component (cf. PSR J0218+4232; Navarro et al. 1995). We are planning further ATCA observations of this source in order to distinguish between these possibilities.
B1706–16: An unresolved off-pulse source of flux density $`4.0\pm 0.3`$ mJy is seen at the pulsar’s position. This source is less than 15% linearly or circularly polarized, but so is the pulsar itself \[Gould & Lyne 1998\]. While the pulse profile is quite narrow and shows no evidence for an interpulse, the fact that the VLA gating is set on-line means that, as for PSR B0136+57, we are unable to rule out a component of the pulse profile as the source of this detection. As for PSR B1634–45 above, we plan to re-observe this pulsar with the ATCA in order to clarify this situation.
B1718–35: This pulsar is in a complicated region, and suffers significant confusion from the nearby star-forming region NGC 6334 (e.g. Brooks & Whiteoak 2000). Gating shows the pulsar to be located at the center of a $`4\mathrm{}`$ radio nebula, G351.70+0.66, which is also clearly visible in data from the MGPS \[Green et al. 1999\]. This region has a distinct counterpart at 60 $`\mu `$m in IRAS data, and is probably thermal.
B1727–33: Two ultra-compact H ii regions are in the field, one of which, G354.19–0.06 \[Becker et al. 1994\], has a flux density of 0.3 Jy and is located $`10^{}`$ from the pulsar position. The pulsar was not detected, its catalogued flux of 2.9 mJy corresponding to a signal-to-noise of only $`2\sigma `$.
B1730–37: The sensitivity of the observations was reduced by the presence of PMN J1733–3722 (flux density 0.6 Jy), just $`2^{}`$ away.
B1754–24: The $`1\sigma `$ uncertainty in the right ascension of this pulsar was previously $`14^{}`$; as discussed above, we have now greatly improved on this position. This more precise position puts the pulsar along the same line of sight as the large diffuse H ii region G5.33+0.08 \[Lockman, Pisano & Howard 1996\]. Emission from the latter is clearly seen in NVSS data and in the 90 cm image of Frail, Kassim & Weiler \[Frail, Kassim & Weiler 1994\], but is largely resolved out by our observations.
## 4 Discussion
As discussed in §1, the parameters of the current survey were chosen to improve on the sensitivity of previous surveys, in particular that carried out by FS97. These factors are summarised in Fig 2, where the sensitivity of FS97’s search is compared to that presented here. Only for PWN with radii between $`0\stackrel{}{.}4`$ and $`1^{\prime \prime }`$ is FS97’s search more sensitive than ours. Between $`1^{\prime \prime }`$ and $`10^{\prime \prime }`$, the current results are considerably (up to 100 times) more sensitive, while at scales smaller than $`0\stackrel{}{.}4`$ and larger than $`10^{\prime \prime }`$, our data probe a parameter space to which FS97 were not sensitive at all.
The results reported in §3 showed that most of the sources in our sample had no detectable PWN associated with them. The exceptions were PSRs B0136+57, B1634–45 and B1706–16, for which unresolved off-pulse sources were detected at the pulsar position. From the current data, we are unable to conclusively determine whether this emission corresponds to emission from the pulsar or from a compact PWN. While we plan to investigate these sources further, for the purposes of the present discussion we assume these observations to be non-detections, but with a sensitivity corresponding to the flux density of the off-pulse source (rather than to the noise in the surrounding area of the image).
To quantify the significance of our non-detections, we follow FS97 in characterising a radio PWN’s integrated luminosity, $`L_R`$, by
$$L_R=ϵ\dot{E}\mathrm{erg}\mathrm{s}^1,$$
(1)
where $`\dot{E}`$ erg s<sup>-1</sup> is the associated pulsar’s spin-down luminosity, and $`ϵ`$ is the fraction of $`\dot{E}`$ which goes into radio emission.
Assuming a typical PWN spectral index $`\alpha =0.3`$ and integrating from 10 MHz to 100 GHz, the corresponding 1.4 GHz flux density is
$$S_{1.4}=2.1\times 10^5\frac{ϵ\dot{E}_{34}}{d^2}\mathrm{mJy},$$
(2)
where $`d`$ kpc is the distance to the pulsar and $`\dot{E}=10^{34}\dot{E}_{34}`$ erg s<sup>-1</sup>. We generally use distances from the pulsar catalogue \[Taylor et al. 1995\], derived either from a pulsar’s dispersion measure \[Taylor & Cordes 1993\] or from its kinematic distance based on H i measurements (e.g. Frail & Weisberg 1990).
If we do not detect a PWN in our observations, we can potentially put an upper limit on $`S_{1.4}`$, and hence on $`L_R`$ and $`ϵ`$. However, as demonstrated in Fig 2, these limits depend on the angular size we expect for the PWN. As previously discussed by FS97 and SGJ99, PWN can in general be divided into two distinct classes (assuming that the pulsar is not inside a SNR, which appears to be the case for all the sources in our sample): those which are confined by the gas pressure of the ambient ISM (“static PWN”), and those confined by ram-pressure resulting from motion of the pulsar through the ISM (“bow-shock PWN”).
Let us first consider the case of a static PWN. The bubble in the ISM driven by the pulsar wind will expand supersonically into the ambient medium, producing a PWN of radius \[Arons 1983\]
$$R_{\mathrm{static}}=0.14\times \left(\frac{\dot{E}_{34}t_3^3}{n}\right)^{1/5}\mathrm{pc},$$
(3)
where $`n`$ cm<sup>-3</sup> is the density of the ambient medium (assumed to be pure hydrogen), and $`t_3`$ kyr is the period for which the pulsar has been interacting with the ISM. In further discussion we assume that this age is given by the pulsar’s characteristic age, $`\tau _cP/2\dot{P}`$.
However a PWN is only static while $`\dot{R}_{\mathrm{static}}>V_{\mathrm{PSR}}`$, where $`V_{\mathrm{PSR}}`$ km s<sup>-1</sup> is the pulsar’s space velocity. Re-arranging Equation (3) of FS97 (and correcting for a missing factor of $`4\pi `$ in their results), we find that a PWN will be static for velocities and ambient densities for which:
$$nV_{\mathrm{PSR}}^54\times 10^9\dot{E}_{34}/t_3^2.$$
(4)
When this condition is not met, the pulsar has “overtaken” its own static PWN, and a bow-shock PWN results. The resulting PWN is much smaller than the static PWN, and has a size determined by a balance between the pressure of the relativistic pulsar wind and the ram-pressure resulting from the pulsar’s motion,
$$R_{\text{bow-shock}}=0.63\left(\frac{\dot{E}_{34}}{nV_{\mathrm{PSR}}^2}\right)^{1/2}\mathrm{pc},$$
(5)
where we have assumed that all the spin-down luminosity of the pulsar goes into the wind, and that the radius of a bow-shock PWN is 1.5 times the radius at which the ram and wind pressures balance \[van Buren & McCray 1988\].
Thus for a given $`n`$ and $`V_{\mathrm{PSR}}`$, we can use Equation (4) to determine whether a PWN has a static or a bow-shock morphology, and from this use either Equation (3) or (5) to determine its radius, and hence its angular extent, $`\theta _{\mathrm{PWN}}`$.
For a pulsar-gated observation with a 1.4 GHz RMS sensitivity of $`\sigma `$ mJy beam<sup>-1</sup> at a resolution $`\theta _\alpha ^{\prime \prime }\times \theta _\delta ^{\prime \prime }`$, suppose no off-pulse source is detected at the 5$`\sigma `$ level. There are three possible reasons for this non-detection: the PWN is unresolved and below the point-source sensitivity of the observations, the PWN is resolved and below the surface brightness sensitivity limit, or the PWN is larger than the maximum spatial scale to which we are sensitive. We now consider what limits on $`ϵ`$ can be derived for each situation.
If $`\theta _{\mathrm{PWN}}`$ is smaller than the resolution limit, then $`S_{1.4}<5\sigma `$ and from Equation (2) the corresponding upper limit on $`ϵ`$ is
$$ϵ<ϵ_0=2.4\times 10^5\sigma \left(\dot{E}_{34}/d^2\right)^1.$$
(6)
Note that if an off-pulse point source of 1.4 GHz flux density $`S_{\mathrm{pt}}`$ is detected at the position of the pulsar (as was the case for three of the pulsars in our sample), then $`5\sigma `$ should be replaced by $`S_{\mathrm{pt}}`$ in this expression.
If the PWN is extended, but smaller than the largest spatial scale to which our interferometer is sensitive, then
$$ϵ<\frac{ϵ_0\theta _{\mathrm{PWN}}^2}{\theta _\alpha \theta _\delta },$$
(7)
where we have assumed a Gaussian profile for the PWN. Finally, if the PWN is larger than our observations can detect, we can put no limit on $`ϵ`$.
FS97 have argued, on the basis of their non-detections, that the lack of observable radio PWN around most pulsars implies $`ϵ10^6`$, two orders of magnitude less than for pulsars with detected radio PWN. However, in making this calculation they assumed that the ambient density was $`n=1`$ cm<sup>-3</sup>. Through Equation (4), when combined with a reasonable velocity, this implied that all the pulsars they observed were powering bow-shock PWN, and the consequent high ram pressure ensured that these sources would largely be unresolved by their $`0\stackrel{}{.}8`$ resolution.
However their assumed ambient density is not representative of our current best understanding of the ISM. While the relative filling fractions are uncertain, available evidence supports a multi-phase ISM of which 90% by volume is a combination of a warm medium of density $`n=0.3`$ cm<sup>-3</sup> and a hot ionised component of density $`n=0.003`$ cm<sup>-3</sup> (see Ferrière 1998a for a recent discussion and overview). In the hot low-density component, we still expect most PWN to be bow shocks, but the ram pressure is greatly reduced and PWN will consequently be much more extended. Thus many PWN in this low density medium will be larger than the resolution limit of FS97, and through Equation (7), the limit on $`ϵ`$ much less stringent than claimed.
In order to better consider the limits on $`ϵ`$, we therefore carry out the following calculation for 27 of the pulsars in Table 1 (PSR B1706–44 is excluded as it has a candidate PWN), and additionally for the 4 pulsars discussed by SGJ99. For each pulsar in our sample, we adopt possible densities of $`n=0.3`$ cm<sup>-3</sup> and $`n=0.003`$ cm<sup>-3</sup>. Ten of the pulsars in our sample have measured proper motions, and we use the corresponding 3D space velocities determined by Cordes & Chernoff \[Cordes & Chernoff 1998\]; one other pulsar (PSR B1055–52) has had a scintillation velocity determined for it \[Johnston, Nicastro & Koribalski 1998\]. For the remaining sources we set $`V_{\mathrm{PSR}}=380`$ km s<sup>-1</sup>, corresponding to the mean pulsar velocity of the distribution of Cordes & Chernoff \[Cordes & Chernoff 1998\]. Using Equation (4) we determine whether each corresponding PWN is bow-shock or static, and then consequently determine $`\theta _{\mathrm{PWN}}`$. Upper limits on $`ϵ`$ are then determined from Equations (6) and (7), and the results given in Table 2. For the 16 sources in our sample which were also observed by FS97, we can make a similar calculation based on their results (converting their data to 1.4 GHz assuming a spectral index $`\alpha =0.3`$); these revised values of $`ϵ_{\mathrm{max}}`$ are also listed in Table 2. Note that many PWN are unresolved for either choice of ambient density, and so have the same value of $`ϵ_{\mathrm{max}}`$ in both cases.
We can compare our upper limits to values of $`ϵ`$ for known PWN. In Table 3 of FS97, $`ϵ`$ is listed for the six pulsars then known to have detected radio PWN, to which we add PWN which have since been associated with PSR J0537–6910 ($`ϵ=5\times 10^4`$; Lazendic & Dickel 1998), PSR B0906–49 ($`ϵ=2\times 10^6`$; GSFJ98) and PSR J1811–1926 ($`ϵ<2\times 10^3`$; Morsi & Reich 1987; Torii et al. 1999). For these nine sources, six measurements lie in the reasonably narrow range $`(15)\times 10^4`$. Those sources lying outside this range are PSR B0833–45, for which it is unclear just what part of the surrounding SNR is pulsar-powered, PSR B0906–49, which GSFJ98 argue is substantially different from other PWN, and PSR J1811–1926, where the radio PWN is faint and has poorly constrained properties. Thus the data available suggest that a “typical” detectable PWN has $`ϵ10^4`$.
FS97 argued that typical non-detections corresponded to $`ϵ_{\mathrm{max}}2\times 10^6`$, significantly less than for detected PWN. However, it can be seen from the results in Table 2 that for $`n=0.003`$ cm<sup>-3</sup> some PWN become too large to be detected by their data, while for the remaining pulsars the limit rises to $`ϵ_{\mathrm{max}}3\times 10^4`$. Thus we argue that the observations of FS97 could have missed “normal” PWN around most of their sample if most of these sources are in low density regions, and hence their observations do not constrain pulsars which lack detectable PWN to be any different in their wind properties from those pulsars with observed PWN.
On the other hand, the current observations can potentially provide constraining limits on $`ϵ`$. We first consider the six young and energetic pulsars in our sample (i.e. the first six pulsars in Table 2). These pulsars are defined approximately by $`\dot{E}_{34}50`$ and $`t_350`$, properties similar to those pulsars with detectable radio PWN. For either assumed ISM density, our data constrain these pulsars to have upper limits on $`ϵ`$ in the range $`(0.0020.2)\times 10^4`$, significantly less than for pulsars with observed PWN. While this seems to imply genuinely low values of $`ϵ`$, for $`n=0.003`$ cm<sup>-3</sup>, conditions are such that these pulsars have only recently overtaken their static nebulae. Lowering their velocities slightly, arguing that their actual ages are less than their characteristic ages, or accepting that realistically, the transition from static to bow-shock PWN does not happen instantaneously, it seems likely that these pulsars are still producing static nebulae, whose extents are much larger than for bow-shocks. In this case, the corresponding limits become $`ϵ_{\mathrm{max}}10^4`$, and are not constraining.
We thus argue that our non-detections of PWN can be explained even if all young and energetic pulsars have similar wind properties. The difference between detectable and non-detectable PWN seems to be that detected PWN are either in dense regions of the ISM or in SNRs, in which there is sufficient external pressure to confine the pulsar wind and produce an observable PWN. However, pulsars with no PWN are in the low density phase of the ISM and so produce unobservable “ghost remnants” \[Blandford et al. 1973\]. With the exception of PSR B0906–49 (GSFJ98), pulsars with observed PWN also have associated SNRs, while all those young pulsars without PWN also have no associated SNRs. We indeed expect SNRs to be undetectable in the hot component of the ISM \[Kafatos et al. 1980, Gaensler & Johnston 1995\], consistent with our conclusion above that it is a low ambient density which causes a PWN around a young pulsar to be undetectable.
A notable exception is PSR B1757–24, which is associated with both a SNR and a radio PWN despite having $`n=0.003`$ cm<sup>-3</sup> \[Frail & Kulkarni 1991, Manchester et al. 1991\]. In this case, the pulsar is inferred to have a transverse space velocity $``$1500 km s<sup>-1</sup> \[Frail, Kassim & Weiler 1994\]; this not only supplies the necessary ram pressure to produce an observable PWN, but has caused the pulsar to overtake the shell of the associated SNR, re-energising the remnant with its passage. If it were not for the extreme velocity of the pulsar, neither the PWN nor the SNR would be detectable, as expected in a low density region.
The relative numbers of young pulsars with detected and undetected SNRs/PWN suggests an approximate filling fraction $``$50% for the low density component of the ISM. This is somewhat more than recent estimates of 15–20% \[Ferrière 1998b\], but can be explained by the fact that we expect young pulsars to be preferentially located in low density regions produced by the powerful winds of their progenitors and by previous supernovae in the region.
The majority of pulsars in our sample are considerably less energetic ($`\dot{E}_{34}<50`$) and older ($`t_3100`$) than pulsars around which radio PWN have been observed. Values of $`ϵ_{\mathrm{max}}`$ for these pulsars are plotted in Fig 3, from which it can be seen that for either choice of ambient density, the distribution peaks around $`ϵ_{\mathrm{max}}10^5`$. These pulsars have all long since overtaken their static PWN, and will have bow shock PWN for any sensible choice of $`n`$ and $`V_{\mathrm{PSR}}`$. The resulting size of such a PWN is not a strong function of $`n`$ or $`V_{\mathrm{PSR}}`$; to produce values of $`ϵ_{\mathrm{max}}>10^4`$, consistent with detected PWN, requires uniformly low space velocities, $`V_{\mathrm{PSR}}150`$ km s<sup>-1</sup>, for these pulsars. However, only $``$5% of pulsars are thought to be traveling at less than 150 km s<sup>-1</sup> \[Cordes & Chernoff 1998\]. While the ages and distances we have used for these pulsars have their associated uncertainties, for bow-shock PWN values of $`ϵ_{\mathrm{max}}`$ are independent of age, and distances would have to uniformly be increased by a factor of three to shift the peak in $`ϵ_{\mathrm{max}}`$ to a value in agreement with that seen for detected PWN. The lower value of $`ϵ`$ we have derived for older pulsars is thus a result quite robust to the assumptions and uncertainties involved in its derivation, and we therefore argue that these pulsars have winds which are genuinely at least an order of magnitude less efficient at producing radio emission than the winds of young and energetic pulsars.
FS97 consider various reasons why older pulsars might appear to have lower values of $`ϵ`$. Possibilities include:
1. that the PWN are resolved out by the observations;
2. that an increasing fraction of $`\dot{E}`$ goes into pulsed X-rays and $`\gamma `$-rays (see Thompson et al. 1994);
3. that their winds are dominated by Poynting flux rather than relativistic particles;
4. that the injection spectrum of particles in the pulsar wind has shifted to higher energies.
Our observations can conclusively rule out alternative (i), as we can detect PWN produced for almost all feasible values of $`n`$ and $`V_{\mathrm{PSR}}`$. While we cannot distinguish between the remaining three possibilities, we note that of detected PWN, that associated with the oldest pulsar, PSR B0906–49, also has the lowest value of $`ϵ`$ and the steepest spectral index (GSFJ98). Since the spectral index of a PWN is directly related to the spectrum of injected particles \[Pacini & Salvati 1973\], this result tentatively suggests that the efficiency of the wind in producing radio emission is related to the injection spectrum, and that alternative (iv) might then best explain the observations.
## Conclusions
We have used pulsar-gating at 1.4 GHz to search for radio PWN around 27 pulsars. Our search was up to 100 times more sensitive than the only other comparable survey, and was carried out on spatial scales corresponding to a much wider range of ambient densities and pulsar velocities. Including data from previous work by SGJ99, non-detections towards 28 pulsars, plus inconclusive results in three other cases, have allowed us to determine upper limits on $`ϵ`$, the fraction of a pulsar’s spin-down luminosity which goes into producing radio emission from a PWN.
We find that the data are consistent with virtually all young energetic pulsars having $`ϵ10^4`$. The lack of PWN around $``$50% of young pulsars can be explained if they are in low ambient densities (0.003 cm<sup>-3</sup>), consistent with the absence of associated supernova remnants around these sources.
For older pulsars, any reasonable choice of ambient density and pulsar velocity results in upper limits on the wind efficiency $`ϵ<10^5`$, ten times less than for young pulsars. Thus pulsars seem to become less efficient at producing radio wind nebulae as they age; we speculate that this result is due to the spectrum of injected relativistic particles steepening in older pulsars. This possibility can be tested through X-ray observations towards such pulsars – it is likely that Chandra will make many new detections of X-ray PWN, and through consequent imaging spectroscopy, we may finally be able to probe the winds around older pulsars.
Of the $``$55 non-recycled pulsars with $`\dot{E}>3\times 10^{34}`$ erg s<sup>-1</sup>, almost all have now been searched for associated radio PWN down to a good sensitivity. If those sources with no detectable PWN are indeed in low density regions of the ISM, it seems unlikely that we will ever find radio PWN around them with current telescopes. For example, if PSR B1046–58 is powering a static PWN with $`ϵ=10^4`$, the resulting radio nebula would be $`20^{}`$ across with a flux density at 1.4 GHz of 0.5 Jy. To detect this source would require $`\sigma =0.3`$ mJy arcmin<sup>-2</sup>, which is generally below the confusion limit for instruments capable of imaging sources this large. We might have to wait for the large increase in sensitivity promised by the Square Kilometre Array in order to make further progress.
## Acknowledgements
We are particularly grateful to Walter Brisken for his help with pulsar-gating at the VLA, Barry Clark and Miller Goss for their generous re-scheduling of a failed observing run, and Andrew Lyne for supplying timing solutions for many of the pulsars we observed. We also thank Phillip Hicks for his assistance with the observations, Froney Crawford for information on the pulse profile of PSR B1634–45, and Jon Arons, Jim Cordes, Vicky Kaspi, Michael Pivovaroff and Eric van der Swaluw for useful discussions and suggestions. This research has made use of the data base of published pulse profiles maintained by the European Pulsar Network, available at http://www.mpifr-bonn.mpg.de/pulsar/data, NASA’s Astrophysics Data System Abstract Service and of the SIMBAD database, operated at CDS, Strasbourg, France. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. The Australia Telescope is funded by the Commonwealth of Australia for operation as a National Facility managed by CSIRO. B.M.G. acknowledges the support of NASA through Hubble Fellowship grant HF-01107.01-98A awarded by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS 5–26555. B.W.S. is supported by NWO Spinoza grant 08-0 to E.P.J. van den Heuvel.
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# How behavior of systems with sparse spectrum can be predicted on a quantum computer
## 1 Introduction and background
A behavior of typical quantum system cannot be analyzed at hand or by classical computer because huge dimension of its Hamiltonian makes it impossible to solve Shröedinger equation even numerically. Nevertheless one can expect that this analysis would be easier in the framework of quantum computing. The idea of this approach is to force one quantum system to simulate a behavior of other more simple primary system with some profit in time. In the particular case where the primary system is a classical computer with oracle this is a problem of quantum speedup for classical computations: given a classical algorithm with oracle, is there a quantum algorithm computing the same function faster using the same oracle quantumly? Examples: Shor factoring algorithm (look at \[Sh\]), Grover search algorithm (look at \[Gr\]).
In general case the problem of prediction has the following form: given a low of evolution for some system (classical or quantum) is there a device (wizard) of not exponential size which can predict the behavior of this system? To predict the behavior means that given a time instant, wizard returns a state of initial system in this instant earlier than this state appears in its natural evolution. Note that in general case of quantum evolution a wizard may use quantumly not only Hamiltonian (like in case of quantum computations with oracle) but some hidden information about primary system. What kind of such information may be useful for prediction? This is the information about eigenvalues of primary Hamiltonian. This article analyses the possibility of quantum predictions in terms of spectrum of primary system.
Quantum method of finding eigenvalues was first presented by Abrams and Lloyd in their work \[AL\]. Their method, generalizing a quantum part of Shor factoring algorithm, uses Quantum Fourier transform (QFT) and requires of order $`N`$ implementations of initial Hamiltonian where $`N`$ is the dimension of the main space. Thus, their algorithm requires exponential time and it cannot be used for quantum speedup. However, for the systems whose spectrum is sparse the idea of this approach with QFT may be used for predictions. I proceed with the exact definitions.
## 2 Definitions
Let $`H`$ be some Hamiltonian which induces unitary transformations $`U`$ in some $`N`$ dimensional space $`Y`$ of states of $`n`$ qubits by the usual rule in quantum mechanics: $`U(\tau )=\mathrm{exp}(iH\tau /h)`$, $`N=2^n`$. In order to apply notions from the algorithm theory, such as complexity, etc., we assume that the time of Hamiltonian action is always $`\tau =1`$. We may of course choose another value for $`\tau `$ but this would only involve the change in time scale without any subsequences for the algorithm complexity. Call a system primary if its evolution is determined by unitary transformation $`U`$.
We assume that the complexity of computation is the number of $`U`$ applications. Time of all other transformations we use here is assumed to be negligible comparatively with the time of computation. Quantum Fourier transform (QFT) requires the time of order $`n^2`$ (look at \[Sh\]). Here we shall use this transform for $`\mathrm{log}_2(M)`$ gubits in the following form:
$$\mathrm{QFT}_M:|s\frac{1}{\sqrt{M}}\underset{l=0}{\overset{M1}{}}\mathrm{exp}(2\pi isl/M)|l$$
Let eigenvalues of $`H`$ have the form $`2\pi \omega _kh,k=0,1,\mathrm{},N1`$, when frequencies $`\omega _k`$ are real numbers from the segment $`[0,1)`$. Then the eigenvalues of $`U`$ will be $`\mathrm{exp}(2\pi i\omega _k)`$. This is not loss in generality because we always can choose other unit $`\tau `$ for Hamiltonian action time. Let us give the precise definitions of the ”knowledge of eigenvalues” of $`H`$. It means that there exists another Hamiltonian $`W_H`$ called a wizard for $`H`$ which acts in $`N^2`$ dimensional Hilbert space and returns an approximation of frequencies $`\omega _k`$ in within $`1/N`$ given eigenvector $`\mathrm{\Phi }_k`$ of $`H`$. Introduce the following notations. Every frequency $`\omega _k`$ has binary notation $`0.ϵ_1ϵ_2\mathrm{}`$. Let $`pn`$ be some integer, $`M=2^p`$. Denote the string of ones and zeroes $`ϵ_1ϵ_2\mathrm{}e_p`$ by $`\overline{ϵ}_k^p`$. $`ϵ_k^p`$ may be considered as integer if we cancel all first zeroes. And vise versa, every integer $`l`$ less than $`M`$ can be written in form $`\overline{ϵ}_k^p`$ if we add suitable number of zeroes in front. Then the number $`\stackrel{~}{\omega }_k^M=0.\overline{ϵ}_k^p`$ will be approximation of $`w_k`$ in within $`1/M`$. We shall define this number $`0.\overline{ϵ}_k^p`$ from $`[0,1)`$ by $`(0.l)_p`$. Wizard action in these notations will be
$$W_U|\mathrm{\Phi }_k,b|\mathrm{\Phi }_k,b\overline{ϵ}_k^p$$
where $``$ is the biwise addition modulo $`2`$.
At last introduce auxiliary transformations $`\mathrm{WH}`$ and $`\mathrm{U}_{seq}^p`$. Let the memory is divided into the main part $`x`$ of $`n`$ qubits and ancilla $`a`$ of $`pn`$ qubits: $`|x,a`$, $`M=2^p`$. Put:
1). $`\mathrm{WH}|x,a=\frac{1}{\sqrt{M}}\underset{s=0}{\overset{M1}{}}(1)^{as}|x,s`$. This is Walsh - Hadamard transform applied to ancilla where $``$ denotes dot product modulo 2.
2). $`\mathrm{U}_{seq}^p|x,a=|\mathrm{U}^ax,a`$. This is the result of $`a`$ sequential applications of $`U`$ to the main register.
Walsh-Hadamard transform can be fulfilled in a standard model of quantum computer. To fulfill $`\mathrm{U}_{seq}^p`$ it would suffice to use the following oracle $`\mathrm{U}_{cond}`$ (conditional application of $`\mathrm{U}`$) depending on $`\mathrm{U}`$: $`\mathrm{U}_{cond}|x,\alpha |x^{},\alpha `$, where $`x^{}=x`$ if $`\alpha =0`$ and $`x^{}=\mathrm{U}x`$ if $`\alpha =1`$, $`\alpha `$ is one qubit register. An application of $`\mathrm{U}_{cond}`$ cannot be reduced to the simple using of $`\mathrm{U}`$ as an oracle because a conditional application of $`\mathrm{U}`$ is quantumly controlled by the second register. One proposal about its practical implementation can be found in the section 4.5.
## 3 Wizard transformation
### 3.1 How wizard predicts evolution
Assume temporarily that we have exact equations $`\omega _k=\stackrel{~}{\omega }_k^N`$. Given a wizard transformation $`W_U`$ how can we predict evolution of initial system? Let $`|\xi `$ denote initial state as a contents of $`n`$ qubits main register. Let
$$\xi =\underset{k=0}{\overset{N1}{}}x_k|\mathrm{\Phi }_k$$
(1)
be the expansion of our state in basis of eigenvectors of $`U`$. Add $`n`$ qubits ancillary register initialized by zeroes and obtain the state $`|\xi ,0^n`$.
Now apply the wizard transformation $`W_U`$ to the main register ($`p=n`$). It gives the state
$$\xi ^{}=\underset{k=0}{\overset{N1}{}}x_k|\mathrm{\Phi }_k,\overline{ϵ}_k^p.$$
Given the number $`t`$ we can turn each state in the first register by angle determined by the second register and $`t`$: $`2\pi \omega _kt`$. This gives the state
$$\xi ^{\prime \prime }=\underset{k=0}{\overset{N1}{}}x_k\mathrm{exp}(2\pi i\omega _kt)|\mathrm{\Phi }_k,\overline{ϵ}_k^p.$$
At last apply wizard again obtaining
$$\xi ^{\prime \prime \prime }=\underset{k=0}{\overset{N1}{}}x_k\mathrm{exp}(2\pi i\omega _kt)|\mathrm{\Phi }_k,0^n,$$
and now cancellation ancilla gives exactly the state $`U^t|\xi `$ which is the state in time instant $`t`$ of initial system. Thus given a wizard we can predict the behavior of initial system provided the action of wizard takes smaller time than $`t/2`$.
### 3.2 Simulation a wizard
In this subsection I describe the procedure of simulation a wizard action on zero ancilla. The particular case $`p=n`$ of this procedure was proposed in other purpose in the paper \[AL\].
We start from the state of the form (1) with ancilla attached initialized by zeroes. Wizard action on the initial state with zero ancilla is defined by the following:
$$\mathrm{QFT}_M\mathrm{U}_{seq}^p\mathrm{WH}|\xi ,0^p.$$
(2)
Why this procedure must work? Let us again assume temporarily that $`\omega _k=\stackrel{~}{\omega }_k^M`$ exactly. Then we have $`\mathrm{WH}|\xi ,0^p=\chi _0=\frac{1}{\sqrt{M}}\underset{k=0}{\overset{N1}{}}\underset{s=0}{\overset{M1}{}}x_k|\mathrm{\Phi }_k,s`$. Application of $`\mathrm{U}_{sec}^p`$ gives: $`\chi _1=\mathrm{U}_{sec}^p\chi _0=\frac{1}{\sqrt{M}}\underset{k=0}{\overset{N1}{}}\underset{s=0}{\overset{M1}{}}x_k\mathrm{exp}(2\pi i\omega _ks)|\mathrm{\Phi }_k,s`$. Now $`\mathrm{QFT}_M\chi _1=\underset{k=0}{\overset{N1}{}}x_k|\mathrm{\Phi }_k,\overline{ϵ}_k^p`$. This is just the wizard action. But in general case the result of QFT transform has not so simple form. It means that we cannot merely clean ancilla by repetition of wizard action in predicting algorithm even if we have a wizard action on all words. Hence the precision of a wizard simulation by (2) must be elaborated in more details.
### 3.3 Accuracy of a wizard simulation
Let $`\{\stackrel{~}{\omega }_{k,i}\}`$ be some set of integers. Denote $`L_\epsilon (\stackrel{~}{\omega }_{k,i})=\{i:|(0.\stackrel{~}{\omega }_{k,i})_p\omega _k|\epsilon `$ or
$`|(0.\stackrel{~}{\omega }_{k,i})_p\omega _k1|\epsilon \}`$. Let $`\xi `$ be a state of the form (1).
Definition A transformation $`W`$ of the form
$$W:|\xi ,0^p\underset{k=0}{\overset{N1}{}}\underset{i=0}{\overset{M1}{}}\lambda _{i,k}|\mathrm{\Phi }_k,\stackrel{~}{\omega }_{k,i}$$
is called a transformation of $`W_{\delta ,\epsilon }`$ type if $`\underset{k=0}{\overset{N1}{}}\underset{iL_\epsilon (\stackrel{~}{\omega }_{k,i})}{}|\lambda _{i,k}|^21\delta `$ for any $`\xi `$. Thus, $`\delta `$ is an error probability from the quantum superposition, and $`\epsilon `$ is a precision of eigenvalues approximations.
###### Lemma 1
The transformation $`\mathrm{QFT}_M\mathrm{U}_{seq}^p\mathrm{WH}`$ belongs to the type $`W_{\frac{1}{K},\frac{K}{M}}`$ for any $`K\{1,2,\mathrm{},M\}`$.
Proof
Denote $`\mathrm{QFT}_M\mathrm{U}_{seq}^p\mathrm{WH}|\xi ,0^p`$ by $`\chi _{p,2}`$. Then we have
$$\chi _{p,2}=\frac{1}{M}\underset{k=0}{\overset{N1}{}}\underset{l=0}{\overset{M1}{}}x_kH_{l,k}|\mathrm{\Phi }_k,l$$
(3)
where $`H_{l,k}=\underset{s=0}{\overset{M1}{}}\mathrm{exp}(2\pi is(\omega _k(0.l)_p))`$. Put $`\mathrm{\Delta }=\omega _k(0.l)_p`$. Then by summing the progression we obtain
$$H_{l,k}=\frac{1\mathrm{exp}(2\pi iM\mathrm{\Delta })}{1\mathrm{exp}(2\pi i\mathrm{\Delta })}.$$
(4)
Let $`\mathrm{\Delta }=\frac{K}{M}+\delta ,K=K(l)`$, where $`0\delta <1/M,K`$ integer. Then each $`H_{l,k}`$ depends on $`K:H_{l,k}=H_{l,k}(K)`$. Fix some value for $`K`$: $`K_0`$. It means that the accuracy of eigenvalues approximations will be $`K_0/M`$. Estimate the sum of squared amplitudes for states $`|\mathrm{\Phi }_k,l`$ for which $`K:K_0KMK_0`$. Using that module of denominator in (4) is separated from zero by $`\mathrm{min}\{\frac{\pi K}{M},\frac{\pi (MK)}{M}\}`$, we have:
$$\begin{array}{cc}& \underset{k=0}{\overset{N1}{}}\underset{K=K_0}{\overset{MK_0}{}}|x_k|^2\frac{1}{M^2}|H_{l,k}(K)|^2\\ & 4\underset{k=0}{\overset{N1}{}}\underset{K=K_0}{\overset{MK_0}{}}|x_k|^2/\pi ^2K^2\frac{4}{\pi ^2}\underset{K=K_0}{\overset{MK_0}{}}\frac{1}{K^2}\frac{1}{K_0}.\end{array}$$
This is exactly what is needed. Lemma 1 is proved.
What would happen if we apply other sequence of transformations $`\mathrm{QFT}_M^1(\mathrm{U}_{seq}^p)^1\mathrm{WH}`$ instead of $`\mathrm{QFT}_M\mathrm{U}_{seq}^p\mathrm{WH}`$ for revealing frequencies?
###### Lemma 2
Let for some state $`\chi `$
$$\begin{array}{ccc}\chi _2& =\mathrm{QFT}_M\mathrm{U}_{seq}^p\mathrm{WH}|\chi ,0^p& =\underset{k}{}\underset{l=0}{\overset{M1}{}}y_{k,l}|\mathrm{\Phi }_k,l\\ \chi _2^{}& =\mathrm{QFT}_M^1(\mathrm{U}_{seq}^p)^1\mathrm{WH}|\chi ,0^p& =\underset{k}{}\underset{l=0}{\overset{M1}{}}y_{k,l}^{}|\mathrm{\Phi }_k,l\end{array}$$
Put $`\delta _{k,l}=\stackrel{~}{\omega }_k^N(0.l)_p`$, $`\mathrm{\Delta }=\omega _k(0.l)_p`$, $`\chi _2^{\prime \prime }=\underset{k}{}\underset{l=0}{\overset{M1}{}}y_{k,l}^{}\mathrm{exp}(2\pi i(M1)\delta _{k,l})|\mathrm{\Phi }_k,l`$. Let $`M/N<ϵ`$. Then
$$\chi _2^{\prime \prime }\chi _2<7ϵ$$
Proof
We have $`y_{k,l}=x_kH_{k,l}`$ (look at (3),(4)), where $`H_{k,l}=H_{k,l}(\mathrm{\Delta })`$. Then $`y_{k,l}^{}=x_kH_{k,l}(\mathrm{\Delta })`$. Now $`H_{k,l}(\mathrm{\Delta })=H_{k,l}(\mathrm{\Delta })\mathrm{exp}(2\pi i(M1)\mathrm{\Delta })`$, $`|\delta _{k,l}\mathrm{\Delta }|1/N`$, $`2\pi (M1)|\delta _{k,l}\mathrm{\Delta }|7M/N`$ and Lemma 2 follows.
### 3.4 Complexity of a wizard
It is readily seen that the above procedure for $`p=n`$ requires $`N`$ applications of the initial transformation $`U`$, hence it cannot be used on purpose to predict its evolution. It turns that in general case the following theorem takes place.
###### Theorem 1
It is impossible to simulate a wizard action with precision up to $`O(1/N)`$ using less than $`\mathrm{\Omega }(N)`$ conditional application of initial transformation $`\mathrm{U}_{cond}`$.
Proof
A lower bound for the time of quantum computation of the PARITY function is $`N/2`$ by the results \[FGGS\] and \[BBCMW\]. These constructions can be extended to the case when we can use more general oracle $`\mathrm{U}_{cond}`$ instead of $`\mathrm{U}`$ when evaluating PARITY. Assuming that a wizard action can be simulated using less than $`\mathrm{\Omega }(N)`$ application of $`\mathrm{U}_{cond}`$, in view of the previous subsection we would be able to compute the function PARITY on a quantum computer in time less than $`N/2`$ (adding appropriate constant to the number $`n`$ in order to enhance accuracy) which leads to the contradiction. Theorem 1 is proved.
Call a unitary transformation $`U`$ classical if it maps basic states to basic states. One can ask is there any way to predict a behavior of classical system on a quantum computer? This statement is independent of a form of spectrum. It turns that in general case a behavior of the bulk of classical systems in a short time frame is impossible on quantum computer even on one step ahead (look at \[Oz\]). Here short time means approximately $`O(N^{1/7})`$, where $`N`$ \- the number of all states.
But for one class of systems the quantum prediction is possible. This is the class of systems with sparse spectrum.
## 4 Prediction of the evolution of the systems with sparse spectrum
### 4.1 Case of precise eigenvalues
What would happen if we use less number of qubits $`p<n`$ instead of $`n`$ in the ancilla? Simulation of the wizard will be shorter, but it will be simulated with the corresponding loss in precision and we obtain no prediction.
Nevertheless, if the spectrum is sparse then prediction is possible. Indeed, suppose that we have classical algorithm $`h`$ enhancing the accuracy of eigenvalues approximation.
Namely, let $`p<n`$, $`M=2^p`$ and let $`h:\{0,1\}^p\{0,1\}^n`$ be integer function mapping rough frequency approximation up to $`1/M`$ to the more precise approximation up to $`1/N`$.
Given the initial state $`\xi `$ and time instant $`t`$ we shall now describe the procedure of quick finding $`U^t\xi `$. The idea is simple: Repeat the procedure from above with only $`p`$ ancillary qubits instead of $`n`$, and enhance accuracy by $`h`$.
At first we assume for the simplicity that we have exact equalities $`\omega _k=\stackrel{~}{\omega }_k^M`$. Here is the detailed description of the algorithm.
First we apply Walsh-Hadamard transform to the ancilla and obtain
$$\chi _0=\frac{1}{\sqrt{M}}\underset{k=0}{\overset{N1}{}}\underset{s=0}{\overset{M1}{}}x_k|\mathrm{\Phi }_k,s$$
Then we apply $`\mathrm{U}_{seq}^p`$. Since $`sp`$ this operation requires $`p`$ conditional applications of $`U`$. This gives the state
$$\chi _1=\frac{1}{\sqrt{M}}\underset{k=0}{\overset{N1}{}}\underset{s=0}{\overset{M1}{}}x_kU^s|\mathrm{\Phi }_k,s=\frac{1}{\sqrt{M}}\underset{k=0}{\overset{N1}{}}\underset{s=0}{\overset{M1}{}}\mathrm{exp}(2\pi i\omega _ks)x_k|\mathrm{\Phi }_k,s.$$
Now application of $`\mathrm{QFT}_M`$ to the second register yields:
$$\chi _2=\underset{k=0}{\overset{N1}{}}x_k|\varphi _k,\stackrel{~}{\omega }_k$$
where $`0.\stackrel{~}{w}_k`$ is approximation of eigenvalue $`\omega _k`$ in within $`1/M`$.
This is the point when we will use that the spectrum is sparse. Add one more register with $`n`$ qubits initialized by zeroes: $`\underset{k=0}{\overset{N1}{}}x_k|\mathrm{\Phi }_k,\stackrel{~}{\omega }_k,0^n`$. Apply the unitary version of the algorithm $`h`$ enhancing accuracy to two ancillary registers: $`|a,b|a,bh(a)`$. This gives the state $`\underset{k=0}{\overset{N1}{}}x_k|\mathrm{\Phi }_k,\stackrel{~}{\omega }_k,h(\stackrel{~}{\omega }_k)`$. Now we turn every state of the form $`|\mathrm{\Phi }_k,\stackrel{~}{\omega }_k,l`$ to the angle $`2\pi 0.lt`$ and obtain the state $`\underset{k=0}{\overset{N1}{}}x_k\mathrm{exp}(2\pi i0.lt)|\mathrm{\Phi }_k,\stackrel{~}{\omega }_k,l`$. Here $`t`$ is a given time instant. Now apply unitary version of $`h`$ which cleans up $`n`$ last ancillary qubits and then discard them:
$$\underset{k=0}{\overset{N1}{}}x_k\mathrm{exp}(2\pi \omega _kt)|\mathrm{\Phi }_k,\stackrel{~}{\omega }_k,$$
and again $`\mathrm{QFT}_M`$ to the $`p`$ ancillary qubits:
$$\frac{1}{\sqrt{M}}\underset{k=0}{\overset{N1}{}}\underset{s=0}{\overset{M1}{}}x_k\mathrm{exp}(2\pi i\omega _kt)\mathrm{exp}(2\pi i\stackrel{~}{\omega }_ks)|\mathrm{\Phi }_k,s,$$
then again $`\mathrm{U}_{seq}^p`$ and Walsh-Hadamard transform to the ancilla. The result will be:
$`\underset{k=0}{\overset{N1}{}}x_k\mathrm{exp}(2\pi i\omega _kt)|\mathrm{\Phi }_k,0^p`$. The wanted state is in the main register now.
Note that all the work with enhancing an accuracy here is not mandatory because $`\omega _k=\stackrel{~}{\omega }_k^M`$ exactly. Why this scheme doesn’t work for the case when $`\stackrel{~}{\omega }_k^M`$ are only approximations of the true eigenvalues $`\omega _k`$ ? The point is that $`\chi _2`$ will not have so simple form and in addition multiplication frequencies by $`t`$ will cause big error for $`t=O(N)`$. Thus for the general case more refined algorithm is needed.
### 4.2 General case
We have numbers $`n`$ and $`p`$. Choose some $`q:pq<n`$ and put $`L=2^q`$. Without loss in generality we can extend $`h`$ to the mapping $`h:\{0,1\}^q\{0,1\}^n`$ so that if $`\stackrel{~}{\omega }_k^L=(0.x)_q`$ then $`\stackrel{~}{\omega }_k^N=(0.h(x))_n`$.
Our initial state has the form
$$\xi =\underset{k=0}{\overset{N1}{}}x_k|\mathrm{\Phi }_k,0^q$$
1. Apply $`\mathrm{QFT}_L\mathrm{U}_{seq}^q\mathrm{WH}`$ to the initial state to get $`\chi _{q,2}`$.
2. Add one more register with $`n`$ qubits initialized by zeroes:
$`\underset{k}{}\underset{l=0}{\overset{L1}{}}y_{k,l}|\mathrm{\Phi }_k,l,0^p`$ and apply the unitary version of the algorithm $`h`$ enhancing accuracy to two ancillary registers: $`\chi _3=\underset{k}{}\underset{l=0}{\overset{L1}{}}y_{k,l}|\mathrm{\Phi }_k,l,h(l)`$.
3. Turn every state of the form $`|\mathrm{\Phi }_k,l,h(l)`$ to the angle $`2\pi 0.h(l)t`$ and obtain the state $`\chi _4=\underset{k}{}\underset{l=0}{\overset{L1}{}}y_{k,l}\mathrm{exp}(2\pi i0.h(l)t)|\mathrm{\Phi }_k,l,h(l)`$. Here $`t`$ is given time instant.
4. Put $`\delta _{k,l}^{}=(0.h(l))_n(0.l)_q`$. Turn every state of the form $`|\mathrm{\Phi }_k,l,h(l)`$ to the angle $`2\pi (L1)\delta _{k,l}^{}`$. Denote the result by $`\chi _5`$.
5. Now apply unitary version of $`h`$ which cleans up $`n`$ last ancillary qubits and then discard them:
$`\chi _6=\underset{k}{}\underset{l=0}{\overset{L1}{}}y_{k,l}\mathrm{exp}(2\pi i0.h(l)t)\mathrm{exp}(2\pi i(L1)\delta _{k,l}^{})|\mathrm{\Phi }_k,l`$.
6. Apply to $`\chi _6`$ the transformation $`\mathrm{WH}\mathrm{U}_{sec}^q\mathrm{QFT}_L`$, then observe and discard ancilla.
The time of this algorithm is $`2L+2t_{\mathrm{QFT}_L}+2t_h+w`$ where $`t_{\mathrm{QFT}_L}`$ is the time of Fourier transform, $`t_h`$ is the time of enhancing accuracy and $`w`$ is the time of rotations.
###### Lemma 3
Denote the result of algorithm 1-6 by $`\mathrm{U}_{q,t}`$. For any $`\delta >0`$ for any $`p,q,n,t`$ altering so that $`q=p+c,L=2^q`$, $`t:0<tCN`$, where $`c=\mathrm{log}_2(\frac{4}{\delta })`$, $`pc`$, $`C=\frac{\delta }{14}`$ we shall have
$$\mathrm{U}_{q,t}\mathrm{U}^t<\delta $$
Proof
Introduce the simplifying notations:
$`\chi =\underset{k}{}\underset{l=0}{\overset{L1}{}}y_{k,l}\mathrm{exp}(2\pi i\omega _kt)|\mathrm{\Phi }_k,l,h(l)`$, $`\stackrel{~}{\chi }=\mathrm{QFT}_L^1(\mathrm{U}_{seq}^q)^1\mathrm{WH}\mathrm{exp}(2\pi i\omega _kt)\underset{k}{}x_k|\mathrm{\Phi }_k,0^q`$,
$`\chi ^{}=\underset{k}{}\underset{l=0}{\overset{L1}{}}y_{k,l}\mathrm{exp}(2\pi i\omega _kt)\mathrm{exp}(2\pi i(L1)\delta _{k,l}^{})|\mathrm{\Phi }_k,l`$.
By Lemma 1 $`\chi _4\chi <\frac{\delta }{2}`$. The passages $`\chi _4\chi _6`$ and $`\chi \chi ^{}`$ are fulfilled by the same unitary transformation, which preserves lengths. Consequently, $`\chi _6\chi ^{}<\frac{\delta }{2}`$. By Lemma 2 $`\stackrel{~}{\chi }\chi ^{}<\frac{\delta }{2}`$. Then triangle inequality gives $`\chi _6\stackrel{~}{\chi }<\delta `$. The passages $`\chi _6U_{q,t}`$ and $`\stackrel{~}{\chi }U^t`$ are fulfilled by the same unitary transformation $`\mathrm{WH}\mathrm{U}_{sec}^q\mathrm{QFT}_L`$. Hence $`U^tU_{q,t}<\delta `$. Lemma 3 is proved.
By Lemma 3 this algorithm gives the prediction of state in time instant $`O(N)`$ in time $`O(M)`$ if classical algorithm enhancing accuracy obtain eigenvalues in negligible time. Thus if it is possible to enhance accuracy of eigenvalues from $`ϵ_1`$ to $`ϵ`$ the speedup will be $`\frac{ϵ_1(1\rho )}{56ϵ}`$, where $`\rho `$ is a fidelity, $`0<\rho <1`$. The result can be formulated as
###### Theorem 2
Given a Hamiltonian of system with sparse spectrum and algorithm enhancing the accuracy of eigenvalues from $`\epsilon _1`$ to $`\epsilon `$ and a fidelity $`\rho `$, a state of the system in the time instant $`\frac{1\rho }{14\epsilon }`$ can be obtained in time $`\frac{4}{(1\rho )\epsilon _1}`$ with this fidelity.
### 4.3 Generalizations
#### 4.3.1 Sparse areas of spectrum
A spectrum of real system like molecule typically contains strips where spectrum is continuos separated by gaps where energy levels are absent at all. Let $`w`$ be the maximal width of strips and $`g`$ be the minimal width of gapes.
Regard the times of evolution less than $`1/w`$ (here we assume the system of units where Plank constant is unit). For these times the width of strips is negligible and we can assume that $`\epsilon _1=1/g`$ in Theorem 2. So in this situation $`\epsilon =w`$ and we obtain the following generalization of Theorem 2.
If $`t=o(1/w)`$ we can predict the state $`U(t)`$ of primary system in time $`t_{pre}=O(t\frac{w}{g})`$.
Yet more generalization can be obtained if we consider the spectrum which is not sparse at all but has sparse areas. Here if the initial state of primary system is concentrated in sparse area, the prediction is possible with such probability which is equal to the degree of concentration.
#### 4.3.2 Travelling in a time: predictions and restoration of a history
A wizard transformation may be used not only for predictions a future but for restoration of a history as well. If we replace the time $`t`$ by $`t`$ in predicting procedure then we obtain the state of primary system in time instant $`t`$ which means the restoration of a history. Perform procedure from the section 4.2 with $`t`$ instead of $`t`$, where $`0<t<CN`$. This makes possible to obtain the state of primary system in time of order $`N`$ in the past. Thus we conclude the following generalization of Theorem 2.
###### Theorem 3
Given a Hamiltonian of some system, its states in time instant $`T=\frac{1\rho }{14\epsilon }`$ in future as well as in time instant $`T`$ in the past can be obtained in time $`\frac{4}{(1\rho )\epsilon _1}`$ with fidelity $`\rho `$ provided each eigenvalue can be quickly calculated with precision up to $`\epsilon `$ given its approximation with accuracy $`\epsilon _1`$.
Note that for restoration of a history we don’t need a sparse spectrum, may be $`\epsilon _1=\epsilon `$. This is surprising because the algorithm uses only operator $`\mathrm{U}`$ when the natural way to obtain a state from the past is to use an inverse transformation $`\mathrm{U}^1`$ and there is no evident way to simulate the action of the inverse transformation by means of $`\mathrm{U}`$.
### 4.4 Examples
#### 4.4.1 Shor factoring algorithm
Formula (2) for a wizard simulation can be considered as a natural generalization for the quantum part of Shor factoring algorithm (\[Sh\]). Let a unitary operator $`\mathrm{U}`$ in (2) has the form $`\mathrm{U}|x|ax(\mathrm{mod}q)`$, where $`ax`$ is a multiplication of integers, $`(a,q)=1`$. Then eigenvectors of $`\mathrm{U}`$ have the form $`\frac{1}{\sqrt{r}}\underset{j=0}{\overset{r1}{}}\mathrm{exp}(2\pi ikj/r)|a^j(\mathrm{mod}q)`$ and its eigenvalues are $`\mathrm{exp}(2\pi ij/r)`$ where $`r`$ is a period of $`a`$ (minimal integer such that $`a^r1(\mathrm{mod}q)`$). If we apply an operator (2) with this $`\mathrm{U}`$ to zero initial state and observe the second register then by Lemma 1 we obtain an approximation of a number $`j/r`$ up to $`O(1/N)`$ with high probability. Using this procedure sequentially we can restore the value $`r`$ of a period. This procedure was used for factoring by Shor.
Here we use a wizard simulation (2) as a technical element of the algorithm 4.2. Alternatively, one can try to use Kitaev’s method instead of QFT in (2). Namely, consider one controlling qubit and use Hadamard transform after conditional application of $`\mathrm{U}`$ and repeat this with numerous controlling qubits. Then by some auxiliary transformation it is possible to extract eigenvalues in ancillary register (look at \[Ki\] for details). But then we must get rid of revealed eigenvalues like in the point 6 of the algorithm and it faces a little difficulty when using Hadamard transform instead of QFT in (2), because we have not $`\mathrm{U}_{cond}^1`$ and must always manage with only $`\mathrm{U}_{cond}`$.
#### 4.4.2 Grover search algorithm
Consider Grover algorithm for the fast quantum search: $`U(t)|\stackrel{~}{0}=(I_{\stackrel{~}{0}}I_a)^t|\stackrel{~}{0}`$, where $`I_b`$ denotes inversion of the state $`b`$: $`I_b|b=|b,I_b|s=|s`$ for $`b|s=0`$. It was shown in \[Gr2\] that if $`|\stackrel{~}{0}|a|=O(1/\sqrt{N})`$ then $`U(t_1)|a`$ for $`t_1=O(\sqrt{N})`$ independently of $`\stackrel{~}{0}`$, where $`N`$ is the dimension of main space. What will happen if we apply Theorem 2 for $`U`$ as a primary evolution? The problem in fact will be two dimensional and the minimal gap between eigenvalues will be of order $`1/\sqrt{N}`$. Hence, by Theorem 2 we can predict the states $`U(t)`$ for $`t=O(N)`$ in time of order $`\sqrt{N}`$. As for $`t=O(\sqrt{N})`$ we obtain no additional speedup.
### 4.5 About practical implementations
The two main components of the algorithm are conditional iterations of a primary transformation $`\mathrm{U}`$ and QFT.
The main difficulty for the practical implementation of this method is in conditional iterations of $`\mathrm{U}`$. Given only a primary device realizing $`\mathrm{U}`$ one cannot immediately fulfil this iterations because it requires quantum control on the number of iterations. The solution may be following. Decomposition of a primary device into elementary parts that can be included to the quantumly controlled circuit and realize the conditional iterations of $`\mathrm{U}`$ for all parts simultaneously. All controlling qubits should be used in entangled state as a Shröedinger cat. Given a quantum gate array computing $`\mathrm{U}`$ we can easily construct a new gate array computing $`\mathrm{U}_{cond}`$. To put it otherwise a conditional application of $`\mathrm{U}`$ is possible through the control on microscopic level (this bears a resemblance with the control in living cells).
This scheme can be reformulated by means of analogous quantum computing if we consider QFT on an ancillary register as a passage to the canonically conjugate magnitude. Say, if we use a value of coordinates in operations with a register then canonically conjugate will be the corresponding impulse. The passage from the coordinate representation of a wave function to the impulse representation in one-dimensional case can be defined as
$$\varphi (p)=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{exp}(ipx/h)\psi (x)𝑑x$$
where a probability to obtain impulse in a segment $`(p,p+dp)`$ is $`|\varphi (p)|^2dp/2\pi h`$. Assume that we have one particle which can be located in $`M`$ points of the form $`x=0,\frac{1}{M},\frac{2}{M},\mathrm{},\frac{M1}{M}`$. Then the coordinate quantum space for one particle will be $`M`$ dimensional. Consider the corresponding integral sum for $`\varphi (p)`$ in the system of units where Plank constant $`h`$ is one. This integral sum will be just the sum in the definition of QFT where $`x`$ plays a role of $`s/M`$ and $`p/2\pi `$ plays a role of $`l`$.
Then the main algorithm acquires the following general form.
a). Primary evolution quantumly controlled by a magnitude containing in the properly prepared ancillary register.
b). Simple actions depending on the canonically conjugate magnitude.
c). Repetition of a).
Procedure of such a kind can predict a behavior of a primary system with sparse spectrum and restore the state of arbitrary system in the past. This scheme seems to be very simple and would be interesting to find its natural physical analog. The good starting point here may be the comparison between spectral features of the known systems and their functions and complexity.
## 5 Conclusion
Formulate again the main result: the method is presented which makes possible to obtain states of a primary quantum system earlier than they appear naturally in its evolution provided the spectrum of system is sparse. This speedup will be the more if the gaps between continuous strips of spectrum increase comparatively to the width of strips. This method can be applied also for restoration of states of the arbitrary primary system in the past. This method yields a speedup and it can make possible to fit into the time when coherent states exist and thus fight decoherence in quantum computations.
## 6 Acknowledgements
I am grateful to Lov Grover for his kind invitation to Bell Labs, discussions and attention to my work. I also thank Peter Hoyer and David DiVincenzo for fast replies and useful comments.
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# Microlensing Constraints on the Frequency of Jupiter Mass Planets
## 1. Introduction
A Galactic microlensing event occurs when a massive, compact object (the lens) passes near to our line-of-sight to a more distant star (the source). If the lens, observer, and source are perfectly aligned, then the lens images the source into a ring, called the Einstein ring, which has angular radius of<sup>1</sup><sup>1</sup>1For the scaling relation on the far right of equations (1), (2), and (3), we have assumed $`D_\mathrm{S}=8\mathrm{kpc}`$ and $`D_\mathrm{L}=6.5\mathrm{kpc}`$, typical distances to the lens and source for microlensing events toward the bulge.
$$\theta _\mathrm{E}\left[\frac{4GM}{c^2}\frac{D_{\mathrm{LS}}}{D_\mathrm{L}D_\mathrm{S}}\right]^{1/2}480\mu \mathrm{as}\left(\frac{M}{M_{}}\right)^{1/2},$$
(1)
where $`M`$ is the mass of the lens, and $`D_{\mathrm{LS}}`$, $`D_\mathrm{S}`$, $`D_\mathrm{L}`$ are the lens-source, observer-source, and observer-lens distances, respectively. This corresponds to a physical distance at the lens plane of
$$r_\mathrm{E}=\theta _\mathrm{E}D_\mathrm{L}3\mathrm{AU}\left(\frac{M}{M_{}}\right)^{1/2}.$$
(2)
If the lens is not perfectly aligned with the line-of-sight, then the lens splits the source into two images. The separation of these images is $`𝒪(\theta _\mathrm{E})`$ and hence unresolvable. However, the source is also magnified by the lens, by an amount that depends on the angular separation between the lens and source in units of $`\theta _\mathrm{E}`$. Since the lens, observer, and source are all in relative motion, this magnification is a function of time: a ‘microlensing event.’ The time scale for such an event is
$$t_\mathrm{E}\frac{\theta _\mathrm{E}}{\mu _{\mathrm{rel}}}40\mathrm{days}\left(\frac{M}{M_{}}\right)^{1/2},$$
(3)
where $`\mu _{\mathrm{rel}}`$ is the relative lens-source proper motion.
If the primary lens has a planetary companion, and the position of this companion happens to be in the path of one of the two images created during the primary event, then the planet will perturb the light from this image, creating a deviation from the primary light curve (see Figure 1). The duration of this perturbation is $`\sqrt{q}t_\mathrm{E}`$, where $`q`$ is the mass ratio between the planet and primary. Hence, for a Jupiter/Sun mass ratio ($`q10^3`$), the perturbation time scale is $`𝒪(\mathrm{day})`$. These short-duration deviations are the signatures of planets orbiting the primary lenses. Note that since the perturbation time scale is considerably less than $`t_\mathrm{E}`$, the majority of the light curve will be indistinguishable from a single lens.
Three parameters determine the magnitude of the perturbation, and hence define the observables. These are mass ratio $`q`$, the instantaneous angular separation $`d`$ between the planet and primary in units of $`r_\mathrm{E}`$, and the angle $`\alpha `$ between the projected planet/star axis and the path of the source. As $`q`$ decreases, the perturbation time scale decreases, although the magnitude of the deviation does not necessarily decrease. Thus very small mass ratio planets ($`q<10^5`$) can be detected using microlensing, although the detection probability is small. The lower limit to the detectable $`q`$ is set practically by the sampling of the primary event, and ultimately by the finite size of the source stars (Bennett & Rhie 1996). A microlensing event is generally alerted only if the minimum angular impact parameter in units of $`\theta _\mathrm{E}`$ satisfies $`u_01`$, which corresponds to image positions between $`(0.61.6)\theta _\mathrm{E}`$. Since the planet must be near one of these images to create a perturbation, microlensing is most sensitive to planets with separations $`0.6<d<1.6`$, the ‘lensing zone.’ The angle $`\alpha `$, which is of no physical interest, is uniformly distributed. Only certain values of $`\alpha `$ will create detectable deviations. Thus integration over $`\alpha `$ defines a geometric detection probability.
Microlensing as a method to detect extrasolar planets was first suggested by Mao & Paczyński (1991), and was expanded upon by Gould & Loeb (1992) who demonstrated that if all lenses had a Jupiter analog, than $`20\%`$ of all light curves should exhibit $`>5\%`$ deviations. Since these two seminal papers, many authors have explored the use of microlensing to detect planets. It is not our intention to provide a comprehensive review of this field. However, of particular relevance is the paper by Griest & Safizadeh (1998, GS98) who demonstrated that, for high-magnification events (those with maximum magnification $`A_{\mathrm{max}}>10`$), the detection probability for planets in the lensing zone is $`100\%`$. Thus high-magnification events are an extremely efficient means of detecting extrasolar planets. The results of GS98 also imply that multiple planets in the lensing zone should betray their presence in high-magnification events (Gaudi, Naber & Sackett 1998).
## 2. Limits on Companions in OGLE-1998-BUL-14
The basic requirements to detect planets with microlensing are good temporal sampling and photometric precision. Since the optical depth to microlensing is low, $`𝒪(10^6)`$, the survey teams that discover microlensing events toward the Galactic bulge must monitor of order one million stars on a nightly basis in order to detect any events. Therefore, they generally have insufficient sampling to detect the short-duration perturbations to the primary light curve (see, e.g., Figure 2). However, these survey teams (OGLE, Udalski et al. 1994; MACHO, Alcock et al. 1996; EROS; Glicenstein et al., these proceedings) issue alerts, notification of ongoing events. This allows follow-up collaborations (GMAN, Alcock et al. 1997; PLANET, Albrow et al. 1998; MPS, Rhie et al. 1999; MOA, Yock, these proceedings) to monitor these events frequently with high-quality photometry to search for planetary deviations. In particular, the PLANET collaboration has access to four telescopes located in Chile, South Africa, Western Australia, and Tasmania, and can monitor events nearly round-the-clock, weather permitting.
Figure 2 shows PLANET photometry of an event alerted by the OGLE collaboration, OGLE-1998-BUL-14. This was a high-magnification event ($`A_{\mathrm{max}}16)`$ with $`t_\mathrm{E}40\mathrm{days}`$, making it an excellent candidate to search for planetary deviations. PLANET obtained a total of 600 data points for this event: 461 $`I`$-band and 139 $`V`$-band. The median sampling interval is about 1 hour, or $`10^3t_\mathrm{E}`$, with very few gaps greater than 1 day. The $`1\sigma `$ scatter in $`I`$ over the peak of the event (where the sensitivity to planets is the highest) is $`1.5\%`$. The dense sampling and excellent photometry means that our efficiency to detect massive companions should be quite high. In fact, examination of the residuals from a single lens model (Fig. 2) reveal no obvious deviations of any kind.
To be more quantitative, we simultaneously search for binary-lens fits and calculate the detection efficiency $`ϵ(d,q)`$ of OGLE-1998-BUL-14 as a function of separation and mass ratio using a method proposed by Gaudi & Sackett (2000). For details on the implementation for this event, see Albrow et al. (2000).
We find no binary-lens models in the parameter ranges $`0d4`$ and $`10^5q1`$ that provide significantly ($`\mathrm{\Delta }\chi ^2>10`$) better fits to the OGLE-1998-BUL-14 dataset. We therefore conclude that the light curve of OGLE-1998-BUL-14 is consistent with a single lens.
In Figure 4 we show the detection efficiency $`ϵ`$ of our OGLE-1998-BUL-14 dataset to companions as a function of the mass $`M_\mathrm{p}`$ and orbital radius $`a`$ of the companion. Parameter combinations shaded in black are excluded at the 95% significance level. Stellar companions to the primary lens of OGLE-1998-BUL-14 with separations between $`2\mathrm{AU}`$ and $`11\mathrm{AU}`$ are excluded. Companions with mass $`10M_\mathrm{J}`$ are excluded between $`3\mathrm{AU}`$ and $`7\mathrm{AU}`$. Although we cannot exclude a Jupiter-mass companion at any separation, we had a $`80\%`$ chance of detecting such a companion at $`3\mathrm{AU}`$. The detection efficiency for OGLE-1998-BUL-14 is $`>25\%`$ at $`a=3\mathrm{AU}`$ for all companion masses $`M_\mathrm{p}>0.03M_\mathrm{J}`$. We find that we had a $`60\%`$ chance of detecting a companion with the mass and separation of Jupiter ($`M_\mathrm{p}=M_\mathrm{J}`$ and $`a=5.2\mathrm{AU}`$), and a $`5\%`$ chance of detecting a companion with the mass and separation of Saturn ($`M_\mathrm{p}=0.3M_\mathrm{J}`$ and $`a=9.5\mathrm{AU}`$) in the light curve of OGLE-1998-BUL-14.
Thus, although Jupiter analogs cannot be ruled out in OGLE-1998-BUL-14, the detection efficiencies are high enough that non-detections in several events with similar quality will be sufficient to place meaningful constraints on their abundance.
How do the OGLE-1998-BUL-14 efficiencies compare to planet detection via other methods? In Figure 4 we show the radial velocity detection limit on $`M_\mathrm{p}\mathrm{sin}i`$ for a solar mass primary as a function of the semi-major axis for a velocity amplitude of $`K=20\mathrm{m}\mathrm{s}^1`$, which is the limit found for the majority of the stars in the Lick Planet Search (Cumming, Marcy & Butler 1999). Although we show this limit for the full range of $`a`$, in reality the detection sensitivity extends only to $`a<5\mathrm{AU}`$ due to the finite duration of radial-velocity planet searches and the fact that one needs to observe a significant fraction of an orbital period. In addition, we plot in Figure 4 the $`M_\mathrm{p}\mathrm{sin}i`$ and $`a`$ for planetary candidates detected in the Lick survey. Radial velocity searches clearly probe a different region of parameter space than microlensing, in particular, smaller separations. Note, however, that our OGLE-1998-BUL-14 data set gives us a $`>75\%`$ chance of detecting analogs to two of these extrasolar planets: the third companion to Upsilon And and the companion to 14 Her. Although the efficiency is low, we do have sensitivity to planets with masses as small as $`0.01M_\mathrm{J}`$, considerably smaller than can be detected via radial velocity methods. For comparison, we also show in Figure 10 the astrometric detection limit on $`M_\mathrm{p}`$ for a $`M_{}`$ primary at $`10\mathrm{pc}`$, for an astrometric accuracy of $`\sigma _\mathrm{A}=1\mathrm{mas}`$. For an astrometric campaign of $`11`$ years, this limit extends to $`5\mathrm{AU}`$. Such an astrometric campaign ($`\sigma _\mathrm{A}=1\mathrm{mas},P=11`$ years), would be sensitive to companions similar to those excluded in our analysis of OGLE-1998-BUL-14. The proposed Space Interferometry Mission (SIM) promises $`4\mu \mathrm{as}`$ astrometric accuracy, which would permit the detection of considerably smaller mass companions.
## 3. Combined Limits from the 1998-1999 PLANET Seasons
Clearly one cannot make any statements about the population of the extrasolar planets as a whole based on one event. Fortunately, PLANET has monitored, over the last five years, more than 100 events, a subset of which have temporal sampling and photometric accuracy similar to that of OGLE-1998-BUL-14. Here we present a preliminary analysis of these events.
We select 23 high-quality light curves from the 1998-1999 PLANET seasons and analyze these in the same manner as OGLE-1998-BUL-14 (Albrow et al. 2000). Included in this sample are 5 high-magnification ($`A_{\mathrm{max}}>10`$) events: OGLE-1998-BUL-14, MACHO-1998-BUL-35, OGLE-1999-BLG-5, OGLE-1999-BUL-35 and OGLE-1999-BUL-36.
We find that all 23 events are consistent with a single lens to within our detection threshold. Using this null result, along with the detection efficiency $`ϵ_i(d,q)`$ for each event $`i`$, we find a statistical upper limit to the fraction of these lenses that have a companion of a given separation and mass ratio. If the fraction of lenses with a companion as a function of $`d`$ and $`q`$ is $`f(d,q)`$, then probability that $`N`$ events with individual efficiencies $`ϵ_i(d,q)`$ would give a null result (no detections) is,
$$P=\mathrm{\Pi }_{i=1}^N[1f(d,q)ϵ_i(d,q)].$$
(4)
The 95% confidence level (c.l.) upper limit to $`f(d,q)`$ is found by setting $`P=5\%`$.
In Figure 4 we show the 95% c.l. upper limits to $`f(d,q)`$ for separations $`0d4`$ and $`q=10^2,10^{2.5},`$ and $`10^3`$. We convert these to limits on the fraction of lenses with companions of a given mass and physical separation by assuming that all the primaries have mass $`M=M_{}`$ and distance $`D_\mathrm{L}=6\mathrm{kpc}`$. We find that $`<33\%`$ of these lenses have Jupiter mass planets with separations of 1.5-3 AU. Similarly, $`<33\%`$ have planets of mass $`M_\mathrm{p}3M_\mathrm{J}`$ with separations of 1-4 AU. Although we cannot place an interesting limit on Jupiter analogs, we do find that $`<50\%`$ of lenses have $`3M_\mathrm{J}`$ planets at the separation of Jupiter ($`5.2\mathrm{AU}`$).
## 4. Conclusions
Microlensing offers a unique and complementary method of detecting extrasolar planets. Although many light curves have been monitored in the hopes of detecting the short-duration signature of planetary companions to the primary lenses, no convincing planetary detections have yet been made, despite the fact that data of sufficient quality are being acquired to detect such companions. These null results indicate that Jupiter-mass companions with separations in the ‘lensing zone,’ $`1.53\mathrm{AU}`$, occur is less than 1/3 of systems.
The potential for this field is enormous. Current microlensing searches for planets will continue to monitor events alerted toward the bulge, and either push these limits down to levels probed by radial velocity surveys ($`5\%`$), or finally detect planets, and measure the frequency of companions at separations more relevant to our solar system. Next generation microlensing planet searches have the promise of obtaining a robust statistical estimate of the fraction of stars with planets of mass as low as that of the Earth.
### Acknowledgments.
We thank the MACHO, OGLE and EROS collaborations for providing real-time alerts. We are especially grateful to the observatories that support our science (Canopus, CTIO, Perth and SAAO) via the generous allocations of time that make this work possible. This work was supported by grants AST 97-27520 and AST 95-30619 from the NSF, by grant NAG5-7589 from NASA, by a grant from the Dutch ASTRON foundation through ASTRON 781.76.018, by a Marie Curie Fellowship from the European Union, by “coup de pouce 1999” award from Ministère de l’Éducation nationale, de la Rechereche et de la Technologie, and by a Presidential Fellowship from the Ohio State University.
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Alcock, C., et al. 1996, ApJ, 463, L67
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# Sequences In non-commutative 𝐿^𝑝-spaces
## 1. Introduction
In , Kadec and Pełczyński proved that if $`1p<\mathrm{}`$ then every bounded sequence $`\{f_n\}_{n=1}^{\mathrm{}}`$ in $`L^p[0,1]`$ has a subsequence that can be decomposed into two extreme sequences $`\{g_k\}_{k=1}^{\mathrm{}}`$ and $`\{h_k\}_{k=1}^{\mathrm{}}`$, where the $`h_k`$’s are pairwise disjoint and the $`g_k`$’s are $`L_p`$-equi-integrable that is $`\underset{m(A)0}{lim}\mathrm{sup}_k\chi _Ag_k_p0`$ and $`h_kg_k`$ for every $`k1`$. This result was used to study different structures of subspaces of $`L^p[0,1]`$. Later, the same decomposition property was proved for larger classes of Banach function spaces (see for Orlicz spaces with $`\mathrm{\Delta }_2`$-condition and $`q`$-concave lattices, for some symmetric spaces). There are however Banach lattices with sequences for which the above decomposition is not possible. Examples of reflexive, $`p`$-convex Banach lattices without the subsequence decomposition can be found in a paper of Figiel et al . Subsequently, Weis characterized, in terms of uniform order continuity conditions and ultrapowers, all Banach lattices where such property is possible. For the case of rearrangement invariant function spaces, the order-continuous spaces, in which the above decomposition works, were fully characterized as those that have the Fatou property (equivalently, those that contains no subspace isomorphic to $`c_0`$). In , a version of Kadec-Pełczynski decomposition was considered for preduals of semi-finite von Neumann algebras.
It is the intention of the present paper to give an extension of the Kadec-Pełczyński decomposition stated above to the case of general non-commutative symmetric spaces of measurable operators. Let $``$ be a von Neumann algebra, equipped with a distinguished faithful, normal, semi-finite trace $`\tau `$ and $`E`$ be a rearrangement invariant Banach function space on $`[0,1]`$ or the half line $`(0,\mathrm{})`$. We define equi-integrability in the non-commutative setting as generalizations of Akemann’s characterization of weak compactness on preduals of von Neumann algebras. Using such notion, we provide an analogue of the Kadec-Pełczyński decomposition for non-commutative spaces. Namely, we proved that if $`E`$ is order continuous and satisifies the Fatou property then the corresponding symmetric space of measurable operators $`E(,\tau )`$ has the subsequence splitting property. Our approach allows ones to consider more general spaces such as quasi-Banach rearrangement invariant spaces that are $`\alpha `$-convex with constant $`1`$ and satisfy non trivial $`q`$-lower estimate with constant $`1`$. In particular splitting of bounded sequences is valid in non-commutative $`L^p`$-spaces for $`0<p<\mathrm{}`$. It should be noted that Sukochev obtain a similar result for the case of finite von Neumann algebras. We also remark that since $`c_0`$ fails the subsequence splitting property, our result for the case where $`E`$ is a Banach space case is the best possible.
As application of the main result, we study the structure of subspaces of $`L^p(,\tau )`$ for $`1<p<\mathrm{}`$ which generalizes the case $`p=1`$ treated in .
We refer to and for general information concerning von Neumann algebras as well as non-commutative integration, to and for Banach lattice theory.
## 2. Definitions and preliminary results
Throughout, $`H`$ is a given Hilbert space and $`(H)`$ denotes a semi-finite von Neumann algebra with a distinguished normal, faithful semi-finite trace $`\tau `$. The identity in $``$ will be denoted by $`\mathrm{𝟏}`$ and $`_p`$ will stand for the set of all (self adjoint) projections in $``$. A closed and densely defined operator $`a`$ on $`H`$ is said to be affiliated with $``$ if $`ua=au`$ for all unitary operator $`u`$ in the commutant $`^{}`$ of $``$.
A closed and densely defined operator $`x`$, affiliated with $``$, is called $`\tau `$-measurable if for every $`\epsilon >0`$, there exists an orthogonal projection $`p`$ such that $`p(H)\text{dom}(x)`$, $`\tau (\mathrm{𝟏}p)<\epsilon `$ and $`xp`$. The set of all $`\tau `$-measurable operators will be denoted by $`\stackrel{~}{}`$. The set $`\stackrel{~}{}`$ is a $``$-algebra with respect to the strong sum, the strong product and the adjoint operation. Given a self-adjoint operator $`x`$ in $`\stackrel{~}{}`$, we denote by $`e^x()`$ its spectral measure. Recall that $`e^{|x|}(B)`$ for all Borel sets $`B`$ and $`x\stackrel{~}{}`$. For fixed $`x\stackrel{~}{}`$ and $`t0`$, we define
$$\mu _t(x)=\mathrm{inf}\left\{s0:\tau (e^{|x|}(s,\mathrm{}))t\right\}.$$
The function $`\mu _{(.)}(x):[0,\mathrm{})[0,\mathrm{}]`$ is called the generalized singular value function (or decreasing rearrangement) of $`x`$. We note that $`\mu _t(x)<\mathrm{}`$ for every $`t>0`$. For a complete study of $`\mu _{(.)}`$, we refer the reader to . The topology defined by the metric on $`\stackrel{~}{}`$ obtained by setting
$$d(x,y)=\mathrm{inf}\left\{t0:\mu _t(xy)t\right\},\text{for}x,y\stackrel{~}{},$$
is called the measure topology. It is well-known that a net $`(x_\alpha )_{\alpha I}`$ in $`\stackrel{~}{}`$ converge to $`x\stackrel{~}{}`$ in measure topology if and only if for every $`\epsilon >0`$, $`\delta >0`$, there exists $`\alpha _0I`$ such that whenever $`\alpha \alpha _0`$, there exists a projection $`p_p`$ such that
$$(x_\alpha x)p_{}<\epsilon \text{and}\tau (\mathrm{𝟏}p)<\delta .$$
It was shown in that $`(\stackrel{~}{},d)`$ is a complete metric space.
Recall that if we consider $`=L^{\mathrm{}}(^+,m)`$, where $`m`$ is the Lebesgue measure on $`^+`$ then $``$ is an abelian von Neumann algebra acting on $`L^2(^+,m)`$ via multiplication, with the trace being the integration with respect to $`m`$, then $`\stackrel{~}{}=L^0(^+,m)`$ (the usual space of all measurable functions on $`^+`$) and the generalized singular value $`\mu (f)`$ is precisely the decreasing rearrangement of the function $`|f|`$ (usually denoted by $`f^{}`$ in Banach lattice theory).
###### Definition 2.1.
A symmetric quasi-Banach function space on $`^+`$ is a quasi-Banach lattice $`E`$ of measurable functions with the following properties:
* $`E`$ is an order ideal in $`L^0(^+,m)`$;
* $`E`$ is rearrangemant invariant in the sense of (p .114);
* $`E`$ contains all finitely supported simple functions.
Unless stated otherwise, $`E`$ will always denote a symmetric quasi-Banach function space on $`^+`$. We define the symmetric space of measurable operators $`E(,\tau )`$ by setting
$$E(,\tau ):=\{x\stackrel{~}{}:\mu (x)E\}$$
and
$$x_{E(,\tau )}=\mu (x)_E\text{for all}xE(,\tau ).$$
It is shown in (Lemma 4.1) that if $`E`$ is $`\alpha `$-convex (for some $`0<\alpha 1`$) with constant $`1`$, then $`_{E(,\tau )}`$ is an $`\alpha `$-norm that is for every $`x,yE(,\tau )`$,
$$x+y_{E(,\tau )}^\alpha x_{E(,\tau )}^\alpha +y_{E(,\tau )}^\alpha .$$
Equipped with $`_{E(,\tau )}`$, the space $`E(,\tau )`$ is a $`\alpha `$-Banach space. The space $`E(,\tau )`$ is often referred to as the non-commutative analogue of the function space $`E`$. We remark that if $`0<p<\mathrm{}`$ and $`E=L^p(^+,m)`$ then $`E(,\tau )`$ coincides with the usual non-commutative $`L^p`$-space associated to the semi-finite von Neumann algebra $``$. Also if $`E=L^{\mathrm{}}(^+,m)`$, then $`L^{\mathrm{}}(,\tau )`$ is the von Neumann algebra $``$. We refer to , and for some background on the space $`E(,\tau )`$.
###### Definition 2.2.
A quasi-Banach function space $`E`$ is said to satisfy a lower $`q`$-estimate if there exists a positive constant $`C>0`$ such that for all finite sequences $`\{x_n\}`$ of mututally disjoint elements in $`E`$,
$$\left(x_n^q\right)^{\frac{1}{q}}Cx_n.$$
The least such constant $`C`$ is called the constant of the lower $`q`$-estimate. Recall that if $`E`$ is a quasi-Banach function space and $`1<p<\mathrm{}`$,
$$E^{(p)}=\left\{xL^0(^+,m);|x|^pE\right\}withx_{E^{(p)}}=|x|^p_E^{\frac{1}{p}}.$$
We will need the following known result. A proof can be found in .
###### Proposition 2.3.
Assume that $`E`$ is order-continuous and $`\alpha `$-convex with constant $`1`$ for some $`0<\alpha 1`$.
* If $`xE(,\tau )`$ and $`ef`$ are projections in $``$ then $`exe_{E(,\tau )}fxf_{E(,\tau )}`$;
* If $`xE(,\tau )`$ and $`e_\beta _\beta 0`$ is a net of projections in $``$ then $`xe_\beta _{E(,\tau )}_\beta 0`$.
The following definition isolates the main topic of this paper.
###### Definition 2.4.
Let $`E`$ be a quasi-Banach function space on $`^+`$ and $`K`$ be a bounded subset of $`E(,\tau )`$. We will say that $`K`$ is $`E`$-equi-integrable if $`\underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}e_nxe_n_{E(,\tau )}=0`$ for every decreasing sequence $`\{e_n\}_{n=1}^{\mathrm{}}`$ of projections with $`e_n_n0`$.
We remark, from Proposition 2.3, that since $`\{e_n\}_{n=1}^{\mathrm{}}`$ is decreasing, so is the sequence $`\left\{\mathrm{sup}_{xK}e_nxe_n_{E(,\tau )}\right\}_{n=1}^{\mathrm{}}`$ and therefore the limit in the definition above always exists. This definition was motivated by the commutative case on one hand and the characterization of weakly compact subsets of $`L^1(,\tau )`$ by Akemann (see also p.150) on the other. Using this terminology, Akemann’s characterization can be stated as in the commutative case: relatively weakly compact subsets of $`L^1(,\tau )`$ are exactly the equi-integrable sets. Such characterization is not valid in general. For $`1<p<\mathrm{}`$, any set of normalized disjoint sequence cannot be $`L^p`$-integrable but since $`L^p`$ is reflexive, such set is relatively weakly compact.
On the next proposition, we will show that for the general case, one implication always holds.
###### Proposition 2.5.
Assume that $`E`$ is an order-continuous symmetric Banach function space and $`K`$ is a $`E`$-equi-integrable set in $`E(,\tau )`$ then $`K`$ is relatively weakly compact.
The proposition will be proved in several steps. Recall that if $`E`$ is a symmetric Banach function space, then $`E(,\tau )`$ is a subset of $`L^1(,\tau )+`$ and therefore if $`p`$ is a projection in $`L^1(,\tau )`$ and $`K`$ is a subset of $`E(,\tau )`$, then $`pK`$ and $`Kp`$ are subsets of $`L^1(,\tau )`$.
###### Lemma 2.6.
Let $`p`$ be a projection in $`L^1(,\tau )`$ and $`K`$ be a $`E`$-equi-integrable subset of $`E(,\tau )`$. The sets $`pKp`$ and $`pK(1p)`$ are relatively weakly compact in $`L^1(,\tau )`$.
To see this lemma, it is enough to check that these sets are $`L^1`$-equi-integrable. Let $`T:E(,\tau )L^1(,\tau )`$ be the linear map defined by $`xTx=pxp`$. This map is well-defined and one can deduce from the closed graph theorem that it is bounded. Let $`\{e_n\}_{n=1}^{\mathrm{}}`$ be a sequence of projections with $`e_n_n0`$. For each $`n1`$, set $`f_n`$ to be the right support projection of $`e_np`$. By the definition of support projections, $`f_np`$. So $`\{f_n\}_{n=1}^{\mathrm{}}`$ is a sequence of finite projections. We also note that (see for instance the proof of \[19, Proposition 1.6 p.292\]),
$$f_n=e_n(\mathrm{𝟏}p)(\mathrm{𝟏}p)$$
and by Kaplansky formula (see for instance \[12, Theorem 6.1.6 p.403\]),
$$f_ne_ne_n(\mathrm{𝟏}p).$$
Since $`\tau (f_n)=\tau (e_ne_n(\mathrm{𝟏}p))\tau (p)`$ and $`\{e_ne_n(\mathrm{𝟏}p)\}_{n=1}^{\mathrm{}}`$ converges to zero, $`\{f_n\}_{n=1}^{\mathrm{}}`$ converges to zero. Now since the $`f_n`$’s are finite projections, we conclude that if $`g_n=_{kn}f_k`$, then $`\{g_n\}_{n=1}^{\mathrm{}}`$ converges to zero. Therefore, for every $`xK`$,
$`e_npxpe_n_1`$ $`=e_np(g_nxg_n)pe_n_1`$
$`p(g_nxg_n)p_1`$
$`Tg_nxg_n_{E(,\tau )}.`$
Since $`K`$ is $`E`$-equi-integrable, one obtain that
$$\underset{n\mathrm{}}{lim}\mathrm{sup}_{ypKp}e_nye_n_1T\underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}g_nxg_n_{E(,\tau )}=0$$
which concludes that $`pKp`$ is relatively weakly compact in $`L^1(,\tau )`$.
For $`pK(\mathrm{𝟏}p)`$, set $`S:E(,\tau )L^1(,\tau )`$ be the map defined by $`xSx=px(\mathrm{𝟏}p)`$. As above, $`S`$ is bounded. Let $`\{e_n\}_{n=1}^{\mathrm{}}`$ and $`\{g_n\}_{n=1}^{\mathrm{}}`$ be sequences of projections as discribed above. For each $`n1`$, let $`s_n`$ be the left support projection of $`(\mathrm{𝟏}p)e_n`$. Then $`s_n=e_npp`$ for every $`n1`$ and the sequence $`\{s_n\}_{n=1}^{\mathrm{}}`$ is decreasing. It is claimed that $`s_n_n0`$.
For this, it is enough to check that $`e_np_np`$. In fact, $`e_npe_npe_np`$ and the sequence defined by the right hand side of the equivalence converges to $`p`$ which implies that
$$\underset{n\mathrm{}}{lim}\tau (e_npe_n)=\underset{n\mathrm{}}{lim}\tau (pe_np)=\tau (p)$$
which gives
$$\underset{n\mathrm{}}{lim}\tau (e_nppe_n)=0.$$
But since $`(e_nppe_n)^2=(e_nppe_n)+e_np+pe_n`$, we can conclude that
$$\underset{n\mathrm{}}{lim}e_nppe_n_2=0.$$
From this, we get (by passing to a subsequence if necessary) that $`\{e_nppe_n\}_{n=1}^{\mathrm{}}`$ converges to zero in measure. Similarly, $`\{e_npppe_n\}_{n=1}^{\mathrm{}}`$ converges to zero in measure so $`e_np_np`$ hence $`s_n_n0`$.
To conclude the proof of Lemma 2.6, note that $`g_ns_n`$ so $`g_ns_n=g_n+s_n`$. In particular, $`g_ns_n_n0`$ and we get that
$$\begin{array}{cc}\hfill \underset{n\mathrm{}}{lim}\mathrm{sup}_{ypK(\mathrm{𝟏}p)}e_nye_n_1& =\underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}e_npx(\mathrm{𝟏}p)e_n_1\hfill \\ & =\underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}e_np(g_ns_n)x(g_ns_n)(\mathrm{𝟏}p)e_n_1\hfill \\ & \underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}p(g_ns_n)x(g_ns_n)(\mathrm{𝟏}p)_1\hfill \\ & S\underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}(g_ns_n)x(g_ns_n)_{E(,\tau )}=0.\hfill \end{array}$$
The proof is complete.
###### Lemma 2.7.
Let $`p`$ and $`K`$ be as in Lemma 2.6. Then $`pK`$ is relatively weakly compact in $`E(,\tau )`$.
Note first that $`pK`$ is $`E`$-equi-integrable. This can be seen by applying the series of argument used in Lemma 2.6, considering the operators $`T`$ and $`S`$ as maps from $`E(,\tau )`$ into $`E(,\tau )`$.
Let $`\{px_n\}_{n=1}^{\mathrm{}}`$ be a bounded sequence in $`pK`$. From Lemma reffinitecompact1, we can assume that $`\{px_n\}_{n=1}^{\mathrm{}}`$ is weakly convergent in $`L^1(,\tau )`$. Fix $`\phi E^{}(,\tau )_+`$ and let $`\phi =_0^{\mathrm{}}t𝑑e_t`$ be its spectral decomposition. For each $`k1`$, set $`q_k:=e^\phi ((0,k))`$. We remark that $`\phi q_k=q_k\phi `$. For $`m,n`$,
$`\phi ,px_npx_m`$ $`=\phi \phi q_k,px_npx_m+\phi q_k,px_npx_m`$
$`=\phi (\mathrm{𝟏}q_k),px_npx_m+\phi q_k,px_npx_m`$
$`=(\mathrm{𝟏}q_k)\phi (\mathrm{𝟏}q_k),px_npx_m+\phi q_k,px_npx_m.`$
This gives
$$\left|\phi ,px_npx_m\right|\left|\tau \left(\phi (1q_k)(px_npx_m)(1q_k)\right)\right|+\left|\phi q_k,px_npx_m\right|.$$
Since $`\phi q_k`$ belongs to $``$,
$$\underset{n,m\mathrm{}}{lim\; sup}\left|\phi q_k,px_npx_m\right|2\phi _{E^{}(,\tau )}\mathrm{sup}_{aK}(\mathrm{𝟏}q_k)pa(\mathrm{𝟏}q_k)_{E(,\tau )}.$$
Now since $`\mathrm{𝟏}q_k_k0`$ and $`pK`$ is $`E`$-equi-integrable, we obtain that
$$\underset{n,m\mathrm{}}{lim}\phi q_k,px_npx_m=0.$$
The lemma is proved.
To complete the proof of Proposition 2.5, let $`\{p_k\}_{k=1}^{\mathrm{}}`$ be a sequence of projections that increases to $`\mathrm{𝟏}`$ and $`\tau (p_k)<\mathrm{}`$ and fix $`\epsilon >0`$. Choose $`k_01`$ such that
$$\mathrm{sup}_{aK}(\mathrm{𝟏}p_{k_0})a(\mathrm{𝟏}p_{k_0})_{E(,\tau )}\epsilon .$$
We have $`K=p_{k_0}K+(\mathrm{𝟏}p_{k_0})Kp_{k_0}+(\mathrm{𝟏}p_{k_0})K(\mathrm{𝟏}p_{k_0})`$ which implies that
$$Kp_{k_0}K+(\mathrm{𝟏}p_{k_0})Kp_{k_0}+\epsilon B_{E(,\tau )}$$
where $`B_{E(,\tau )}`$ denotes the closed unit ball of $`E(,\tau )`$. From Lemma 2.7, the sets $`p_{k_0}K`$ and $`(\mathrm{𝟏}p_{k_0})Kp_{k_0}`$ are relatively weakly compact which concludes that $`K`$ is relatively weakly compact. The proof is complete. ∎
###### Remark 2.8.
If $`\tau (\mathrm{𝟏})<\mathrm{}`$, the proof above can be considerably shortened. In this case, $`E(,\tau )L^1(,\tau )`$ so if $`K`$ is $`E`$-equi-integrable, then it is relatively weakly compact in $`L^1(,\tau )`$ and on can argue as in the last part of Lemma 2.7 to conclude that $`K`$ is relatively weakly compact in $`E(,\tau )`$.
The following proposition should be compared with (Theorem 5.1 and Theorem 5.2). It generalizes well known property of equi-integrable sets in function spaces to the non-commutative settings.
###### Proposition 2.9.
Let $`E`$ be a symmetric quasi-Banach space function space and $`K`$ be a $`E`$-equi-integrable subset of $`E(,\tau )`$. For each sequence $`\{x_n\}_{n=1}^{\mathrm{}}`$ in $`K`$ and $`x\overline{K}`$, the following are equivalent:
* $`\underset{n\mathrm{}}{lim}x_nx_{E(,\tau )}=0`$;
* $`\{x_n\}_{n=1}^{\mathrm{}}`$ converges to $`x`$ in measure (as $`n\mathrm{}`$).
###### Proof.
The implication $`(a)(b)`$ is trivial. For $`(b)(a)`$, we will assume that $`x=0`$. Recall that there exists $`0<\alpha 1`$, such that $`L^\alpha (,\tau )E(,\tau )+L^\alpha (,\tau )`$ with $`x_{+L^\alpha (,\tau )}x_{E(,\tau )}2x_{L^\alpha (,\tau )}`$ for every $`xL^\alpha (,\tau )`$. The proposition will be proved by showing the following lemma:
###### Lemma 2.10.
For every $`p_p`$ with $`\tau (p)<\mathrm{}`$, $`\underset{n\mathrm{}}{lim}x_np_{E(,\tau )}=\underset{n\mathrm{}}{lim}px_n_{E(,\tau )}=0`$.
To see this lemma, fix $`\epsilon >0`$ and let $`C=\mathrm{max}\{1,\tau (p)\}`$. Since $`K`$ is equi-integrable, there exists $`\delta >0`$ such that whenever $`q_p`$ satisfies $`\tau (q)<\delta `$, then for every $`n`$, $`qx_nq_{E(,\tau )}\epsilon /(2)^{1/\alpha }`$. Since both $`\{x_n\}_{n=1}^{\mathrm{}}`$ and $`\{x_n^{}\}_{n=1}^{\mathrm{}}`$ converge to zero in measure, one can choose $`n_01`$ such that for each $`nn_0`$, there exists a projection $`p_n_p`$ with $`\tau (\mathrm{𝟏}p_n)<\delta `$,
$$x_np_n_{}<\frac{\epsilon }{2[4C]^{\frac{1}{\alpha }}}$$
and
$$x_n^{}p_n_{}<\frac{\epsilon }{2[4C]^{\frac{1}{\alpha }}}.$$
For $`nn_0`$,
$`x_np_{E(,\tau )}^\alpha `$ $`x_np_np_{E(,\tau )}^\alpha +p_nx_n(\mathrm{𝟏}p_n)p_{E(,\tau )}^\alpha +(\mathrm{𝟏}p_n)x_n(\mathrm{𝟏}p_n)p_{E(,\tau )}^\alpha `$
$`2^\alpha \mathrm{max}\{x_np_n_{}^\alpha ,x_np_np_{L^\alpha (,\tau )}^\alpha \}`$
$`+2^\alpha \mathrm{max}\{x_n^{}p_n_{}^\alpha ,p(\mathrm{𝟏}p_n)x_n^{}p_n_{L^\alpha (,\tau )}^\alpha \}+(\mathrm{𝟏}p_n)x_n(\mathrm{𝟏}p_n)_{E(,\tau )}^\alpha `$
$`2.2^\alpha \mathrm{max}\{\epsilon ^\alpha /2^\alpha 4C,(\epsilon ^\alpha /2^\alpha 4C)\tau (p)\}+\epsilon ^\alpha /2`$
$`\epsilon ^\alpha .`$
A similar estimate works for $`\{x_n^{}p\}_{n=1}^{\mathrm{}}`$. The lemma is verified.
To complete the proof of Proposition 2.9, choose a mutually disjoint family $`\{e_i\}_{iI}`$ of projections in $``$ with $`_{iI}e_i=\mathrm{𝟏}`$ for the strong operator topology and $`\tau (e_i)<\mathrm{}`$ for all $`iI`$. Using a similar argument as in , one can get an at most countable subset $`\{e_k\}_{k=1}^{\mathrm{}}`$ of $`\{e_i\}_{iI}`$ such that for each $`e_i`$ outside of $`\{e_k\}_{k=1}^{\mathrm{}}`$, $`e_ix_n=x_ne_i=0`$ for every $`n`$. Let $`e=_ke_k`$. Replacing $``$ by $`ee`$ and $`\tau `$ by its restriction on $`ee`$, we may assume that $`e=\mathrm{𝟏}`$. Let $`p_n=_{kn}e_k`$. It is clear that $`p_n_n0`$ and $`\tau (\mathrm{𝟏}p_n)<\mathrm{}`$ for every $`n`$. Fix $`\epsilon >0`$ and choose $`n_01`$ such that
$$\mathrm{sup}_np_{n_0}x_np_{n_0}_{E(,\tau )}\epsilon .$$
We get that
$$\underset{n\mathrm{}}{lim\; sup}x_n_{E(,\tau )}^\alpha \underset{n\mathrm{}}{lim}x_n(\mathrm{𝟏}p_{n_0})_{E(,\tau )}^\alpha +\underset{n\mathrm{}}{lim}x_n^{}(\mathrm{𝟏}p_{n_0}_{E(,\tau )}^\alpha +\epsilon =\epsilon $$
and since $`\epsilon `$ is arbitrary, the proof is complete. ∎
For the rest of this section, we collect some results that will be useful in the later sections of the paper.
The inequality given below can be viewed as the analogue of the well-known fact on normal functional on von Neumann algebra, $`|\phi (a)|^2\phi |\phi |(aa^{})`$ whenever $`a`$ and $`\phi _{}`$ \[19, Proposition 4.6 p. 146\], to the general case of symmetric spaces of measurable operators.
###### Proposition 2.11.
Let $`xE(,\tau )`$ and $`y`$ then
$$xy_{E(,\tau )}|x|y_{E(,\tau )}x_{E(,\tau )}^{\frac{1}{2}}y^{}|x|y_{E(,\tau )}^{\frac{1}{2}}.$$
###### Proof.
Let $`x=u|x|`$ be the polar decomposition of $`x`$. Then $`xy_{E(,\tau )}=u|x|y_{E(,\tau )}u_{\mathrm{}}|x|y_{E(,\tau )}`$. Also $`|x|y_{E(,\tau )}=|x|^{\frac{1}{2}}|x|^{\frac{1}{2}}y_{E(,\tau )}`$ and using Hölder’s inequality,
$`|x|y_{E(,\tau )}`$ $`|x|^{\frac{1}{2}}_{E^{(2)}(,\tau )}|x|^{\frac{1}{2}}y_{E^{(2)}(,\tau )}`$
$`=x_{E(,\tau )}^{\frac{1}{2}}y^{}|x|y_{E(,\tau )}^{\frac{1}{2}}.`$
###### Remark 2.12.
Let $`K`$ be a bounded subset of $`E(,\tau )`$. If we set $`|K|:=\{|a|:aK\}`$, then it is clear from Proposition 2.11 that if $`|K|`$ is $`E`$-equi-integrable then for every decreasing projections $`e_n_n0`$, $`\underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}xe_n_{E(,\tau )}=\underset{n\mathrm{}}{lim}\mathrm{sup}_{xK}e_nx_{E(,\tau )}=0`$. In particular if $`|K|`$ is $`E`$-equi-integrable then so is $`K`$.
###### Proposition 2.13.
Assume that $`E`$ is $`\alpha `$-convex with constant 1 for some $`0<\alpha 1`$. Let $`\{p_n\}_{n=1}^{\mathrm{}}`$ be a sequence of decreasing projections in $``$ and $`K`$ be a bounded subset of $`E(,\tau )`$ such that:
* $`p_n_n0`$;
* For each $`n1`$, the sets $`(\mathrm{𝟏}p_n)K`$ and $`|K(\mathrm{𝟏}p_n)|`$ are $`E`$-equi-integrable.
Then $`K`$ is $`E`$-equi-integrable if and only if $`\underset{n\mathrm{}}{lim}\mathrm{sup}_{aK}p_nap_n_{E(,\tau )}=0.`$
###### Proof.
We will show the non trivial implication. Fix $`f_k_k0`$, a sequence in $`_p`$. We need to show that $`\underset{k\mathrm{}}{lim}\mathrm{sup}_{aK}f_kaf_k_{E(,\tau )}=0`$.
We will assume without loss of generality that $`K`$ is a subset of the unit ball of $`E(,\tau )`$. For every $`aK`$,
$`f_kaf_k`$ $`=f_k(\mathrm{𝟏}p_n)af_k+f_kp_naf_k`$
$`=f_k(\mathrm{𝟏}p_n)af_k+f_kp_na(\mathrm{𝟏}p_n)f_k+f_kp_nap_nf_k.`$
Since $`E(,\tau )`$ is $`\alpha `$-convex, we get:
$`f_kaf_k_{E(,\tau )}^\alpha `$ $`f_k(\mathrm{𝟏}p_n)af_k_{E(,\tau )}^\alpha +f_kp_na(\mathrm{𝟏}p_n)f_k_{E(,\tau )}^\alpha +f_kp_nap_nf_k_{E(,\tau )}^\alpha `$
$`f_k(\mathrm{𝟏}p_n)af_k_{E(,\tau )}^\alpha +a(\mathrm{𝟏}p_n)f_k_{E(,\tau )}^\alpha +p_nap_n_{E(,\tau )}^\alpha .`$
Using Proposition 2.11 on the second term, we have
$`f_kaf_k_{E(,\tau )}^\alpha `$ $`f_k(\mathrm{𝟏}p_n)af_k_{E(,\tau )}^\alpha +a(\mathrm{𝟏}p_n)_{E(,\tau )}^{\frac{\alpha }{2}}f_k|a(\mathrm{𝟏}p_n)|f_k_{E(,\tau )}^{\frac{\alpha }{2}}`$
$`+p_nap_n_{E(,\tau )}^\alpha .`$
Let $`\epsilon >0`$, choose $`n_0`$ large enough so that $`\mathrm{sup}_{aK}p_{n_0}ap_{n_0}_{E(,\tau )}<\epsilon `$. We conclude that
$`\underset{k\mathrm{}}{lim}\mathrm{sup}_{aK}f_kaf_k_{E(,\tau )}^\alpha `$ $`\underset{k\mathrm{}}{lim}\mathrm{sup}_{aK}f_k(\mathrm{𝟏}p_{n_0})af_k_{E(,\tau )}^\alpha `$
$`+\underset{k\mathrm{}}{lim}\mathrm{sup}_{aK}f_k|a(\mathrm{𝟏}p_{n_0})|f_k_{E(,\tau )}^{\frac{\alpha }{2}}+\epsilon ^\alpha .`$
By (ii), the first two terms converge to zero so $`\underset{k\mathrm{}}{lim}\mathrm{sup}_{aK}f_kaf_k_{E(,\tau )}\epsilon `$ and since $`\epsilon `$ is arbitrary, the proof is complete. ∎
The next proposition can be found in (Proposition 2.5).
###### Proposition 2.14.
Assume that $`E`$ is $`\alpha `$-convex with constant $`1`$ for some $`0<\alpha 1`$ and satisfies a lower $`q`$-estimate with constant $`1`$ for some finite $`q\alpha `$. If $`k=2q/\alpha `$, then for all $`yE(,\tau )`$, for all projections $`e,f`$ with $`e+f=1`$ and $`\tau (e)<\mathrm{}`$, it follows that
$$eye_{E(,\tau )}^k+eyf_{E(,\tau )}^k+fye_{E(,\tau )}^k+fyf_{E(,\tau )}^ky_{E(,\tau )}^k.$$
The proof of the next lemma is just a notational adjustment of the proof of a lemma from the $`L^1`$-case so we will leave the details to the interested readers.
###### Lemma 2.15.
Let $`\{x_k\}_{k=1}^{\mathrm{}}`$ be a bounded sequence in $`E(,\tau )`$ and $`\{e_n\}_{n=1}^{\mathrm{}}`$ be a decreasing sequence of projections in $``$ such that $`e_n_n0`$. Assume that $`\underset{n\mathrm{}}{lim}\mathrm{sup}_ke_nx_ke_n_{E(,\tau )}=\gamma >0`$ then there exists a subsequence $`\{x_{k_n}\}_{n=1}^{\mathrm{}}`$ so that $`\underset{n\mathrm{}}{lim}e_nx_{k_n}e_n_{E(,\tau )}=\gamma `$.
## 3. Kadec-Pełczyński theorem for symmetric spaces
The main result of the present article is the following theorem.
###### Theorem 3.1.
Let $`E`$ be an order continuous symmetric quasi-Banach function space in $`^+`$ that is $`\alpha `$-convex with constant $`1`$ for some $`0<\alpha 1`$ and suppose that $`E`$ satisfies a lower $`q`$-estimate with constant $`1`$ for some $`q\alpha `$.
Let $`\{x_n\}_{n=1}^{\mathrm{}}`$ be a bounded sequence in $`E(,\tau )`$ then there exists a subsequence $`\{x_{n_k}\}_{k=1}^{\mathrm{}}`$ of $`\{x_n\}_{n=1}^{\mathrm{}}`$, bounded sequences $`\{y_k\}_{k=1}^{\mathrm{}}`$ and $`\{z_k\}_{k=1}^{\mathrm{}}`$ in $`E(,\tau )`$ and a decreasing sequence of projections $`p_k_n0`$ in $``$ such that:
* $`x_{n_k}=y_k+z_k`$ for all $`k1`$;
* $`\{y_k:k1\}`$ is $`E`$-equi-integrable and $`p_ky_kp_k=0`$ for all $`k1`$;
* $`\{z_k\}_{k=1}^{\mathrm{}}`$ is such that $`p_kz_kp_k=z_k`$ for all $`k1`$.
The proof will be divided into several steps. Without loss of generality, we will assume that the sequence $`\{x_k\}_{k=1}^{\mathrm{}}`$ is a subset of the unit ball of $`E(,\tau )`$. Since we are dealing with sequences, we can and do assume without loss of generality that $``$ is countably decomposable (see and the proof of Proposition 2.9 above for the details of such reduction).
Set $`𝒟_1:=\{\{e_n\}_{n=1}^{\mathrm{}}_p:e_n_n0\text{and}\tau (e_1)<\mathrm{}\}`$ and consider
$$\delta :=\mathrm{sup}\left\{\underset{n\mathrm{}}{lim}\mathrm{sup}_ke_n|x_k|e_n_{E(,\tau )}:\{e_n\}_{n=1}^{\mathrm{}}𝒟_1\right\}.$$
As in , one can show that $`\delta `$ is attained and using Lemma 2.15, one can choose a subsequence of $`\{x_n\}_{n=1}^{\mathrm{}}`$ (which we will denote again by $`\{x_n\}_{n=1}^{\mathrm{}}`$) and $`\{e_n\}_{n=1}^{\mathrm{}}𝒟_1`$ such that
(3.1)
$$\delta =\underset{n\mathrm{}}{lim}e_n|x_n|e_n_{E(,\tau )}$$
and
(3.2)
$$\delta ^{}=\underset{n\mathrm{}}{lim}e_n|x_n^{}|e_n_{E(,\tau )}=\mathrm{sup}\{\underset{n\mathrm{}}{lim}\mathrm{sup}_kq_n|x_k^{}|q_n_{E(,\tau )};\{q_n\}_{n=1}^{\mathrm{}}𝒟_1\}.$$
For each $`n1`$, set $`v_n:=x_ne_nx_ne_n`$ and let $`V:=\{v_n:n1\}`$.
###### Lemma 3.2.
There exists a sequence of projections $`\{g_n\}_{n=1}^{\mathrm{}}`$ in $``$ with:
* For every $`n1`$, $`g_n\mathrm{𝟏}e_n`$;
* $`\tau (g_n)<\mathrm{}`$, in particular $`g_n`$ is a finite projection;
* $`g_n^n\mathrm{𝟏}`$.
###### Proof.
The lemma can be obtained inductively. Since $``$ is countably decomposable, there exists $`\phi _0`$ a faithful normal state in $`_{}`$. Since $`\mathrm{𝟏}e_n`$ is a semifinite projection, there exists a sequence of projections $`\{g_j^{(n)}\}_{j=1}^{\mathrm{}}`$ with $`\tau (g_j^{(n)})<\mathrm{}`$ for every $`j1`$ and $`g_j^{(n)}^j\mathrm{𝟏}e_n`$. One can choose $`j_n1`$ such that $`\phi _0(1e_n)\phi _0(g_{j_n}^{(n)})<1/n`$. Set
$$\{\begin{array}{cc}g_n:=g_{j_1}^{(1)}\hfill & \text{for }n=1\hfill \\ g_n=g_{j_n}^{(n)}g_{n1}\hfill & \text{for }n>1.\hfill \end{array}$$
It is easy to verify that $`\{g_n\}_{n=1}^{\mathrm{}}`$ satisfies the requirements of the lemma. ∎
For each $`n1`$, let $`p_n=\mathrm{𝟏}g_n`$. Clearly $`p_n_n0`$, $`\mathrm{𝟏}p_n`$ is a finite projection and $`p_ne_n`$ for each $`n1`$.
###### Lemma 3.3.
For each $`n1`$, the sets $`\left|V(\mathrm{𝟏}p_n)\right|`$ and $`\left|(\mathrm{𝟏}p_n)V\right|`$ are $`E`$-equi-integrable.
###### Proof.
Let us show that for every $`n1`$, $`\left|V(1p_n)\right|`$ is an $`E`$-equi-integrable set. Assume that there exists $`k_01`$ such that $`|V(\mathrm{𝟏}p_{k_0})|`$ is not $`E`$-equi-integrable. By definition, there exists a decreasing sequence of projections $`q_n_n0`$ such that $`\underset{n\mathrm{}}{lim}\mathrm{sup}_{a|V(\mathrm{𝟏}p_{k_0})|}q_naq_n_{E(,\tau )}>0`$, that is
$$\underset{n\mathrm{}}{lim}\mathrm{sup}_mq_n|v_m(\mathrm{𝟏}p_{k_0})|q_n_{E(,\tau )}>0.$$
Choose a strictly increasing sequence $`\{m_n\}_{n=1}^{\mathrm{}}`$ of $``$ such that
$$\underset{n\mathrm{}}{lim}q_n|v_{m_n}(1p_{k_0})|q_n_{E(,\tau )}>0.$$
Let $`u_{n,k_0}`$ be a bounded operator such that $`|v_{m_n}(\mathrm{𝟏}p_{k_0})|=u_{n,k_0}v_{m_n}(\mathrm{𝟏}p_{k_0})`$. We get that
$`q_n|v_{m_n}(\mathrm{𝟏}p_{k_0})|q_n_{E(,\tau )}`$ $`=q_nu_{n,k_0}v_{m_n}(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}`$
$`=q_nu_{n,k_0}[x_{m_n}e_{m_n}x_{m_n}e_{m_n}](\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}.`$
We recall that $`e_{k_0}p_{k_0}`$ and since $`\{e_n\}_{n=1}^{\mathrm{}}`$ is decreasing, for $`m_nk_0`$, $`e_{m_n}p_{k_0}`$ and therefore $`e_{m_n}(\mathrm{𝟏}p_{k_0})=0`$ and since $`q_nu_{n,k_0}_{\mathrm{}}1`$, we obtain that for $`n`$ large enough,
$`q_n|v_{m_n}(\mathrm{𝟏}p_{k_0})|q_n_{E(,\tau )}`$ $`=q_nu_{n,k_0}(x_{m_n})(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}`$
$`x_{m_n}(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}.`$
Using Proposition 2.11, with $`x=x_{m_n}`$ and $`y=(\mathrm{𝟏}p_{k_0})q_n`$, we get
$`q_n|v_{m_n}(\mathrm{𝟏}p_{k_0})|p_n_{E(,\tau )}`$ $`x_{m_n}_{E(,\tau )}^{\frac{1}{2}}q_n(1p_{k_0})|x_{m_n}|(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}^{\frac{1}{2}}`$
$`q_n(\mathrm{𝟏}p_{k_0})|x_{m_n}|(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}^{\frac{1}{2}}.`$
This implies that
$$\underset{n\mathrm{}}{lim\; sup}q_n(\mathrm{𝟏}p_{k_0})|x_{m_n}|(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}>0.$$
Let $`s_n`$ be the left support projection of $`(\mathrm{𝟏}p_{k_0})q_n`$ (this is equal to the right support projection of $`q_n(\mathrm{𝟏}p_{k_0})`$). We have
$`q_n(\mathrm{𝟏}p_{k_0})|x_{m_n}|(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}`$ $`=q_n(\mathrm{𝟏}p_{k_0})s_n|x_{m_n}|s_n(\mathrm{𝟏}p_{k_0})q_n_{E(,\tau )}`$
$`s_n|x_{m_n}|s_n_{E(,\tau )}.`$
By the definition of support projection, $`s_n(\mathrm{𝟏}p_{k_0})`$ for every $`n1`$, so $`\{s_n\}_{n=1}^{\mathrm{}}`$ is a sequence of finite projections. As in proof of Lemma 2.6, we note that $`s_n=q_np_{k_0}p_{k_0}`$ and as before, $`s_nq_nq_np_{k_0}`$. Now since $`q_n_n0`$, $`q_nq_np_{k_0}_n0`$ hence $`\tau (s_n)=\tau (q_nq_np_{k_0})`$ converges to zero which implies that $`s_n_n0`$. Therefore, $`\{s_n\}_{n=1}^{\mathrm{}}𝒟_1`$.
In summary, we get $`\{s_n\}_{n=1}^{\mathrm{}}𝒟_1`$ with $`s_n1p_{k_0}`$ for each $`n1`$ and for some $`\gamma >0`$,
(3.3)
$$\underset{n\mathrm{}}{lim\; sup}s_n|x_{m_n}|s_n=\gamma .$$
Let $`f_n:=s_ne_{m_n}`$.
For each $`m_nk_0`$, $`s_n\mathrm{𝟏}p_{k_0}\mathrm{𝟏}e_{k_0}`$ so $`s_ne_{m_n}`$ hence $`f_n=s_n+e_{m_n}`$. In particular $`\{f_n\}_{n=1}^{\mathrm{}}𝒟_1`$.
Using Proposition 2.14 (it applies since $`\tau (f_n)<\mathrm{})`$,
$`f_n|x_{m_n}|f_n_{E(,\tau )}^{\frac{2q}{\alpha }}`$ $`s_n|x_{m_n}|s_n_{E(,\tau )}^{\frac{2q}{\alpha }}+e_{m_n}|x_{m_n}|s_n_{E(,\tau )}^{\frac{2q}{\alpha }}`$
$`+s_n|x_{m_n}|e_{m_n}_{E(,\tau )}^{\frac{2q}{\alpha }}+e_{m_n}|x_{m_n}|e_{m_n}_{E(,\tau )}^{\frac{2q}{\alpha }}`$
$`s_n|x_{m_n}|s_n_{E(,\tau )}^{\frac{2q}{\alpha }}+e_{m_n}|x_{m_n}|e_{m_n}_{E(,\tau )}^{\frac{2q}{\alpha }}.`$
Taking the limit as $`n`$ tends to $`\mathrm{}`$, one gets from (3.1) and (3.3) that $`\delta ^{\frac{2q}{\alpha }}\gamma ^{\frac{2q}{\alpha }}+\delta ^{\frac{2q}{\alpha }}`$. This is a contradiction since $`\gamma >0`$.
We conclude that for every $`n1`$, the set $`|V(\mathrm{𝟏}p_n)|`$ is an $`E`$-equi-integrable set.
For the case of $`|(\mathrm{𝟏}p_n)V|`$, it is enough to repeat the argument above for $`V^{}(\mathrm{𝟏}p_n)`$ using the definition of $`\delta ^{}`$ (instead of $`\delta `$). Details are left to the reader. This ends the proof of the lemma. ∎
We will proceed to the proof of Theorem 3.1. Consider two cases.
Case 1: Assume that $`V`$ is $`E`$-equi-integrable.
It is enough to set $`y_n=x_{m_n}e_{m_n}x_{m_n}e_{m_n}`$ and $`z_n=e_{m_n}x_{m_n}e_{m_n}`$.
Case 2: Assume that $`V`$ is not $`E`$-equi-integrable.
Proposition 2.13 and Lemma 3.3 imply that there exists $`\nu >0`$ such that
$$\underset{n\mathrm{}}{lim}\mathrm{sup}_{vV}p_nvp_n_{E(,\tau )}=\nu >0.$$
Choose a subsequence $`\{v_{n_k}\}_{k=1}^{\mathrm{}}`$ such that
(3.4)
$$\underset{k\mathrm{}}{lim}p_kv_{n_k}p_k_{E(,\tau )}=\nu >0.$$
For each $`k1`$, let $`w_k:=v_{n_k}p_kv_{n_k}p_k`$ and set
$$W:=\{w_k:k1\}.$$
###### Lemma 3.4.
The set $`W`$ is $`E`$-equi-integrable.
###### Proof.
We note first that if $`kn`$, then $`(\mathrm{𝟏}p_n)w_k=(\mathrm{𝟏}p_n)v_k`$ and $`w_k(1p_n)=v_k(1p_n)`$ so for fixed $`n1`$, $`(\mathrm{𝟏}p_n)W=\{(\mathrm{𝟏}p_n)w_k:k<n\}\{(\mathrm{𝟏}p_n)v_k:kn\}`$. Similarly, $`W(\mathrm{𝟏}p_n)=\{w_k(\mathrm{𝟏}p_n):k<n\}\{v_k(\mathrm{𝟏}p_n):kn\}.`$
Lemma 3.3 implies that for every $`n1`$, both $`|W(\mathrm{𝟏}p_n)|`$ and $`(\mathrm{𝟏}p_n)W`$ are $`E`$-equi-integrable sets. Therefore, if $`W`$ is not $`E`$-equi-integrable, there would be a subsequence $`\{w_{k(j)}\}_{j=1}^{\mathrm{}}`$ of $`\{w_k\}_{k=1}^{\mathrm{}}`$ and $`\epsilon >0`$ such that
(3.5)
$$\underset{j\mathrm{}}{lim}p_jw_{k(j)}p_j_{E(,\tau )}=\epsilon .$$
Using Proposition 2.11 on $`v_{n_{k(j)}}`$ and $`p_j=(p_jp_{k(j)})+p_{k(j)}`$, we obtain:
$`p_jv_{n_{k(j)}}p_j_{E(,\tau )}^{\frac{2p}{\alpha }}`$ $`(p_jp_{k(j)})v_{n_{k(j)}}(p_jp_{k(j)})_{E(,\tau )}^{\frac{2q}{\alpha }}+p_{k(j)}v_{n_{k(j)}}(p_jp_{k(j)})_{E(,\tau )}^{\frac{2q}{\alpha }}+`$
$`(p_jp_{k(j)})v_{n_{k(j)}}p_{k(j)}_{E(,\tau )}^{\frac{2q}{\alpha }}+(p_jp_{k(j)})v_{n_{k(j)}}(p_jp_{k(j)})_{E(,\tau )}^{\frac{2q}{\alpha }}.`$
Taking into account the identities, $`(p_jp_{k(j)})v_{n_{k(j)}}(p_jp_{k(j)})=(p_jp_{k(j)})w_{k(j)}(p_jp_{k(j)})`$, $`(p_jp_{k(j)})v_{n_{k(j)}}p_{k(j)}=p_jw_{k(j)}p_{k(j)}`$ and $`p_{k(j)}v_{n_{k(j)}}(p_jp_{k(j)})=p_{k(j)}w_{k(j)}p_j`$, one can deduce that,
$`p_jv_{n_{k(j)}}p_j_{E(,\tau )}^{\frac{2q}{\alpha }}`$ $`(p_jp_{k(j)})w_{k(j)}(p_jp_{k(j)})_{E(,\tau )}^{\frac{2q}{\alpha }}+p_jw_{k(j)}p_{k(j)}_{E(,\tau )}^{\frac{2q}{\alpha }}+`$
$`p_{k(j)}w_{k(j)}p_j_{E(,\tau )}^{\frac{2q}{\alpha }}+p_{k(j)}v_{n_{k(j)}}p_{k(j)}_{E(,\tau )}^{\frac{2q}{\alpha }}.`$
Let $`C(q,\alpha )`$ be the norm of the identity map from $`\mathrm{}_3^{\frac{2q}{\alpha }}`$ onto $`\mathrm{}_3^\alpha `$, where $`\mathrm{}_3^{\frac{2q}{\alpha }}`$ (resp. $`\mathrm{}_3^\alpha `$) denotes the $`3`$-dimensional $`\mathrm{}^{\frac{2q}{\alpha }}`$-space (resp. $`\mathrm{}^\alpha `$-space). We have
$`p_jv_{n_{k(j)}}p_j_{E(,\tau )}^{\frac{2q}{\alpha }}`$ $`C(q,\alpha )^{\frac{2q}{\alpha }}[((p_jp_{k(j)})w_{k(j)}(p_jp_{k(j)})_{E(,\tau )}^\alpha +`$
$`p_jw_{k(j)}p_{k(j)}_{E(,\tau )}^\alpha +p_{k(j)}w_{k(j)}p_j_{E(,\tau )}^\alpha )^{\frac{1}{\alpha }}]^{\frac{2q}{\alpha }}+`$
$`p_{k(j)}v_{n_{k(j)}}p_{k(j)}_{E(,\tau )}^{\frac{2q}{\alpha }}.`$
We remark that $`p_jw_{k(j)}p_j=(p_jp_{k(j)})w_{k(j)}(p_jp_{k(j)})+p_jw_{k(j)}p_{k(j)}+p_{k(j)}w_{k(j)}p_j`$ and since $`E(,\tau )`$ is $`\alpha `$-convex (with constant $`1`$), the above inequality implies
$$p_jv_{n_{k(j)}}p_j_{E(,\tau )}^{\frac{2q}{\alpha }}C(q,\alpha )^{\frac{2q}{\alpha }}p_jw_{k(j)}p_j_{E(,\tau )}^{\frac{2q}{\alpha }}+p_{k(j)}v_{n_{k(j)}}p_{k(j)}_{E(,\tau )}^{\frac{2q}{\alpha }},$$
and taking the limit as $`j\mathrm{}`$, we get from (3.4) and (3.5) that
$$\nu ^{\frac{2q}{\alpha }}C(q,\alpha )^{\frac{2q}{\alpha }}\epsilon ^{\frac{2q}{\alpha }}+\nu ^{\frac{2q}{\alpha }}.$$
This is a contradiction since $`\epsilon >0`$, so $`W`$ is a $`E`$-equi-integrable set. The lemma is proved.
To complete the proof of Theorem 3.1, we note that $`W=\{x_{n_k}p_kx_{n_k}p_k:k1\}`$ so if we set $`y_k=x_{n_k}p_kx_{n_k}p_k`$ and $`z_k=p_kx_{n_k}p_k`$. The proof of is complete. ∎
###### Remarks 3.5.
(1) If $`=(\mathrm{}^2)`$ with the usual trace, then every projection of finite trace is a finite rank projection so in the proof above, $`\delta =\delta ^{}=0`$. In the particular case of unitary matrix space $`C_E`$ where $`E`$ is a symmetric sequence space, one proceed directly to Case 2 by setting $`W:=\left\{x_np_nx_np_n;n1\right\}`$ where $`\{p_n\}_{n=1}^{\mathrm{}}`$ is an arbitrary sequence of projections satisfying: $`p_n_n0`$ and for every $`n1`$, $`1p_n`$ is a finite projection.
(2) If $``$ is a finite von Neumann algebra with a normalized finite trace $`\tau `$ and $`E`$ is a symmetric space on $`[0,1]`$ satisfying the assumptions of Theorem 3.1, it is enough to take $`p_n=e_n`$ (i.e $`g_n=\mathrm{𝟏}e_n`$ on Lemma 3.2) and conclude immediately as in Lemma 3.3 that $`V`$ is $`E`$-equi-integrable.
(3) In the proof above, it is clear that the projections $`\{p_k\}_{k=1}^{\mathrm{}}`$ are such that either $`\tau (p_1)<\mathrm{}`$ or $`\tau (\mathrm{𝟏}p_k)<\mathrm{}`$ for all $`k1`$. In fact, the argument above shows that if $`\{e_n\}_{n=1}^{\mathrm{}}`$ is a sequence in $`𝒟_1`$ that attained the quantities $`\delta `$ and $`\delta ^{}`$, then any sequence of projections satisfying $`p_n_n0`$, $`e_np_n`$ for each $`n1`$ and $`\tau (\mathrm{𝟏}p_n)<\mathrm{}`$ for each $`n1`$, would satisfy the conclusion of Theorem 3.1.
The following extension shows that if one considers finitely many bounded sequences in $`E(,\tau )`$, one can choose a single sequence of projections that works for each sequence.
###### Corollary 3.6.
If $``$ and $`E`$ are as in Theorem 3.1 and $`\left\{x^{(1)}\right\}_{n=1}^{\mathrm{}},\left\{x^{(2)}\right\}_{n=1}^{\mathrm{}},\mathrm{},\left\{x^{(j_0)}\right\}_{n=1}^{\mathrm{}}`$ be finitely many bounded sequences in $`E(,\tau )`$. Then there exist a strictly increasing sequence $`\{n_k\}_{k=1}^{\mathrm{}}`$ of $``$ and a sequence of decreasing projections $`p_k_k0`$ in $``$ such that for each $`1jj_0`$, the set $`\{x_{n_k}^{(j)}p_kx_{n_k}^{(j)}p_k:k1\}`$ is $`E`$-equi-integrable.
###### Proof.
For $`1jj_0`$, we set, as in the proof of Theorem 3.1,
$$\delta _j:=\mathrm{sup}\left\{\underset{n\mathrm{}}{lim}\mathrm{sup}_ke_n|x_k^{(j)}|e_n_{E(,\tau )}:\{e_n\}_{n=1}^{\mathrm{}}𝒟_1\right\}.$$
One can choose a strictly increasing sequence $`\{n_k\}_{k=1}^{\mathrm{}}`$ in $``$ such that for each $`1jj_0`$, there exists a sequence $`\{e_k^{(j)}\}_{k=1}^{\mathrm{}}𝒟_1`$ with
$$\delta _j=\underset{k\mathrm{}}{lim}e_k^{(j)}|x_{n_k}^{(j)}|e_k^{(j)}_{E(,\tau )}$$
and
$$\delta _j^{}=\underset{k\mathrm{}}{lim}e_k^{(j)}|x_{n_k}^{(j)}|e_k^{(j)}_{E(,\tau )}=\mathrm{sup}\{\underset{n\mathrm{}}{lim}\mathrm{sup}_kq_n|x_{n_k}^{(j)}{}_{}{}^{}|q_n_{E(,\tau )};\{q_n\}_{n=1}^{\mathrm{}}𝒟_1\}.$$
For every $`k1`$, set $`e_k:=_{1jj_0}e_k^{(j)}`$. Since $`\tau (e_k)_{j=1}^{j_0}\tau (e_k^{(j)})`$, it is clear that the sequence $`\{e_k\}_{k=1}^{\mathrm{}}`$ belongs to $`𝒟_1`$ and each of the $`\delta _j`$’s and $`\delta _j^{}`$’s are attained at $`\{e_k\}_{k=1}^{\mathrm{}}`$. One can complete the proof by procceding as in the proof of Theorem 3.1, simmultaneously on the finite set of sequences and the fixed $`\{e_k\}_{k=1}^{\mathrm{}}`$. ∎
Our next result shows that the decreasing projections in the decomposition can be replaced by mutually disjoint projections.
###### Theorem 3.7.
Let $`E`$ be an order continuous quasi-Banach function space as in Theorem 3.1. Let $`\{x_n\}_{n=1}^{\mathrm{}}`$ be a bounded sequence in $`E(,\tau )`$ then there exists a subsequence $`\{x_{n_k}\}_{k=1}^{\mathrm{}}`$ of $`\{x_n\}_{n=1}^{\mathrm{}}`$, bounded sequences $`\{\phi _k\}_{k=1}^{\mathrm{}}`$ and $`\{\zeta _k\}_{k=1}^{\mathrm{}}`$ in $`E(,\tau )`$ and mutually disjoint sequence of projections $`\{e_k\}_{k=1}^{\mathrm{}}`$ such that:
* $`x_{n_k}=\phi _k+\zeta _k`$ for all $`k1`$;
* $`\{\phi _k:k1\}`$ is $`E`$-equi-integrable and $`e_k\phi _ke_k=0`$ for all $`k1`$;
* $`\{\zeta _k\}_{k=1}^{\mathrm{}}`$ is such that $`e_k\zeta _ke_k=\zeta _k`$ for all $`k1`$.
###### Proof.
Let $`\{x_n\}_{n=1}^{\mathrm{}}`$ be a bounded sequence in $`E(,\tau )`$ and suppose (by taking a subsequence if necessary), $`x_n=y_n+z_n`$ with $`p_ny_np_n=0`$, the set $`\{y_n:n1\}`$ is $`E`$-equi-integrable and $`p_nz_np_n=z_n`$ for all $`n1`$, be the decomposition of $`\{x_n\}_{n=1}^{\mathrm{}}`$ as in Theorem 3.1.
Let $`n_1=1`$. Since $`p_n_n0`$ and
$$p_1z_1p_1(p_1p_n)z_1(p_1p_n)=p_nz_1p_1+p_1z_1p_np_nz_1p_n,$$
Proposition reforder-continuity(part(ii)) shows that
$$\underset{n\mathrm{}}{lim}p_1z_1p_1(p_1p_n)z_1(p_1p_n)_{E(,\tau )}=0.$$
Choose $`n_2>n_1=1`$ such that
$$p_1z_1p_1(p_1p_{n_2})z_1(p_1p_{n_2})_{E(,\tau )}<\frac{1}{2}.$$
Inductively, one can construct $`n_1<n_2<\mathrm{}<n_k<\mathrm{}`$ such that
$$p_{n_k}z_{n_k}p_{n_k}(p_{n_k}p_{n_{k+1}})z_{n_k}(p_{n_k}p_{n_{k+1}})_{E(,\tau )}<\frac{1}{2^k}.$$
Since $`z_n=p_nz_np_n`$ for every $`n1`$, one gets
$$z_{n_k}(p_{n_k}p_{n_{k+1}})z_{n_k}(p_{n_k}p_{n_{k+1}})_{E(,\tau )}<\frac{1}{2^k}.$$
For every $`k1`$, set
$`e_k:`$ $`=p_{n_k}p_{n_{k+1}}`$
$`\zeta _k:`$ $`=(p_{n_k}p_{n_{k+1}})z_{n_k}(p_{n_k}p_{n_{k+1}})`$
$`\phi _k:`$ $`=y_{n_k}+[z_{n_k}e_kz_{n_k}e_k].`$
Since $`\{y_{n_k}:k1\}`$ is a $`E`$-equi-integrable set and $`\underset{k\mathrm{}}{lim}z_{n_k}e_kz_{n_k}e_k_{E(,\tau )}=0`$, it is clear that $`\{\phi _k:k1\}`$ is $`E`$-equi-integrable. Also $`\{e_k\}_{k=1}^{\mathrm{}}`$ is mutually disjoint. The proof is complete. ∎
###### Corollary 3.8.
Let $`E`$ be an order-continuous symmetric Banach function space on $`^+`$ with the Fatou propery. Let $`\{x_n\}_{n=1}`$ be a bounded sequence in $`E(,\tau )`$ then there exists a subsequence $`\{x_{n_k}\}_{k=1}^{\mathrm{}}`$ of $`\{x_n\}_{n=1}^{\mathrm{}}`$, bounded sequences $`\{\phi _k\}_{k=1}^{\mathrm{}}`$ and $`\{\zeta _k\}_{k=1}^{\mathrm{}}`$ in $`E(,\tau )`$ and mutually disjoint sequence of projections $`\{e_k\}_{k=1}^{\mathrm{}}`$ such that:
* $`x_{n_k}=\phi _k+\zeta _k`$ for all $`k1`$;
* $`\{\phi _k:k1\}`$ is $`E`$-equi-integrable and $`e_k\phi _ke_k=0`$ for all $`k1`$;
* $`\{\zeta _k\}_{k=1}^{\mathrm{}}`$ is such that $`e_k\zeta _ke_k=\zeta _k`$ for all $`k1`$.
###### Proof.
Assume that $`E`$ has the Fatou property (equivalenty $`E`$ does not contain $`c_0`$). Since $`E`$ is symmetric, $`Ec_0`$ is equivalent to $`E`$ not containing $`\mathrm{}_{\mathrm{}}^n`$ uniformly, and therefore $`E`$ satisfies the $`q`$-lower estimate for some $`q`$ and one can renorm $`E`$ so that it satisfies the lower $`q`$-estimate of constant $`1`$. All of these facts can be found in . ∎
The proof of Theorem 3.1 can be adjusted to obtain decompositions where the projections are taken only on one side, that is, the following result follows:
###### Corollary 3.9.
Let $`E`$ be an order continuous quasi-Banach function space in $`^+`$ that is $`\alpha `$-convex with constant $`1`$ for some $`0<\alpha 1`$ and suppose that $`E`$ satisfies a lower $`q`$-estimate with constant $`1`$ for some $`q\alpha `$. Let $`\{x_n\}_{n=1}^{\mathrm{}}`$ be a bounded sequence in $`E(,\tau )`$ then there exist a subsequence $`\{x_{n_k}\}_{k=1}^{\mathrm{}}`$ of $`\{x_n\}_{n=1}^{\mathrm{}}`$, bounded sequences $`\{y_k\}_{k=1}^{\mathrm{}}`$ and $`\{z_k\}_{k=1}^{\mathrm{}}`$ in $`E(,\tau )`$ and decreasing projections $`e_k_n0`$ in $``$ such that:
* $`x_{n_k}=y_k+z_k`$ for all $`k1`$;
* $`e_ky_k=0`$ for all $`k1`$ and $`\underset{n\mathrm{}}{lim}\mathrm{sup}_{k1}f_ny_k_{E(,\tau )}=0`$ for every $`f_n_n0`$.
* $`\{z_k\}_{k=1}^{\mathrm{}}`$ is such that $`e_kz_k=z_k`$ for all $`k1`$.
###### Definition 3.10.
A subspace $`X`$ of $`L^p(,\tau )`$ is called strongly embedded into $`L^p(,\tau )`$ if the $`L^p`$ and the measure topologies on $`X`$ coincide.
The following result is a direct application of Proposition 2.9 and Theorem 3.7.
###### Theorem 3.11.
Let $`1p<\mathrm{}`$. Every subspace of $`L^p(,\tau )`$ either contains almost isometric copies of $`\mathrm{}^p`$ or is strongly embedded in $`L^p(,\tau )`$.
The next corollary should be compared with \[18, Theorem 2.4\].
###### Corollary 3.12.
Assume that $``$ is finite and $`p>2`$. Every subspace of $`L^p(,\tau )`$ either contains almost isometric copies of $`\mathrm{}^p`$ or is isomorphic to a Hilbert space.
For the commutative case, the space $`\mathrm{}^p`$ can not be strongly embedded in $`L^p[0,1]`$ for $`0<p<2`$. This is due to Kalton for $`0<p<1`$ and Rosenthal for the case $`1p<2`$ (see also for another approach). A non-commutative analogue should be of interest.
Problem: Let $``$ be a semifinite von Neumann algebra and $`0<p<2`$. Does $`\mathrm{}^p`$ strongly embed into $`L^p(,\tau )`$?
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# New Neighbors from 2MASS: Activity and Kinematics at the Bottom of the Main Sequence
## 1 Introduction
Catalogs of nearby stars (Gliese & Jahreiß, 1991; Kirkpatrick, Henry, & Simons, 1995; Reid, Hawley, & Gizis, 1995) and high proper motion stars (Luyten, 1979) are grossly deficient in very low mass (VLM) dwarfs. With spectral types of M7 and later, these objects, sometimes called “ultracool M dwarfs,”<sup>1</sup><sup>1</sup>1 All spectral types in this paper are on the Kirkpatrick, Henry, & McCarthy (1991) M dwarf and Kirkpatrick et al. (1999b) L dwarf systems. L dwarfs are cooler than “ultracool” M dwarfs. are so optically faint that even nearby ones eluded searches based on the older (pre-1980s) sky surveys. These dwarfs have particular importance because they lie at or below the hydrogen burning limit – and have proven not only to be estimators of the numbers of dark brown dwarfs, but also present interesting astrophysical challenges in their own right.
The new generation of sky surveys allows this deficiency to be addressed and large samples of nearby VLM dwarfs to be identified. The Two Micron All-Sky Survey<sup>2</sup><sup>2</sup>22MASS data and documentation are available at http://www.ipac.caltech.edu/2mass (Skrutskie et al., in prep.; hereafter 2MASS) provides reliable photometry in the JHK<sub>s</sub> passbands, close to the peak of emission for these cool dwarfs. Furthermore, the Second Palomar Sky Survey (Reid et al., 1991), hereafter POSS II) provides B<sub>J</sub>,R<sub>F</sub>, and I<sub>N</sub> photographic photometry in the northern hemisphere. In the southern hemisphere, the UK Schmidt and ESO sky survey plates provide $`B_J`$ and $`R_F`$ magnitudes. In sum, it is becoming possible to identify both the least luminous stars and young massive brown dwarfs by their optical and near-infrared colors alone over most of the sky.
We present first results of a search using near-infrared and optical sky survey data aimed at completing the nearby star catalog for the ultracool M dwarfs. We discuss the sample selection and spectroscopic followup in Section 2. Although the sample discussed in this paper includes only a small fraction of the total population of nearby ultracool dwarfs, it represents a fourfold increase in the number of such sources known. We discuss some preliminary results concerning the statistical properties of these sources in the latter sections of this paper. We discuss some stars of special interest in Section 3 and the 2MASS colors of ultracool M dwarfs in Section 4. The local space density of VLM dwarfs is discussed in Section 5, their activity and kinematics are discussed in Section 6, and finally our conclusions and future prospects are discussed in 7.
## 2 Data
### 2.1 Sample Selection
Our results are based on three observational samples. For our initial observing run, in July 1998, we used both photometric and proper motion criteria to define a sample of candidate VLM dwarfs. Based on the results from this run and further experience analyzing 2MASS data, we were able to improve our selection criteria for our December 1998 and subsequent observing runs. There are thus three samples with different properties, and when necessary we distinguish them as “Sample A” (July 1998), “Sample B” (December 1998), and “Sample C” (June 1999) respectively. Samples B and C have nearly identical selection criteria, and when combined are refered to as Sample BC. In all samples, objects within 20 degrees of the Galactic plane were excluded. Unless otherwise stated, all the analysis in Section 5 and 6 is based on Sample BC.
Sample A was based on 2MASS data processed by July 1998. All these data were obtained at the Mt. Hopkins 2MASS telescope and a total of 363 square degrees were searched. All objects classified as extended by the 2MASS pipeline were eliminated (Jarrett et al. (2000)). The 2MASS data were correlated with the USNO’s PMM scans of the POSS II plates, using a preliminary version of the software used by the 2MASS Rare Objects Core Project (this software will be described in more detail in a future publication by Monet et al.). This provided three additional colors: B<sub>J</sub>, R<sub>F</sub>, and I<sub>N</sub>. Zero-point calibrations for the POSS II scans were not available, but we have since found that rough zero-points (good to $`\pm 0.5`$ magnitudes) are $`B=B_{Jinst}+2`$, $`R_C=R_{Finst}1`$, and $`I_C=I_{Ninst}1`$. These do not account for plate-to-plate variations or the significant color terms expected in $`B_J`$ and $`R_F`$.
All objects that met the following criteria were observed:
1. $`9.0`$ K$`{}_{s}{}^{}13.0`$
2. $`0.95`$ J-K$`{}_{s}{}^{}<1.30`$
3. no B<sub>J</sub> detection
4. Significant ($`>2\sigma `$) proper motion
Objects were also observed if they satisfied:
1. $`9`$ K$`{}_{s}{}^{}12.0`$
2. $`0.95`$ J-K$`{}_{s}{}^{}<1.30`$
3. B<sub>J</sub> detection
4. Significant ($`>2\sigma `$) proper motion
The proper motion criterion requires elaboration. The magnitude of the observed positional offset between the 2MASS and POSS II source is compared to the distribution of all 2MASS-POSS II correlations. Only objects with significant proper motion were selected. Thus, while we selected sources that are moving with high confidence (roughly $`2\sigma `$), the actual cutoff in terms of arcseconds per year depends on the epoch difference between the F plate and 2MASS, which varies between 0 and $`10`$ years.
As shown in Section 2.2, these criteria led to the identification of a number of new nearby M dwarfs, but they are flawed in some respects. First, the POSS II and UKST IIIa J plates are sufficiently sensitive that many nearby (bright) VLM dwarfs are detected. Moreover, the blue cutoff in J-K<sub>s</sub> allowed mid-M dwarfs to enter the sample. We selected this blue cutoff because the prototypical M7 dwarf VB 8 has J-K$`=0.95`$ (Leggett, 1992). However, M dwarfs between spectral type M0 and M6.5 all have J-K$`{}_{s}{}^{}0.9`$, and, with uncertainties of $`\sigma _{JK}0.04`$ mag, significant numbers are scattered to J-K$`{}_{s}{}^{}>0.95`$. Consequently, our initial attempts to select ultracool M dwarfs produced a sample which is heavily contaminated by distant mid-M dwarfs. As it turned out, there were no 2MASS sources meeting our color and magnitude cuts that were not paired with a POSS II source in Sample A.
Based on experience gained from this analysis, we revised our sample selection for the December 1998 observing run. Because a larger area was available, we focused on brighter VLM dwarfs which should lie within $`20`$ parsecs. The Sample B criteria are:
1. K$`{}_{s}{}^{}12.0`$
2. J-K$`{}_{s}{}^{}1.00`$
3. R<sub>F</sub> \- K$`{}_{s}{}^{}>3.5`$ or I<sub>N</sub> \- K$`{}_{s}{}^{}2.0`$
4. $`\delta <+30\mathrm{deg}`$
5. $`\alpha <13^h00^m`$ or $`\alpha >20^h00^m`$
6. J-H $`\frac{4}{3}`$ H-K $`+0.25`$
No selection based on proper motion was applied, and therefore Sample B is kinematically unbiased. A total of 2977 sq. degrees were searched. The J-K<sub>s</sub> cutoff excludes early and mid-M dwarfs, but also some of the bluer M7 dwarfs. Nearby bright L dwarfs, however, are included in these samples, since we impose no red cutoff. Both samples are magnitude selected, and therefore are biased towards overluminous stars and unresolved near-equal luminosity binaries, although binaries with separations of a few arcseconds may be excluded by the extended source provision. The position selection is due to the requirement that the objects be observable from Las Campanas in December. Ultracool M dwarfs have R-K$`{}_{s}{}^{}>5.5`$ and I-K$`>3.4`$ (Leggett, 1992), but we used a more liberal selection to allow for uncertainties in the calibration of the photographic magnitudes. In Section 4, we show that R-K$`{}_{s}{}^{}>4.9`$ includes all M8 and later dwarfs. The J-H,H-K cut excludes M giants.
Sample C was selected for our June 1999 Kitt Peak observing run. The selection was identical to the Sample B selection, except that a different area was covered and only objects with $`RK_s>4.9`$ were selected. Processed data which lay within the following limits were selected:
1. $`\alpha >11^h00^m`$ and $`\delta >+6\mathrm{°}`$
2. $`16^h20^m<\alpha <23^h35^m`$ and $`36<\delta <+6\mathrm{°}`$
### 2.2 Spectroscopy and Data Analysis
Sample A was observed on UT dates July 30 – August 1 1998 using the Double Spectrograph and the Hale 200-in. telescope during an observing run that was primarily devoted to our ongoing spectroscopic survey of Luyten high proper motion stars (Gizis & Reid, 1997). The wavelength coverage included 6290 to 8800 Å at a resolution of 3 Å.
Sample B was observed on December 2 – 7 1998 using the Modular Spectrograph on the Las Campanas 100-in telescope. The “Tek 5” chip, a 2048-square CCD with $`24\mu `$ pixels, was used with a 600 l/mm grating blazed at 7500 Å. The useful wavelength range of the spectra was 6100 - 9400 Å at a resolution of 6 Å. A few targets (including all three objects with J-K$`{}_{s}{}^{}>1.3`$, which had been previously identified as L dwarf candidates) were observed during Keck observing runs (see Kirkpatrick et al. 1999b, hereafter K99, and Kirkpatrick et al. 2000, hereafter K00) using LRIS (Oke et al., 1995). The resolution was 9 Å with wavelength coverage from 6300 to 10100 Å. Observations of the flare dwarf 2MASSI J0149089+295613 have already been described in Liebert et al. (1999). The H$`\alpha `$ activity levels adopted here are the average quiescent values. The known object BRI 1222-1222 was not observed, and we rely on the spectroscopic observations reported by Kirkpatrick, Henry, & Simons (1995) and Tinney & Reid (1998, hereafter TR). 2MASSW J0354013+231633 has been previously published as 2MASP J0354012+231635 (Kirkpatrick, Beichman, & Skrutskie, 1997) but the spectral observations presented here are independent.
Sample C was observed June 22 – 23 1999 using the R.C. spectrograph and the Kitt Peak 4m telescope. The wavelength coverage was 6140 Å to 9200 Å using the 2048 CCD, but the extreme ends of the spectra were out of focus. A few objects were observed at Palomar 200-in. in May 1999. The four new L dwarfs identified with the Kitt Peak data were reobserved at higher signal-to-noise using Keck in July 1998. The spectral measurements used for classification on the K99 system are given in Table 2.
All spectra were extracted and flux calibrated using IRAF. M dwarf spectral types were measured by overplotting dwarfs of known spectral type, and should be good to $`\pm 0.5`$ subclasses. All L dwarfs have Keck observations and were classified as in K99 and K00. A few M dwarfs in Table 1 have classifications that differ by 0.5 subclasses from previously published values – we leave our values unalterred as representative of our uncertainties. H$`\alpha `$ fluxes were measured assuming the data were photometric (there were no clouds for our observations in December 1998). We assume that BC$`{}_{K}{}^{}=3.2`$ as derived by Tinney, Mould, & Reid (1993) for ultracool M dwarfs. We believe the H$`\alpha `$ fluxes should be viewed with caution since slit losses and the high airmass of many of the observations will increase the uncertainties (the spectrograph was not adjusted to the parallactic angle). However, our derived H$`\alpha `$ to bolometric luminosity ratios are consistent with those of TR, and in any case the H$`\alpha `$ emission strength is variable in these dwarfs. In Table 3, we list our measured and derived parameters for the ultracool dwarfs in Sample BC.
Proper motions were estimated by measuring positions off the DSS (POSS I/UKST J) and XDSS (POSS II/UKST R) images of the photographic sky surveys. When a second epoch photographic sky survey was not available, we used the 2MASS images instead. All motions reported are relative to other stars in the field, but the correction to absolute motions is negligible compared to other sources of error in the kinematics. The proper motion reported for LHS2397a is from Luyten (1979) and BRI1222-1222 is from Tinney (1996). Note that all of our targets are visible on the DSS images, but most were previously unrecognized. The targets which had no POSS II pairings in our initial processing all proved to have high proper motions. They are visible on the XDSS but lie outside the 8 arcsecond search radius employed in cross-referencing against the photographic data. Since the 2MASS positions are highly accurate, and both the 2MASS images and DSS images are (or will shortly be) easily accessible electronically, we are not presenting finding charts.
Using the available parallaxes for late-type dwarfs (Monet et al., 1992; Tinney et al., 1995; Tinney, 1996; Kirkpatrick et al., 1999b), we find (Figure 1) the linear fit $`M_K=7.593+2.25\times `$ J-K<sub>s</sub> This fit is only valid for M7 and later dwarfs over the color range $`1.0`$ J-K$`{}_{S}{}^{}1.6`$, and should be modified as more L dwarf parallaxes are measured at USNO. The observed scatter is $`\sigma =0.36`$ magnitudes. The distances and tangential velocities derived using this estimate are listed in Table 3. We caution, however, that the distances derived for many of the M7 and M7.5 dwarfs may be underestimated in this paper. As can be seen in Figure 1, the main sequence bends sharply at spectral type M7. Our selection of only targets with J-K$`{}_{s}{}^{}1.0`$ tends to select M6’s and M7’s with overestimated colors, while a spectral classification error of only 0.5 subclasses from M6.5 to M7.0 leads to a large error in $`M_K`$.
## 3 Stars of Special Interest
A few of the targets deserve special comment. Our search has identified one very nearby star and a number of very high proper motion dwarfs. The M8.0 dwarf 2MASSW J0027559+221932 has a photometric parallax that places it within ten parsecs; given the uncertainties, it may lie within the eight parsec sample. We note also that a number of dwarfs in Table 3 have motions greater than 1 arcsecond per year, but do not appear in the LHS Catalog even though they are visible on the POSS plates.
Seven L dwarfs are part of our Sample BC. The L5 dwarf 2MASSW J1507476-162738, also in Reid et al. (2000), was selected for this project but does not lie in the Sample BC area. Additional observations and discussion of 2MASSI J0746425+200032, 2MASSW J0036159+182110 and 2MASSW J1439283+192915 are given in Reid et al. (2000). 2MASSW J1439283+192915 is in the original K99 paper while 2MASSW J0036159+182110 was observed at Keck for the K00 paper. 2MASSW J1300425+191235 is particularly surprising: it has the J-K<sub>s</sub> color of an M8 dwarf but has an L1 spectrum. It will be discussed further in an additional paper, but we note here that the estimated distance and tangential velocity is based on the J-K<sub>s</sub> color and should be viewed with great caution. Nevertheless, it apparently has a high velocity, and is likely to be old. While we report photometric distance estimates only, most of these dwarfs are on the USNO parallax program and accurate distances should be forthcoming.
Since our selection is based upon photometry only, we are sensitive to wide binary pairs where the 2MASS observations of the secondaries are unaffected by the primaries. Two M dwarf secondaries that do not meet the spatial restrictions and were specially observed have been reported in Gizis et al. (1999b). One of the Sample C ultracool M dwarfs also appears to be a secondary. The M8.0 dwarf 2MASSW J2331016-040618 is 447 arcseconds west and 65 arcseconds south of the F8 dwarf HD 221356. Our photometric parallax of 26.3 parsecs for the M dwarf is consistent with Hipparcos trigonometric parallax of 26.24 parcsecs (Perryman et al., 1997), as are the observed proper motions. This is apparently a wide binary system with separation of 0.057 parsecs. The F8 primary may provide a useful age and composition constraint on the M dwarf.
Two of our sources have been previously identified as candidate Hyades members. LP 475-855 has been discussed as an Hyades candidate by Eggen (1993) although it was rejected as too bright by Leggett, Harris, & Dahn (1994). Our photometry supports the latter conclusion, although conceivably it could be a foreground (escaping?) Hyades member if it is an unresolved equal-mass binary. Our initial observation of this star found it in a flare state, with H$`\alpha `$ equivalent width of 40 Å. A 25 December 1998 Keck spectrum found H$`\alpha `$ of only 7 Å, which is the value we report in Table 1. LP 415-20 (Bryja 262) has been extensively discussed as a Hyades member and has been classified as an M6.5 dwarf (Bryja, Humphreys, & Jones, 1994). The difference in spectral types is within our uncertainties, and we note that VO is visible in Bryja’s plot of the spectrum. The poor distance estimate for this object is consistent with our belief that our M7 distances based on J-K<sub>s</sub> colors are unreliable and should be viewed with caution.
## 4 Colors
In Figure 2, we plot the R-K<sub>s</sub>,J-K<sub>s</sub> diagram for Sample B. We concluded that adjusting our color criterion to R<sub>F</sub>-K$`{}_{s}{}^{}>4.9`$ would increase our selection efficiency without losing M8 and later dwarfs, and we adopted this selection criterion for Sample C. This observation is consistent with our estimate that our simple $`R_F`$ zero-point is good to $`\pm 0.5`$ magnitudes.
Our observations show a good correlation between the far-red spectral type and the 2MASS near-IR colors. In Figure 3, we plot the observed J-K<sub>s</sub> color distribution as a function of spectral type. As expected, the M7.0 and M7.5 distributions are truncated by our requirement that J-K$`{}_{s}{}^{}1.0`$. The histograms suggest that of order one M8-M8.5 dwarf may be expected to be missed due to this requirement, and it appears that it is very unlikely than any (normal) M9 dwarfs are missed. 2MASS photometry for these bright sources is expected to be good to 0.03 magnitudes.
In Figure 4, we plot the 2MASS near-infrared color-color diagram for our sample. It is evident that the M/L dwarf sequence lies well below our imposed J-H,H-K cut, and therefore we are not missing ultracool M dwarfs due to this criteria. We fit the relation $`JK_s=(0.146\pm 0.117)+(1.238\pm 0.263)\times HK_s`$ assuming the errors in each color are 0.042 and including the M7.0 to M9.5 dwarfs. This relation may be convenient as a representative sequence in the 2MASS color system. It is interesting to note that the dwarfs around (H-K,J-H) $`=(0.45,0.62)`$ are nearly all classified as M7-M7.5 dwarfs, while the dwarfs at (H-K,J-H) $`=(0.42,0.7)`$ but with similar J-K<sub>s</sub> colors are nearly all classified as M8-M8.5 dwarfs. This may reflect some relation between the red-optical region (dominated by TiO and VO, and influenced by dust) and the IR colors (dominated by H<sub>2</sub>O and H<sub>2</sub>, and also influenced by dust) – or it is due to some subtle bias in our classifications or photometry (for example, perhaps we tend to select M7 dwarfs whose H-K<sub>s</sub> color has been overestimated). Note that one of the outliers below the normal J-H,H-K<sub>s</sub> relation is the peculiar L dwarf 2MASSW J1300425+191235.
## 5 Luminosity Function
Our Sample BC is the first large sample of bright, photometrically selected ultracool M dwarfs. Using our data and derived distances, we can estimate the luminosity function using Schmidt’s (1968) $`V/V_{max}`$ technique. The space density is
$$\mathrm{\Phi }=\frac{1}{V_{max}}$$
$$V_{max}=\frac{\mathrm{\Omega }}{3}\left(10.0^{(K_{lim}M_K+5.0)/5.0}\right)^3$$
In our case, K$`{}_{lim}{}^{}=12.0`$ and $`\mathrm{\Omega }=6040`$ sq. degrees. The corresponding variance is
$$\sigma _\mathrm{\Phi }^2=\frac{1}{V_{max}^2}$$
The space densities are given in Table 4. We derive a space density of $`0.0045\pm 0.0008`$ ultracool M dwarfs per cubic parsec. According to our adopted color-magnitude relation, this is for dwarfs in the range $`9.8<M_K<10.8`$. Therefore, the corresponding luminosity function bin is $`\mathrm{\Phi }(M_K=10.3)=0.0048\pm 0.0009`$ dwarfs per cubic parsec per K magnitude. Since Tinney, Mould, & Reid (1993) have shown that $`BC_K=3.2`$ for these late dwarfs, this may be represented in bolometric magnitudes as $`\mathrm{\Phi }(M_{bol}=13.5)=0.0048\pm 0.0009`$ dwarfs per cubic parsec per bolometric magnitude. We note, however, that this value excludes the M7 dwarfs which will also contribute near M$`{}_{K}{}^{}=9.8`$, so our value is a lower limit. The space density for the early L dwarfs is half that of the ultracool M dwarfs, although we caution that the distance estimate for 2M1300 may be incorrect.
Schmidt’s statistic measures the uniformity of the density distribution of a sample, effectively providing an estimate of sample completeness. For a uniform sample, $`<V/V_{max}>=0.5`$ with an uncertainty of $`\frac{1}{\sqrt{12N}}`$ where N is the number of stars observed. Both our L dwarf and M8.0-M8.5 sample lie within 1 $`\sigma `$ of this value, suggesting that we are complete. The value for M9.0 to M9.5 dwarfs is more problematic, indicating that we have either excluded a few nearby, very bright M9 dwarfs, or that there happen to be no such very nearby dwarfs in our survey volume.
Our space density for the ultracool M dwarfs is consistent with Tinney’s (1993) value of $`\mathrm{\Phi }(M_{bol}=13.5)=0.0076\pm 0.0031`$ dwarfs per cubic parsec per bolometric magnitude, which was based on selection with R<sub>F</sub> and I<sub>N</sub> photographic magnitudes but K followup of all VLM dwarfs to improve photometric parallaxes. Only 6 dwarfs contributed to this bin, accounting for Tinney’s larger uncertainty relative to our sample. Delfosse et al. (1999) have analyzed the DENIS Mini-survey and found 19 M8 and later dwarfs, including 3 L dwarfs. They do not estimate the M dwarf space density, but use the three L dwarfs to estimate $`\mathrm{\Phi }(M_{bol}=15.3)0.011\pm 0.006`$ dwarfs per cubic parsec per bolometric magnitude. We note that their estimated $`M_K`$ for the ultracool M dwarfs are inconsistent with our adopted values (and Figure 1) since they consider their M8-M9 dwarfs to have $`M_K>11`$.
Malmquist bias will affect our sample. Stobie, Ishida & Peacock (1989) have shown that the luminosity function will be overestimated by:
$$\frac{\mathrm{\Delta }\mathrm{\Phi }}{\mathrm{\Phi }}=0.5\sigma ^2\left[\left(0.6\mathrm{ln}10\right)^21.2\mathrm{ln}10\frac{\mathrm{\Phi }^{}}{\mathrm{\Phi }}+\frac{\mathrm{\Phi }^{}}{\mathrm{\Phi }^{\prime \prime }}\right]$$
A model luminosity function can be used to derive the first and second derivatives. The scatter of parallax stars about the linear fit adopted here is $`\sigma =0.36`$. Adopting this value for $`\sigma `$, and making the assumption that the luminosity function is flat, we find that $`\frac{\mathrm{\Delta }\mathrm{\Phi }}{\mathrm{\Phi }}=0.21`$: i.e., the values we derive overestimate the true space densities by $`20\%`$. Since this is only a preliminary sample, we defer further analysis of the Malquist bias, as well as the effect of unresolved binaries, until additional data are available. In the long term, trigonometric parallaxes and searches for companions will allow this issue to be addressed directly.
Our derived space density can best be compared to the luminosity function of nearby stars. Figure 5 plots the Reid & Gizis (1997) ($`\delta >30\mathrm{°}`$) luminosity function for stars within eight parsecs with our M8.0-L4.5 data point added to it. Our results suggest that the dropoff seen in for the faintest ($`M_K>10`$) dwarfs in the eight parsec sample is in part due to incompleteness. <sup>3</sup><sup>3</sup>3Note that, as seen in Reid & Gizis’s Fig. 2, the oft-used 5.2 parsec sample shows the same feature, albeit with less significance due to the very small volume. Applying standard mass-luminosity relations to $`\mathrm{\Phi }(M_K)`$ derived from the 8-parsec sample implies a turnover in the mass function close to the hydrogen-burning limit. There is, however, no reason to expect the star formation process to be cognizant of the mass limit for hydrogen burning. Reid et al. (1999a) have modelled a sample of twenty 2MASS and and three DENIS L dwarfs and conclude that the substellar mass function is consistent with an extension of the power-law matched to data for stars with masses between 0.1 and $`1M_{}`$. The higher space densities measured for ultracool dwarfs in this paper suggest a greater degree of continuity across the stellar/substellar boundary. Continuation of the present survey should identify the missing dwarfs within eight parsecs.
If compared to the classical “photometric” luminosity functions (Stobie, Ishida & Peacock, 1989; Tinney, 1993) which have a peak at $`M_{bol}=10`$ and a dropoff to $`M_{bol}=12`$, then our data would imply a rise in the luminosity function at the stellar/substellar boundary. However, this peak and dropoff are an artifact of the data analysis due to the incorrect assumption of a linear color-magnitude relation (Reid & Gizis, 1997) and/or other systematic errors such as unresolved binaries (Kroupa, Tout & Gilmore, 1993). We believe that the nearby-star sample is a better comparison for our sample, and we emphasize again that the M dwarfs here identified require follow-up trigonometric parallax determinations and high-resolution imaging and radial velocity searches for companions to produce a definitive luminosity function.
## 6 Activity and Kinematics
### 6.1 Review
The BC sample was not selected on the basis of proper motions, and therefore is (relatively) unbiased in terms of kinematics.<sup>4</sup><sup>4</sup>4 The existence of age-luminosity, age-metallicity, metallicity-luminosity, age-activity, and age-kinematics correlations implies that there may be kinematic and/or activity bias due to our Malquist-type luminosity bias. If the distances are underestimated due to bias, then the estimated tangential velocities will also be biased. It is useful to review the properties of nearby disk stars and the already-known properties of ultracool M dwarfs before discussing our kinematic and activity measurements.
Stars are born with low space velocity dispersions and and high chromospheric activity levels. Over time, the space velocity of the stars increases as they interact with the Galactic disk. Using BAFGK dwarfs with known ages, Wielen (1974) showed that the total space velocity increases from $`\sigma _{tot}=19`$ km/s at mean ages of 0.4 to $`0.9\times 10^9`$ yr to 34 km/s at $`2\times 10^9`$ yr to $`48`$ km/s at $`5\times 10^9`$ yrs. The high chromospheric activity levels of young stars with convective envelopes is attributed to a dynamo which is driven by rotation. As the star ages, angular momentum loss through the the stellar wind spins down the star, causing the chromospheric activity to in turn decrease. Wielen (1974) did indeed find that Ca II emission line strength in M dwarfs was related to kinematics, with older stars showing less activity and higher space velocities. As progressively less massive (later spectral type) stars are considered, the observed frequency of high H$`\alpha `$ activity increases. This is not due only to the fact that H$`\alpha `$ emission is more detectable against the cool photosphere – Hawley, Gizis, & Reid (1996, hereafter HGR) showed that the percentage of highly active M dwarfs increases with cooler spectral types even when H$`\alpha `$ activity is compared to the star’s bolometric luminosity. The increased lifetime of activity is confirmed by observations of open clusters Hawley, Reid, & Tourtellot (1999). The connection to rotation in early M dwarfs is confirmed by Delfosse et al. (1998), who have found that the incidence of rapid rotators is higher among cooler M dwarfs, and that the rapid rotators are active. There is some evidence that the rotation-activity relation is breaking down (Hawley, Reid, & Gizis, 1999). In summary, mid-M dwarfs maintain H $`\alpha `$ emission for billions of years as they slowly spin down. Kinematics are a good age indicator, but only in a statistical sense – individual stars can not be accurately dated by the velocities.
For the ultracool M dwarfs, there is considerable evidence that the standard stellar age-activity and rotation-activity relations no longer apply. Basri & Marcy (1995) found that the M9.5 dwarf BRI 0021-0214 had very rapid rotation ($`v\mathrm{sin}i=40`$ km/s) but little H$`\alpha `$ activity. TR have shown that the lithium M9 brown dwarf LP 944-20 (Tinney, 1998) is also a member of this class of “inactive, rapid rotators” as are the two of the DENIS L dwarfs (Martín et al., 1997). TR argue that observations of these field objects as well as open clusters indicate that the violation of the age-activity connection is primarily correlated with mass (the physical mechanisms remain unknown). Basri (1999) reports that rapid rotation is more common among objects of lower luminosity, and proposes that the H$`\alpha `$ activity is powered by a turbulent dynamo that is quenched at high rotation rates. There is some evidence that even among the “inactive, rapid rotators” $`\mathrm{log}\left(\frac{L_{H\alpha }}{L_{bol}}\right)`$ is related to age, as it decreases from $`4.6`$ for the Pleiades to $`5.2`$ for $`0.51.0\times 10^9`$ yr brown dwarfs. This trend may be true for younger ages, since the very low mass ($`0.010.06M_{}`$), very young ($`<10`$ Myr) M8.5 brown dwarf $`\rho `$ Oph 162349.8-242601 has $`EW_{H\alpha }>50`$Å (Luhman, Liebert & Rieke, 1997; Martín, Basri, & Zapatero Osorio, 1999), while the young ($`1`$ Myr), possible brown dwarfs ($`M0.07M_{}`$) V410 Tau X3 (M8.5) and X6 (M6) have $`EW_{H\alpha }15`$ Å (Martín, Basri, & Zapatero Osorio, 1999). Both imply higher activity levels than their older Pleiades and field counterparts, though it should be noted that they are also lower mass. In contrast to the “inactive, rapid rotators’,” some ultracool M dwarfs do show H$`\alpha `$ emission, but they have lower rotation rates ($`20`$ km/s, TR). However, even a rotation rate of $`5`$ km/s is adequate to maintain H$`\alpha `$ emission in mid-M dwarfs (Bopp & Fekel, 1977), so the limits on the “low” rotation rates of these ultracool M dwarfs are not surprising by comparison. A so-far unique field object is the M9.5 dwarf PC 0025+0447 (Schneider et al., 1991; Martín, Basri, & Zapatero Osorio, 1999), whose quiescent H$`\alpha `$ emission ($`300`$ Å) is comparable to highly active mid-M dwarfs in terms of $`\frac{L_{H\alpha }}{L_{bol}}`$. The nature of this object and its emission is uncertain — Martín, Basri, & Zapatero Osorio (1999) argue that this object is a very young brown dwarf, suggesting that the USNO parallax indicating ordinary ultracool M dwarf luminosity is incorrect. Some of the known M8-M9 dwarfs are definitely brown dwarfs. The “inactive, rapid rotator” M9 dwarf LP 944-20 has lithium absorption and a luminosity that indicates it has a mass between 0.056 and 0.064 $`M_{}`$ (Tinney, 1998). The Pleiades brown dwarfs Teide 1 and Calar 3 have spectral types of M8, lithium absorption, and anomalous VO and Na features due to low-surface gravity (Martín, Rebolo, & Zapatero-Osorio, 1996).
### 6.2 The Properties of M and L dwarfs
The BC sample of ultracool M dwarfs provides the first opportunity for a thorough investigation of the distribution of activity in these VLM dwarfs. In Figure 6, we compare the percentage of ultracool M dwarfs observed in emission to the HGR statistics for nearby K7 to M6.5 dwarfs. Note that that $`80`$ to $`300`$ dwarfs contribute to each of the HGR bins up to spectral type M4.5. Only a few objects contribute to the M6.0 and M6.5 bins, which also are probably kinematically biased against young, active stars due to incompleteness in the pCNS3.<sup>5</sup><sup>5</sup>5HGR’s statistics for M7 and later dwarfs are sparse, but further observations have revealed that they are incorrect. Their Table 5 should show that 2 of 2 M7 dwarfs, 2 of 2 M8 dwarfs, and 2 of 3 M9.0-M9.5 dwarfs show emission. We extend our M dwarf sample to even cooler dwarfs by using the K99 and K00 data, who report the strength of H$`\alpha `$ emission in their L dwarfs. Their sample is photometrically-selected and kinematically unbiased. The sample shows a steady decline in H$`\alpha `$ emission frequency from 60% (80% if a marginal detection is included) for type L0 down to only 8% (25% if two marginal detections are included) for L4 dwarfs. None of the twenty dwarfs with spectral type L5.0 or later show definite emission (two have marginal detections). For the earlier dwarfs, emission of 1 Å equivalent width would have been detected in almost all the objects; however, the upper limit on the equivalent widths for the latest L dwarfs was typically somwhat larger. The fact that there is so little photospheric continuum for the latest L dwarfs around the H$`\alpha `$ feature should compensate for the lower sensitivity in terms in equivalent widths.
The data indicate that the frequency of emission increases with later spectral type (cooler temperatures), until at spectral type M7 all of our targets show detectable emission. Indeed, we are not aware of any inactive M7 dwarfs (HGR; Gizis & Reid 1997). This indicates that the dwarfs can maintain detectable levels of activity for the lifetime of the Galactic disk. Later than M7, the H$`\alpha `$ emission frequency begins to decrease, with our sample of ultracool M dwarfs merging cleanly with the L dwarfs. This coincides with the breakdown of the rotation-activity relation already noted for M9 and L dwarfs and reflects the apparent relative inability of the ultracool M dwarfs to heat the chromosphere discussed by TR. The percentage of emission at spectral type M6 is particularly uncertain, as the HGR sample may be biased toward higher velocity, hence older, stars at such low luminosities. Sixteen of our nineteen M6-M6.5 dwarfs show emission, but they have been effectively selected on the basis of unusually red $`JK_s`$ colors, and may be biased in some way. In any case, there is little doubt that some high velocity, presumably very old, M6 dwarfs are no longer active.
Since the H$`\alpha `$ line is seen against an increasingly faint photosphere for these ultracool dwarfs, the $`\mathrm{log}\left(\frac{L_{H\alpha }}{L_{bol}}\right)`$ ratio is more indicative of the true level of activity. We plot our M dwarfs, the HGR early M dwarfs, and the K99 L dwarf data in Figure 7. For the K99 data, we have measured the continuum level off the observed spectra and assumed $`BC_K=3.33`$ as measured by Tinney et al. (1993) for GD 165B to convert the K99 equivalent widths to flux ratios. Note that the increase in maximum observed activity levels from K7 to the peak at M3-M5 reflects the increased lifetime of emission for the lower mass stars. Earlier M dwarfs with activity levels near -3 are known in young clusters, but do not have long enough lifetimes to appear in the local sample. The lower envelope of data points is set by the fact that the minimum observable H$`\alpha `$ emission of 1Å equivalent width corresponds to a lower luminosity fraction in cooler dwarfs. The addition of our data to the HGR data clearly indicates that beyond M7 the level of activity is indeed declining. This decline continues for lower (L dwarf) temperatures. The decline is quite steep – in only three subclasses (M8 to L1) the activity drops by one full dex. Martín et al. (1999) note a “slight trend toward decreasing emission in the L dwarfs” – our conclusions differ due to our larger sample and their use of equivalent widths only.
While the quiescent H$`\alpha `$ chromospheric activity is declining, our data suggest that flare activity is common in the ultracool M dwarfs, and may be a significant contributor to the activity energy budget. We summarize evidence for variability in Table 5. Since these events represent a strong enhancement of the H$`\alpha `$ line strength, we suggest that they may be flares. At least a few dwarfs apparently maintain strong quiescent emission – our H$`\alpha `$ line strength of 29 Å for LP412-31 is identical with the value observed by Martín, Rebolo, & Zapatero-Osorio (1996). Other ultracool M dwarfs not in our sample have been seen to flare – Reid et al. (1999b) recently observed a flare on the “inactive” dwarf BRI 0021-0214, while Martín, Rebolo, & Magazzu (1994) observed the H-alpha EW of LHS 2065 on two consecutive nights as 7.5 Å then 20.3 Å. Martín et al. (1999) have observed flares in a number of ultracool M dwarfs. Flaring activity has thus been observed in ultracool M dwarfs with both very low and high levels of quiescent emission.
Assuming that these strong variations are due to flares, we estimate the flare rate from our own statistics. At least four of the fifty-three ultracool M targets were flaring the first time we observed them for this program — implying that ultracool M dwarfs spend $`7\%`$ of the time in a flare state. This is consistent with the observed flare rates of the “inactive” M9.5 dwarf BRI 0021-0214 (Reid et al., 1999b) and the monitoring of 2M0149 (Liebert et al., 1999). This flare rate is a lower limit, since some of our other targets may also be flaring, but lacking additional spectra we cannot tell whether they are merely active as for LP 412-31, and since we have no way of identifying weaker flares. The H$`\alpha `$ equivalent widths appear to be enhanced by a factor of $`10`$ in the observed flares, implying that perhaps half of the H$`\alpha `$ luminosity is emitted during flares.
We now consider the relationship between kinematics and activity for the M8 - M9.5 dwarfs. In Figure 8, we plot the observed relation between the tangential velocity and H$`\alpha `$ emission. There is a striking relationship between activity and velocity, in a sense opposite to that observed in more massive M dwarfs. All the dwarfs with strong H$`\alpha `$ emission have large velocities ($`v_{tan}>20`$ km/s). There is also a striking population of low-velocity, low-activity ultracool M dwarfs. With the exception of the high-velocity, inactive M9 dwarf 2MASSW J0109216+294925, the least active stars appear to be drawn from a lower velocity population. It is difficult to fairly characterize the tangential velocity dispersion ($`\sigma _{tan}^2=\sqrt{\sigma _{ra}^2+\sigma _{dec}^2}`$) of these populations, but the inactive, low-velocity population in the lower left of Figure 8 may be characterized by $`\sigma _{tan}15`$ km/s. While the low-velocity, low-activity population seems to to have EW$`{}_{H\alpha }{}^{}7`$Å and $`v_{tan}25`$ km/s, we can calculate a velocity dispersion only for a purely activity selected sample. As an illustration, the dwarfs with $`EW_{H\alpha }<3`$Å (excluding 2MASSW J0109216+294925) have $`\sigma _{tan}=13`$ km/s; in contrast, those with more emission have $`\sigma _{tan}=38`$ km/s. Using the approximation that $`\sigma _{tot}=\sqrt{\frac{3}{2}}\sigma _{tan}`$, the implied total space dispersions of the two populations are 16 km/s and 47 km/s. Comparing to Wielen (1974), the active M dwarfs are apparently drawn from a $`5\times 10^9`$ yr population, but the inactive M dwarfs are consistent with a $`0.5`$ Gyr population. This estimate is crude at best, but it seems clear that the overall ultracool M dwarf population is drawn from a long-lived, presumably stellar population, while the group of less active, low-velocity stars represent a younger ($`1`$ Gyr) population .
Despite the smaller sample size, the properties of L dwarfs are of considerable interest. The Sample BC L0-L4 dwarf velocities are typical of an old disk population. Even excluding 2M1300 due to unusual color, we find $`\sigma _{tan}=56`$ km/s, while including it we find $`\sigma _{tan}=70`$ km/s. Only two of the Sample BC L dwarfs show emission – the low velocity 2M1108 has the strongest emission at 7.8Å, while the high velocity 2M1506 has weak 1Å emission. While the velocities of only two L dwarfs are not definitive, the velocity distribution of the inactive L dwarfs suggest they are mostly old. Adding the information provided by the work of K99 and K00 to our data provides strong clues, that just as in the ultracool M dwarfs, the traditional activity-age relationship is broken, perhaps even reversed. L dwarfs that show lithium absorption are necessarily younger and lower mass than L dwarfs of the same spectral type which have destroyed lithium. Thus, using the traditional stellar age-activity relation, one would expect them to be more chromospherically active. Consider the L1 to L4.5 dwarfs, where lithium is detectable even at the low resolution of the K99/K00 Keck LRIS observations. Only one L dwarf, Kelu-1, shows both H$`\alpha `$ emission and lithium absorption.<sup>6</sup><sup>6</sup>6It is interesting to note that Basri (1999) finds that Kelu-1 is rotating extremely rapidly: 80 km/s. Eleven other such L dwarfs show H$`\alpha `$ emission but do not have lithium absorption. Twelve L dwarfs show lithium absorption but do not have H$`\alpha `$ emission (four of these have marginal H$`\alpha `$ detections or noise consistent with emission of less than 2Å). While many L1-L4 dwarfs have neither H$`\alpha `$ emission nor lithium absorption, it seems clear that the chromospherically active L dwarfs are drawn from an older, more massive population than the lithium L dwarfs. Beyond L4.5, there are no definite cases of H$`\alpha `$ emission, although lithium absorption is present for $`50\%`$ of the L dwarfs.
Brown dwarfs have also been identified in nearby young clusters. In Figure 9, we compare the activity of our field dwarfs to young brown dwarfs. Shown are young brown dwarfs from the $`\sigma `$ Ori cluster (Bejar, Zapatero Osorio & Rebolo, 1999; Zapatero Osorio et al., 1999) and the young brown in the $`\rho `$ Oph cloud (Luhman, Liebert & Rieke, 1997) – both these clusters are less than 10 Myr old. Also shown are confirmed Plieades brown dwarfs (Martín, Rebolo, & Zapatero-Osorio, 1996; Zapatero Osorio et al., 1997; Martín et al., 1998a, b) with age $`10^8`$ years. In order to suggest the age evolution of the field M dwarfs, those with $`v_{tan}<20`$ km/s have been marked as open symbols. These low-velocity dwarfs are likely to be younger than the other field dwarfs. It is evident that the Plieades M8 and later brown dwarfs are not more active that the typical field dwarfs, although the M6-M7 brown dwarfs appear to be more active. Like the young field L dwarfs that have lithium aborption, the Pleiades and $`\sigma `$ Ori L dwarfs do not show emission, even though some older field L dwarfs do. In the case of the $`\rho `$ Oph brown dwarf, Luhman, Liebert & Rieke (1997) have shown that the emission is probably due to accretion from the circumstellar disk or envelope detected in the mid-IR. Similar accretion may account for the $`\sigma `$ Ori strong emitters, while the absence of accretion would account for the weak emission in the other half of the $`\sigma `$ Ori brown dwarf sample. If the emission is chromospheric, then some other factor (such as rotation) is needed to explain the large spread in activity levels. It should be noted that these young brown dwarfs are probably hotter at a given spectral type (Luhman, Liebert & Rieke, 1997), which suggests that if temperature were plotted the brown dwarfs would appear even less active compared to field dwarfs.
We thus summarize the observations:
1. Although the fraction of dwarfs showing H$`\alpha `$ emission reaches 100% at spectral type M7, the fraction that show chromospheric activity drops rapidly for later spectral types.
2. The fraction of energy in chromospheric H$`\alpha `$ for those dwarfs that are active drops rapidly as a function of spectral type beyond M6.
3. Low velocity, kinematically young M8.0-M9.5 dwarfs have weaker activity than many higher velocity, old M8.0-M9.5 dwarfs
4. Flaring is common among the M7-M9.5 dwarfs
5. The early L (L0-L4) dwarfs in Sample BC have old kinematics.
6. Early L (L1-L4.5) dwarfs with H$`\alpha `$ emission are old and massive enough to have burned lithium
7. L1-L4.5 dwarfs with lithium are unlikely to have H$`\alpha `$ emission.
8. None of the L dwarfs later than L4.5 have H$`\alpha `$ emission but half have lithium.
9. The two known L dwarfs in young clusters do not show H$`\alpha `$ emission
10. Young Pleiades M8-M9 dwarfs are less active than the higher velocity field M dwarfs.
### 6.3 Discussion
How can these observations be explained? We believe they imply that the maximum activity level is a strong function of temperature beyond spectral type M7, with lower temperature objects able to maintain less emission. Additionally, beyond spectral type M7 substellar dwarfs tend to have less activity than stellar dwarfs.
In Figure 10, we plot evolutionary sequences from Burrows et al. (1993) and Baraffe et al. (1998). In the Baraffe et al. (1998) models the hydrogen burning limit is at 0.072 $`M_{}`$ and the lithium burning limit is at $`0.055M_{}`$. Also shown is an estimated temperature scale from Reid et al. (1999a) based on the arguments made in K99. The models indicate that it takes $`10^9`$ years for stars near the hydrogen burning limit to settle into the M8 and cooler temperatures (Figure 10). M8 and cooler temperatures are possible at younger ages, but they are substellar objects which continue to cool with time. Thus, the low velocity, low activity population is likely to be a population of substellar/transition objects. By the time they are older than 1 Gyr, they appear as L dwarfs or even cooler T dwarfs. Thus, the comparison of low velocity M8-M9.5 dwarfs to high velocity M8-M9.5 dwarfs is the same as a comparison of younger, lower mass ($`0.07M_{}`$) objects to older, higher mass objects ($`0.08M_{}`$) at the same temperature. The same occurs when comparing the K99/K00 lithium L1-L4 ($`0.055M_{}`$) dwarfs to the K99/K00 non-lithium L1-L4 ($`0.07M_{}`$) dwarfs. These mass estimates are only meant as illustrative values – mass estimates are subject to many uncertainties and a range of masses and ages will be sampled.
Both the field M8-M9 dwarf and field L dwarf observations show that at a given spectral type, the less massive dwarfs are less active, even though they are younger. It is perhaps worth noting that theoretical models suggest that the lower mass objects will be slightly more luminous with a smaller surface gravity (Burrows et al., 1993; Baraffe et al., 1998). The comparison to cluster brown dwarfs is consistent with this trend (Figure 9). The Plieades brown dwarfs which have cooled to these ultracool M and L temperatures are less active than old field ultracool M and L dwarfs — but more active than the low velocity ultracool M dwarfs. There are thus strong suggestions, as already noted in Section 6.1, that activity levels do decrease with age in brown dwarfs, and that therefore activity levels are dependent upon temperature, mass, and age. The importance of accretion needs to be investigated for the youngest ($`<10`$ Myr) ages. The dominant effect is temperature, as both young and old objects show the rapid fall in activity levels beyond spectral type M7.
At the same time, the models suggest that stars, or at least very-long lived hydrogen burning transition objects, are likely to exist down to L0 - L4 temperatures. This is completely consistent with our empirical finding that the early L dwarfs have old kinematics. We note that Kirkpatrick et al. (1999a) find a temperature of $`1900\pm 100`$ K for the L4 dwarf GD 165B and constrain the age to be greater than 1.2 Gyr using updated models and the white dwarf primary’s cooling age and argue it is just below the substellar limit, near the transition region between stars and brown dwarfs (formally, they actually derive the minimum stellar mass using the models). L dwarfs with lithium must be below the lithium burning limit ($`0.055M_{}`$; Chabrier & Baraffe 1997) and younger than 1 Gyr (Figure 10). The inactivity of these lithium L dwarfs demonstrates that the lowest mass objects cannot sustain significant activity at a temperature (or luminosity) that is adequate for sustaining some activity in older but more massive dwarfs.
The decline in the frequency of activity may be associated with two effects. First, as later spectral types are considered, a larger fraction of very low mass (substellar) objects contribute, and these are more likely to be inactive in field samples. Second, activity among the L dwarfs may die out in time, since the observed high velocity L dwarfs are inactive – although we cannot tell if they were ever active. It would be of great interest to find whether or not the early K99/K00 L dwarfs which are active have low or high velocities. The high-velocity, low-activity M9.5 dwarf 2MASSW J0109216+294925 may be a young brown dwarf that happens to have high velocity, an old object which has never been active, or an old stellar M9 dwarf whose chromospheric activity has declined with age. In any case, it is worthy of additional study. We note that the high observed flare rate implies that the rotation rate may decrease with time, even among the “inactive, rapid rotators” if the flaring is associated with mass loss and/or a stellar wind. We speculate this may provide a mechanism for the evolution of activity. Additional observations are needed to determine what the rotational velocities are as a function of mass, spectral type, and age.
What fraction of the ultracool M dwarfs are likely to be substellar? While we cannot identify which individual objects are brown dwarfs, we can identify a number of probably young objects. Three of our M8.0-M9.5 dwarfs show no H$`\alpha `$ emission and very low velocity; another three have equally low velocities and $`EW_{H\alpha }<3`$Å. Out of a total population of 32 M8.0-M9.5 dwarfs, our data suggest that $`1020\%`$ are brown dwarfs. These objects should be more likely to have lithium absorption (like LP 944-20), but most of the brown dwarfs will be massive enough to burn lithium. Indeed, we note that none show the distinctive signs of low surface gravity that characterize the Pleiades M8 brown dwarfs Teide 1 and Calar 3 (Martín, Rebolo, & Zapatero-Osorio, 1996), so none of our targets are very young ($`10^8`$ years). Approximately ten objects belong to the low velocity, low activity group in the lower left of Figure 8 – that is a third of the sample, but some fraction of these will stabilize as hydrogen burning L dwarfs, in order to account for the observed popoulation of high velocity early L dwarfs. These fractions will be somewhat overestimated for the Galactic disk population, since the old, large scale height population will be underrepresented locally. Other effects may also be important, such as the fact that we have estimated distances for all dwarfs using one color-absolute magnitude relation. Adding kinematic ages, as in this study, provides an additional constraint on the modelling necessary to determine the field substellar mass function (Reid et al., 1999a).
The nature of the ultracool M dwarfs has been debated for some time in the literature. While convential wisdom suggested that most if not all field ultracool M dwarfs are stellar, many suggestions that most ultracool M dwarfs are substellar have been made, most of which have been discredited. We remark that Bessell (1991) noted that the ultracool M dwarfs are expected to be a mixture of young brown dwarfs and older stars – and he also noted that there was a paucity of high proper motion ultracool M dwarfs expected from the stellar population in the LHS catalogue (Luyten, 1979). Our study has identified a number of high proper M dwarfs which appear on the red POSS plates but were overlooked for the LHS catalogue – evidently, the faintest of these targets on the blue plate precluded their detection by Luyten and contributed to the effect noted by Bessell. Our results show that most ultracool M dwarfs are old, and hence stellar, but perhaps 10-20% are a younger population of brown dwarfs.
We end our discussion with a few caveats due to our photometric selection. Both the relative numbers and kinematics of “active” and “inactive” M dwarfs will be changed if one group is preferentially brighter at $`M_K`$ for its $`JK`$ color. Indeed, the inactive brown dwarf LP 944-20 lies one magnitude below the active M dwarf LHS 2397a in the $`JK`$, $`M_K`$ HR diagram. Preliminary USNO parallaxes show instrinsic dispersion in the HR diagram for the late M and L dwarfs (Dahn, private communication). If inactive ultracool M dwarfs are subluminous compared to our adopted relation, we will have overestimated their velocities; correspondingly, if the more active dwarfs are “superluminous,” we have underestimated their velocities. Fortunately, this would only strengthen our evidence that active ultracool M dwarfs are older. The intrinsic dispersion presumably depends upon such ill-understood factors as metallicity, age, surface gravity, and dust formation. Another possible bias on activity levels is that we favor the inclusion of unresolved binaries. Amongst the earlier M dwarfs, very short period systems have enhanced chromospheric activity due to tidal effects maintaining high rotation rates(Young, Sadjadi, & Harlan, 1987) — however, even if this mechanism works in the ultracool M dwarfs, which seems unlikely if the rotation-activity relation has broken down, only $`5\%`$ of earlier type M dwarfs show emission due to this effect, so it should be negligible.
## 7 Summary
We show that a sample of bright, nearby ultracool M and L dwarfs can be selected without proper motion bias using 2MASS and PMM scans. Our initial samples include high proper motion objects, visible on the POSS plates, that should be added to an updated version of the LHS Catalogue, and one M8.0 dwarf with a photometric parallax that places it within 10 parsecs. We intend to continue this study in order to complete the nearby star catalog for the lowest mass stars.
Using our initial sample, we estimate the space density of dwarfs near the hydrogen-burning limit. We show that the dropoff near the hydrogen burning limit in the five and eight parsec nearby star samples is likely to be due to incompleteness. This is more consistent with a smooth relation across the hydrogen burning limit. Trigonometric parallaxes and searches for companions will help improve the space density estimate.
Most importantly, we use our spectroscopic observations of our well-defined sample to explore the relationships between age, kinematics, and chromospheric activity for the ultracool M and L dwarfs. We show that the observations can be understood if activity is primarily related to temperature and secondarily mass and age, and that lower mass (substellar) objects have weaker chromosperes. Thus, the classical relation that strong H$`\alpha `$ emission implies youth is not valid for these dwarfs. Instead, strong H$`\alpha `$ emitters in the field are likely to be old ($`1`$ Gyr) stars, while weaker emitters are often young ($`<1`$ Gyr), lower-mass brown dwarfs. This does not exclude the idea that for a given dwarf, H$`\alpha `$ activity declines with age – but spectral type (temperature) is the observable in the field. The local population of ultracool M dwarfs apparently consists both of the most massive (lithium burning) brown dwarfs and the lowest mass (hydrogen burning) stars, with the substellar objects making up a significant fraction of the sample. The early L (L0-L4) dwarfs are consistent with an old, at least partially stellar population. The evidence thus suggests, as do some models, that early L dwarfs can be stable hydrogen-burning stars. Expansion of the sample with follow-up observations should clarify the relative contribution of stars and brown dwarfs to these temperature ranges.
We thank Suzanne Hawley and Mike Skrutskie for useful discussions. We thank the staffs of Las Campanas, Palomar, Keck, and Kitt Peak observatories for their assistance in the observations and the people at University of Massachusetts and IPAC for their efforts in making 2MASS a reality. JEG and JDK acknowledge the support of the Jet Propulsion Laboratory, California Institute of Technology, which is operated under contract with NASA. This work was funded in part by NASA grant AST-9317456 and JPL contract 960847. INR, JDK, and JL acknowledge funding through a NASA/JPL grant to 2MASS Core Project science This publication makes use of data products from 2MASS, which is a joint project of the University of Massachusetts and IPAC, funded by NASA and NSF. Some of the data presented herein were obtained at the W.M. Keck Observatory, which is operated as a scientific partnership among Caltech, the University of California and NASA. JEG accessed the DSS as a Guest User, Canadian Astronomy Data Centre, which is operated by the Herzberg Institute of Astrophysics, National Research Council of Canada. This research has made use of the Simbad database, operated at CDS, Strasbourg, France.
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# Toroidal Lie algebras and Bogoyavlensky’s 2+1-dimensional equation
## 1 Introduction
In the early 1980s, Date et al. discovered the remarkable fact that the Lie algebra $`𝔤𝔩(\mathrm{})`$ acts on the solutions of the KP(Kadomtsev-Petviashvili) hierarchy. The link between infinite-dimensional Lie algebras and soliton equations has been widely extended. In particular, Kac and Wakimoto developed a scheme to construct hierarchies from vertex operator representations of general affine Lie algebras. On these backgrounds, Billig and Iohara et al. derived new hierarchies of Hirota bilinear equations from $`2`$-toroidal Lie algebra $`𝔤^{\mathrm{tor}}`$, a central extension of double loop algebra $`𝔤[s^{\pm 1},t^{\pm 1}]`$ of a simple Lie algebra $`𝔤`$. Applying the construction to the case $`𝔤^{\mathrm{tor}}=𝔰𝔩_{\mathrm{}}^{\mathrm{tor}}`$ (and their principal vertex realization), we have an extension of the $`\mathrm{}`$-reduced KP hierarchy of Hirota bilinear form.
For the simplicity of exposition, we here deal with the case $`𝔤=𝔰𝔩_2`$. The typical Hirota bilinear equations of lower degrees in hierarchy of , are the following
$`\left(D_x^44D_xD_t\right)\tau \tau =0,`$ (1.1)
$`\left(D_yD_x^3+2D_yD_t6D_zD_x\right)\tau \tau =0.`$ (1.2)
If we count the degrees of Hirota’s $`D`$-operators $`D_y,D_x,D_z,D_t`$ as $`0,1,2,3`$ respectively, the Hirota bilinear equations above have degrees $`4`$ and $`3`$ respectively. The first one (1.1) is a famous bilinear form of the KdV equation $`4u_t=6uu_x+u_{xxx},`$ where $`u\stackrel{\mathrm{def}}{=}2(\mathrm{log}\tau )_{xx}.`$ The second one (1.2) can be written as a non-linear equation known as Bogoyavlensky’s $`2+1`$-dimensional equation
$`u_z={\displaystyle \frac{1}{4}}u_{xxy}+vu_x+uu_y,`$ (1.3)
where we set $`v\stackrel{\mathrm{def}}{=}(\mathrm{log}\tau )_{xy}`$ and used the KdV equation to eliminate the terms including derivative $`_t.`$ Equation (1.3) has the following equivalent Lax form
$`{\displaystyle \frac{P}{z}}=[P_y+C,P]\text{where}P\stackrel{\mathrm{def}}{=}_x^2+u,C\stackrel{\mathrm{def}}{=}v_x+{\displaystyle \frac{3}{4}}u_y.`$ (1.4)
The main purpose of this paper, motivated by the fact above, is to establish a Lax formalism of the Hirota bilinear equations arising from $`𝔰𝔩_{\mathrm{}}^{\mathrm{tor}}`$ for any $`\mathrm{}2`$. For this purpose, we shall introduce the $`\mathrm{}`$-Bogoyavlensky hierarchy. We shall also construct a family of special solutions that generalize the breaking solitons studied by Bogoyavlensky.
We shall develop our discussions as follows. In section 2, we discuss the hierarchy on a viewpoint of representation theory. First we give definition of the Lie algebra $`\stackrel{~}{𝔤}^{\mathrm{tor}}`$, and then explain a general result on their representations(Lemma 1). Using this lemma, we construct a representation of $`\stackrel{~}{𝔰𝔩}_{\mathrm{}}^{\mathrm{tor}}`$, based on the $`\mathrm{}`$-reduction procedure of $`𝔤𝔩(\mathrm{})`$ in . We shall show that every vector $`\tau `$ in the $`SL_{\mathrm{}}^{\mathrm{tor}}`$-orbit of the vacuum vector satisfies a hierarchy of Hirota bilinear equations, which is essentially the same as those derived in ,, although the generating functions are apparently different. In section 3, we first give a review on the Lax formalism of the KP hierarchy. Next we present a heuristic introduction to Bogoyavlensky’s hierarchy. These arguments lead us to formulate the $`\mathrm{}`$-Bogoyavlensky hierarchy for each $`\mathrm{}2,`$ where the so-called Wilson-Sato operator plays the significant role. In section 4, we give a residue formula for the formal Baker-Akhiezer functions of the hierarchy by the calculus of pseudo-differential operators. The formula leads to exactly the same system of Hirota bilinear equations derived in section 2. Note that the most part of sections 3,4 can be read independently from the preceding section. Section 5 is devoted to investigate the special solutions of the hierarchy, namely the Wronskian solutions and the $`N`$-soliton solutions. As an appendix, we give a table of Hirota equations of low degree in section 7.
It should be mentioned that the Bogoyavlensky’s equation is a dimensional reduction, to $`2+1`$ dimensions, of the four dimensional self-dual Yang-Mills equation. Mason and Sparling pointed out that many integrable systems in $`2+1`$ and lower dimensions (in particular, $`1+1`$ dimensional soliton equations such as the KdV equation) can be thus derived from the self-dual Yang-Mills equation. In fact, Schiff rediscovered Bogoyavlensky’s equation, having been unaware of Bogoyavlensky’s work. Schiff’s work was further extended by Yu et al. .
Throughout the paper, the base field is the field $``$ of complex numbers. The symbol $``$ stands for the set of integers. The binomial coefficients $`\left(\genfrac{}{}{0pt}{}{n}{k}\right)`$ is defined as $`n(n1)\mathrm{}(nk+1)/k!`$ for $`k0.`$ For a vector space $`V`$, we denote by $`V[z^{\pm 1}],V[z^{\pm 1},w^{\pm 1}]`$ the spaces $`V[z,z^1],V[z,z^1,w,w^1]`$ of Laurent polynomials with the coefficients in $`V`$ respectively. For an integer $`\mathrm{}2,`$ we denote by $`𝔤𝔩_{\mathrm{}}`$ (resp. $`𝔰𝔩_{\mathrm{}}`$) the Lie algebras of the all (resp. traceless) $`\mathrm{}\times \mathrm{}`$ matrices.
Note added: After the first version of this paper had written, we learned that F. Calogero had studied equation (1.3) in 1975. We thank Kouichi Toda and his collaborators for the information.
## 2 Deriving hierarchy from representation theory
### 2.1 Definitions of the toroidal Lie algebra
Let $`𝔤`$ be a finite-dimensional simple Lie algebra over $`.`$ Let $`R`$ be the ring of Laurent polynomials of two variables $`[s^{\pm 1},t^{\pm 1}].`$ The module of Kähler differentials $`\mathrm{\Omega }_R`$ of $`R`$ is defined with the canonical derivation $`d:R\mathrm{\Omega }_R`$. As an $`R`$-module, $`\mathrm{\Omega }_R`$ is freely generated by $`ds`$ and $`dt.`$ Let $`\overline{}:\mathrm{\Omega }_R\mathrm{\Omega }_R/dR`$ be the canonical projection. Let $`𝒦`$ denote $`\mathrm{\Omega }_R/dR.`$ Let $`(|)`$ be the normalized Killing form () on $`𝔤.`$ We define the Lie algebra structure on $`𝔤^{\mathrm{tor}}\stackrel{\mathrm{def}}{=}𝔤R𝒦`$ by
$`[Xf,Yg]=[X,Y]fg+(X|Y)\overline{(df)g},[𝒦,𝔤^{\mathrm{tor}}]=0.`$ (2.1)
In the paper , Kassel proved that the bracket defines a universal central extension of $`𝔤R`$ as a Lie algebra over $`,`$ for a large class of commutative algebra $`R.`$ One can find in a simpler proof (a natural extension of Wilson’s discussions ) of the Kassel’s result that works for the cases when the ground field of $`R`$ has characteristic $`0.`$
Let $`𝒟=R_{\mathrm{log}s}R_{\mathrm{log}t}`$ be the Lie algebra of derivations on $`R.`$ A derivation $`\delta 𝒟`$ can be naturally extended to a derivation on the Lie algebra $`𝔤R`$ by
$`\delta \left(Xf\right)\stackrel{\mathrm{def}}{=}X\delta f.`$ (2.2)
It is known that a derivation $`\delta 𝒟`$ acting on the Lie algebra $`𝔤R`$ has a natural extension to $`𝔤^{\mathrm{tor}}`$. The action on the center $`𝒦`$ is given explicitly as follows . First we shall define the action of $`𝒟`$ on $`\mathrm{\Omega }_R`$ by:
$`f_{\mathrm{log}u}\left(gd\mathrm{log}v\right)\stackrel{\mathrm{def}}{=}\left(f_{\mathrm{log}u}g\right)d\mathrm{log}v+\delta _{u,v}gdf\text{for}u,v=s,t.`$ (2.3)
The action preserves the exact forms $`dR,`$ namely we have $`f_{\mathrm{log}u}(dg)=d\left(f_{\mathrm{log}u}g\right)`$ for $`u=s,t,`$ and so induces an action, also denoted by the same notation, on $`𝒦=\mathrm{\Omega }_R/dR.`$
Let $`\stackrel{ˇ}{𝒟}`$ denote the Lie subalgebra $`R_{\mathrm{log}t}`$ of $`𝒟.`$ We shall add the derivations $`\stackrel{ˇ}{𝒟}`$ to $`𝔤^{\mathrm{tor}}`$ to get the Lie algebra
$`\stackrel{~}{𝔤}^{\mathrm{tor}}=𝔤^{\mathrm{tor}}\stackrel{ˇ}{𝒟}.`$
Here the bracket in $`𝔤^{\mathrm{tor}}`$ is given by (2.1). We define the bracket between $`\stackrel{ˇ}{𝒟}`$ and $`𝔤^{\mathrm{tor}}`$ by the action (2.2), and (2.3), $`[\delta ,a]=\delta (a)(\delta \stackrel{ˇ}{𝒟},a𝔤^{\mathrm{tor}}).`$ The bracket in $`\stackrel{ˇ}{𝒟}`$ is defined by
$`[f_{\mathrm{log}t},g_{\mathrm{log}t}]=\left(f(_{\mathrm{log}t}g)g(_{\mathrm{log}t}f)\right)_{\mathrm{log}t}\overline{(_{\mathrm{log}t}g)d(_{\mathrm{log}t}f)}.`$
Note that if $`f,g[t^{\pm 1}]`$ in the above formula, then it is equivalent to the relation of the Virasoro algebra. Remark that $`\stackrel{~}{𝔤}^{\mathrm{tor}}`$ is not the semidirect product of $`𝔤^{\mathrm{tor}}`$ with the action of $`\stackrel{ˇ}{𝒟}`$ by (2.2),(2.3).
We have, for $`u=s,t`$, the Lie subalgebras
$`\widehat{𝔤}_u\stackrel{\mathrm{def}}{=}𝔤[u^{\pm 1}]\overline{d\mathrm{log}u},`$
with the brackets given by
$`[Xu^m,Yu^n]=u^{m+n}[X,Y]+m\delta _{m+n,0}(X|Y)K_u,`$
which are isomorphic to the affine Lie algebra $`\widehat{𝔤}`$ with the canonical central element $`K_u\stackrel{\mathrm{def}}{=}\overline{d\mathrm{log}u}.`$
We prepare the generating series of $`\stackrel{~}{𝔤}^{\mathrm{tor}}`$ as follows:
$`A_m(z)\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n}{}}As^nt^mz^{n1},D_m^t(z)\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n}{}}s^nt^m_{\mathrm{log}t}z^{n1},`$
$`K_m^s(z)\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n}{}}\overline{s^nt^md\mathrm{log}s}z^n,K_m^t(z)\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n}{}}\overline{s^nt^md\mathrm{log}t}z^{n1},`$
where $`m`$ and $`A𝔤.`$ The relation $`\overline{d(s^nt^m)}=0`$ can be neatly expressed by these generating series as
$`{\displaystyle \frac{}{z}}K_m^s(z)+mK_m^t(z)=0.`$
### 2.2 A Heisenberg subalgebra $`𝔡`$ of $`\stackrel{~}{𝔤}^{\mathrm{tor}}`$
The Lie algebra $`\stackrel{~}{𝔤}^{\mathrm{tor}}`$ has a Heisenberg subalgebra $`𝔡`$ with basis $`\{s^n_{\mathrm{log}t},\overline{s^nd\mathrm{log}t},d\mathrm{log}s(n)\}.`$ In fact, if we set $`\varphi _n\stackrel{\mathrm{def}}{=}\overline{s^nd\mathrm{log}t},`$ $`\varphi _n^{}\stackrel{\mathrm{def}}{=}s^n_{\mathrm{log}t},`$ and $`\mathrm{}\stackrel{\mathrm{def}}{=}\overline{d\mathrm{log}s},`$ then we have
$`[\varphi _m^{},\varphi _n]=m\delta _{m+n,0}\mathrm{},[\varphi _m,\varphi _n]=[\varphi _m^{},\varphi _n^{}]=[\mathrm{},𝔡]=0.`$
The Heisenberg algebra $`𝔡`$ is degenerate in the sense that $`[\varphi _0^{},\varphi _0]=0.`$ We also note that $`𝔡`$ commute with $`\widehat{𝔤}_s.`$
Next we realize the action of $`𝔡`$ on the space of polynomials
$`F_𝔡\stackrel{\mathrm{def}}{=}[y_{\mathrm{}},y_2\mathrm{},y_3\mathrm{},\mathrm{};y_{\mathrm{}}^{},y_2\mathrm{}^{},y_3\mathrm{}^{},\mathrm{}][e^{y_0},e^{y_0}],\mathrm{𝐯𝐚𝐜}_𝔡\stackrel{\mathrm{def}}{=}1F_𝔡`$
where the action of $`𝔡`$ is defined by:
$`\varphi _m\{\begin{array}{cc}\frac{}{y_m\mathrm{}^{}}\hfill & m>0\hfill \\ my_m\mathrm{}\hfill & m0\hfill \end{array},\varphi _m^{}\{\begin{array}{cc}\frac{}{y_m\mathrm{}}\hfill & m0\hfill \\ my_m\mathrm{}^{}\hfill & m<0\hfill \end{array},\mathrm{}=\mathrm{id}.`$
For convenience’ sake, we adopted rather peculiar notation for the variables $`y_n\mathrm{}`$ etc.
We define the generating series
$`\varphi (z)\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n}{}}\varphi _nz^{n1},\varphi ^{}(z)\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n}{}}\varphi _n^{}z^{n1},`$
and for each $`m`$, the following vertex operators
$`V_m(z)\stackrel{\mathrm{def}}{=}\mathrm{exp}\left(m{\displaystyle \underset{n>0}{}}y_n\mathrm{}z^n\right)e^{my_0}\mathrm{exp}\left(m{\displaystyle \underset{n>0}{}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{}{y_n\mathrm{}^{}}}z^n\right).`$ (2.4)
We also introduce the normal product $`:\varphi (z)V_m(z):\stackrel{\mathrm{def}}{=}\varphi (z)_{<0}V_m(z)+V_m(z)\varphi (z)_{>0}`$ where $`\varphi (z)_{<0}\stackrel{\mathrm{def}}{=}_{n<0}\varphi _nz^{n1}`$ and $`\varphi (z)_{>0}\stackrel{\mathrm{def}}{=}_{n>0}\varphi _nz^{n1}.`$
We shall say a representation $`V`$ of the affine Lie algebra $`\widehat{𝔤}_s`$ is restricted if and only if for all $`vV`$ and $`X𝔤`$, there exists $`N`$ such that $`Xs^nv=0`$ for $`n>N.`$ Now we can state the following lemma due to Iohara et al. and Berman-Billig .
###### Lemma 1
Let $`(V,\pi )`$ be a restricted representation of $`\widehat{𝔤}_s`$ such that $`\overline{d\mathrm{log}s}\mathrm{id}_V.`$ Then we can define the representation $`\pi ^{\mathrm{tor}}`$ of $`\stackrel{~}{𝔤}^{\mathrm{tor}}`$ on $`V_𝔡`$ such that
$`A_m(z)`$ $`A^\pi (z)V_m(z),K_m^s(z)`$ $`V_m(z),`$
$`D_m^t(z)`$ $`\varphi ^{}(z)V_m(z),K_m^t(z)`$ $`:\varphi (z)V_m(z):,`$
where $`A𝔤,m`$ and $`A^\pi (z)\stackrel{\mathrm{def}}{=}_n\pi (As^n)z^{n1}.`$
### 2.3 Review on a representation of $`𝔤𝔩(\mathrm{})`$ and its $`\mathrm{}`$-reduction
In the subsection, we collect basic results on representation of $`𝔤𝔩(\mathrm{})`$ which is used in the sequel. One can find in the original paper the proofs we omit. We shall mainly follow the notation used in book .
Let the associative $``$-algebra $`𝒜`$ be generated by $`\psi _i,\psi _i^{}(i+\frac{1}{2})`$ with the relations:
$`\psi _i\psi _j+\psi _j\psi _i=\psi _i^{}\psi _j^{}+\psi _j^{}\psi _i^{}=0,\psi _i^{}\psi _j+\psi _j\psi _i^{}=\delta _{i+j,0}.`$
Consider a left $`𝒜`$-module with a cyclic vector $`\mathrm{𝐯𝐚𝐜}_𝒜`$ satisfying
$`\psi _i\mathrm{𝐯𝐚𝐜}_𝒜=\psi _i^{}\mathrm{𝐯𝐚𝐜}_𝒜=0\text{for}i>0.`$ (2.5)
This $`𝒜`$-module $`𝒜\mathrm{𝐯𝐚𝐜}_𝒜`$ is called the fermionic Fock space, which we denote by $`_𝒜.`$
We define the following infinite formal Laurent series of the variable $`\lambda `$
$`\psi (\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{i+\frac{1}{2}}{}}\psi _i\lambda ^{i\frac{1}{2}},\psi ^{}(\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{i+\frac{1}{2}}{}}\psi _i^{}\lambda ^{i\frac{1}{2}},`$
and the normal ordering of the following quadratic expression to be
$`:\psi _i\psi _j^{}\stackrel{\mathrm{def}}{=}\{\begin{array}{cc}\psi _i\psi _j^{}\hfill & \text{if}i<0\text{or}j>0\hfill \\ \psi _j^{}\psi _i\hfill & \text{if}i>0\text{or}j<0\hfill \end{array}.`$
Consider the following set of infinite complex matrices
$`\left\{A=(a_{ij})_{i,j+\frac{1}{2}}\right|a_{ij}=0\text{for}|ij|0\}.`$
It forms a Lie algebra over $``$ with the usual bracket $`[A,B]_0\stackrel{\mathrm{def}}{=}ABBA`$, which we shall denote by $`\overline{𝔤𝔩}(\mathrm{})`$. We shall define the 2-cocycle $`\omega `$ on $`\overline{𝔤𝔩}(\mathrm{})`$ by :
$`\omega (E_{ij},E_{i^{}j^{}})\stackrel{\mathrm{def}}{=}\delta _{ji^{}}\delta _{ij^{}}(\theta (j)\theta (i))`$
where $`E_{ij}\stackrel{\mathrm{def}}{=}\left(\delta _{ii^{}}\delta _{jj^{}}\right)_{i^{},j^{}+\frac{1}{2}}`$ is the matrix unit and $`\theta `$ is defined by $`\theta (i)\stackrel{\mathrm{def}}{=}1`$ if $`i>0`$ and $`\theta (i)\stackrel{\mathrm{def}}{=}0`$ if $`i<0.`$ Let us introduce the central extension $`𝔤𝔩(\mathrm{})\stackrel{\mathrm{def}}{=}\overline{𝔤𝔩}(\mathrm{})c`$ with the central element $`c`$ and the bracket:
$`[A,B]\stackrel{\mathrm{def}}{=}[A,B]_0+\omega (A,B)c,A,B\overline{𝔤𝔩}(\mathrm{}).`$
###### Lemma 2
The map $`E_{ij}:\psi _i\psi _j^{}:,c\mathrm{id}__𝒜`$ defines the representation $`\pi _{\mathrm{}}`$ of $`𝔤𝔩(\mathrm{})`$ on the fermionic Fock space $`_𝒜.`$
In terms of the generating series defined by
$`E(\lambda ,\mu )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{i,j+\frac{1}{2}}{}}E_{ij}\lambda ^{i\frac{1}{2}}\mu ^{j\frac{1}{2}},`$
the lemma has the equivalent expression
$`E(\lambda ,\mu ):\psi (\lambda )\psi ^{}(\mu ):.`$
Fix an integer $`\mathrm{}2.`$ Let $`A`$ be a matrix $`(a_{ij})_{i,j+\frac{1}{2}}`$ in $`𝔤𝔩(\mathrm{}).`$ We assume the condition:
$`a_{i,j}=a_{i+\mathrm{},j+\mathrm{}}\text{for all}i,j+\frac{1}{2}.`$ (2.6)
Then we introduce, for each $`n`$, the $`\mathrm{}\times \mathrm{}`$ complex matrix by
$`A_n\stackrel{\mathrm{def}}{=}\left(a_{i+\frac{1}{2},j+\frac{1}{2}+n\mathrm{}}\right)_{0i,j\mathrm{}1}`$
and define an $`\mathrm{}\times \mathrm{}`$ matrix of Laurent polynomial of the variable $`s`$
$`{\displaystyle \underset{n}{}}A_ns^n.`$
Conversely, let $`_nA_ns^n`$ be a Laurent polynomial in $`𝔤𝔩_{\mathrm{}}[s^{\pm 1}].`$ Then we associate the $`\mathrm{}`$-periodic infinite matrix in the following block form
$`{\displaystyle \underset{n}{}}A_ns^n\left(\begin{array}{ccccc}\mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& A_0& A_1& A_2& \mathrm{}\\ \mathrm{}& A_1& A_0& A_1& \mathrm{}\\ \mathrm{}& A_2& A_1& A_0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right)𝔤𝔩(\mathrm{}).`$
Let us denote by $`\iota _{\mathrm{}}`$ the above map, which is clearly injective. Let $`X,Y𝔤𝔩_{\mathrm{}}`$ and consider $`Xs^m,Ys^n`$ $`(m,n)`$ as elements of $`𝔤𝔩(\mathrm{}).`$ Then the commutation relation is given by
$`[Xs^m,Ys^n]=[X,Y]s^{m+n}+m\delta _{m+n,0}\mathrm{tr}(XY)c.`$
So, in particular, if we extend $`\iota _{\mathrm{}}`$ to $`\widehat{𝔰𝔩}_{\mathrm{}}=𝔰𝔩_{\mathrm{}}[s^{\pm 1}]K_s`$ by $`\iota _{\mathrm{}}(K_s)\stackrel{\mathrm{def}}{=}c`$, then it gives the embedding $`\iota _{\mathrm{}}`$ of Lie algebras $`\widehat{𝔰𝔩}_{\mathrm{}}`$ into $`𝔤𝔩(\mathrm{}).`$
We shall describe a base of $`\iota _{\mathrm{}}(\widehat{𝔰𝔩}_{\mathrm{}})\widehat{𝔰𝔩}_{\mathrm{}}`$ for the later use.
###### Definition 1
For each $`\mathrm{}`$-th root of unity $`\zeta 1`$, we define $`E_n^\zeta 𝔤𝔩(\mathrm{})(n)`$ by the following series
$`E(\lambda ,\zeta \lambda )={\displaystyle \underset{n}{}}E_n^\zeta \lambda ^{n1}.`$
It is easy to see that $`E_n^\zeta `$ satisfies $`\mathrm{}`$-periodic condition (2.6). If we shall regard $`E_n^\zeta `$ as an element of $`𝔤𝔩_{\mathrm{}}[s^{\pm 1}],`$ then $`E_n^\zeta `$ is traceless. Note that $`E_n^\zeta `$ is homogeneous of the principal degree $`n`$, here by the principal degree of the matrix $`s^ne_{ij}`$ in $`𝔤𝔩_{\mathrm{}}`$ we mean $`ji+n\mathrm{}.`$
Example 1 For $`\mathrm{}=2,\zeta =1`$, we have
$`E_{2i+1}^\zeta =\left(\begin{array}{cc}0& s^i\\ s^{i+1}& 0\end{array}\right),E_{2i}^\zeta =\left(\begin{array}{cc}s^i& 0\\ 0& s^i\end{array}\right)(i).`$
If we define $`\mathrm{\Lambda }_n\stackrel{\mathrm{def}}{=}_{i+\frac{1}{2}}E_{i,i+n}(n\{0\}),`$ then we have the following Heisenberg relations
$`[\mathrm{\Lambda }_m,\mathrm{\Lambda }_n]=m\delta _{m+n,0}c.`$
###### Proposition 1
The matrices $`E_n^\zeta (\zeta ^{\mathrm{}}=1,\zeta 0,n),\mathrm{\Lambda }_n(n,n0mod\mathrm{})`$ form a basis of $`𝔰𝔩_{\mathrm{}}[s^{\pm 1}].`$ Therefore the representation $`\pi _{\mathrm{}}\iota _{\mathrm{}}`$ of $`\widehat{𝔰𝔩}_{\mathrm{}}`$ is given by
$`{\displaystyle \underset{n}{}}E_n^\zeta \lambda ^{n1}:\psi (\lambda )\psi ^{}(\zeta \lambda ):(\zeta ^{\mathrm{}}=1,\zeta 1),`$ (2.7)
$`\mathrm{\Lambda }_n{\displaystyle \underset{i+\frac{1}{2}}{}}:\psi _i\psi _{i+n}^{}:(n,n0mod\mathrm{}),K_s\mathrm{id}.`$
###### Proposition 2
If we identify $`E_m^\zeta ,\mathrm{\Lambda }_m`$ with the elements of $`𝔤𝔩_{\mathrm{}}[s^{\pm 1}],`$ then we have
$`s^mE_n^\zeta =E_{n+m\mathrm{}}^\zeta ,s^m\mathrm{\Lambda }_n=\mathrm{\Lambda }_{n+m\mathrm{}}.`$ (2.8)
Note that the centralizer of $`\mathrm{\Lambda }_{\mathrm{}}`$ in $`𝔤𝔩(\mathrm{})`$ coincide with $`\iota _{\mathrm{}}(\widehat{𝔰𝔩}_{\mathrm{}})+1.`$
###### Lemma 3
Let $`\mathrm{\Omega }_𝒜(\lambda )\stackrel{\mathrm{def}}{=}\psi (\lambda )\psi ^{}(\lambda )`$ be the operator on $`_𝒜^{\mathrm{tor}}\stackrel{\mathrm{def}}{=}_𝒜_𝒜.`$ Then we have
$`[{\displaystyle _{\lambda =\mathrm{}}}\lambda ^n\mathrm{}\mathrm{\Omega }_𝒜(\lambda )𝑑\lambda ,\widehat{𝔰𝔩}_{\mathrm{}}]=0\text{for}n.`$
Proof. We use the expression
$`{\displaystyle _{\lambda =\mathrm{}}}\lambda ^n\mathrm{}\mathrm{\Omega }_𝒜(\lambda )𝑑\lambda ={\displaystyle \underset{i+\frac{1}{2}}{}}\psi _i\psi _{i+n\mathrm{}}^{}.`$ (2.9)
For $`A=_{i,j+\frac{1}{2}}a_{ij}:\psi _i\psi _j^{}:`$ that satisfy (2.6), we have
$`[A1+1A,{\displaystyle \underset{k}{}}\psi _k\psi _{k+n\mathrm{}}^{}]`$
$`=`$ $`{\displaystyle \underset{k}{}}[A,\psi _k]\psi _{k+n\mathrm{}}^{}+{\displaystyle \underset{k}{}}\psi _k[A,\psi _{k+n\mathrm{}}^{}]`$
$`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle \underset{j}{}}a_{jk}\psi _j\psi _{k+n\mathrm{}}^{}+{\displaystyle \underset{k}{}}\psi _k{\displaystyle \underset{j}{}}(a_{k+n\mathrm{},j})\psi _j^{}`$
$`=`$ $`{\displaystyle \underset{k,j}{}}\left(a_{j,kn\mathrm{}}a_{j+n\mathrm{},k}\right)\psi _j\psi _k^{}=0,`$
where we used (2.6) in the last line. Q.E.D.
### 2.4 Equation for the $`SL_{\mathrm{}}^{\mathrm{tor}}`$-orbit of the vacuum
Combining the results in the preceding subsection and Lemma 1, we have the representation of $`\stackrel{~}{𝔰𝔩}_{\mathrm{}}^{\mathrm{tor}}`$ on the space $`_𝒜^{\mathrm{tor}}\stackrel{\mathrm{def}}{=}_𝒜_𝔡`$ defined as $`\pi _{\mathrm{},𝒜}^{\mathrm{tor}}\stackrel{\mathrm{def}}{=}(\iota _{\mathrm{}}\pi _{\mathrm{}})^{\mathrm{tor}}.`$ Now we are in a position to define the following important operator:
$`\mathrm{\Omega }_𝒜^{\mathrm{tor}}(\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{m}{}}\psi (\lambda )V_m(\lambda ^{\mathrm{}})\psi ^{}(\lambda )V_m(\lambda ^{\mathrm{}}).`$
The operator satisfy the following property.
###### Lemma 4
We have
$`[{\displaystyle _{\lambda =\mathrm{}}}\lambda ^j\mathrm{}\mathrm{\Omega }_𝒜^{\mathrm{tor}}(\lambda )𝑑\lambda ,𝔰𝔩_{\mathrm{}}^{\mathrm{tor}}]=0\text{for any}j.`$
Here the contour integral is understood symbolically, namely, just to extract the coefficient of $`\lambda ^1:\lambda ^n𝑑\lambda /(2\pi i)=\delta _{n,1}`$.
Proof. We first note that the operator $`\mathrm{\Omega }_𝒜^{\mathrm{tor}}(\lambda )`$ is the product of $`\mathrm{\Omega }_𝒜(\lambda )=\psi (\lambda )\psi ^{}(\lambda )`$ and $`_mV_m(\lambda ^{\mathrm{}})V_m(\lambda ^{\mathrm{}})`$ that commute with each other. In view of lemma 3 and the fact that the latter operator is a series of $`\lambda ^{\mathrm{}},`$ one can see the lemma above is valid. Q.E.D.
Let $`SL_{\mathrm{}}^{\mathrm{tor}}`$ denote a group of invertible linear transformations on $`_𝒜^{\mathrm{tor}}`$ generated by the exponential action of the elements in $`𝔰𝔩_{\mathrm{}}R`$ acting locally nilpotently.
###### Corollary 1
Let $`\tau _𝒜^{\mathrm{tor}}`$ be in the $`SL_{\mathrm{}}^{\mathrm{tor}}`$-orbit of $`\mathrm{𝐯𝐚𝐜}_𝒜^{\mathrm{tor}}\stackrel{\mathrm{def}}{=}\mathrm{𝐯𝐚𝐜}_𝒜\mathrm{𝐯𝐚𝐜}_𝔡.`$ Then we have
$`{\displaystyle _{\lambda =\mathrm{}}}\lambda ^j\mathrm{}\mathrm{\Omega }_𝒜^{\mathrm{tor}}(\lambda )\left(\tau \tau \right)𝑑\lambda =0\text{for}j0.`$ (2.10)
Proof. By the preceding Lemma, it suffices to show the equations above for $`\tau =\mathrm{𝐯𝐚𝐜}_𝒜^{\mathrm{tor}}.`$ In fact, by the expressions (2.9) and (2.4), one can directly verify this. Q.E.D.
### 2.5 Bosonization
Here we present a summary of the boson-fermion correspondence, that says the space $`_𝒜`$ can be identified with the space of polynomials
$`_{}\stackrel{\mathrm{def}}{=}[x_1,x_2,x_3,\mathrm{}][e^{x_0},e^{x_0}],\mathrm{𝐯𝐚𝐜}_{}\stackrel{\mathrm{def}}{=}1_{}.`$
To state the correspondence precisely, we introduce the charged vacua
$`\mathrm{𝐯𝐚𝐜}_𝒜^{[n]}\stackrel{\mathrm{def}}{=}\{\begin{array}{cc}\psi _{n+\frac{1}{2}}^{}\mathrm{}\psi _{\frac{1}{2}}^{}\mathrm{𝐯𝐚𝐜}_𝒜\hfill & n<0\hfill \\ \mathrm{𝐯𝐚𝐜}_𝒜\hfill & n=0\hfill \\ \psi _{n+\frac{1}{2}}\mathrm{}\psi _{\frac{1}{2}}\mathrm{𝐯𝐚𝐜}_𝒜\hfill & n>0\hfill \end{array}.`$
###### Lemma 5
There exists an unique linear isomorphism $`\sigma :_𝒜_{}`$ such that
$`\sigma (\mathrm{𝐯𝐚𝐜}_𝒜^{[n]})=\mathrm{𝐯𝐚𝐜}_{}^{[n]},\sigma \mathrm{\Lambda }_m\sigma ^1=\{\begin{array}{cc}\frac{}{x_m}\hfill & m>0\hfill \\ mx_m\hfill & m<0\hfill \end{array},`$
where we set $`\mathrm{𝐯𝐚𝐜}_{}^{[n]}\stackrel{\mathrm{def}}{=}e^{nx_0}(n).`$
We introduce the vertex operators by
$`X(\lambda )`$ $`\stackrel{\mathrm{def}}{=}`$ $`\mathrm{exp}\left(\xi (𝒙,\lambda )\right)e^{x_0}\lambda ^{_{x_0}}\mathrm{exp}\left(\xi (\stackrel{~}{}_𝒙,\lambda ^1)\right),`$
$`X^{}(\lambda )`$ $`\stackrel{\mathrm{def}}{=}`$ $`\mathrm{exp}\left(\xi (𝒙,\lambda )\right)e^{x_0}\lambda ^{_{x_0}}\mathrm{exp}\left(\xi (\stackrel{~}{}_𝒙,\lambda ^1)\right),`$
$`X(\lambda ,\mu )`$ $`\stackrel{\mathrm{def}}{=}`$ $`\mathrm{exp}\left(\xi (𝒙,\lambda )\xi (𝒙,\mu )\right)\lambda ^{_{x_0}}\mu ^{_{x_0}}\mathrm{exp}\left(\xi (\stackrel{~}{}_𝒙,\lambda ^1)+\xi (\stackrel{~}{}_𝒙,\mu ^1)\right)`$
where $`\stackrel{~}{}_𝒙`$ stands for $`(_{x_1},\frac{1}{2}_{x_2},\frac{1}{3}_{x_3},\mathrm{})`$ and we set
$`\xi (𝒙,\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}x_n\lambda ^n`$ (2.11)
Then we have the following correspondences of operators:
$`\sigma \psi (\lambda )\sigma ^1=X(\lambda ),\sigma \psi ^{}(\lambda )\sigma ^1=X^{}(\lambda ),`$
$`\sigma (:\psi (\lambda )\psi ^{}(\mu ):)\sigma ^1={\displaystyle \frac{1}{\lambda \mu }}(X(\lambda ,\mu )1),`$
where $`\frac{1}{\lambda \mu }`$ stands for the series $`_{n=0}^{\mathrm{}}\lambda ^{n1}\mu ^n.`$
### 2.6 Hirota bilinear equation arising from toroidal Lie algebras
If we translate equation (2.10) into bosonic language, then it comes out a hierarchy of Hirota bilinear equations. So let us introduce the following bosonic counterpart of operator $`\mathrm{\Omega }_𝒜^{\mathrm{tor}}(\lambda )`$
$`\mathrm{\Omega }_{}^{\mathrm{tor}}(\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{m}{}}X(\lambda )V_m(\lambda ^{\mathrm{}})X^{}(\lambda )V_m(\lambda ^{\mathrm{}})=(\sigma \mathrm{id})\mathrm{\Omega }_𝒜^{\mathrm{tor}}(\lambda )(\sigma ^1\mathrm{id}).`$
We define $`_{}^{\mathrm{tor}}\stackrel{\mathrm{def}}{=}_{}_𝔡`$ and $`\mathrm{𝐯𝐚𝐜}_{}^{\mathrm{tor}}\stackrel{\mathrm{def}}{=}\mathrm{𝐯𝐚𝐜}_{}\mathrm{𝐯𝐚𝐜}_𝔡`$. To state next theorem, we define the elementary Schur polynomials $`p_n(𝒙)(n=0,1,\mathrm{})`$ by the generating series
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}p_n(𝒙)\lambda ^n=\mathrm{exp}\left({\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}x_k\lambda ^k\right),`$
here $`𝒙`$ stands for $`(x_1,x_2,x_3,\mathrm{}).`$
Notation We shall use the notation $`\stackrel{ˇ}{𝒚}\stackrel{\mathrm{def}}{=}(y_{\mathrm{}},y_2\mathrm{},\mathrm{})`$, $`𝒚^{}\stackrel{\mathrm{def}}{=}(y_{\mathrm{}}^{},y_2\mathrm{}^{},\mathrm{})`$, $`y\stackrel{\mathrm{def}}{=}y_0`$, and $`𝒚=(y,\stackrel{ˇ}{𝒚})`$ in the sequel of this paper.
###### Theorem 1
Let $`\tau =\tau (𝐱,𝐲)`$ be in the $`SL_{\mathrm{}}^{\mathrm{tor}}`$-orbit of $`\mathrm{𝐯𝐚𝐜}_{}^{\mathrm{tor}}_{}^{\mathrm{tor}}.`$ Then $`\tau `$ satisfies the bilinear equation:
$`{\displaystyle \underset{m,k=0}{\overset{\mathrm{}}{}}}p_m(2𝒂)p_{m+k\mathrm{}+j\mathrm{}+1}(\stackrel{~}{D_𝒙})p_k((\kappa D_y)𝒃)e^{𝒂,D_𝒙+𝒃,D_{\stackrel{ˇ}{𝒚}}}\tau (𝒙,𝒚)\tau (𝒙,𝒚)=0`$
for $`j0`$ where $`𝐚=(a_1,a_2,\mathrm{}),𝐛=(b_{\mathrm{}},b_2\mathrm{},\mathrm{})`$ and $`\kappa `$ are indeterminates and
$`\stackrel{~}{D_𝒙}\stackrel{\mathrm{def}}{=}(D_{x_1},{\displaystyle \frac{D_{x_2}}{2}},{\displaystyle \frac{D_{x_3}}{3}},\mathrm{}),𝒂,D_𝒙\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_nD_{x_n},𝒃,D_{\stackrel{ˇ}{𝒚}}\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n\mathrm{}D_{y_n\mathrm{}}.`$ (2.12)
Proof. By the bosonic realization, we shall identify $`_{}^{\mathrm{tor}}_{}^{\mathrm{tor}}`$ with the following space of polynomials
$`[e^{\pm x_0^{(i)}},e^{\pm y_0^{(i)}},x_n^{(i)},y_n\mathrm{}^{(i)},y_n\mathrm{}^{(i)}(i=1,2,n=1,2,\mathrm{})].`$
To convert equation (2.10) into bosonic language, we shall introduce new variables,
$`x_n`$ $`\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{2}}(x_n^{(1)}x_n^{(2)}),`$ $`\overline{x}_n`$ $`\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{2}}(x_n^{(1)}+x_n^{(2)})`$ $`\text{for}n0,`$
$`y_n\mathrm{}`$ $`\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{2}}(y_n\mathrm{}^{(1)}y_n\mathrm{}^{(2)}),`$ $`\overline{y_n\mathrm{}}`$ $`\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{2}}(y_n\mathrm{}^{(1)}+y_n\mathrm{}^{(2)})`$ $`\text{for}n0,`$
$`y_n\mathrm{}^{}`$ $`\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{2}}(y_n\mathrm{}^{(1)}y_n\mathrm{}^{(2)}),`$ $`\overline{y_n\mathrm{}^{}}`$ $`\stackrel{\mathrm{def}}{=}{\displaystyle \frac{1}{2}}(y_n\mathrm{}^{(1)}+y_n\mathrm{}^{(2)})`$ $`\text{for}n1.`$
Then we obtain
$`X(\lambda )X^{}(\lambda )=\mathrm{exp}\left(2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}x_n\lambda ^n\right)e^{2x_0}\lambda ^{_{x_0}}\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{}{x_n}}\lambda ^n\right)`$
and
$`{\displaystyle \underset{m}{}}`$ $`V_m(\lambda ^{\mathrm{}})V_m(\lambda ^{\mathrm{}})`$
$`=`$ $`{\displaystyle \underset{m}{}}`$ $`\mathrm{exp}\left(2m{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}y_n\mathrm{}\lambda ^n\mathrm{}\right)e^{2my_0}\mathrm{exp}\left(m{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle \frac{}{y_n\mathrm{}^{}}}\lambda ^n\mathrm{}\right).`$
Using a lemma by Billig (, Proposition 3, see also ), we obtain the Hirota bilinear equation. Here we made a reduction $`D_{y_n\mathrm{}^{}}=0,`$ since, by the construction of the representation, $`\tau `$ in the orbit does not depend on the variables $`y_n\mathrm{}^{}.`$ Q.E.D.
Remark.
1. By the construction of the representation, it is clear that our $`\tau =\tau (𝒙,𝒚)`$ does not depend on the variables $`x_{\mathrm{}},x_2\mathrm{},\mathrm{}.`$
2. If we put $`b_{\mathrm{}}=b_2\mathrm{}=\mathrm{}=0`$, then equation (2.12) is reduced to
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}p_m(2𝒂)p_{m+j\mathrm{}+1}(\stackrel{~}{D_𝒙})e^{𝒂,D_𝒙}\tau \tau =0(j0)`$
This is known as bilinear equation for the $`\mathrm{}`$-reduced KP hierarchy .
Example 2 Let $`\mathrm{}=2`$. From the coefficient of $`a_3`$ for $`j=0`$ , we have a equation:
$`\left(D_x^44D_xD_t\right)\tau \tau =0,\text{where}x=x_1,t=x_3.`$ (2.13)
This can be transformed into a non-linear equation
$`2\rho _{xt}={\displaystyle \frac{1}{2}}\rho _{xxxx}+3\rho _{xx}^2`$ (2.14)
where we set $`\rho =\mathrm{log}\tau .`$ Differentiating by $`x`$, we have the original KdV equation:
$`u_t={\displaystyle \frac{1}{4}}u_{xxx}+{\displaystyle \frac{3}{2}}uu_x,u\stackrel{\mathrm{def}}{=}2\rho _{xx}.`$ (2.15)
Example 3 Let $`\mathrm{}=2`$. The coefficient of $`b_2`$ in (2.12) for $`j=0`$ gives an equation
$`\left({\displaystyle \frac{1}{6}}D_{y_0}D_{x_1}^3+{\displaystyle \frac{1}{3}}D_{y_0}D_{x_3}D_{x_1}D_{y_2}\right)\tau \tau =0.`$ (2.16)
Equation (2.16) times $`6`$ is nothing but (1.3) with $`y_0=y,x_1=x,y_2=z,x_3=t.`$
We shall use the following formulae of the logarithmic transformations:
$`{\displaystyle \frac{D_{x}^{}{}_{}{}^{3}D_y\tau \tau }{\tau ^2}}=12\rho _{xx}\rho _{xy}+2\rho _{xxxy},{\displaystyle \frac{D_xD_y\tau \tau }{\tau ^2}}=2\rho _{xy},`$ (2.17)
where $`\rho \stackrel{\mathrm{def}}{=}\mathrm{log}\tau .`$ These allows us to write equation (2.16) as a non-linear equation:
$`2\rho _{xy_2}=2\rho _{xx}\rho _{xy}+{\displaystyle \frac{1}{3}}\rho _{xxxy}+{\displaystyle \frac{2}{3}}\rho _{yt}.`$ (2.18)
Example 4 Let us derive Bogoyavlensky’s $`2+1`$-dimensional equation by our hierarchy of Hirota equations. We differentiate (2.18) by $`x`$ to give:
$`2\rho _{xxy_2}`$ $`=`$ $`2(\rho _{xx}\rho _{xy})_x+{\displaystyle \frac{1}{3}}\rho _{xxxxy}+{\displaystyle \frac{2}{3}}\rho _{xyt}`$ (2.19)
$`=`$ $`2\rho _{xxx}\rho _{xy}+2\rho _{xx}\rho _{xxy}+{\displaystyle \frac{1}{3}}\rho _{xxxxy}+{\displaystyle \frac{1}{3}}\left({\displaystyle \frac{1}{2}}\rho _{xxxx}+3\rho _{xx}^2\right)_y`$ (2.20)
$`=`$ $`{\displaystyle \frac{1}{2}}\rho _{xxxxy}+2\rho _{xxx}\rho _{xy}+4\rho _{xx}\rho _{xxy}.`$ (2.21)
We have used equation (2.15) in the second line. Putting $`u=2\rho _{xx}`$, it reads
$`u_{y_2}={\displaystyle \frac{1}{4}}u_{xxy}+{\displaystyle \frac{1}{2}}(_x^1u_y)u_x+uu_y,`$ (2.22)
where we interpret $`_x^1u_y`$ as $`2\rho _{xy}.`$ After the simple scale transformation $`_x=2_x^{},_{y_2}=_{y_2^{}},u=4u^{}`$, (2.22) coincide with (1.22) in Chap. II in .
Example 5 Let $`\mathrm{}=3`$. There is a non-trivial equation of degree $`4`$ unique up to scaler, which is not in the $`3`$-reduced KP hierarchy (the Boussinesq hierarchy), i.e.
$`\left(4D_{x_1}D_{y_3}D_{y_0}D_{x_4}D_{y_0}D_{x_1}^2D_{x_2}\right)\tau \tau =0.`$ (2.23)
For instance, this arises from the coefficient of $`b_3`$ in (2.12) for $`j=0`$. Introducing $`\rho =\mathrm{log}\tau `$, we can rewrite (2.23) as a non-linear equation:
$`4\rho _{x_1y_3}\rho _{y_0x_4}\rho _{y_0x_1x_1x_2}4\rho _{y_0x_1}\rho _{x_1x_2}2\rho _{y_0x_2}\rho _{x_1x_1}=0.`$ (2.24)
## 3 Lax formalism of Bogoyavlensky hierarchy
### 3.1 Formal pseudo-differential operators
The Lax formalism of the KP hierarchy is described in the language of formal pseudo-differential operators (PsDO for short) on a line. An affine coordinate $`x`$ of this line is to be identified with the first variable $`x_1`$ in the bosonization of free fermions in the previous section. Let $`_x`$ denote the derivation $`/x`$. In an abstract setting, one can take an arbitrary ring with a derivation $``$ (namely, a differential ring); we shall rather naively consider, e.g., the ring of formal power series of $`x`$.
A formal PsDO is a formal linear combination, $`A=_na_n_x^n`$, of integer powers of $`_x`$ with coefficients $`a_n=a_n(x)`$ that depend on $`x`$. The index $`n`$ ranges over all integers with an upper bound $`N`$. The least upper bound is called the order of this PsDO. The first non-vanishing coefficient $`a_N`$ is called the leading coefficient. If the leading coefficient is equal to $`1`$, the PsDO is said to be monic. It is convenient to use the following notation:
$`(A)_k`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \underset{nk}{}}a_n_x^n,`$
$`(A)_k`$ $`\stackrel{\mathrm{def}}{=}`$ $`{\displaystyle \underset{nk}{}}a_n_x^n,`$
$`(A)_k`$ $`\stackrel{\mathrm{def}}{=}`$ $`a_k.`$
Addition and multiplication (or composition) of two PsDO’s are defined as follows. Addition of two PsDO’s is an obvious operation, namely, the term wise sum of the coefficients. Multiplication is defined by extrapolating the Leibniz rule
$`_x^nf={\displaystyle \underset{k0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{k}}\right)f^{(k)}_x^{nk},`$
to the case where $`n`$ is negative. Here “$``$” stands for composition of two operators, and $`f^{(k)}`$ the $`k`$-th derivative $`^kf/x^k`$ of $`f`$. More explicitly, the product $`C=AB`$ of two PsDO’s $`A=_na_n_x^n`$ and $`B=_nb_n_x^n`$ is given by
$`C={\displaystyle \underset{m,n,k}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{m}{k}}\right)a_m^{(k)}b_n_x^{m+nk}.`$
Note that the $`n`$-th order coefficient $`c_n=(C)_n`$ is the sum of a finite number of terms. We shall frequently write $`AB`$ rather than $`AB`$ if it does not cause confusion.
Another important operation is the formal adjoint $`AA^{}`$:
$`A={\displaystyle \underset{n}{}}a_n_x^nA^{}={\displaystyle \underset{n}{}}(_x)^na_n.`$
This is an anti-homomorphism, namely, for any pair $`A,B`$ of PsDO’s,
$`(AB)^{}=B^{}A^{}.`$
Any PsDO $`A=_{nN}a_n_x^n`$ with an invertible leading coefficient $`a_N`$ has an inverse PsDO. In particular, a monic PsDO is invertible.
### 3.2 KP hierarchy and its $`\mathrm{}`$-reductions
The standard Lax formalism of the KP hierarchy uses a monic first order PsDO (the Lax operator) of the form
$`L\stackrel{\mathrm{def}}{=}_x+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}g_{n+1}_x^n.`$
The hierarchy is defined by the Lax equations
$`{\displaystyle \frac{L}{x_n}}=[B_n,L]`$
with an infinite number of time variables $`𝒙=(x_1,x_2,\mathrm{})`$. The Zakharov-Shabat operators $`B_n`$ are given by
$`B_n\stackrel{\mathrm{def}}{=}(L^n)_0.`$
We identify $`x_1`$ with $`x`$, namely $`x_1=x`$. This is consistent with the first Lax equation
$`{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{x_1}}=[_x,\mathrm{\Psi }(\lambda )]={\displaystyle \frac{\mathrm{\Psi }(\lambda )}{x}}.`$
These Lax equations are associated with the linear equations
$`L\mathrm{\Psi }(\lambda )=\lambda \mathrm{\Psi }(\lambda ),{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{x_n}}=B_n\mathrm{\Psi }(\lambda ).`$
Here $`\mathrm{\Psi }(\lambda )`$ (the “wave function”) is understood to be a function of both $`𝒙`$ and $`\lambda `$, $`\mathrm{\Psi }(\lambda )=\mathrm{\Psi }(𝒙,\lambda )`$, though we shall frequently omit writing $`𝒙`$ explicitly. The Lax equations ensure the existence of a non-vanishing solution of the above linear equations. Of particular importance is a solution (called the formal Baker-Akhiezer function) of the form
$`\mathrm{\Psi }(\lambda )\stackrel{\mathrm{def}}{=}\left(1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}w_n\lambda ^n\right)e^{\xi (\lambda )},`$
where the first factor is generally a formal Laurent series with variable coefficients $`w_n=w_n(𝒙)`$, and $`\xi (\lambda )`$ is given by
$`\xi (\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}x_n\lambda ^n.`$
The exponential factor is the same thing that we have encountered in the bosonization of free fermions. One can make sense of the action of PsDO’s on $`\mathrm{\Psi }(\lambda )`$ by simply extrapolating the derivation rule
$`_x^ne^{\xi (\lambda )}=\lambda ^ne^{\xi (\lambda )}`$
to negative powers of $`_x`$.
Let us now introduce the Wilson-Sato operator
$`W\stackrel{\mathrm{def}}{=}1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}w_n_x^n.`$
This enables us to express the formal Baker-Akhiezer function as:
$`\mathrm{\Psi }(\lambda )=We^{\xi (\lambda )}.`$ (3.1)
The linear equations for $`\mathrm{\Psi }(\lambda )`$ can now be converted to the equations
$`L=W_xW^1,{\displaystyle \frac{W}{x_n}}=B_nWW_x^n`$ (3.2)
for $`W`$. The first equation shows the relation between the Lax operator $`L`$ and the Wilson-Sato operator $`W`$ (note that $`W`$ is monic, hence invertible.) In particular, the Zakharov-Shabat operators can be written
$`B_n=\left(W_x^nW^1\right)_0.`$
If one substitutes the Zakharov-Shabat operators in the second equation for $`W`$ by this expression, the outcome is the equation
$`{\displaystyle \frac{W}{x_n}}=\left(W_x^nW^1\right)_0WW_x^n=\left(W_x^nW^1\right)_1W.`$
Note that this is a nonlinear system of evolution equations for the coefficients $`w_n`$. This is a master equation of the KP hierarchy; the Lax equations can be derived from this nonlinear system via aforementioned relation (3.2) between $`L`$ and $`W`$. It is this nonlinear system that Sato linearized on an infinite dimensional Grassmann manifold (Sato’s universal Grassmannian) .
The $`\mathrm{}`$-reduced hierarchy (the $`\mathrm{}`$-KdV hierarchy) is defined by the constraint
$`(L^{\mathrm{}})_1=0.`$ (3.3)
The constraint means that $`L^{\mathrm{}}`$ be a differential operator, i.e.,
$`P\stackrel{\mathrm{def}}{=}L^{\mathrm{}}=_x^{\mathrm{}}+u_2_x^\mathrm{}2+\mathrm{}+u_{\mathrm{}}.`$
This constraint is preserved under time evolutions, because the associated Lax equations
$`{\displaystyle \frac{P}{x_n}}=[(P^{n/\mathrm{}})_0,P],`$
for $`P`$ become a closed system of unconstrained evolution equations for the coefficients $`u_2,\mathrm{},u_{\mathrm{}}`$ of $`P`$. The time evolutions in $`x_{\mathrm{}},x_2\mathrm{},\mathrm{}`$ are trivial,
$`{\displaystyle \frac{P}{x_n\mathrm{}}}=[P^n,P]=0,`$
and, in turn, give rise to an eigenvalue problem,
$`P\mathrm{\Psi }(\lambda )=\lambda ^{\mathrm{}}\mathrm{\Psi }(\lambda ).`$
The other time evolutions can be interpreted as isospectral deformations of this eigenvalue problem. In the case where $`\mathrm{}=2`$, this is exactly the well known characterization of the KdV hierarchy.
### 3.3 Heuristic Introduction of Bogoyavlensky hierarchy
We here present a heuristic consideration that leads to a hierarchy of higher evolution equations for the Bogoyavlensky equation.
The point of departure is the Lax representation
$`{\displaystyle \frac{P}{z}}=[P_y+C,P]=P{\displaystyle \frac{P}{y}}+[C,P]`$
mentioned in Introduction (1.4). Here $`P`$ is the Lax operator $`_x^2+p`$ of the KdV equation and $`C`$ is a first order differential operator. This Lax equation is associated with the linear equations
$`P\mathrm{\Psi }(\lambda )=\lambda ^2\mathrm{\Psi }(\lambda ),{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{z}}=(P_y+C)\mathrm{\Psi }(\lambda ).`$
Note that the second linear equation can be rewritten
$`{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{z}}=(\lambda ^2_y+A)\mathrm{\Psi }(\lambda ),`$
where
$`A\stackrel{\mathrm{def}}{=}C{\displaystyle \frac{P}{y}}.`$
Furthermore, if we assume the existence of the Wilson-Sato operator $`W`$ with which $`\mathrm{\Psi }(\lambda )`$ takes the form
$`\mathrm{\Psi }(\lambda )=We^{\xi (\lambda )},\xi (\lambda )=x\lambda +\mathrm{},`$
the time evolution of $`W`$ is determined by the equation
$`{\displaystyle \frac{W}{z}}`$ $`=`$ $`P{\displaystyle \frac{W}{y}}+CW`$
$`=`$ $`P[_y,W]+CW`$
$`=`$ $`(P_y+C)WW_x^2_y.`$
We have used the identity $`PW=W_x^2`$ in the last two lines.
If one compares the last equation with the equations in the KP hierarchy, one will soon notice that the role of $`_x^n`$ is now played by $`_x^2_y`$. The (two-dimensional!) differential operator $`P_y+C`$ should be a counterpart of $`B_n`$’s. A natural generalization of this observation is to consider the higher order operators $`_x^{2n}_y`$, $`n=1,2,\mathrm{}`$. The associated evolution equations, with “time variables” $`y_2(=z),y_4,\mathrm{}`$, will accordingly take such a form as
$`{\displaystyle \frac{W}{y_{2n}}}=D_{2n}WW_x^{2n}_y`$
where $`D_{2n}`$ is a two-dimensional differential operator in $`x`$ and $`y`$. In fact, one can specify $`D_{2n}`$ in more detail: This equation can be rewritten
$`{\displaystyle \frac{W}{y_{2n}}}=(D_{2n}P^n_y)W+P^n{\displaystyle \frac{W}{y}},`$
and this implies that $`D_{2n}P^n_y`$ is a differential operator in $`x`$ only, which we call $`C_{2n}`$:
$`D_{2n}=P^n_y+C_{2n}.`$
We are thus led to the evolution equations
$`{\displaystyle \frac{W}{y_{2n}}}`$ $`=`$ $`(P^n_y+C_{2n})WW_x^{2n}_y`$
$`=`$ $`P^n{\displaystyle \frac{W}{y}}+C_{2n}W.`$
Just like the Zakharov-Shabat operators $`B_n`$ of the KP hierarchy, the differential operator $`C_{2n}`$ is uniquely determined by the evolution equation itself:
$`C_{2n}=\left({\displaystyle \frac{W}{y_{2n}}}W^1P^n{\displaystyle \frac{W}{y}}W^1\right)_0=\left(P^n{\displaystyle \frac{W}{y}}W^1\right)_0.`$
Therefore the evolution equation of $`W`$ can be rewritten
$`{\displaystyle \frac{W}{y_{2n}}}=\left(P^n{\displaystyle \frac{W}{y}}W^1\right)_1W.`$
These equations, along with the aforementioned evolution equations in $`x_{2n+1}`$’s, give a nonlinear system of evolution equations for the coefficients $`w_n`$ of $`W`$. We call this system the Bogoyavlensky hierarchy or, more precisely, the Bogoyavlensky-KdV hierarchy.
The spatial variable $`y`$ may be identified with the zero-th time variable $`y_0`$, because the corresponding operator $`C_0`$ is equal to zero so that the evolution equation for $`W`$ reduces to
$`{\displaystyle \frac{W}{y_0}}={\displaystyle \frac{W}{y}}.`$
It is straightforward to write down the associated linear equations. They take two equivalent (but apparently different) forms. One expression is a direct consequence of the evolution equation of $`W`$:
$`{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{y_{2n}}}=(P^n_y+C_{2n})\mathrm{\Psi }(\lambda ).`$
Another expression is the following:
$`{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{y_{2n}}}=(\lambda ^{2n}_y+A_{2n})\mathrm{\Psi }(\lambda ),`$
where
$`A_{2n}\stackrel{\mathrm{def}}{=}C_{2n}{\displaystyle \frac{P^n}{y}}.`$
### 3.4 Lax equations of Bogoyavlensky hierarchy
Let us derive Lax equations of the Bogoyavlensky hierarchy. To this end, we introduce the PsDO
$`Q\stackrel{\mathrm{def}}{=}{\displaystyle \frac{W}{y}}W^1`$
as the second Lax operator. This operator, just like $`P=W_x^2W^1`$, is related to conjugation by the Wilson-Sato operator:
$`_yQ=W_yW^1.`$
The commutation relation $`[_y,_x^2]=0`$ thereby implies the commutation relation
$`[Q_y,P]=0`$
or, equivalently,
$`{\displaystyle \frac{P}{y}}=[Q,P].`$
Note that this relation itself takes the form of a Lax equation.
These operators have already appeared in the aforementioned differential operators $`C_{2n}`$ and $`A_{2n}`$. $`C_{2n}`$ can be indeed written
$`C_{2n}=(P^nQ)_0.`$
Moreover, using the Lax-type equation
$`{\displaystyle \frac{P^n}{y}}=[Q,P^n]`$
(which is a consequence of the commutation relations of $`P`$ and $`Q`$), one can confirm that $`A_{2n}`$ can be written
$`A_{2n}=(QP^n)_0.`$
Having these operators, we can now write down Lax equations of the Bogoyavlensky hierarchy:
###### Proposition 3
The operators $`P`$ and $`Q`$ satisfy the equations
$`[_{x_{2n+1}}B_{2n+1},P]=0,`$ $`[_{x_{2n+1}}B_{2n+1},Q_y]=0,`$
$`[_{y_{2n}}P^n_yC_{2n},P]=0,`$ $`[_{y_{2n}}P^n_yC_{2n},Q_y]=0,`$
or, equivalently,
$`{\displaystyle \frac{P}{x_{2n+1}}}`$ $`=`$ $`[B_{2n+1},P],`$
$`{\displaystyle \frac{Q}{x_{2n+1}}}`$ $`=`$ $`{\displaystyle \frac{B_{2n+1}}{y}}+[B_{2n+1},Q],`$
$`{\displaystyle \frac{P}{y_{2n}}}`$ $`=`$ $`P^n{\displaystyle \frac{P}{y}}+[C_{2n},P],`$
$`{\displaystyle \frac{Q}{y_{2n}}}`$ $`=`$ $`P^n{\displaystyle \frac{Q}{y}}+{\displaystyle \frac{C_{2n}}{y}}+[C_{2n},Q].`$
Proof. Let us prove the equations for the $`y_{2n}`$-derivatives; the proof of the other equations are parallel and even simpler. First, differentiating $`P=W_x^2W^1`$ gives
$`{\displaystyle \frac{P}{y_{2n}}}`$ $`=`$ $`{\displaystyle \frac{W}{y_{2n}}}_x^2W^1W_x^2W^1{\displaystyle \frac{W}{y_{2n}}}W^1`$
$`=`$ $`[{\displaystyle \frac{W}{y_{2n}}}W^1,P].`$
We can now use the evolution equations
$`{\displaystyle \frac{W}{y_{2n}}}=P^nQW+C_{2n}W`$
and the commutation relation $`[Q,P]=P/y`$, and find that
$`{\displaystyle \frac{P}{y_{2n}}}`$ $`=`$ $`[P^nQ+C_{2n},P]`$
$`=`$ $`P^n{\displaystyle \frac{Q}{y}}+[C_{2n},P].`$
Similarly, from the identity $`Q_y=W_yW^1`$,
$`{\displaystyle \frac{Q}{y_{2n}}}`$ $`=`$ $`[{\displaystyle \frac{W}{y_{2n}}}W^1,Q_y]`$
$`=`$ $`[P^nQ+C_{2n},Q_y]`$
$`=`$ $`[P^n,Q]Q+{\displaystyle \frac{P^nQ}{y}}+{\displaystyle \frac{C_{2n}}{y}}+[C_{2n},Q]`$
$`=`$ $`[P^n,Q]Q+{\displaystyle \frac{P^n}{y}}Q+P^n{\displaystyle \frac{Q}{y}}+{\displaystyle \frac{C_{2n}}{y}}+[C_{2n},Q].`$
The first two terms in the last line cancel each other, again because of the commutation relation of $`Q`$ and $`P`$. This completes the proof. Q.E.D.
### 3.5 Formulation of $`\mathrm{}`$-Bogoyavlensky hierarchy
We are now in a position to formulate the $`\mathrm{}`$-Bogoyavlensky hierarchy for a general value of $`\mathrm{}2`$. This hierarchy is an extension of the $`\mathrm{}`$-reduced KP hierarchy with additional independent variables $`𝒚=(y,\stackrel{ˇ}{𝒚})`$, where $`y=y_0`$ and $`\stackrel{ˇ}{𝒚}=(y_{\mathrm{}},y_2\mathrm{},\mathrm{})`$. The coefficients $`w_n`$ of the Wilson-Sato operator
$`W\stackrel{\mathrm{def}}{=}1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}w_n_x^n`$
are understood to be functions of these variables also, $`w_n=w_n(𝒙,𝒚).`$ Assume that $`W`$ satisfy the $`\mathrm{}`$-reduced constraints, cf.(3.3):
$`\left(W_x^{\mathrm{}}W^1\right)_1=0,`$
and define the differential operator
$`P\stackrel{\mathrm{def}}{=}W_x^{\mathrm{}}W^1.`$
* The most fundamental part of this hierarchy is the evolution equations
$`{\displaystyle \frac{W}{x_n}}`$ $`=`$ $`\left(W_x^nW^1\right)_1W,`$
$`{\displaystyle \frac{W}{y_n\mathrm{}}}`$ $`=`$ $`\left(W_x^n\mathrm{}W^1{\displaystyle \frac{W}{y}}W^1\right)_1W.`$
of $`W`$. These equations can also be written
$`{\displaystyle \frac{W}{x_n}}`$ $`=`$ $`B_nWW_x^n,`$ (3.4)
$`{\displaystyle \frac{W}{y_n\mathrm{}}}`$ $`=`$ $`W_x^n\mathrm{}W^1{\displaystyle \frac{W}{y}}+C_n\mathrm{}W`$
$`=`$ $`(W_x^n\mathrm{}W^1_y+C_n\mathrm{})WW_x^n\mathrm{}_y,`$
where $`B_n`$ and $`C_n\mathrm{}`$ are differential operators defined by
$`B_n`$ $`\stackrel{\mathrm{def}}{=}`$ $`\left(W_x^nW^1\right)_0,`$
$`C_n\mathrm{}`$ $`\stackrel{\mathrm{def}}{=}`$ $`\left(W_x^n\mathrm{}W^1{\displaystyle \frac{W}{y}}W^1\right)_0.`$ (3.6)
* The evolution equations of $`W`$ in $`y_n\mathrm{}`$’s can also be rewritten
$`{\displaystyle \frac{W}{y_n\mathrm{}}}={\displaystyle \frac{W}{y}}_x^n\mathrm{}+A_n\mathrm{}W,`$
where
$`A_n\mathrm{}\stackrel{\mathrm{def}}{=}(QP^n)_0=C_n\mathrm{}{\displaystyle \frac{P^n}{y}}.`$ (3.7)
* Lax equations of this hierarchy are described by $`P`$ and the following PsDO:
$`Q\stackrel{\mathrm{def}}{=}{\displaystyle \frac{W}{y}}W^1.`$
They satisfy the commutation relation
$`[Q,P]={\displaystyle \frac{P}{y}}.`$
$`B_n`$ and $`C_n\mathrm{}`$ can now be expressed as
$`B_n=(P^{n/\mathrm{}})_0,C_n\mathrm{}=(P^nQ)_0.`$
$`P`$ and $`Q`$ satisfy the Lax equations:
$`{\displaystyle \frac{P}{x_n}}`$ $`=`$ $`[B_n,P],`$
$`{\displaystyle \frac{Q}{x_n}}`$ $`=`$ $`{\displaystyle \frac{B_n}{y}}+[B_n,Q],`$
$`{\displaystyle \frac{P}{y_n\mathrm{}}}`$ $`=`$ $`P^n{\displaystyle \frac{P}{y}}+[C_n\mathrm{},P],`$
$`{\displaystyle \frac{Q}{y_n\mathrm{}}}`$ $`=`$ $`P^n{\displaystyle \frac{Q}{y}}+{\displaystyle \frac{C_n\mathrm{}}{y}}+[C_n\mathrm{},Q].`$
* The formal Baker-Akhiezer function
$`\mathrm{\Psi }(\lambda )\stackrel{\mathrm{def}}{=}We^{\xi (\lambda )}`$ (3.8)
satisfies the linear equations
$`P\mathrm{\Psi }(\lambda )=\lambda ^{\mathrm{}}\mathrm{\Psi }(\lambda ),`$ $`Q\mathrm{\Psi }(\lambda )={\displaystyle \frac{\mathrm{\Psi }(\lambda )}{y}},`$
$`{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{x_n}}=B_n\mathrm{\Psi }(\lambda ),`$ $`{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{y_n\mathrm{}}}=\left(P^n_y+C_n\mathrm{}\right)\mathrm{\Psi }(\lambda ).`$ (3.9)
The last equation can also be written
$`{\displaystyle \frac{\mathrm{\Psi }(\lambda )}{y_n\mathrm{}}}=\left(\lambda ^n\mathrm{}_y+A_n\mathrm{}\right)\mathrm{\Psi }(\lambda ).`$ (3.10)
Example 6 Let us consider the case $`\mathrm{}=2`$. The Lax operator has the form $`P=_x^2+u`$ where $`u=2w_{1x}.`$ We can write down $`C_2=A_2+\frac{P}{y}`$ as follows:
$`C_2=c_1_x+c_2,c_1=w_{1y},c_2=w_{1y}w_1w_{2y}2w_{1xy}.`$ (3.11)
The evolution equation with respect to $`y_2`$ reads
$`{\displaystyle \frac{P}{y_2}}=P{\displaystyle \frac{P}{y}}+[C_2,P].`$ (3.12)
One can calculate the right-hand side of (3.12) as follows:
$`P{\displaystyle \frac{P}{y}}+[C_2,P]`$ (3.13)
$`=`$ $`(2u_{xy}c_{1xx}2c_{2x})_x+u_{xxy}+uu_y+c_1u_xc_{2xx},`$ (3.14)
where the coefficients of $`_x^2`$ cancel trivially because $`u=2w_{1x}`$.
Here we use the following relation which is a direct consequence of the $`2`$-reduced condition
$`w_{1xx}+2w_{2x}2w_1w_{1x}=0.`$ (3.15)
Using this, we now have $`c_{2x}=\frac{3}{2}w_{1xxy}.`$ This finally leads to
$`{\displaystyle \frac{u}{y_2}}={\displaystyle \frac{1}{4}}u_{xxy}+vu_x+uu_y\text{where}v\stackrel{\mathrm{def}}{=}w_{1y}.`$ (3.16)
Caution! The expression $`w_{2x}`$, for example, does not mean $`_x^2w`$ but $`_xw_2`$.
Example 7 We shall demonstrate the third order flow $`y_3`$ in case of $`\mathrm{}=3`$. The Lax operator $`P=_x^3+U_x+V`$ is given as:
$`U=3w_{1x},V=3w_{1xx}3w_{2x}+3w_1w_{1x}.`$ (3.17)
Operator $`C_3=c_1_x^2+c_2_x+c_3`$ is given explicitly as:
$`c_1`$ $`=`$ $`w_{1y},c_2=w_{1y}w_1w_{2y}3w_{1xy},`$ (3.18)
$`c_3`$ $`=`$ $`w_{3y}w_{1y}w_1^2+w_1(w_{2y}+3w_{1xy})+w_{1y}(w_2+5w_{1x})3w_{1xxy}3w_{2xy}.`$ (3.19)
Now we recall the $`3`$-reduced condition. In particular, the coefficient of $`_x^1`$ in $`W_x^3W^1`$ should vanish identically. This condition is written as:
$`w_{3x}=\frac{1}{3}w_{1xxx}w_{2xx}+w_{1x}^2+w_{1x}w_2+w_1w_{1xx}+w_1w_{2x}w_{1x}w_1^2.`$ (3.20)
Using (3.17),(3.18), and (3.19), it is simple to see that the right-hand side of evolution equation,
$`_{y_3}P=P{\displaystyle \frac{P}{y}}+[C_3,P],`$ (3.21)
is of order at most two. To verify that the right-hand side is actually of first order, we eliminate $`w_3`$ by using equation (3.20). Thus, it should be emphasized that the $`3`$-reduced condition is crucial to introduce flow $`y_3`$ by (3.21).
After eliminating $`w_3`$, we have the following explicit forms of evolution equations:
$`_{y_3}U`$ $`=`$ $`6w_{1y}w_{2xx}w_{1xxxxy}+9w_{1xy}w_{2x}2w_{2xxxy}12w_{1y}w_{1x}^215w_{1xy}w_1w_{1x}`$
$`+6w_{1x}w_{2xy}9w_{1y}w_1w_{1xx}+12w_{1xy}w_{1xx}+9w_{1xxy}w_{1x}+2w_1w_{1xxxy}`$
$`+5w_{1y}w_{1xxx}+3w_{1xx}w_{2y},`$
$`_{y_3}V`$ $`=`$ $`9w_{1xy}w_{1x}^218w_{1xy}w_1w_{1xx}+3w_{1xxx}w_{2y}+3w_{1x}w_{2xxy}3w_{1x}^2w_{2y}`$
$`15w_{1y}w_{1x}w_{1xx}+9w_{1xy}w_{2xx}+9w_{1xxy}w_{2x}12w_{1xxy}w_1w_{1x}`$
$`+3w_{1y}w_{2xxx}6w_{1y}w_1w_{1xxx}+2w_{1y}w_{1xxxx}+7w_{1xy}w_{1xxx}+9w_{1xxy}w_{1xx}`$
$`+3w_{1xxxy}w_{1x}+9w_{1xx}w_{2xy}3w_{2y}w_1w_{1xx}+12w_{1y}w_1w_{1x}^2+3w_{1y}w_1^2w_{1xx}`$
$`9w_{2x}w_1w_{1xy}+9w_{2x}w_{2xy}+3w_{2y}w_{2xx}9w_{2x}w_{1y}w_{1x}3w_{1y}w_1w_{2xx}`$
$`9w_1w_{1x}w_{2xy}+9w_1^2w_{1x}w_{1xy}w_{2xxxxy}\frac{1}{3}w_{1xxxxxy}+w_1w_{1xxxxy}.`$
### 3.6 Explicit forms of operators $`A_n`$
We give a supplementary discussion on the forms of differential operators $`A_n`$. The differential operators $`A_n\mathrm{}`$ and $`C_n\mathrm{}`$ are written in terms of $`W`$ (3.7),(3.6) and hence in the coefficients $`w_i`$’s of $`W.`$ In particular, we have
$`A_n={\displaystyle \underset{i=1}{\overset{n}{}}}a_i^{(n)}_x^{ni}=\left({\displaystyle \frac{W}{y}}_x^nW^1\right)_0.`$
Here we shall recursively determine the coefficients $`a_i^{(n)}.`$ If we cast $`\mathrm{\Psi }(\lambda )=We^{\xi (\lambda )}`$ into linear equation (3.10), then we obtain
$`{\displaystyle \frac{W}{y_n}}={\displaystyle \underset{i=1}{\overset{n}{}}}a_i^{(n)}(_x+\lambda )^{ni}W+\lambda ^n{\displaystyle \frac{W}{y}}.`$
Taking the polynomial part of the equation above with respect to $`\lambda ,`$ we obtain
$`{\displaystyle \underset{i=1}{\overset{n}{}}}a_i^{(n)}(_x+\lambda )^{ni}\left({\displaystyle \underset{j=0}{\overset{n}{}}}w_j\lambda ^j\right)+{\displaystyle \underset{j=1}{\overset{n}{}}}(_yw_j)\lambda ^{nj}=0,w_0\stackrel{\mathrm{def}}{=}1.`$
Equating coefficients of $`\lambda ^{n1}`$ in the equation above, we obtain $`a_1^{(n)}=_yw_1.`$ Moreover, if $`a_1^{(1)},\mathrm{},a_{i1}^{(n)}`$ are already known, then $`a_i^{(n)}`$ can be uniquely determined by the coefficients of $`\lambda ^{ni}`$. In particular, we have $`a_2^{(n)}=w_1_yw_1_yw_2.`$ Note that $`a_i^{(n)}`$ is a differential polynomial of $`w_1,\mathrm{},w_i`$ with respect to $`_x`$ and $`_y.`$
Example 8 For any $`n`$ we have $`a_1^{(n)}=w_{1y},a_2^{(n)}=w_{1y}w_1w_{2y}.`$ We shall give the formulae for $`n5`$:
$`a_3^{(3)}`$ $`=`$ $`w_{3y}w_{1y}w_1^2+w_1w_{2y}+w_{1y}w_2+2w_{1y}w_{1x},`$
$`a_3^{(4)}`$ $`=`$ $`w_1w_{2y}w_{1y}w_1^2+3w_{1y}w_{1x}w_{3y}+w_{1y}w_2,`$
$`a_4^{(4)}`$ $`=`$ $`w_{4y}w_1^2w_{2y}+w_{1y}w_1^35w_{1x}w_{1y}w_1+w_1w_{3y}2w_2w_{1y}w_1`$
$`+w_2w_{2y}+2w_{1x}w_{2y}+w_{1y}w_3+3w_{1y}w_{2x}+3w_{1y}w_{1xx},`$
$`a_3^{(5)}`$ $`=`$ $`w_{1y}w_1^2+w_1w_{2y}w_{3y}+4w_{1y}w_{1x}+w_{1y}w_2,`$
$`a_4^{(5)}`$ $`=`$ $`6w_{1y}w_{1xx}+w_{1y}w_1^3w_1^2w_{2y}+w_1w_{3y}7w_{1x}w_{1y}w_12w_2w_{1y}w_1`$
$`+3w_{1x}w_{2y}+w_2w_{2y}w_{4y}+4w_{1y}w_{2x}+w_{1y}w_3,`$
$`a_5^{(5)}`$ $`=`$ $`w_{5y}+w_{1y}w_4+w_3w_{2y}+3w_{1xx}w_{2y}+3w_{2x}w_{2y}+4w_{1y}w_{1xxx}+3w_2w_{1y}w_1^2`$ (3.23)
$`2w_2w_1w_{2y}6w_2w_{1y}w_{1x}+9w_{1x}w_{1y}w_1^25w_{1x}w_1w_{2y}2w_3w_{1y}w_19w_{1xx}w_{1y}w_1`$
$`7w_{2x}w_{1y}w_1+6w_{1y}w_{2xx}+4w_{1y}w_{3x}+w_2w_{3y}w_{1y}w_2^2+2w_{1x}w_{3y}8w_{1y}w_{1x}^2`$
$`w_{1y}w_1^4+w_1^3w_{2y}w_1^2w_{3y}+w_1w_{4y}.`$
## 4 Reproducing bilinear equations in Lax formalism
### 4.1 Dual formal Baker-Akhiezer function
We define the dual $`\mathrm{\Psi }^{}(\lambda )`$ of the formal Baker-Akhiezer function $`\mathrm{\Psi }(\lambda )`$ as follows:
$`\mathrm{\Psi }^{}(\lambda )\stackrel{\mathrm{def}}{=}W^1e^{\xi (\lambda )}.`$
###### Proposition 4
$`\mathrm{\Psi }^{}(\lambda )`$ satisfies the linear equations
$`P^{}\mathrm{\Psi }(\lambda )=\lambda ^{\mathrm{}}\mathrm{\Psi }^{}(\lambda ),`$ $`Q^{}\mathrm{\Psi }(\lambda )={\displaystyle \frac{\mathrm{\Psi }^{}(\lambda )}{y}},`$
$`{\displaystyle \frac{\mathrm{\Psi }^{}(\lambda )}{x_n}}=B_n^{}\mathrm{\Psi }^{}(\lambda ),`$ $`{\displaystyle \frac{\mathrm{\Psi }^{}(\lambda )}{y_n\mathrm{}}}=\left(P^n_yA_n\mathrm{}\right)\mathrm{\Psi }^{}(\lambda ).`$
The last equation can also be written
$`{\displaystyle \frac{\mathrm{\Psi }^{}(\lambda )}{y_n\mathrm{}}}=\left(\lambda ^n\mathrm{}_yC_n\mathrm{}^{}\right)\mathrm{\Psi }^{}(\lambda ).`$
Proof. The first and second equations are immediate from the operator relations
$`P^{}W^1`$ $`=`$ $`W^1_x,`$
$`Q^{}W^1`$ $`=`$ $`W^1{\displaystyle \frac{W}{y}}W^1={\displaystyle \frac{W^1}{y}}.`$
The third equation can be derived as follows:
$`{\displaystyle \frac{W^1}{x_n}}`$ $`=`$ $`W^1{\displaystyle \frac{W^{}}{x_n}}W^1`$
$`=`$ $`W^1\left(B_nWW_x^n\right)^{}W^1`$
$`=`$ $`B_n^{}W^1+W^1(_x)^n.`$
Similarly,
$`{\displaystyle \frac{W^1}{y_n\mathrm{}}}`$ $`=`$ $`W^1{\displaystyle \frac{W^{}}{y_n\mathrm{}}}W^1`$
$`=`$ $`W^1\left(P^n{\displaystyle \frac{W}{y_n\mathrm{}}}+C_n\mathrm{}W\right)^{}W^1`$
$`=`$ $`{\displaystyle \frac{W^1}{y}}(_x)^n\mathrm{}C_n\mathrm{}W^1,`$
which implies the fourth equation. The last equation can be readily derived from the fourth equation. Q.E.D.
### 4.2 Residue formula of product of PsDO’s
The following lemma is a clue for deriving the bilinear equation that $`\mathrm{\Psi }(\lambda )`$ and $`\mathrm{\Psi }^{}(\lambda )`$ satisfy.
###### Lemma 6
For any pair $`A,B`$ of PsDO’s and non-negative integer $`j`$,
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\lambda =\mathrm{}}}\left(_x^jAe^{x\lambda }\right)\left(Be^{x\lambda }\right)𝑑\lambda =(1)^{j+1}(BA^{})_{j1}.`$
In particular, if $`A`$ and $`B`$ are PsDO’s of zero-th order of the form $`A=a_0+O(_x^1)`$ and $`B=b_0+O(_x^1)`$,
$`BA^{}=b_0a_0(j0){\displaystyle _{\lambda =\mathrm{}}}\left(_x^jAe^{x\lambda }\right)\left(Be^{x\lambda }\right)𝑑\lambda =0.`$
Proof. Let $`A`$ and $`B`$ be written $`A=_na_n_x^n`$ and $`B=_nb_n_x^n`$. Since
$`_x^jA={\displaystyle \underset{m}{}}{\displaystyle \underset{k0}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{j}{k}}\right)a_m^{(k)}_x^{j+mk},`$
the contour integral can be calculated as:
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\lambda =\mathrm{}}}\left(_x^kAe^{x\lambda }\right)\left(Be^{x\lambda }\right)𝑑\lambda `$
$`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\lambda =\mathrm{}}}{\displaystyle \underset{j,k}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{j}{k}}\right)a_m^{(k)}\lambda ^{j+mk}{\displaystyle \underset{n}{}}b_n(\lambda )^nd\lambda `$
$`=`$ $`{\displaystyle \underset{kmn1=j}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{j}{k}}\right)(1)^na_m^{(k)}b_n.`$
On the other hand,
$`BA^{}`$ $`=`$ $`{\displaystyle \underset{m,n}{}}b_n_x^n(_x)^ma_m`$
$`=`$ $`{\displaystyle \underset{m,n,k}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{m+n}{k}}\right)(1)^ma_m^{(k)}_x^{m+nk}`$
$`=`$ $`{\displaystyle \underset{m,n,k}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{kmn1}{k}}\right)(1)^{km}a_m^{(k)}b_n_x^{m+nk}`$
$`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \underset{kmn1=j}{}}\left({\displaystyle \genfrac{}{}{0pt}{}{j}{k}}\right)(1)^na_m^{(k)}b_n(_x)^{j1}.`$
Therefore the coefficient of $`_x^{j1}`$ coincides with the countour integral. The lemma is thus proven. Q.E.D.
### 4.3 Bilinearization of linear equations
Let us slightly change the notation, namely, we make explicit the functional dependence on $`𝒙`$ and $`𝒚,`$ e.g., $`\mathrm{\Psi }(𝒙,𝒚,\lambda )`$, $`\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )`$ and $`\xi (𝒙,\lambda )`$ stand for $`\mathrm{\Psi }(\lambda )`$, $`\mathrm{\Psi }^{}(\lambda )`$ and $`\xi (\lambda )`$ in the previous notation.
We now apply Lemma 6 to the case where $`A=W`$ and $`B=W^1`$. This eventually leads to a bilinear equation for $`\mathrm{\Psi }(𝒙,𝒚,\lambda )`$ and $`\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )`$ as follows:
1. By the lemma applied to $`A=W`$ and $`B=W^1`$ gives the bilinear identities
$`{\displaystyle _{\lambda =\mathrm{}}}\left(_x^j\mathrm{\Psi }(𝒙,𝒚,\lambda )\right)\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )𝑑\lambda =0`$
for $`j0`$ (note that $`x_2\lambda ^2+x_3\lambda ^3+\mathrm{}`$ in $`e^{\pm \xi (𝒙,\lambda )}`$ cancel out). Thus, we have
$`{\displaystyle _{\lambda =\mathrm{}}}\left(R\mathrm{\Psi }(𝒙,𝒚,\lambda )\right)\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )𝑑\lambda =0`$
for any differential operator $`R=_{n0}r_n(𝒙,𝒚)_x^n`$ in $`x.`$
2. Iteration of the evolution equation of $`\mathrm{\Psi }(𝒙,𝒚,\lambda )`$ gives rise to higher order equations of the form
$`_{x_1}^{\alpha _1}_{x_2}^{\alpha _2}\mathrm{}\mathrm{\Psi }(𝒙,𝒚,\lambda )=B_{\alpha _1,\alpha _2,\mathrm{}}\mathrm{\Psi }(𝒙,𝒚,\lambda )`$
for $`\alpha _1,\alpha _2,\mathrm{}0`$, $`B_{\alpha _1,\alpha _2,\mathrm{}}`$ being a differential operator in $`x`$.
3. Combining these equations with the last bilinear identity, we obtain the bilinear equations
$`{\displaystyle _{\lambda =\mathrm{}}}\left(_{x_1}^{\alpha _1}_{x_2}^{\alpha _2}\mathrm{}\mathrm{\Psi }(𝒙,𝒚,\lambda )\right)\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )𝑑\lambda =0`$
for $`\alpha _1,\alpha _2,\mathrm{}0`$.
4. These bilinear equations can be cast into a generating function: Introduce an infinite number of new variables $`𝒂=(a_1,a_2,\mathrm{})`$ , and sum up the bilinear equations over $`\alpha _1,\alpha _2,\mathrm{}0`$ with the weight $`a_1^{\alpha _1}a_2^{\alpha _2}\mathrm{}/\alpha _1!\alpha _2!\mathrm{}`$. The outcome is a single bilinear equation of the form
$`{\displaystyle _{\lambda =\mathrm{}}}\mathrm{\Psi }(𝒙+𝒂,𝒚,\lambda )\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )𝑑\lambda =0.`$
5. Finally, rewriting $`𝒙+𝒂`$ to $`𝒙^{}`$, we obtain the bilinear identity
$`{\displaystyle _{\lambda =\mathrm{}}}\mathrm{\Psi }(𝒙^{},𝒚,\lambda )\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )𝑑\lambda =0.`$
Remark. Note that the status of $`x_{\mathrm{}},x_2\mathrm{},\mathrm{}`$ are somewhat different from others: Since $`w_n`$ do not depend on these variables, they can appear in $`e^{\pm \xi (𝒙,\lambda )}`$ only. The net effect of these variables is thereby insertion of the exponential factor
$`\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(x_n\mathrm{}^{}x_n\mathrm{})\lambda ^n\mathrm{}\right)`$
in the contour integral of the bilinear equations. Therefore the last bilinear equation can be further reduced to the bilinear equations
$`{\displaystyle _{\lambda =\mathrm{}}}\lambda ^j\mathrm{}\mathrm{\Psi }(𝒙^{},𝒚,\lambda )\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )|_{x_{\mathrm{}}^{}=x_2\mathrm{}^{}=\mathrm{}=x_{\mathrm{}}=x_2\mathrm{}=\mathrm{}=0}d\lambda =0`$
for $`j0`$.
This is not the end of the story. What we have done is simply to bilinearize the KP-like part of the $`\mathrm{}`$-Bogoyavlensky hierarchy. The other evolution equations $`y_n\mathrm{}`$’s still remain to be bilinearized.
Bilinearization of the evolution equations in $`y_n\mathrm{}`$’s can be achieved by the following steps, which are almost parallel to the previous calculations:
1. We rewrite the evolution equations as
$`(_{y_n\mathrm{}}\lambda ^n\mathrm{}_y)\mathrm{\Psi }(𝒙,𝒚,\lambda )=A_n\mathrm{}\mathrm{\Psi }(𝒙,𝒚,\lambda )`$
and again do iteration. This yields the higher order equations
$`(_y_{\mathrm{}}\lambda ^{\mathrm{}}_y)^{\beta _1}(_{y_2\mathrm{}}\lambda ^2\mathrm{}_y)^{\beta _2}\mathrm{}\mathrm{\Psi }(𝒙,𝒚,\lambda )`$
$`=`$ $`A_{\beta _1,\beta _2,\mathrm{}}\mathrm{\Psi }(𝒙,𝒚,\lambda )`$
for $`\beta _1,\beta _2,\mathrm{}0`$, $`A_{\beta _1,\beta _2,\mathrm{}}`$ being a differential operator in $`x`$.
2. Combining these equations with the bilinear equations that we have derived above, we now obtain the bilinear equations
$`{\displaystyle _{\lambda =\mathrm{}}}`$ $`\left((_y_{\mathrm{}}\lambda ^{\mathrm{}}_y)^{\beta _1}(_{y_2\mathrm{}}\lambda ^2\mathrm{}_y)^{\beta _2}\mathrm{}\mathrm{\Psi }(𝒙^{},𝒚,\lambda )\right)`$
$`\times `$ $`\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )d\lambda =0.`$
3. Introduce new variables $`𝒃=(b_{\mathrm{}},b_2\mathrm{},\mathrm{})`$ and sum up the bilinear identities over $`\beta _1,\beta _2,\mathrm{}0`$ with the weight $`b_{\mathrm{}}^{\beta _1}b_2\mathrm{}^{\beta _2}\mathrm{}/\beta _1!\beta _2!\mathrm{}`$. This sum contains power series of $`_{y_n\mathrm{}}\lambda ^n\mathrm{}_y`$, which is essentially an exponential:
$`{\displaystyle \underset{\beta 0}{}}{\displaystyle \frac{b^\beta }{\beta !}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(_{y_n\mathrm{}}\lambda ^n\mathrm{}_y)^{\beta _n}=\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n\mathrm{}(_{y_n\mathrm{}}\lambda ^n\mathrm{}_y)\right).`$
The action of this exponential operator can be easily understood by the obvious identity
$`\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n\mathrm{}(_{y_n\mathrm{}}\lambda ^n\mathrm{}_y)\right)f(y,\stackrel{ˇ}{𝒚})=f(y{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n\mathrm{}\lambda ^n\mathrm{},\stackrel{ˇ}{𝒚}+𝒃)`$
that holds true for any function $`f(y,\stackrel{ˇ}{𝒚})`$. We thus obtain the bilinear identity
$`{\displaystyle _{\lambda =\mathrm{}}}\mathrm{\Psi }(𝒙^{},y\eta (𝒃,\lambda ),\stackrel{ˇ}{𝒚}+𝒃,\lambda )\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )𝑑\lambda =0,`$
where for indeterminates $`𝒃=(b_{\mathrm{}},b_2\mathrm{},\mathrm{})`$ we set
$`\eta (𝒃,\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n\mathrm{}\lambda ^n\mathrm{}.`$ (4.1)
4. These calculations can be repeated, now starting with the linear equations for $`\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )`$. This yields a bilinear equation with the variables $`𝒚=(y,\stackrel{ˇ}{𝒚})`$ in $`\mathrm{\Psi }^{}`$ being shifted as
$`(y,\stackrel{ˇ}{𝒚})(y\eta (𝒄,\lambda ),\stackrel{ˇ}{𝒚}+𝒄)`$
where $`𝒄=(c_{\mathrm{}},c_2\mathrm{},\mathrm{})`$ are newly introduced variables.
We thus eventually find the following result:
###### Theorem 2
The formal Baker-Akhiezer functions $`\mathrm{\Psi }(𝐱,𝐲,\lambda )`$ and $`\mathrm{\Psi }^{}(𝐱,𝐲,\lambda )`$ satisfy the bilinear equation
$`{\displaystyle _{\lambda =\mathrm{}}}`$ $`\mathrm{\Psi }(𝒙^{},y\eta (𝒃,\lambda ),\stackrel{ˇ}{𝒚}+𝒃,\lambda )`$ (4.2)
$`\times `$ $`\mathrm{\Psi }^{}(𝒙,y\eta (𝒄,\lambda ),\stackrel{ˇ}{𝒚}+𝒄,\lambda )d\lambda =0.`$
Here $`𝐱`$, $`𝐱^{}`$, $`y`$, $`𝐛`$ and $`𝐜`$ are understood to be independent variables.
Remark.
1. Conversely, one can reproduce the evolution equations of $`W`$ etc. from these bilinear equations.
2. The preceding remark on the status of $`x_{\mathrm{}},x_2\mathrm{},\mathrm{}`$ also applies to this bilinear equation. Thus the bilinear equation holds true if any nonnegative power of $`\lambda ^{\mathrm{}}`$ is inserted. Accordingly, by taking a linear combination of those equations, one eventually obtains the slightly generalized (but actually equivalent) form
$`{\displaystyle _{\lambda =\mathrm{}}}`$ $`f(\lambda ^{\mathrm{}})\mathrm{\Psi }(𝒙^{},y\eta (𝒃,\lambda ),\stackrel{ˇ}{𝒚}+𝒃,\lambda )`$
$`\times `$ $`\mathrm{\Psi }^{}(𝒙,y\eta (𝒄,\lambda ),\stackrel{ˇ}{𝒚}+𝒄,\lambda )d\lambda =0.`$
of the bilinear equation that holds for any power series $`f(\lambda ^{\mathrm{}})=_{n0}f_n\lambda ^n\mathrm{}`$.
### 4.4 $`\tau `$ function and Hirota bilinear equations
We can define the $`\tau `$ function in the same way as the case of the KP hierarchy. Namely, the $`\tau `$ function is a function $`\tau =\tau (𝒙,𝒚)`$ that satisfies the equations
$`\mathrm{\Psi }(𝒙,𝒚,\lambda )`$ $`=`$ $`{\displaystyle \frac{\tau (𝒙ϵ(\lambda ),𝒚)}{\tau (𝒙,𝒚)}}e^{\xi (𝒙,\lambda )},`$ (4.3)
$`\mathrm{\Psi }^{}(𝒙,𝒚,\lambda )`$ $`=`$ $`{\displaystyle \frac{\tau (𝒙+ϵ(\lambda ),𝒚)}{\tau (𝒙,𝒚)}}e^{\xi (𝒙,\lambda )},`$
where
$`ϵ(\lambda )\stackrel{\mathrm{def}}{=}({\displaystyle \frac{1}{\lambda }},{\displaystyle \frac{1}{2\lambda ^2}},\mathrm{},{\displaystyle \frac{1}{n\lambda ^n}},\mathrm{}).`$
Note that this definition allows the indeterminacy
$`\tau (𝒙,𝒚)f(𝒚)\tau (𝒙,𝒚)`$
with $`f(𝒚)`$ being an arbitrary function of $`y`$ and $`\stackrel{ˇ}{𝒚}`$ only. Equation (4.2) now turns into the equation
$`{\displaystyle _{\lambda =\mathrm{}}}`$ $`\mathrm{exp}\left(\xi (𝒙^{}𝒙,\lambda )\right)`$
$`\times `$ $`{\displaystyle \frac{\tau (𝒙^{}ϵ(\lambda ),y\eta (𝒃,\lambda ),\stackrel{ˇ}{𝒚}+𝒃)}{\tau (𝒙^{},y\eta (𝒃,\lambda ),\stackrel{ˇ}{𝒚}+𝒃)}}{\displaystyle \frac{\tau (𝒙+ϵ(\lambda ),y\eta (𝒄,\lambda ),\stackrel{ˇ}{𝒚}+𝒄)}{\tau (𝒙,y\eta (𝒄,\lambda ),\stackrel{ˇ}{𝒚}+𝒄)}}d\lambda =0.`$
for the $`\tau `$ function. Notice here that the two factors in the denominator are power series of $`\lambda ^{\mathrm{}}`$. According to the second remark of the last theorem, one can insert any power series $`f(\lambda ^{\mathrm{}})`$ in the foregoing equation. If we choose
$`f(\lambda ^{\mathrm{}})=\tau (𝒙^{},y\eta (𝒃,\lambda ),\stackrel{ˇ}{𝒚}+𝒃)\tau (𝒙,y\eta (𝒄,\lambda ),\stackrel{ˇ}{𝒚}+𝒄),`$
the two factors in the denominator are cancel out. The outcome is the bilinear equation
$`{\displaystyle _{\lambda =\mathrm{}}}`$ $`\mathrm{exp}\left(\xi (𝒙^{}𝒙,\lambda )\right)`$
$`\times `$ $`\tau (𝒙^{}ϵ(\lambda ),y\eta (𝒃,\lambda ),\stackrel{ˇ}{𝒚}+𝒃)\tau (𝒙+ϵ(\lambda ),y\eta (𝒄,\lambda ),\stackrel{ˇ}{𝒚}+𝒄)d\lambda =0.`$
for the $`\tau `$ function.
To rewrite this bilinear equation into the Hirota form, we replace $`𝒙^{}𝒙+𝒂`$, $`𝒙𝒙𝒂`$, and choose $`𝒄`$ to be equal to $`𝒃`$. The bilinear equation can be thereby converted into a Hirota form:
$`{\displaystyle _{\lambda =\mathrm{}}}`$ $`\mathrm{exp}\left(2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_n\lambda ^n\right)\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{D_{x_n}}{n}}\lambda ^n\right)\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n\mathrm{}D_y\right)`$
$`\times `$ $`e^{𝒂,D_𝒙+𝒃,D_{\stackrel{ˇ}{𝒚}}}\tau (𝒙,𝒚)\tau (𝒙,𝒚)d\lambda =0.`$
The exponential including $`\lambda `$ can be expanded into a series of Schur functions:
$`\mathrm{exp}\left(2{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}a_n\lambda ^n\right)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}p_n(2𝒂)\lambda ^n,`$
$`\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{D_{x_n}}{n}}\lambda ^n\right)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}p_n(\stackrel{~}{D_𝒙})\lambda ^n,`$
$`\mathrm{exp}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_n\mathrm{}\lambda ^n\mathrm{}D_y\right)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}p_n(𝒃D_y)\lambda ^n\mathrm{}.`$
By plugging these identities into the last bilinear equations and do contour integrals, we obtain bilinear equations of the form
$`{\displaystyle \underset{m,k=0}{\overset{\mathrm{}}{}}}`$ $`p_m(2𝒂)p_{m+k\mathrm{}+1}(\stackrel{~}{D_𝒙})p_k(𝒃D_y)`$
$`\times `$ $`e^{𝒂,D_𝒙+𝒃,D_{\stackrel{ˇ}{𝒚}}}\tau (𝒙,𝒚)\tau (𝒙,𝒚)=0.`$
These bilinear equations almost coincide with those that have been derived from the representation theory of toroidal algebras. Although they do not agree, the discrepancy is rather superficial: If one shifts $`𝒂`$ to $`𝒂+\kappa 𝒃/2`$ and move the extra exponential factor $`\mathrm{exp}\left(\kappa _nb_n\mathrm{}\lambda ^n\mathrm{}\right)`$ to the third exponential including $`D_y`$, the outcome exactly reproduces the previous result(cf. Remark in the preceding subsection).
## 5 Special solutions of hierarchy
### 5.1 Construction of the solutions of the Wronskian type
We shall apply the method developed by Date in () to construct special solutions of $`\mathrm{}`$-Bogoyavlensky hierarchy. The method corresponds to consider Baker-Akhiezer functions on rational singular curves with $`N`$ nodes. Such algebro-geometric backgrounds are given by Krichever in (see also of Appendix 1 for an exposition) and Manin in Chap. III. One can also consult instructive lectures given by Mumford in (IIIb, section 5) that treat the subject.
As the data for the solution constructed below, let us consider the following:
* $`\alpha _i=\alpha _i(𝒚),c_i=c_i(𝒚)(i=1,\mathrm{},N)`$ are arbitrary (local) solution of the equations
$`{\displaystyle \frac{\alpha _i}{y_n\mathrm{}}}\alpha _i^n\mathrm{}{\displaystyle \frac{\alpha _i}{y}}=0,{\displaystyle \frac{c_i}{y_n\mathrm{}}}\alpha _i^n\mathrm{}{\displaystyle \frac{c_i}{y}}=0n=1,2\mathrm{},`$ (5.1)
such that we have $`\alpha _i\alpha _j`$ for $`ij,`$
* $`\zeta _i(i=1,\mathrm{},N)`$ are $`\mathrm{}`$-th roots of unity not equal to $`1.`$
* $`r_i(i=1,\mathrm{},N)`$ are arbitrary complex numbers.
We put $`\beta _i(𝒚)\stackrel{\mathrm{def}}{=}\zeta _i\alpha _i(𝒚).`$ Then we have the equations
$`{\displaystyle \frac{\beta _i}{y_n\mathrm{}}}\beta _i^n\mathrm{}{\displaystyle \frac{\beta _i}{y}}=0,n=1,2\mathrm{}.`$
Note that equation (5.1) have constant solutions, namely, for distinct complex numbers $`\alpha _i(i=1,\mathrm{},N)`$ and arbitrary complex numbers $`c_i(i=1,\mathrm{},N),`$ $`\alpha _i(𝒚)=\alpha _i,c_i(𝒚)=c_i`$ trivially satisfy equations (5.1).
Let us construct a special solution of $`\mathrm{}`$-Bogoyavlensky hierarchy, which we shall seek in the form
$`\mathrm{\Psi }(\lambda ,r)=\left(\lambda ^N+w_1\lambda ^{N1}+\mathrm{}+w_N\right)e^{\xi _r(\lambda )},`$ (5.2)
with $`w_n=w_n(𝒙,𝒚)`$ being unknown functions. Here we set $`\xi _r(\lambda )\stackrel{\mathrm{def}}{=}\xi (𝒙,\lambda )+ry+r\eta (\stackrel{ˇ}{𝒚},\lambda )`$ for any constant $`r`$ (cf. (2.11),(4.1)).
We define the $`N\times N`$ matrix
$`\mathrm{\Xi }_N=(𝒇^{(0)},𝒇^{(1)},\mathrm{},𝒇^{(N1)}),𝒇^{(i)}\stackrel{\mathrm{def}}{=}_x^i𝒇`$
where $`𝒇`$ denote the $`N`$-column vector $`{}_{}{}^{t}(f_1,\mathrm{},f_N)`$ and
$`f_i\stackrel{\mathrm{def}}{=}e^{\xi _{r_i}(\beta _i)}+c_ie^{\xi _{r_i}(\alpha _i)}.`$ (5.3)
Note that, for any constant solution $`\{\alpha _i,c_i\}`$ indicated above, we clearly have $`det\mathrm{\Xi }_N0`$.
###### Theorem 3
Let $`\{\alpha _i,\zeta _i,c_i,r_i(i=1,\mathrm{},N)\}`$ be a given set of data described above. Suppose we have $`det\mathrm{\Xi }_N0`$. Then the condition
$`\mathrm{\Psi }(\beta _i,r_i)+c_i\mathrm{\Psi }(\alpha _i,r_i)=0(i=1,\mathrm{},N)`$ (5.4)
uniquely determine the function $`\mathrm{\Psi }(\lambda ,r)`$ of the form in (5.2), and $`\mathrm{\Psi }(\lambda )=\mathrm{\Psi }(\lambda ,r)`$ solves equations (3.10). Explicitly, we have the following expressions
$`w_{Nk}={\displaystyle \frac{det(𝒇^{(0)},\mathrm{},𝒇^{(k1)},𝒇^{(N)},𝒇^{(k+1)},\mathrm{},𝒇^{(N1)})}{det(𝒇^{(0)},\mathrm{},𝒇^{(k1)},𝒇^{(k)},𝒇^{(k+1)},\mathrm{},𝒇^{(N1)})}}(k=1,\mathrm{},N).`$ (5.5)
In particular, the corresponding $`\tau `$ function is given by the formula
$`\tau =det\mathrm{\Xi }_N.`$ (5.6)
Proof. Using the obvious differential equations
$`_x^ne^{\xi _r(\alpha _i)}=\alpha _i^ne^{\xi _r(\alpha _i)},_x^ne^{\xi _r(\beta _i)}=\beta _i^ne^{\xi _r(\beta _i)},`$
the condition (5.4) can be written in the following form
$`W_N𝒇=\mathrm{𝟎}`$ (5.7)
where by $`W_N`$ we denote the differential operator
$`W_N\stackrel{\mathrm{def}}{=}_x^N+w_1_x^{N1}+\mathrm{}+w_N.`$ (5.8)
Since $`det\mathrm{\Xi }_N0,`$ we can solve equation (5.7) by Cramer’s formula, to obtain expression (5.5).
To show the function $`\mathrm{\Psi }(\lambda ,r)`$ satisfy (3.9), we shall show that the PsDO defined by $`W\stackrel{\mathrm{def}}{=}W_N_x^N`$ solves equations (3.4) and (3.5).
It suffices to show equation (3.4) for some differential operator $`B_n`$ in $`x`$, since differential operator $`B_n`$ is uniquely determined by the evolution equation itself. Differentiating (5.7) with respect to $`x_n`$ and using $`_{x_n}𝒇=_x^n𝒇`$, we have
$`\left({\displaystyle \frac{W_N}{x_n}}+W_N_x^n\right)𝒇=\mathrm{𝟎}.`$ (5.9)
There exist differential operators $`B_n`$ and $`R`$ in $`x`$ such that
$`{\displaystyle \frac{W_N}{x_n}}+W_N_x^n=B_nW_N+R,\underset{_x}{ord}(R)N1.`$ (5.10)
¿From (5.7) and (5.9) we have
$`R𝒇=\mathrm{𝟎}(i=1,\mathrm{},N).`$
This implies $`R=0`$ since $`det\mathrm{\Xi }_N0`$. Hence by multiplying $`_x^N`$ from the right to equation (5.10) with $`R=0,`$ we obtain (3.4).
Now we recall $`\alpha _i^{\mathrm{}}=\beta _i^{\mathrm{}}`$. Since the dependence of $`f_i`$ on $`x_n\mathrm{}`$ comes only through the overall factor $`\mathrm{exp}\left(\alpha _i^n\mathrm{}x_n\mathrm{}\right)=\mathrm{exp}\left(\beta _i^n\mathrm{}x_n\mathrm{}\right),`$ the expression (5.5) implies
$`{\displaystyle \frac{W}{x_n\mathrm{}}}=0(n=1,2\mathrm{}).`$
Hence from (3.4) $`W_x^n\mathrm{}W^1`$ is a differential operator, namely we have
$`W_x^n\mathrm{}W^1=W_N_x^n\mathrm{}W_N^1=P^n,`$ (5.11)
where $`P`$ denotes the differential operator $`W_x^{\mathrm{}}W^1.`$
Next we prove (3.5), the evolution equation of $`W`$ with respect to $`y_n\mathrm{}.`$ Using (5.1) and (5.3) one has the equations
$`_{y_n\mathrm{}}𝒇=\left(_x^n\mathrm{}_y\mathrm{\Phi }\right)𝒇\text{for}n=1,2,\mathrm{},`$
$`\text{where}\mathrm{\Phi }\stackrel{\mathrm{def}}{=}\mathrm{diag}(_y(\alpha _1^n\mathrm{}),\mathrm{},_y(\alpha _N^n\mathrm{})).`$
In view of these equations, the differentiation of (5.7) with respect to $`y_n\mathrm{}`$ yields
$`\left({\displaystyle \frac{W_N}{y_n\mathrm{}}}+W_N_x^n\mathrm{}_yW_N\mathrm{\Phi }\right)𝒇=\mathrm{𝟎}.`$ (5.12)
We use (5.11) to rewrite the operator $`W_N_y_x^n\mathrm{}`$ as follows
$`W_N_y_x^n\mathrm{}=W_N_x^n\mathrm{}W_N^1W_N_y=P^nW_N_y=P^n\left({\displaystyle \frac{W_N}{y}}+_yW_N\right)`$ (5.13)
By virtue of (5.13), $`W_N\mathrm{\Phi }=\mathrm{\Phi }W_N`$ and (5.7), we can rewrite (5.12) as
$`\left({\displaystyle \frac{W_N}{y_n\mathrm{}}}P^n{\displaystyle \frac{W_N}{y}}\right)𝒇=\mathrm{𝟎}.`$ (5.14)
Note that the operator in the left hand side of (5.14) is an ordinary differential operator, which does not include $`_y.`$ So we can apply exactly the same argument above to get the unique differential operator $`C_n\mathrm{}`$ in $`x`$ such that
$`{\displaystyle \frac{W_N}{y_n\mathrm{}}}P^n{\displaystyle \frac{W_N}{y}}=C_n\mathrm{}W_N.`$
This completes the proof of (3.5). To see the formula (5.6) is valid, we only have to notice the following identity arising from (4.3) and (3.1):
$`w_1=_x\mathrm{log}\tau .`$ (5.15)
Q.E.D.
### 5.2 $`N`$-soliton solutions
If $`\alpha _i,c_i(i=1,\mathrm{},N)`$ in the preceding subsection are constants, namely do not depend on $`𝒚,`$ then the solution of Wronskian type has clear representation-theoretical meaning. In fact, we shall see that, up to an irrelevant factor, $`\tau _N^W\stackrel{\mathrm{def}}{=}det\mathrm{\Xi }_N`$ coincide with the following $`N`$-soliton solution
$`\tau _N^S\stackrel{\mathrm{def}}{=}\left(1+a_NX(\alpha _N,\beta _N)V_{m_N}(\alpha _N^{\mathrm{}})\right)\mathrm{}\left(1+a_1X(\alpha _1,\beta _1)V_{m_1}(\alpha _1^{\mathrm{}})\right)\mathrm{𝐯𝐚𝐜}_{}^{\mathrm{tor}},`$ (5.16)
where $`a_i`$ are constants we shall specify later. To be rigorous, we have to assume $`|\alpha _N|>\mathrm{}>|\alpha _1|`$ so that the series converges. The fact that $`\tau _N^S`$ is a solution can be readily observed in view of the next lemma (cf. Lemma 3).
###### Lemma 7
Let $`\zeta ^{\mathrm{}}=1,\zeta 1`$ and $`m.`$ Then on the representation $`(_{}^{\mathrm{tor}},\pi _\mathrm{},^{\mathrm{tor}})`$ we have
$`{\displaystyle \underset{n}{}}E_n^\zeta t^m\lambda ^n{\displaystyle \frac{1}{1\zeta }}\left(X(\lambda ,\zeta \lambda )1\right)V_m(\lambda ^{\mathrm{}}).`$ (5.17)
Proof. Let us put
$`\mathrm{\Gamma }(\lambda )\stackrel{\mathrm{def}}{=}{\displaystyle \underset{n}{}}E_n^\zeta t^m\lambda ^n+{\displaystyle \frac{1}{1\zeta }}K_m^s(\lambda ^{\mathrm{}}).`$ (5.18)
Then we have
$`[\varphi _k^{},\mathrm{\Gamma }(\lambda )]`$ $`=`$ $`m\lambda ^k\mathrm{}\mathrm{\Gamma }(\lambda )k`$ (5.19)
and
$`[\mathrm{\Lambda }_k,\mathrm{\Gamma }(\lambda )]`$ $`=`$ $`(1\zeta ^k)\lambda ^k\mathrm{\Gamma }(\lambda )k.`$ (5.20)
(5.19) is a easy consequence of (2.8), and (5.20) follows from the following relation, which can be directly verified
$`[\mathrm{\Lambda }_k,E_n^\zeta t^m]=(1\zeta ^k)\left(E_{n+k}^\zeta t^m+{\displaystyle \frac{\delta _{n+k,0}^{[\mathrm{}]}}{1\zeta }}\overline{s^{(n+k)/\mathrm{}}t^md\mathrm{log}s}\right)k`$ (5.21)
where $`\delta _{m,n}^{[\mathrm{}]}=1`$ if $`mnmod\mathrm{}`$ and $`0`$ otherwise.
By virtue of a standard lemma on vertex operators (see e.g. Lemma 14.5), we have the conclusion that $`\mathrm{\Gamma }(\lambda )`$ is realized on $`_{}^{\mathrm{tor}}`$ by
$`c_{\zeta ,m}X(\lambda ,\zeta \lambda )V_m(\lambda ^{\mathrm{}}),c_{\zeta ,m}\text{is a constant.}`$ (5.22)
To determine the constant $`c_{\zeta ,m}`$, we notice that $`E_0^\zeta t^m\mathrm{𝐯𝐚𝐜}_{}^{\mathrm{tor}}=0.`$ To see this, we only need $`E_0^\zeta \mathrm{𝐯𝐚𝐜}_{}=0,`$ which is obvious from (2.7), (2.5), and the description of the representation in lemma 1. Now it is immediate to see $`c_{\zeta ,m}=(1\zeta )^1,`$ and the proof completes.Q.E.D.
If we define the following matrices
$`E_𝜶=\mathrm{diag}(e^{\xi (\alpha _1)},\mathrm{},e^{\xi (\alpha _N)}),C=\mathrm{diag}(c_1,\mathrm{},c_N),V_𝜶=(\alpha _i^{j1})_{i,j=1,\mathrm{},N}.`$
then $`\mathrm{\Xi }_N=E_𝜷V_𝜷+CE_𝜶V_𝜶`$ and hence we have
$`det\mathrm{\Xi }_N=detE_𝜷detV_𝜷det\left(1+D\mathrm{\Gamma }\right),D\stackrel{\mathrm{def}}{=}CE_𝜶E_𝜷^1,\mathrm{\Gamma }\stackrel{\mathrm{def}}{=}V_𝜶V_𝜷^1.`$ (5.23)
For any matrix $`X=(x_{ij})_{i,j=1,\mathrm{},N}`$ and any subset $`J`$ of $`\{1,2,\mathrm{},N\}`$, we define the principal submatrix $`X_J=(x_{jj^{}})_{j,j^{}J}.`$ We shall use the well-known formula
$`det(1+X)={\displaystyle \underset{J\{1,2,\mathrm{},N\}}{}}detX_J,J\text{runs over all the subsets of}\{1,2,\mathrm{},N\}.`$ (5.24)
###### Lemma 8
Let $`\mathrm{\Gamma }\stackrel{\mathrm{def}}{=}V_𝛂V_𝛃^1=(\gamma _{ij})_{i,jI}.`$ Then we have
$`\gamma _{ij}={\displaystyle \frac{_{k(j)}(\alpha _i\beta _k)}{_{k(j)}(\beta _j\beta _k)}}\text{and}{\displaystyle \frac{det(\mathrm{\Gamma }_J)}{_{jJ}\gamma _{jj}}}={\displaystyle \underset{j,j^{}J,j<j^{}}{}}{\displaystyle \frac{(\alpha _j\alpha _j^{})(\beta _j\beta _j^{})}{(\alpha _j\beta _j^{})(\beta _j\alpha _j^{})}}\text{for}JI.`$
Combining the formulas (5.23),(5.24), and the above lemma (cf. ), we have
$`\tau _N^W=e^{_{i=1}^N\xi _{m_i}(\beta _i)}\mathrm{\Delta }(𝜷){\displaystyle \underset{JI}{}}\left({\displaystyle \underset{jJ}{}}c_j\gamma _{jj}\right)\underset{j<j^{}}{\underset{j,j^{}J}{{\displaystyle }}}{\displaystyle \frac{(\alpha _j\alpha _j^{})(\beta _j\beta _j^{})}{(\alpha _j\beta _j^{})(\beta _j\alpha _j^{})}}e^{_{jJ}(\xi _{m_j}(\alpha _j)\xi _{m_j}(\beta _j))},`$
where $`\mathrm{\Delta }(𝜷)`$ denotes Vandermonde’s determinant $`detV_𝜷=_{i>j}(\beta _i\beta _j).`$ Here we also note that $`\tau _N^S`$ does not depend on $`y_n\mathrm{}^{}`$ from the construction. So if we set $`a_i=c_j\gamma _{jj},`$ we can see that $`\tau _N^W`$ is equal to $`\tau _N^S`$ times an irrelevant factor $`e^{_{i=1}^N\xi _{m_i}(\beta _i)}\mathrm{\Delta }(𝜷).`$
## 6 Conclusion
Inspired by the work of Billig and Iohara et al., we have introduced the $`\mathrm{}`$-Bogoyavlensky hierarchy and clarified its group-theoretic properties, Lax formalism, and special solutions. The group-theoretic characterization of the new hierarchy is based on the $`2`$-toroidal Lie algebra $`𝔰𝔩_{\mathrm{}}^{\mathrm{tor}}`$. We have constructed a representation of this Lie algebra, and derived a system of Hirota bilinear equations for the $`\tau `$ function. The point of departure of the Lax formalism is the observation that the lowest two members of the Hirota bilinear equations of Billig and Iohara et al. are equivalent to Bogoyavlensky’s $`2+1`$ dimensional extension of the KdV equation. Bearing this in mind, we have constructed a system of Lax equations, from which we have been able to reproduce the Hirota bilinear equations. The Lax formalism has also turned out to be useful for understanding special solutions.
Several interesting problems remain open. Firstly, the present construction should be extended to the toroidal Lie algebras associated with simple Lie algebras other than $`𝔰𝔩_{\mathrm{}}`$. This will give an extension of the Drinfeld-Sokolov hierarchies in the spirit of Bogoyavlensky. The same problem can also be raised to another important family of soliton equations, namely those represented by the nonlinear Schrödinger equation. Strachan’s work is remarkable in this respect. He presented a $`2+1`$-dimensional extension and an associated hierarchy of higher equations for soliton equations of this type. Lastly, we would like to mention the problem of constructing special solutions of the algebro-geometric (“finite-band”) type. A precursor of this problem is Cherednik’s work , in which he presented a construction of algebro-geometric solutions to the self-dual Yang-Mills equation. Since our hierarchy is closely related to the self-dual Yang-Mills equation, a similar construction of special solutions will be possible.
## 7 Appendix
We give a list of the Hirota equation of low degree contained in the hierarchy for $`\mathrm{}=2`$. Since we have $`P(D)ff=0`$ for any function $`f`$ and polynomial $`P`$ such that $`P(D)=P(D)`$, we shall list below only the even polynomials in $`D`$. We also drop the terms including $`D_{x_2},D_{x_4},\mathrm{}`$. We shall use the abbreviated notation $`D_n`$ and $`Q_n`$ for $`D_{x_n}`$ and $`D_{y_n}`$ respectively.
| degree 3 | $`D_1^3Q_0+2D_3Q_06D_1Q_2`$ |
| --- | --- |
| degree 4 | $`D_1^44D_1D_4`$ |
| degree 5 | $`\left(D_1^5+20D_1^2D_3+24D_5\right)Q_0120D_1Q_4`$ |
| | $`\left(D_1^5+20D_1^2D_3+24D_5\right)Q_0\left(20D_1^3+40D_3\right)Q_2`$ |
| | $`\left(D_1^510D_1^2D_3+24D_5\right)Q_0+\left(10D_1^340D_3\right)Q_2`$ |
| degree 6 | $`\left(D_1^620D_3^2+10D_1^3D_336D_1D_5\right)Q_0^230\left(D_1^44D_3D_1\right)Q_0Q_2`$ |
| | $`D_1^632D_3^2+4D_1^3D_3`$ |
| | $`D_1^680D_3^220D_1^3D_3+144D_1D_5`$ |
| degree 7 | $`\left(D_1^7+720D_7+70D_1^4D_3+280D_3^2D_1+504D_5D_1^2\right)Q_0`$ |
| | $`\left(420D_1^2D_3+21D_1^5+504D_5\right)Q_2\left(420D_1^3+840D_3\right)Q_4`$ |
| | $`\left(3D_1^7168D_5D_1^2+480D_7\right)Q_0\left(1120D_3280D_1^3\right)Q_4`$ |
| | $`\left(3D_1^7168D_5D_1^2+480D_7\right)Q_0+\left(280D_1^2D_328D_1^5672D_5\right)Q_2`$ |
| | $`\left(5D_1^71440D_770D_1^4D_3+560D_3^2D_1\right)Q_0\left(126D_1^52016D_5\right)Q_2`$ |
| | $`\left(D_1^7+720D_7+70D_1^4D_3+280D_3^2D_1+504D_5D_1^2\right)Q_0`$ |
| | $`\left(840D_1^2D_3+42D_1^5+1008D_5\right)Q_2`$ |
| | $`\left(D_1^7+720D_7+70D_1^4D_3+280D_3^2D_1+504D_5D_1^2\right)Q_05040D_1Q_6`$ |
| | $`\left(D_1^7+720D_7+70D_1^4D_3+280D_3^2D_1+504D_5D_1^2\right)Q_0^35040D_1Q_2^3`$ |
| | $`\left(2520D_1^2D_3+126D_1^5+3024D_5\right)Q_0^2Q_2+\left(2520D_1^3+5040D_3\right)Q_0Q_2^2`$ |
### Acknowledgements
We are grateful to Koji Hasegawa, Kenji Iohara, Tetsuya Kikuchi, Gen Kuroki, Takashi Takebe, Yoshihisa Saito, Ryuichi Sawae, and Takahiro Shiota for fruitful discussions. This work was partly supported by the Grant-in-Aid for Scientific Research (No. 10640165) from the Ministry of Education, Science and Culture.
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# 1 Introduction
## 1 Introduction
The purpose of this Letter is to justify the Abelian dominance in the low-energy region of QCD defined in the maximal Abelian (MA) gauge. This story begins with the idea of ’t Hooft called the Abelian projection. Immediately after this proposal, a hypothesis of Abelian dominance in low-energy physics of QCD was claimed by Ezawa and Iwazaki . By adopting the MA gauge invented by Kronfeld et al. , actually, Abelian dominance was discovered by Suzuki and Yotsuyanagi a decade ago based on numerical simulation on a lattice and has been confirmed by the subsequent simulations, see for reviews. However there is no analytical derivation or proof of Abelian dominance so far. How can one justify or prove the Abelian dominance in low-energy physics in QCD?
In a previous paper , we tried to give an answer by constructing an effective Abelian gauge theory which is considered to be valid in the low-energy region of QCD. We called it the Abelian-projected effective gauge theory (APEGT), although this name is somewhat misleading as will be explained below. Before this work, a number of low-energy effective gauge theories were already proposed based on the idea of Abelian-projection. However, we should keep in mind that these models were constructed by ignoring all the off-diagonal gluon fields from the beginning under the assumption of the Abelian dominance and/or the Abelian electro-magnetic duality, even if they can well describe some features of confinement physics in QCD. In fact, they could not be derived by starting with the QCD Lagrangian. Therefore, one can neither answer how the off-diagonal gluon fields influence the low-energy physics, nor how the Abelian electro-magnetic duality could appear from the non-Abelian gauge theory.
In contrast to these models, the APEGT is a first-principle derivation of effective theory from QCD. It was shown that the off-diagonal gluons do affect the low-energy physics in the sense that off-diagonal gluons renormalize the resulting effective Abelian gauge theory. Moreover, the coupling constant of the effective Abelian gauge theory has the renormalization-scale dependence governed by the renormalization group $`\beta `$-function which is exactly the same as the original QCD, thereby, exhibiting the asymptotic freedom. In this sense, the APEGT reproduces a characteristic feature of the original QCD, asymptotic freedom, even if it is an Abelian gauge theory. In addition, it was demonstrated how the dual Abelian gauge theory (magnetic theory) can in principle be obtained in the low-energy region of QCD. Actually, it is possible to show that monopole condensation leads to a dual Ginzburg-Landau theory supporting the dual superconductor picture of QCD vacuum. A version of the non-Abelian Stokes theorem indicates that the Wilson loop operator can be expressed in terms of diagonal gluon fields, see e.g. . Combining these results, we are able to explain the Abelian dominance in quark confinement, see .
In the derivation of APEGT, however, we have treated the off-diagonal gluons as if they are massive in the MA gauge. This assumption was necessary to justify the procedure of integrating out the off-diagonal gluon fields based on the functional integral, since this integration was interpreted as a step of the Wilsonian renormalization group (RG) of integrating out the massive (high-energy) modes. In view of this, the resultant APEGT is regarded as the low-energy effective theory which is meaningful at least in the length scale $`R>M_A^1`$ with $`M_A`$ being the mass of off-diagonal gluons. In the derivation, moreover, we have integrated out the off-diagonal ghosts and anti-ghosts. This step was also necessary to reproduce the correct coefficient of the $`\beta `$-function.
To really justify the Abelian dominance, therefore, we need to show that the off-diagonal gluons and ghosts become massive in the MA gauge within the same framework as the APEGT. The main purpose of this Letter is to demonstrate that this is indeed the case. Actually, the derivation of off-diagonal gluon mass and ghost mass can be performed within the setting up of the previous paper . In the previous work, we have ignored the ghost self-interactions in the derivation of the APEGT, simply because they were not necessary to obtain the asymptotic freedom. In this Letter we properly take the ghost self-interaction into account. We show that the quartic ghost self-interaction among off-diagonal ghosts leads to two kinds of ghost condensation. As a result, the off-diagonal gluons and off-diagonal ghosts (anti-ghosts) acquire non-zero masses.
It should be remarked that the quartic ghost interaction term is generated by integrating out the off-diagonal gluon fields, even if such an interaction term is absent in the original Lagrangian of QCD. This is due to the nonlinearity of the MA gauge. In general, quartic ghost self-interaction terms are generated in the nonlinear gauge due to radiative corrections. For the theory to be renormalizable, therefore, we need to incorporate quartic ghost self-interaction in the bare Lagrangian via the gauge-fixing and FP ghost term, as pointed out already in Appendix B of . If so, how one can specify the quartic ghost interaction? As a possibility, we introduce it so as to keep the supersymmetry which is quite different from that of supersymmetric theory in theoretical particle physics. It is hidden in the gauge fixing and ghost part of the MA gauge, while there is no supersymmetry in the Yang-Mills Lagrangian, since we are dealing with the usual QCD without supersymmetry. This requirement determines almost uniquely the ghost self-interaction. A special case was already examined in the previous paper in a slightly different context. In this point, this Letter supplements the previous paper by taking into account the ghost self-interactions properly.
## 2 QCD in the modified MA gauge
For the gauge group $`G=SU(N)`$, we consider the Cartan decomposition of the gauge potential into the diagonal and off-diagonal components,
$$𝒜_\mu (x)=𝒜_\mu ^A(x)T^A=a_\mu ^i(x)T^i+A_\mu ^a(x)T^a,$$
(1)
where $`A=1,\mathrm{},N^21`$. Then the maximal Abelian (MA) gauge is defined as follows. We define the functional of off-diagonal gluon fields, $`R[A]:=d^4x\frac{1}{2}A_\mu ^a(x)A^\mu {}_{}{}^{a}(x).`$ The MA gauge is obtained by minimizing the functional $`R[A^U]`$ with respect to the local gauge transformation $`U(x)`$ of $`A_\mu ^a(x)`$. Then we obtain the differential form of the MA gauge,
$$_\mu A^\mu {}_{}{}^{a}gf^{abi}a_\mu ^iA^\mu {}_{}{}^{b}:=D_\mu ^{ab}[a]A^\mu {}_{}{}^{b}=0.$$
(2)
This is nothing but the the background-field gauge with the background field $`a_\mu ^i`$. After the MA gauge is adopted, the original gauge group $`G=SU(N)`$ is broken to the maximal torus group $`H=U(1)^{N1}`$. The MA gauge is a partial gauge fixing which fixes the gauge degrees of freedom for the coset space $`G/H`$.
Following the well-known procedure, the manifest covariant action of QCD in the MA gauge is given by
$`S_{QCD}`$ $`:=`$ $`S_{YM}+S_{GF+FP}+S_F,`$
$`S_{YM}`$ $`=`$ $`{\displaystyle }d^4x{\displaystyle \frac{1}{4}}_{\mu \nu }^A^{\mu \nu }{}_{}{}^{A},`$
$`S_F`$ $`=`$ $`{\displaystyle d^4x\overline{\mathrm{\Psi }}i\gamma ^\mu (_\mu ig𝒜_\mu )\mathrm{\Psi }},`$
$`S_{GF+FP}`$ $`=`$ $`{\displaystyle d^4xi\delta _B[\overline{C}^a(D_\mu [a]A^\mu +\frac{\alpha }{2}B)^a]}.`$ (3)
Here $`S_{GF+FP}`$ is the gauge fixing (GF) and Faddeev-Popov (FP) ghost term where $`\delta _B`$ is the Becchi-Rouet-Stora-Tyupin (BRST) transformation, $`\alpha `$ is the gauge fixing parameter and $`B`$ is the Nakanishi-Lautrup (NL) Lagrange multiplier field. Of course, we can add the gauge fixing term for the residual symmetry $`H`$ which we don’t discuss in this Letter.
In this paper, we adopt the modified MA gauge proposed in ,
$$S_{GF+FP}^{}=d^4xi\delta _B\overline{\delta }_B[\frac{1}{2}A_\mu ^a(x)A^\mu {}_{}{}^{a}(x)\frac{\alpha }{2}iC^a(x)\overline{C}^a(x)],$$
(4)
where $`\delta _B`$ ($`\overline{\delta }_B`$) is the BRST (anti-BRST) transformation. The $`\alpha =2`$ case has been already investigated in . The modified MA gauge is different from the naive MA gauge by the ghost self-interaction, since
$$S_{GF+FP}^{}=d^4xi\delta _B\left[\overline{C}^a\left\{D_\mu [a]A^\mu +\frac{\alpha }{2}B\right\}^ai\frac{\zeta }{2}gf^{abi}\overline{C}^a\overline{C}^bC^ii\frac{\zeta }{4}gf^{abc}C^a\overline{C}^b\overline{C}^c\right],$$
(5)
where we must put $`\zeta =\alpha `$ to recover Eq.(4). Our choice of the Lagrangian (4) (or (5) with $`\zeta =\alpha `$) is invariant under the BRST and anti-BRST transformations. Moreover, it is invariant for arbitrary $`\alpha `$ under the FP ghost conjugation (discrete symmetry) ,
$$C^A\pm \overline{C}^A,\overline{C}^AC^A,B^A\overline{B}^A,\overline{B}^AB^A,𝒜_\mu ^A𝒜_\mu ^A.$$
(6)
Moreover, our choice (4) leads to a renormalizable theory and preserves the hidden supersymmetry as discussed in the final part of this Letter. The Lagrangian (5) should be compared with the Lagrangian in the Lorentz gauge which is invariant under the FP conjugation only in the Landau gauge,
$$_{GF+FP}=i\delta _B[\overline{C}^A(_\mu 𝒜^\mu +\frac{\alpha }{2}B)^A]=+i\delta _B\overline{\delta }_B(\frac{1}{2}𝒜_\mu ^2)+\frac{\alpha }{2}i\overline{\delta }_B(B^AC^A),$$
(7)
although it is also invariant under the BRST and anti-BRST transformations.
By performing the BRST transformation explicitly, we obtain
$`S_{GF+FP}^{}`$ $`=`$ $`{\displaystyle }d^4x\{B^aD_\mu [a]^{ab}A^\mu {}_{}{}^{b}+{\displaystyle \frac{\alpha }{2}}B^aB^a`$ (8)
$`+i\overline{C}^aD_\mu [a]^{ac}D^\mu [a]^{cb}C^big^2f^{adi}f^{cbi}\overline{C}^aC^bA^\mu {}_{}{}^{c}A_{\mu }^{d}`$
$`+i\overline{C}^aD_\mu [a]^{ac}(gf^{cdb}A^\mu {}_{}{}^{d}C_{}^{b})+i\overline{C}^agf^{abi}(D^\mu [a]^{bc}A_\mu ^c)C^i`$
$`+{\displaystyle \frac{\zeta }{8}}g^2f^{abe}f^{cde}\overline{C}^a\overline{C}^bC^cC^d+{\displaystyle \frac{\zeta }{4}}g^2f^{abc}f^{aid}\overline{C}^b\overline{C}^cC^iC^d+{\displaystyle \frac{\zeta }{2}}gf^{abc}iB^bC^a\overline{C}^c`$
$`\zeta gf^{abi}iB^a\overline{C}^bC^i+{\displaystyle \frac{\zeta }{4}}g^2f^{abi}f^{cdi}\overline{C}^a\overline{C}^bC^cC^d\},`$
where the choice (4) specifies the strength of the quartic ghost interactions,
$$\zeta =\alpha .$$
(9)
An implication of this choice will be discussed later. The gauge fixing term (8) is the most general type we consider in the following.
In particular, the $`G=SU(2)`$ case is greatly simplified as
$`S_{GF+FP}^{}`$ $`=`$ $`{\displaystyle }d^4x\{B^aD_\mu [a]^{ab}A^\mu {}_{}{}^{b}+{\displaystyle \frac{\alpha }{2}}B^aB^a`$ (10)
$`+i\overline{C}^aD_\mu [a]^{ac}D^\mu [a]^{cb}C^big^2ϵ^{ad}ϵ^{cb}\overline{C}^aC^bA^\mu {}_{}{}^{c}A_{\mu }^{d}`$
$`+i\overline{C}^agϵ^{ab}(D_\mu [a]^{bc}A_\mu ^c)C^3`$
$`\zeta gϵ^{ab}iB^a\overline{C}^bC^3+{\displaystyle \frac{\zeta }{4}}g^2ϵ^{ab}ϵ^{cd}\overline{C}^a\overline{C}^bC^cC^d\}.`$
Integrating out the NL field $`B^a`$ leads to
$`S_{GF+FP}^{}`$ $`=`$ $`{\displaystyle }d^4x\{{\displaystyle \frac{1}{2\alpha }}(D_\mu [a]^{ab}A^\mu {}_{}{}^{b})^2+(1\zeta /\alpha )i\overline{C}^agϵ^{ab}(D^\mu [a]^{bc}A_\mu ^c)C^3`$ (11)
$`+i\overline{C}^aD_\mu [a]^{ac}D^\mu [a]^{cb}C^big^2ϵ^{ad}ϵ^{cb}\overline{C}^aC^bA^\mu {}_{}{}^{c}A_{\mu }^{d}`$
$`+{\displaystyle \frac{\zeta }{4}}g^2ϵ^{ab}ϵ^{cd}\overline{C}^a\overline{C}^bC^cC^d\}.`$
## 3 Ghost condensation and mass generation due to quartic ghost interaction
The $`\zeta =0`$ case was considered in the previous paper , leaving $`\alpha `$ arbitrary. Even in this case, the quartic ghost self-interaction is generated after integrating out the off-diagonal gluons as mentioned above (see eq.(2.52) and Appendix B of ),
$`z_{4c}g^2ϵ^{ab}ϵ^{cd}\overline{C}^a\overline{C}^bC^cC^d,z_{4c}=4N{\displaystyle \frac{g^2}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mu }{\mu _0}}(N=2),`$ (12)
since the interaction term $`ig^2ϵ^{ad}ϵ^{cb}\overline{C}^aC^bA^\mu {}_{}{}^{c}A_{\mu }^{d}`$ does not vanish even for $`\zeta =0`$ (or $`\alpha =0`$).
If we consider the non-zero $`\zeta `$ case, (12) leads to the renormalization of $`\zeta `$ (together with $`g`$). This is expected from the beginning, since the quartic ghost interaction is a renormalizable interaction. In fact, it has been proven that QCD in the MA gauge is renormalizable by including the quartic ghost interaction.
For simplicity, we first discuss the $`G=SU(2)`$ case.<sup>1</sup><sup>1</sup>1In the SU(2) case, the ghost condensation was seriously discussed by Schaden from a different viewpoint from ours. To incorporate the effect of ghost interaction, we introduce the auxiliary scalar field $`\phi `$ as
$`{\displaystyle \frac{\zeta }{4}}g^2ϵ^{ab}ϵ^{cd}\overline{C}^a\overline{C}^bC^cC^d{\displaystyle \frac{1}{2\zeta g^2}}\phi ^2\phi iϵ^{ab}\overline{C}^aC^b,`$ (13)
where we have used the identity, $`ϵ^{ab}ϵ^{cd}\overline{C}^a\overline{C}^bC^cC^d=2(iϵ^{ab}\overline{C}^aC^b)^2=2(i\overline{C}^aC^a)^2`$ $`=4\overline{C}^1C^1\overline{C}^2C^2.`$ Then the GF+FP term is cast into the form,
$`S_{GF+FP}^{}`$ $`=`$ $`{\displaystyle d^4x\left\{i\overline{C}^a_\mu ^\mu C^a\phi iϵ^{ab}\overline{C}^aC^b\frac{1}{2\zeta g^2}\phi ^2\right\}}+\mathrm{}.`$ (14)
In order to see whether the QCD vacuum chooses a non-trivial $`\phi `$ or not, we consider the effective potential for the $`x`$-independent $`\phi `$ neglecting the kinetic terms. The Coleman-Weinberg type argument (summing up all one-loop ghost diagrams with arbitrary number of external $`\phi `$ fields) <sup>2</sup><sup>2</sup>2Note that the mathematical identity holds $`\mathrm{tr}_{n=1}^{\mathrm{}}\frac{1}{n}\left(\frac{i\phi }{^2}\right)^{2n}=\mathrm{ln}[1+\phi ^2/(^2)^2]=\mathrm{ln}det(_\mu ^\mu \delta ^{ab}\phi ϵ^{ab})\mathrm{ln}det(_\mu ^\mu \delta ^{ab}).`$ or integration over off-diagonal ghosts and anti-ghosts leads to the effective potential $`V(\phi )`$ for $`\phi `$,
$`V(\phi ){\displaystyle d^4x}`$ $`=`$ $`{\displaystyle d^4x\frac{1}{2\zeta g^2}\phi ^2}+i\mathrm{ln}det(_\mu ^\mu \delta ^{ab}\phi ϵ^{ab}).`$ (15)
Hence we obtain
$`V(\phi )`$ $`=`$ $`{\displaystyle \frac{1}{2\zeta g^2}}\phi ^2{\displaystyle \frac{d^4k}{i(2\pi )^4}\mathrm{ln}[(k^2)^2+\phi ^2]}.`$ (16)
The stationary point is given by the zero of the gap equation,
$`V^{}(\phi )`$ $``$ $`\phi \left[{\displaystyle \frac{1}{\zeta g^2}}2{\displaystyle \frac{d^4k}{i(2\pi )^4}\frac{1}{(k^2)^2+\phi ^2}}\right]=0.`$ (17)
Within the minimal subtraction (MS) scheme of the dimensional regularization, the effective potential is obtained as
$$V(\phi )=\frac{1}{2\zeta g^2}\phi ^2+\frac{1}{32\pi ^2}\phi ^2\left[2\mathrm{ln}\left(\frac{1}{4\pi \mu ^2}|\phi |\right)+C\right],$$
(18)
with $`C:=2\gamma 3`$ and the Euler constant $`\gamma =0.5772\mathrm{}`$. As far as $`\zeta 0`$, the gap equation (17) has non-trivial solutions given by $`\phi =\pm \phi _0`$ (besides a trivial one $`\phi =0`$) where
$$v:=\phi _0=4\pi \mu ^2e^{1\gamma }\mathrm{exp}\left[\frac{8\pi ^2}{\zeta g^2(\mu )}\right]>0.$$
(19)
These solutions correspond to global minima of the effective potential. At the global minimum $`\phi =\phi _0`$, $`V(\phi _0)=\frac{1}{32\pi ^2}(\phi _0)^2<0`$. This shows that QCD vacuum prefers a ghost condensate for any value of $`\zeta g^2(0)`$ such that $`\phi _0\zeta g^2iϵ^{ab}\overline{C}^aC^b0.`$
Now we consider the theory around the nontrivial vacuum $`\phi =\phi _0`$. The off-diagonal ghost propagator is modified in the ghost-condensed vacuum into
$$C^a(x)\overline{C}^b(y)=i\frac{d^4k}{i(2\pi )^4}\frac{k^2\delta ^{ab}vϵ^{ab}}{(k^2)^2+v^2}e^{ik(xy)}.$$
(20)
When $`ϵ^{ab}\overline{C}^aC^b=0`$, i.e., $`v=0`$, another ghost condensation $`\overline{C}^aC^a`$ is also zero in the dimensional regularization. However, non-zero condensation $`ϵ^{ab}\overline{C}^aC^b0`$ leads to another condensation $`\overline{C}^aC^a0`$. In the condensed vacuum, the ghost-gluon 4-body interaction, $`ig^2ϵ^{ad}ϵ^{cb}\overline{C}^aC^bA^\mu {}_{}{}^{c}A_{\mu }^{d},`$ leads to a mass term of the off-diagonal gluons,
$$ig^2ϵ^{ad}ϵ^{cb}\overline{C}^aC^bA^\mu {}_{}{}^{c}A_{\mu }^{d}=\frac{1}{2}g^2i\overline{C}^cC^cA^\mu {}_{}{}^{a}A_{\mu }^{a},$$
(21)
where we have used $`\overline{C}^aC^b=\frac{1}{2}(\delta ^{ab}\overline{C}^cC^c+ϵ^{ab}ϵ^{cd}\overline{C}^cC^d).`$ The condensation is given by
$`i\overline{C}^aC^a={\displaystyle \frac{d^4k}{i(2\pi )^4}\frac{2k^2}{(k^2)^2+v^2}}={\displaystyle \frac{v}{16\pi }}>0,`$ (22)
where the signature, i.e., positivity of $`v`$ is determined by analytic continuation to Euclidean region. Thus the off-diagonal gluon acquires the mass given by
$`M_A^2=g^2i\overline{C}^aC^a={\displaystyle \frac{g^2v}{16\pi }}>0.`$ (23)
The dynamically generated mass $`m_A`$ is finite excluding the mass counter term. Note that the introduction of the explicit mass term $`\frac{1}{2}m^2A_\mu ^aA^\mu ^a`$ spoils the renormalizability of the theory.<sup>3</sup><sup>3</sup>3Even if one introduces a bare mass term of the form, $`m^2\left(\frac{1}{2}A_\mu ^aA^\mu {}_{}{}^{a}+i\alpha \overline{C}^aC^a\right)`$, the modified BRST and anti-BRST transformations can be constructed as $`\delta _BB=m^2C,\overline{\delta }_BB=m^2\overline{C}g(\overline{C}\times B),`$ under which the modified Lagrangian is invariant. However, the nilpotency of both transformations is violated as $`\delta _B^2\overline{C}=im^2C,\overline{\delta }_B^2C=im^2\overline{C},`$ leading to the breakdown of physical $`S`$-matrix unitarity, see . Our derivation of off-diagonal gluon mass preserves the renormalizability, see for more details.
Now we proceed to estimate the order of the off-diagonal mass. We impose the renormalization condition at the renormalization point $`M`$, i.e, we define the renormalized coupling $`\zeta g^2`$ by
$$V^{\prime \prime }(\phi )|_{\phi =M^2}=(\zeta g^2)_{(M)}^1.$$
(24)
where $`M`$ is nonzero but arbitrary. The renormalizability implies that arbitrary choice of $`M`$ should not change the physics. Hence we have the $`\mu `$-independence of $`V(\phi )`$ which means that $`(\zeta g^2)_{(M)}^1(4\pi ^2)^1\mathrm{ln}M`$ is a $`M`$-independent constant. Then we have $`(\zeta g^2)_{(M)}=(\zeta g^2)_{(M_0)}[1+(\zeta g^2)_{(M_0)}(4\pi ^2)^1\mathrm{ln}\frac{M}{M_0}]^1`$. Hence $`\zeta g^2`$ satisfies the RG equation, $`M\frac{d}{dM}(\zeta g^2)=\frac{1}{4\pi ^2}(\zeta g^2)^2.`$ From the asymptotic freedom, $`M\frac{dg^2}{dM}=\frac{b_0}{8\pi ^2}g^4,`$ $`\zeta `$ obeys the RG equation, $`M\frac{d\zeta }{dM}=\frac{g^2}{4\pi ^2}\zeta (\zeta b_0/2).`$ Then we find that $`\zeta =0`$ and $`\zeta =b_0/2`$ are the fixed points. The $`M`$ independence of $`V`$ means that $`V`$ satisfies the differential equation,
$$\left[M\frac{}{M}+\beta (g)\frac{}{g}+\gamma _\zeta (g)\frac{}{\zeta }\gamma _\phi (g)\phi \frac{}{\phi }\right]V(\phi )=0,$$
(25)
where $`\gamma _\phi (g)`$ is the anomalous dimension of $`\phi `$ defined by $`\gamma _\phi (g):=\frac{1}{2}M\frac{\mathrm{ln}Z_\phi }{M}`$ and $`\phi _R:=Z_\phi ^{1/2}\phi `$ and $`\gamma _\zeta (g):=M\frac{\zeta }{M}`$. Substituting (18) into (25), we obtain the consistent results, $`\gamma _\phi (g)=\frac{b_0g^2}{16\pi ^2},`$ $`\beta (g)=g\gamma _\phi (g)=\frac{b_0g^3}{16\pi ^2},`$ and $`\gamma _\zeta (g):=\frac{g^2}{4\pi ^2}\zeta (\zeta b_0/2).`$ The non-trivial fixed point $`\zeta =b_0/2`$ yields $`v=4\pi e^{1\gamma }\mu ^2\mathrm{exp}\left(\frac{16\pi ^2}{b_0g^2(\mu )}\right)`$ $`=4\pi e^{1\gamma }\mathrm{\Lambda }_{QCD}^2.`$ Therefore, the condensation $`v`$ is a renormalization-group invariant and the order is given by the QCD scale, $`\mathrm{\Lambda }_{QCD}`$. Hence, the off-diagonal gluon mass is given by $`M_A=\frac{g}{2}e^{(1\gamma )/2}\mathrm{\Lambda }_{QCD}=(\pi \alpha _s)^{1/2}e^{(1\gamma )/2}\mathrm{\Lambda }_{QCD}.`$ This is comparable with the Lattice simulation result, $`M_A1.2\mathrm{GeV}`$, see .
Moreover, the quartic ghost interaction can give a mass for the ghost, since the treatment $`\stackrel{`}{a}`$ la Hartree-Fock approximation leads to $`\frac{\zeta }{4}g^2ϵ^{ab}ϵ^{cd}\overline{C}^a\overline{C}^bC^cC^d=\frac{\zeta }{2}g^2(iϵ^{ab}\overline{C}^aC^b)^2=\frac{\zeta }{2}g^2(i\overline{C}^aC^a)^2\zeta g^2i\overline{C}^aC^ai\overline{C}^bC^b.`$ This implies the off-diagonal ghost mass,
$`M_c^2\zeta g^2i\overline{C}^aC^a=\zeta g^2{\displaystyle \frac{v}{16\pi }}=\zeta M_A^2.`$ (26)
Thus off-diagonal gluons and ghosts can become massive due to ghost self-interactions. Note that $`\overline{C}^aC^a`$ and $`ϵ^{ab}\overline{C}^aC^b`$ are invariant under the residual U(1). Even in the presence of the condensation, the residual U(1) invariance is not broken spontaneously and the diagonal gluon remains massless . These results strongly support the Abelian dominance.
It is possible to extend the above analysis to the SU(3) case . The potential $`V(\phi ^3,\phi ^8)`$ is written in terms of two diagonal combinations, $`\phi ^i\zeta g^2\sqrt{1}f^{iab}\overline{C}^aC^b(i=3,8)`$. In fact, the effective potential for $`SU(3)`$ is given by
$`V(\stackrel{}{\phi })`$ $`=`$ $`{\displaystyle \frac{1}{2\zeta g^2}}\stackrel{}{\phi }\stackrel{}{\phi }{\displaystyle \underset{\alpha =1}{\overset{3}{}}}{\displaystyle \frac{d^4k}{i(2\pi )^4}\mathrm{ln}[(k^2)^2+(\stackrel{}{ϵ}_\alpha \stackrel{}{\phi })^2]},`$ (27)
where $`\stackrel{}{\phi }:=(\phi ^3,\phi ^8)`$ and $`ϵ_\alpha `$ is the root vectors $`ϵ_1=(1,0),ϵ_2=(\frac{1}{2},\frac{\sqrt{3}}{2}),ϵ_3=(\frac{1}{2},\frac{\sqrt{3}}{2}).`$ The schematic plot of the potential is given in Fig.1. It turns out that the potential has the global minima at six points on the three straight lines along the root vectors, i.e., (I) $`\phi ^8=0,\phi ^30`$, (II) $`\phi ^8=\sqrt{3}\phi ^3`$, (III) $`\phi ^8=\sqrt{3}\phi ^3`$. We find that the off-diagonal gluons in the $`SU(3)`$ case have two different masses as follows.<sup>4</sup><sup>4</sup>4This result is obtained up to Weyl symmetry.
$`\mathrm{I}:`$ $`{\displaystyle \frac{1}{\sqrt{2}}}m_{A^1}={\displaystyle \frac{1}{\sqrt{2}}}m_{A^2}=m_{A^4}=m_{A^5}=m_{A^6}=m_{A^7},`$
$`\mathrm{II}:`$ $`m_{A^1}=m_{A^2}={\displaystyle \frac{1}{\sqrt{2}}}m_{A^4}={\displaystyle \frac{1}{\sqrt{2}}}m_{A^5}=m_{A^6}=m_{A^7},`$
$`\mathrm{III}:`$ $`m_{A^1}=m_{A^2}=m_{A^4}=m_{A^5}={\displaystyle \frac{1}{\sqrt{2}}}m_{A^6}={\displaystyle \frac{1}{\sqrt{2}}}m_{A^7}.`$ (28)
The value of larger mass is given by $`m_A^2=\frac{g^2V_0}{16\pi },`$ where $`V_0=4^{1/6}(4\pi \mu ^2)e^{1\gamma }\mathrm{exp}\left(\frac{16\pi ^2}{3\zeta g^2}\right).`$ The $`\mu `$-independence of the potential holds when $`\zeta =b_0/3`$ for $`SU(3)`$. Hence, we obtain $`V_0=4^{1/6}(4\pi )e^{1\gamma }\mathrm{\Lambda }_{QCD}^2.`$ Another way to estimate the order of the off-diagonal mass is based on the identity of the trace anomaly,
$$0|T_\mu ^\mu |0=\frac{\beta (\alpha _s)}{4\alpha _s}0|(_{\mu \nu }^A)^2|0=\frac{11N_c2N_f}{24}0|\frac{\alpha _s}{\pi }(_{\mu \nu }^A)^2|0.$$
(29)
Note that the values of gluon condensate obtained on a lattice are as follows.<sup>5</sup><sup>5</sup>5 The authors would like to thank E.-M. Ilgenfritz for providing this information. In the presence of light quarks, the charmonium sum rules gives $`0|\frac{\alpha _s}{\pi }(_{\mu \nu }^A)^2|01.31.9\times 10^2(\mathrm{GeV})^4.`$ $`0|\frac{\alpha _s}{\pi }(_{\mu \nu }^A)^2|00.152(\mathrm{GeV})^4`$ for $`G=SU(2)`$ and $`0.144(\mathrm{GeV})^4`$ for $`G=SU(3)`$. On the other hand, the vacuum energy of the condensed vacuum (at the global minima) leads to
$$0|T_\mu ^\mu |04V(\phi _0)=\frac{3V_0^2}{16\pi ^2}\text{for SU(3)},\left(\frac{v^2}{8\pi ^2}\text{for SU(2)}\right).$$
(30)
Equating (29) and (30), we obtain $`V_0=(|0|T_\mu ^\mu |0|16\pi ^2/3)^{1/2}3.2(\mathrm{GeV})^2,`$ for $`N_c=3`$ and $`N_f=0`$. Finally, we have $`m_A=(\alpha _sV_0/4)^{1/2},(\alpha _sV_0/8)^{1/2}0.40.5\mathrm{GeV}.`$ These are our predictions. The full details of $`SU(3)`$ case will be given in .
## 4 APEGT of QCD in the modified MA gauge
In order to obtain the “effective” theory which is written in terms of the diagonal fields $`a_\mu ^i,B^i,C^i,\overline{C}^i`$ alone, we intend to integrate out all the off-diagonal fields $`A_\mu ^a,B^a,C^a,\overline{C}^a`$. We call the resultant effective field theory the Abelian-projected effective gauge theory (APEGT). That is to say, the APEGT is defined as
$$\mathrm{exp}(iS_{APEGT})=[dA_\mu ^a][dC^a][d\overline{C}^a][dB^a]\mathrm{exp}(iS_{QCD}).$$
(31)
Hence the vacuum-to-vacuum amplitude (or the partition function) of QCD reads
$`Z_{QCD}`$ $`:=`$ $`{\displaystyle [d𝒜_\mu ^A][dC^A][d\overline{C}^A][dB^A]\mathrm{exp}(iS_{QCD})}`$ (32)
$`=`$ $`{\displaystyle [da_\mu ^i][dC^i][d\overline{C}^i][dB^i]\mathrm{exp}(iS_{APEGT})}.`$
In the naive MA gauge, such an attempt was first performed by Quandt and Reinhardt for $`\alpha =0`$ and subsequently by one of the authors for $`\alpha 0`$, in particular, $`\alpha =1`$ at least for $`G=SU(2)`$. (We have found that the $`\alpha =0`$ case is very special from the viewpoint of renormalizability.) The generalization to $`SU(N)`$ is straightforward .
In the naive MA gauge , the off-diagonal gluons were expected to become massive, while the diagonal gluons were believed to behave in rather complicated way. Recently, the massiveness of off-diagonal gluons has been shown by Monte Carlo simulations on a lattice . An analytical explanation was given at least in the topological sector based on the dimensional reduction of the topological sector to the two-dimensional coset $`G/H`$ nonlinear sigma (NLS) model, see section IV. C of . In this paper we have given another evidence of mass generation of off-diagonal gluons and ghosts. In view of these facts, the integration of massive off-diagonal gluon fields can be interpreted as a step of integration of massive modes in the sense of the Wilsonian renormalization group. In this sense, the APEGT obtained in this way is regarded as the low-energy effective theory describing the physics in the length scale $`R>m_A^1`$ or in the low-energy region $`p<m_A`$.
In order to obtain the explicit form of the APEGT in the modified MA gauge, we repeat the steps performed in to obtain the APEGT. The GF+FP term in the condensed vacuum reads
$`S_{GF+FP}^{}`$ $`=`$ $`{\displaystyle }d^4x[{\displaystyle \frac{1}{2\alpha }}(D_\mu [a]^{ab}A^\mu {}_{}{}^{b})^2+{\displaystyle \frac{1}{2}}M_A^2A^\mu {}_{}{}^{a}A_{\mu }^{a}`$ (33)
$`+`$ $`i\overline{C}^aD_\mu [a]^{ac}D^\mu [a]^{cb}C^bg\phi _0iϵ^{ab}\overline{C}^aC^bg\stackrel{~}{\phi }iϵ^{ab}\overline{C}^aC^bV(\phi _0+\stackrel{~}{\phi })],`$
where we have put $`\phi =\phi _0+\stackrel{~}{\phi }`$. Note that $`V(\phi _0+\stackrel{~}{\phi })=V(\phi _0)+\frac{1}{2}\stackrel{~}{\phi }^2V^{\prime \prime }(\phi _0)+O(\stackrel{~}{\phi }^3)`$ with $`V^{}(\phi _0)=0`$ and $`V^{\prime \prime }(\phi _0)=\frac{1}{8\pi ^2}`$. We perform the integration over (high-energy) massive modes, i.e., off-diagonal gluons $`A_\mu ^a`$ and off-diagonal ghosts $`C^a`$ and anti-ghosts $`\overline{C}^a`$ for the total action $`S_{YM}+S_{GF+FP}^{}`$. In the process of deriving the APEGT, we have introduced the anti-symmetric auxiliary (Abelian) tensor field $`B_{\mu \nu }^i`$ to avoid the quartic self-interactions among the off-diagonal gluons appearing in $`S_{YM}`$ where $`B_{\mu \nu }^i`$ is invariant under the residual gauge transformation $`H=U(1)^{N1}`$. The way of introducing $`B_{\mu \nu }^i`$ is not unique, see and for more details. In the following we discuss one of the original versions . By repeating the procedures in , we can show that the resultant APEGT is written as (up to higher-derivative terms)<sup>6</sup><sup>6</sup>6The higher-derivative terms are suppressed in the low-energy region, since they are of the order $`O(p^2/M_A^2)`$.
$`_{aB}[a,B]`$ $`=`$ $`{\displaystyle \frac{1}{4g^2(\mu )}}f_{\mu \nu }^if^{\mu \nu }{}_{}{}^{i}{\displaystyle \frac{1+z_b}{4}}g^2B_{\mu \nu }^iB^{\mu \nu }{}_{}{}^{i}+{\displaystyle \frac{z_c}{2}}B_{\mu \nu }^i{}_{}{}^{}f_{}^{\mu \nu }{}_{}{}^{i},`$ (34)
where we have defined $`g(\mu ):=Z_a^{1/2}g`$ with $`Z_a:=1z_a+z_d=1+\frac{22}{3}N\frac{g^2}{(4\pi )^2}\mathrm{ln}\frac{\mu }{\mu _0}.`$ Here $`f_{\mu \nu }^i`$ is the Abelian field strength $`f_{\mu \nu }^i:=_\mu a_\nu ^i_\nu a_\mu ^i`$ and $`{}_{}{}^{}f_{\mu \nu }^{i}`$ is the Hodge dual of $`f_{\mu \nu }^i`$, i.e., $`{}_{}{}^{}f_{\mu \nu }^{i}:=\frac{i}{2}ϵ_{\mu \nu \rho \sigma }f^{\rho \sigma }^i`$. This result shows that the off-diagonal gluons can not be ignored and that they influence the APEGT in the form of renormalization of the Abelian sector. In fact, the renormalization factors $`z_a,z_b,z_c,z_d`$ are given by $`z_a=\frac{20}{3}N\frac{g^2}{(4\pi )^2}\mathrm{ln}\frac{\mu }{\mu _0},z_b=+2N\frac{g^2}{(4\pi )^2}\mathrm{ln}\frac{\mu }{\mu _0},z_c=+4N\frac{g^2}{(4\pi )^2}\mathrm{ln}\frac{\mu }{\mu _0},z_d=\frac{2}{3}N\frac{g^2}{(4\pi )^2}\mathrm{ln}\frac{\mu }{\mu _0}`$ where $`\mu `$ is a renormalization scale.
A remarkable fact is that the running of the gauge coupling constant $`g(\mu )`$ is governed by the $`\beta `$-function,
$$\beta (g):=\mu \frac{dg(\mu )}{d\mu }=b_0g^3(\mu ),b_0=\frac{11}{3}N>0,$$
(35)
which is the same as the original Yang-Mills theory. So the APEGT is an effective Abelian gauge theory exhibiting the asymptotic freedom. The coupling between $`B_{\mu \nu }^i`$ and $`{}_{}{}^{}f_{\mu \nu }^{}^i`$ is important to derive the dual Abelian gauge theory which leads to the dual superconductivity. This term is generated through the integration (or radiative corrections) and is absent in the original Lagrangian. In this sense, the APEGT just obtained is non-renormalizable. Nevertheless, the APEGT can be made renormalizable, see for more details. The effect of dynamical quarks can be included into this scheme by integrating out the quark fields. It results in further renormalization leading to the $`\beta `$-function with a different coefficient, $`b_0=\frac{11}{3}N\frac{4}{3}fr_F`$, where $`f`$ is the number of quark flavors and $`r_F`$ is the dimension of fermion representation.
## 5 New extended BRS algebra
It is easy to show that the QCD Lagrangian (3) in the modified MA gauge (10) or (11) has a new global symmetry if it is restricted to $`C^3=0`$ subspace or to the parameter $`\zeta =\alpha `$, that is to say, the Lagrangian is invariant under the two transformations,
$`\delta _+\overline{C}^a(x)=C^a(x),\delta _+(\mathrm{other}\mathrm{fields})=0,`$ (36)
$`\delta _{}C^a(x)=\overline{C}^a(x),\delta _{}(\mathrm{other}\mathrm{fields})=0.`$ (37)
The existence of this symmetry in the Lagrangian in the maximal Abelian gauge was recently noticed by Schaden . After eliminating $`B^a`$ (and putting $`\zeta =\alpha `$), (11) agrees with the Lagrangian examined by Schaden from a quite different viewpoint, the equivariant cohomology . These transformations $`\delta _\pm `$ for the field $`\mathrm{\Phi }`$ are defined by the generators $`Q_\pm `$ as $`\delta _\pm \mathrm{\Phi }=[iQ_\pm ,\mathrm{\Phi }],Q_\pm :=d^3xJ_\pm ^0,`$ where the generators are constructed through the Noether currents,
$`J_+^\mu `$ $`=`$ $`iC^a(D_\mu [a]C)^a=+i\delta _B(C^aA_\mu ^a),`$
$`J_{}^\mu `$ $`=`$ $`+i\overline{C}^a(D_\mu [a]\overline{C})^a=i\overline{\delta }_B(\overline{C}^aA_\mu ^a).`$ (38)
They should be compared with the ghost number,
$`\delta _cC^A(x)`$ $`:=`$ $`[iQ_c,C^A(x)]=C^A(x),`$
$`\delta _c\overline{C}^A(x)`$ $`:=`$ $`[iQ_c,\overline{C}^A(x)]=\overline{C}^A(x),`$ (39)
where $`Q_c`$ is the ghost charge defined by $`Q_c=d^3xJ_c^0,J_c^\mu =i\{(D^\mu [a]\overline{C})^AC^A+\overline{C}^A(D^\mu [a]C)^A\}.`$ Shaden found that there is a SL(2,R) symmetry among $`Q_+,Q_{}`$ and $`Q_c`$, i.e., $`[iQ_c,Q_+]=2Q_+,[iQ_c,Q_{}]=2Q_{},i[Q_+,Q_{}]=Q_c,`$ where the diagonal generator is the ghost number $`Q_c`$.
It is well known that the BRST transformation, anti-BRST transformation and the ghost number generator form the double BRS algebra among three generators, $`Q_B,\overline{Q}_B`$ and $`Q_c`$,
$`[Q_c,Q_c]=0,`$
$`\{Q_B,Q_B\}`$ $`=`$ $`0,\{\overline{Q}_B,\overline{Q}_B\}=0,`$
$`i[Q_c,Q_B]`$ $`=`$ $`Q_B,i[Q_c,\overline{Q}_B]=\overline{Q}_B,`$ (40)
$`\{Q_B,\overline{Q}_B\}=0.`$
By enlarging the double BRS algebra, we find a new extended double BRS algebra among five generators, $`Q_B,\overline{Q}_B,Q_+,Q_{}`$ and $`Q_c`$, supplemented by
$`[Q_B,Q_+]`$ $`=`$ $`0,i[\overline{Q}_B,Q_+]=Q_B,`$
$`i[Q_B,Q_{}]`$ $`=`$ $`\overline{Q}_B,[\overline{Q}_B,Q_{}]=0,`$ (41)
$`i[Q_c,Q_+]`$ $`=`$ $`2Q_+,i[Q_c,Q_{}]=2Q_{},`$ (42)
$`i[Q_+,Q_{}]=Q_c.`$
Note that the new extended BRS algebra closes only on the space of functionals which are invariant under the residual U(1) gauge transformation.
This should be compared with the extended BRS algebra (BRSNO algebra) found by Nakanishi and Ojima in the manifest covariant gauge of the Lorentz type where the additional symmetry is given by
$`\delta _{cc}B^A=ig(C\times C)^A,\delta _{cc}\overline{C}^A=2C^A,\delta _{cc}(\mathrm{other}\mathrm{fields})=0,`$
$`\delta _{\overline{c}\overline{c}}B^A=+ig(\overline{C}\times \overline{C})^A,\delta _{\overline{c}\overline{c}}C^A=+2\overline{C}^A,\delta _{\overline{c}\overline{c}}(\mathrm{other}\mathrm{fields})=0.`$ (43)
Although the BRSNO algebra holds for arbitrary gauge, their generators are conserved only in the Landau gauge $`\alpha =0`$. In the new extended algebra given above, the generators are conserved for an arbitrary gauge parameter $`\alpha `$, but only on the space which is invariant under the residual gauge group.
## 6 Spontaneous breaking of a global symmetry and hidden supersymmetry in MA gauge
The non-zero expectation value $`ϵ^{ab}C^a\overline{C}^b`$ is regarded as the spontaneous breaking of the SL(2,R) symmetry as pointed out by Schaden , since <sup>7</sup><sup>7</sup>7We can consider other types of ghost condensations with the non-zero ghost number, i.e., $`0|[iQ_+,ϵ^{ab}C^a\overline{C}^b]|0=0|ϵ^{ab}C^aC^b|0`$ and $`0|[iQ_{},ϵ^{ab}C^a\overline{C}^b]|0=0|ϵ^{ab}\overline{C}^a\overline{C}^b|0.`$ Three composite operators $`ϵ^{ab}C^a\overline{C}^b,ϵ^{ab}C^aC^b,ϵ^{ab}\overline{C}^a\overline{C}^b`$ are mutually related by the action of the generators $`Q_+`$ or $`Q_{}`$ which are spontaneously broken.
$$0|[iQ_+,ϵ^{ab}\overline{C}^a\overline{C}^b]|0=20|ϵ^{ab}C^a\overline{C}^b|0=0|[iQ_{},ϵ^{ab}C^aC^b]|0.$$
(44)
The non-compact SL(2,R) symmetry is spontaneously broken into the non-compact Abelian subgroup, since the ghost charge $`Q_c`$ is not broken. The massless Nambu-Goldstone (NG) particles associated with this spontaneous symmetry breaking can be confined by the quartet mechanism , i.e., decouple from physical observables, since the current $`J_+^\mu `$ ($`J_{}^\mu `$) is BRST (anti-BRST) exact. Therefore, we need not to worry about the emergence of massless particles.
In the previous paper we have argued that the non-zero mass for the off-diagonal gluons can be understood from the massive spectrum of the coset NLS model in two dimensions, since the GF and FP ghost part for the modified MA gauge in four dimensions is reduced to the coset NLS model in two dimensions by the dimensional reduction $`\stackrel{`}{a}`$ la Parisi and Sourlas . In this Letter we have argued that the quartic ghost interaction is an origin of off-diagonal gluon mass. Now we discuss how two pictures could be related to each other.
It is shown that the action (4) for gauge fixing and FP ghost in the modified MA gauge has the orthosymplectic symmetry $`OSp(4|2)`$ among $`A_\mu ^a,C^a,\overline{C}^a`$ when it is written in the superspace $`X^M:=(x_\mu ,\theta ,\overline{\theta })`$ following the superspace formulation by Bonora and Tonin . This superspace formulation can give a geometric meaning of BRST $`\delta _B`$ and anti-BRST $`\overline{\delta }_B`$ transformations as translations in the Grassmann variables $`\theta `$ and $`\overline{\theta }`$ respectively, $`\delta _B\frac{d}{d\theta }𝑑\theta ,\overline{\delta }_B\frac{d}{d\overline{\theta }}𝑑\overline{\theta },`$ where we have employed the equivalence between the differentiation and integration with respect to the Grassmann variable. Then the GF and FP part in the modified MA gauge is rewritten into the manifest $`OSp(4|2)`$ invariant form,
$$S_{GF+FP}^{}=id^4x𝑑\theta 𝑑\overline{\theta }\mathrm{tr}_{G/H}\left[\frac{1}{2}\eta _{NM}𝒜^N(x,\theta ,\overline{\theta })𝒜^M(x,\theta ,\overline{\theta })\right],$$
(45)
using the Lie-algebra valued superfield (one-form),
$$𝒜_M(X)dX^M=𝒜_\mu (x,\theta ,\overline{\theta })dx^\mu +C(x,\theta ,\overline{\theta })d\theta +\overline{C}(x,\theta ,\overline{\theta })d\overline{\theta },$$
(46)
and a supermetric $`\eta _{NM}=\delta _{\mu \nu }`$ for $`(M,N)=(\mu ,\nu )`$ and $`i\frac{\alpha }{2}`$ for $`(M,N)=(\theta ,\overline{\theta })`$. Thanks to the $`OSp(4|2)`$ invariance of the integrand, it is shown that (45) is reduced to
$$S_{GF+FP}^{}=\pi \alpha d^2z[\frac{1}{2}A_\mu ^a(z)A^\mu {}_{}{}^{a}(z)\frac{\alpha }{2}iC^a(z)\overline{C}^a(z)].$$
(47)
If we restrict the gauge potential to its topological nontrivial piece in the coset $`G/H`$, $`A_\mu ^a\mathrm{tr}\left[T^a\frac{i}{g}U(x)_\mu U(x)^{}\right]:=\frac{1}{g}\mathrm{\Omega }_\mu ^a,`$ the action (47) is nothing but the coset NLS model ,
$$S_{GF+FP}=\frac{\pi \alpha }{2g^2}d^2z\mathrm{\Omega }_\mu ^a(z)\mathrm{\Omega }^\mu {}_{}{}^{a}(z).$$
(48)
Thus, as far as $`\alpha 0`$, the dimensional reduction to the two-dimensional coset NLS model occurs and the massive spectrum in the coset NLS model implies the massive off-diagonal gluon, see section IV.C of .<sup>8</sup><sup>8</sup>8In the $`\alpha =0`$ case, the $`OSp(4|2)`$ invariance is lost and hence the above mechanism of dimensional reduction does not work. On the other hand, the quartic ghost interaction disappears in this case and the ghost condensation generating the off-diagonal gluon mass does not occur and there is no spontaneous breaking of $`SL(2,R)`$ symmetry. In view of these, the case $`\alpha =0`$ is rather special and should be discussed separately. It is also suggestive for the correspondence between two pictures that the symplectic group $`Sp(2)`$ for the Grassmann variables is isomorphic to the $`SL(2,R)`$ mentioned above. The action of the NLS model may have a wrong sign depending on the signature of the parameter $`\alpha `$. This might be related to the fact that the ghost condensate does not vanish even in $`g=0`$ (rather diverges) for $`\alpha <0`$. Note that the dimensional reduction does not imply the equivalence between two Hilbert spaces on which the respective quantum theory is constructed. Thus the mass generation could be related to the spontaneous breaking of $`OSp(4|2)`$ symmetry as claimed in from slightly different viewpoint. Obviously, we need further study on the symmetry breaking.
## 7 Conclusion and discussion
We have shown that the masses of off-diagonal gluons and off-diagonal ghosts are dynamically generated in QCD by adopting the MA gauge. This provides an evidence of the Abelian dominance which is expected to hold in low-energy region of QCD. The MA gauge is a nonlinear gauge and hence the quartic ghost interaction term is inevitably generated by radiative corrections . From the viewpoint of renormalizability of the theory, therefore, we need to add the bare quartic ghost interaction to the original Lagrangian. We have explicitly shown that the quartic ghost interaction leads to ghost–anti-ghost condensations which give the masses of the off-diagonal gluons and ghosts in QCD, although QCD doesn’t have any elementary scalar field.
In this Letter we determined the form of the ghost interaction from the requirement of preserving the hidden supersymmetry (the resulting gauge is called the modified MA gauge). Surprisingly, the resulting Lagrangian in the modified MA gauge exactly coincides with that recently proposed by Schaden (at least for $`SU(2)`$) from quite a different point of view. Therefore, the ghost and anti-ghost condensation can be understood as a spontaneous breaking of the global SL(2,R) symmetry recently claimed by Schaden for the SU(2) case. We have proposed an extended BRS algebra which includes the SL(2,R) algebra. However, it is not clear at present whether the SL(2,R) symmetry can be applied to the gauge group $`SU(N)`$ for $`N>2`$. Finally we argued that the mass generation is also related to the spontaneous breaking of a supersymmetry hidden in the modified MA gauge for arbitrary $`N`$.
In this Letter, although we have pointed out the importance of the quartic interaction term from renormalizability point of view, we have not indicated that the APEGT obtained in our scenario is really renormalizable. The totally renormalizable APEGT can be obtained improving the previous work , see .
Finally, it will be interesting to see how the dynamical mass generation just obtained affects the dual (magnetic) theory. This issue will be discussed from APEGT in a forthcoming paper .
## Acknowledgments
After submitting this paper for publication, the authors were informed by Martin Schaden that he discussed the ghost condensation and its relation to the trace anomaly and obtained the similar results (in the SU(2) case) to ours presented in this paper. This work is supported in part by the Grant-in-Aid for Scientific Research from the Ministry of Education, Science and Culture (10640249).
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# Weak anisotropy and disorder dependence of the in-plane magnetoresistance in high mobility (100) Si-inversion layers
## Abstract
We report studies of the magnetoresistance (MR) in a two-dimensional electron system in (100) Si-inversion layers, for perpendicular and parallel orientations of the current with respect to the magnetic field in the 2D-plane. The magnetoresistance is almost isotropic; this result does not support the suggestion of the orbital origin of the MR in Si-inversion layer. In the hopping regime, however, the MR contains a weak anisotropic component that is non-monotonic in magnetic field. We found that the field, at which the MR saturates, for different samples varies by a factor of two, being lower or higher than the field of complete spin polarization of free carriers. Therefore, the saturation of the MR can not be identified with the spin polarization of free carriers.
For low carrier densities, in the vicinity of what is known as the ‘metal-insulator transition (MIT) in two dimensions’ , a magnetic field $`B_{}`$ applied in the plane of the 2D carriers suppresses the anomalous metallic conduction . The resistivity $`\rho `$ measured at low temperatures ($`TE_F`$) raises roughly proportionally to $`B_{}^2`$ and further saturates (Fermi energy, $`E_F`$, through the paper is in units of K). The observed strong magnetoresistance (MR) and subsequent saturation of the resistivity in high fields is often interpreted as a manifestation of the spin alignment of free carriers by the parallel field . Recently, in Ref. , the strong MR was suggested to be an orbital effect caused by the transition from 2D to 3D behavior, when the magnetic length $`a_B`$ becomes smaller than the thickness of the quasi-2D layer. We note however, that there is a third possibility in which the parallel field also acts on the spins of the localized carriers (or of the donor/acceptor bound states), and causes a corresponding increase in Coulomb scattering. Such a possibility is particularly inherent in models which consider interface localized states.
In order to determine how universal is the magnetoresistance in parallel field, we conducted systematic measurements on a variety of samples with different mobilities. Despite the qualitative similarity of the magnetoresistance in all samples, the characteristic magnetic field at which MR saturates, $`B_{\mathrm{sat}}`$, is found to be non-universal. $`B_{\mathrm{sat}}`$ varies by a factor of 2 from the highest to the lowest mobility sample being, respectively, either lower or higher than the field of complete spin polarization. This result demonstrates that the saturation of the magnetoresistance can not be identified with a field of complete spin polarization of the mobile electrons in Si-MOS samples, the parameter which should depend only on the carrier density but not on a sample.
To probe the relevance of the spin and orbital effects, we have performed measurements of the MR for a magnetic field in the 2D plane, with the bias current $`j`$ directed perpendicular and parallel to the field. In the explored range of densities on both sides of the MIT, and in magnetic fields $`(0÷12)`$ T, we found no strong anisotropy in the MR. This does not support the idea of the orbital origin of the MR in Si-MOS structures, but might be consistent with spin-related mechanisms. We observed only a weak anisotropy in the MR, $`\mathrm{\Delta }\rho =(\rho _{}\rho _{})/2\overline{\rho }5\%`$, which in the hopping regime exhibits a non-monotonic dependence on the magnetic field. To the best of our knowledge, such an anisotropy in the (100) crystal plane has not been observed earlier in Si-MOS systems (a monotonic anisotropy of the MR was reported in Ref. for the anisotropic GaAs/AlGaAs (311) plane). The anisotropic non-monotonic component of the MR indicates an effect of the spin-orbit (SO) coupling on the electron transport in the vicinity of the MIT.
The ac-measurements (3 Hz) of the resistivity were performed at temperatures $`0.270.3`$ K on five (100) Si-MOS samples: Si-9Nj (peak mobility $`\mu ^{peak}=4.3`$m<sup>2</sup>/Vs at $`T=0.3`$ K), Si-153 (3.8 m<sup>2</sup>/Vs), Si-15a (3.2 m<sup>2</sup>/Vs), Si22a (2.7 m<sup>2</sup>/Vs), and Si43a (1.96 m<sup>2</sup>/Vs). The samples were lithographically defined as rectangular Hall bars of the size $`0.8\times 5`$ mm (first four samples) and $`0.256\times 2.5`$ mm (the last one); their long sides (current direction) were aligned along . All samples exhibited qualitatively similar behavior of $`R(B_{})`$. Typical field dependent traces of the resistivity are shown in Fig. 1. Similar to that reported earlier , resistivity grows with magnetic field and then saturates above a certain field, $`B_{\mathrm{sat}}`$. To obtain the data for different orientations of $`j`$ relative to $`B_{}`$, the sample was rotated in situ at low temperature. The carrier density $`n`$ was varied by the gate voltage and its value was determined from the Shubnikov-de Haas (SdH) oscillations, quantum Hall effect , and Hall voltage. For the studied samples, even through the MIT , the density did not deviate more than by 7% from the value extrapolated from the high density regime .
Figure 1 demonstrates the absence of a large anisotropy in the MR, and the presence of a minor anisotropy ($`10\%`$) in high fields and for intermediate densities with $`\rho (jB)`$ systematically exceeding $`\rho (jB)`$ . Both these results disagree with the predictions of Ref. where the MR was associated with a crossover from 2D to 3D behavior with increasing parallel field. In order to quantify the saturation field $`B_{\mathrm{sat}}`$, we used three different empiric definitions for $`B_{\mathrm{sat}}`$ as a field corresponding to the interception of the tangents (as demonstrated in Fig. 1) on the logarithmic (i) or linear (ii) scale of resistance, and (iii) corresponding to an increase of the MR to 97.5% of its maximum value. Three dashed lines in Fig. 1 connect the $`\rho (B_{\mathrm{sat}})`$-points determined according to the definitions (i) to (iii).
The $`B_{\mathrm{sat}}`$-data according to the definition (i) are plotted in Fig. 2 vs electron density. In accord with other results , we find that $`B_{\mathrm{sat}}`$ increases approximately linearly with the density. We therefore assume first that the saturation field corresponds to the complete polarization of the 2D electron system,
$$B_{\mathrm{pol}}=\frac{2E_F}{g^{}\mu _B}=\frac{n}{g^{}m^{}}\frac{h}{e},$$
(1)
where $`g^{}`$ is the effective Landé $`g`$factor, $`m^{}`$ the effective mass and $`\mu _B`$ is the Bohr magneton (valley degeneracy $`g_v=2`$ for (100)Si crystal plane). In Ref. , the $`g^{}m^{}`$ values were measured as a function of the carrier density; it was also found that $`g^{}m^{}`$ is sample independent. Therefore, $`B_{\mathrm{pol}}`$ must also be a universal function of the carrier density. Contrary to this expectation, we found $`B_{\mathrm{sat}}(n)`$ to be noticeably different for different samples (see Fig. 2). This result, evidently, does not support the proposed model where the MR is related to the screening radius (or to the density of states at $`E_F`$) which depends on the free carrier density only.
The $`B_{\mathrm{sat}}`$ values depend on the choice of the empiric definition, however, the difference between different samples persists within any definition of this parameter. We conclude therefore that the magnetic field value, at which the magnetoresistance saturates, does not reflect the alignment of the free carrier spins. In Fig. 2, for comparison, we plotted also the field $`B_{\mathrm{pol}}`$ corresponding to complete spin polarization which was calculated according to Eq. (1) and using the $`g^{}m^{}(n)`$ values determined in Ref. . Clearly, for different samples, the magnetoresistance saturates at fields, which could be either lower or higher than the field of complete spin-polarization, $`B_{\mathrm{pol}}`$. For some samples and over a limited density range, $`B_{\mathrm{sat}}`$ is rather close to $`B_{\mathrm{pol}}`$, as was noted in Ref. .
The $`B_{\mathrm{sat}}(n)`$ dependences shown in Fig. 2 have similar slopes ($`dB_{\mathrm{sat}}/dn5.7`$ T$`/10^{11}\text{cm}^2`$), but each of them extrapolates to zero at a sample-dependent finite offset density, $`n_\mathrm{d}`$. The latter is of the order of the ‘critical density’ $`n_c`$ for the MIT at zero field. More specifically, $`n_\mathrm{d}`$ varies from approximately $`0.4n_c`$ for the highest mobility sample, Si9Nj, to $`n_c`$ for the lowest mobility one, Si-43a . The offset $`n_\mathrm{d}`$ and the inverse peak mobility (a measure of the sample disorder) are in an apparent linear relation (values of $`\mu ^{peak}`$ for each sample are given above). The dotted curve in Fig. 2 represents a schematic dependence $`B_{\mathrm{sat}}(n)`$ with the same slope as the experimental data and with the offset extrapolated linearly to the infinite mobility.
It was suggested in Ref. that the non-zero value of $`n_\mathrm{d}`$ is caused by a spontaneous spin polarization of the electron system. Clearly, such an assumption would lead to a paradox, since more and more disordered samples would exhibit this phenomenon at higher and higher density (or lower inter-electron interaction strength). On the other hand, no disorder- (or sample-) dependence is found in the spin susceptibility ($`\chi g^{}m^{}`$) directly measured in Ref. down to $`n=1\times 10^{11}`$ cm<sup>-2</sup> (i.e., to the critical density of the MIT). It is therefore rather likely that the offset value, $`n_\mathrm{d}`$, and the whole effect of the magnetoresistance saturation in Si-MOS samples are related mainly with a physics of the disorder or of the bound states, rather than with intrinsic properties of free electrons in the 2D system.
The non-universality of $`B_{\mathrm{sat}}`$ reflects a breakdown of the models which assume the magnetoresistance to be related only with the Zeeman energy of free carriers. On the other hand, such a behavior is anticipated for the models which consider the magnetoresistance to be a result of floating up of the band of localized states e.g., such as considered in Refs. .
We now turn to a comparison of the magnetoresistance measured with different current direction; a more accurate analysis reveals their weak non-monotonic difference, $`\mathrm{\Delta }\rho =\rho (jB)\rho (jB)`$. Arrows in Fig. 3 mark two major features, a density-dependent peak ‘A’ and a density-independent broad maximum ‘B’. Both peaks are well pronounced only in the hopping regime (i.e., for $`\rho (B,n)(0.22)h/e^2`$, depending on the magnetic field ) and disappear for densities $`n1\times 10^{11}`$cm<sup>-2</sup>. Whereas the monotonic component of the MR anisotropy (Fig. 3) might in principle be related to the effects of finite thickness of the 2D layer , it is not the case for the peak structure of the MR. We also note that both peaks can not be caused by a perpendicular field component due to a minor misalignment of the sample plane .
At present, there is no explanation for the origin of the density-independent peak ‘B’. As for the peak ‘A’, its characteristic magnetic field $`B_{\mathrm{peak}}`$ increases with density as Fig. 4 shows. A peak of a similar shape in the MR anisotropy was theoretically predicted to occur, due to the interplay between the spin-orbit and Zeeman coupling . In particular, Chen et al. considered the hopping (insulating) conduction regime and found that the peak in the anisotropy should take place at a field $`B_{\mathrm{peak}}^{(i)}=\alpha k_F/g^{}\mu _B`$ (where $`\alpha `$ is the spin-orbit coupling constant, and $`k_F`$ is the Fermi momentum). Raimondi et al. considered the diffusive (metallic) transport regime and found that a peak in anisotropy occurs at a field $`B_{\mathrm{peak}}^{(m)}=2E_F/g^{}\mu _B`$. For comparison, we plotted in Fig. 4 both theoretical dependences $`B_{\mathrm{peak}}(n)`$. In the calculations we did not use adjustable parameters ; for the renormalized $`g^{}`$ and $`g^{}m^{}`$ values we used the experimentally determined data .
There is a similarity between the measured location of the peak ‘A’ and $`B_{\mathrm{peak}}^{(m)}`$ calculated in Ref. , though the calculations are done for the diffusive transport regime whereas in the experiment at such high fields, the transport is hopping . The inconsistency with the $`B_{\mathrm{peak}}^{(i)}`$ calculated in Ref. for the hopping transport cannot be eliminated by selecting any value of $`\alpha `$, the parameter which is supposed to be density independent.
In conclusion, we have shown that at low carrier densities in the vicinity of the MIT, the parallel-field magnetoresistance in (100) Si-inversion layers, is almost independent of the relative orientation of the bias current and magnetic field. This is inconsistent with the orbital origin of the strong MR and supports its spin origin. We have observed a weak ($`5\%`$) anisotropy of the MR in the hopping regime, which is non-monotonic as a function of the magnetic field. In particular, it exhibits a density-dependent sharp peak and a broad maximum. We compared the sharp peak in the anisotropy of the magnetoresistance with the one theoretically predicted and related with the interplay between the spin-orbit and Zeeman coupling. We found the peak shape and the magnetic field position to be consistent with the theoretical predictions. We found that the field at which magnetoresistance of Si-inversion layers saturates is a sample- (and, apparently, a disorder-) dependent parameter. For a variety of samples studied, the ‘saturation’ field may be substantially larger or smaller than the field of complete spin polarization. Therefore, the saturation of the MR can not be identified with the complete spin-alignment of free carriers.
Authors acknowledge discussions with B. L. Altshuler, D. L. Maslov, M. E. Gershenson, H. Kojima, and M. E. Raikh. The work was supported by FWF (project No 13439), INTAS, NATO, NSF, and the Russian programs RFBR, ‘Physics of solid state nanostructures’, ‘Statistical physics’, ‘Integration’ and ‘The State support of the leading scientific schools’.
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# Link Invariants Associated with Gauge Equivalent Solutions of the Yang–Baxter Equation: the One-Parameter Family of Minimal Typical Representations of 𝑈_𝑞[𝑔𝑙(2|1)]
## 1 Introduction
The close connection between solutions of the quantum Yang–Baxter equation (QYBE) and the evaluation of invariants for oriented knots and links through representations of the braid group is well known . The algebraic properties of quantum algebras and superalgebras provide systematic means to construct solutions of the QYBE which can in turn be used to explicitly compute these invariants. The seminal example of this method is the use of the six vertex solution of the QYBE in the evaluation of the Jones polynomial invariant .
The spectral parameter dependent solutions of the trigonometric Yang–Baxter equation (TYBE) arise as evaluation (or loop) representations of *affine* quantum algebras and superalgebras. Braid group representations are obtained from these solutions by taking the limit as the spectral parameter approaches infinity. We refer to these limiting cases as *quantum R matrices* which satisfy the QYBE. As demonstrated by Bracken et al. , gauge equivalent solutions of the TYBE are obtained by considering different gradations of the underlying affine (super)algebraic structure. The significant point is that the explicit braid group representation obtained in the infinite spectral parameter limit depends on the choice of the gradation. This feature has already been observed in . This means that given a single solution of the TYBE, there are many possibilities to in turn obtain a solution of the QYBE.
Our aim in this work is to investigate the construction of link invariants for the case of solutions of the QYBE obtained through non-standard choices of gradation (i.e., not the *homogeneous* gradation) for the affine (super)algebraic structure which underlies the TYBE solutions. We find that in general we are only able to define invariants of *regular isotopy*, which is to say that the results are invariant under the second and third Reidemeister moves. However, we will show that in our examples we can relate the regular isotopy invariants to well known invariants of *ambient isotopy*, viz invariant under all three Reidemeister moves. (Our terminology for regular and ambient isotopy is adopted from .)
Our specific calculations are performed using the solution of the TYBE associated with the one-parameter family of minimal typical representations of the quantum superalgebra $`U_q[gl(2|1)]`$. For the choice of the homogeneous gradation for the untwisted affine extension $`U_q[gl(2|1)^{(1)}]`$, one obtains the Links–Gould invariants recently investigated in detail in . However, for different choices of the gradation the process yields other invariants which can be related back to both the Jones and Alexander–Conway polynomial invariants.
We begin the paper with a description of gauge equivalent solutions of the TYBE without appealing to the gradation structure of quantum affine (super)algebras. For the sake of simplicity, we present these solutions in terms of a simple basis transformation satisfying some particular constraints.
## 2 Gauge equivalent solutions of the TYBE
For an arbitrary vector space $`V`$, let $`R(x)\mathrm{End}(VV)`$ satisfy the TYBE on the tensor product space $`VVV`$:
$`R_{12}(x)R_{13}(xy)R_{23}(y)=R_{23}(y)R_{13}(xy)R_{12}(x).`$ (1)
where $`R_{12}(x)R(x)I`$, etc, and $`x`$ and $`y`$ are arbitrary complex parameters. Consider an invertible matrix $`A(x)\mathrm{End}(V)`$, with the properties:
$`\begin{array}{ccc}\hfill A(x)A(y)& =& A(xy)\hfill \\ \hfill [R(x),A_1(y)A_2(y)]& =& 0,\hfill \end{array}\}`$ (4)
where $`A_1(x)A(x)I`$ and $`A_2(x)IA(x)`$. We immediately deduce that:
$`A(1)`$ $`=`$ $`I,`$
$`A(x)A(y)`$ $`=`$ $`A(y)A(x),`$
$`A^1(x)`$ $`=`$ $`A(\overline{x}),`$
where throughout the paper, we liberally write $`\overline{X}`$ for $`X^1`$, for various $`X`$.
Now set $`(x)=A_1(x)R(x)A_1(\overline{x}).`$ It is an algebra exercise to prove:
$`\begin{array}{ccc}& & _{12}(x)_{13}(xy)_{23}(y)\hfill \\ & & =A_1(x)R_{12}(x)A_1(\overline{x})A_1(xy)R_{13}(xy)A_1(\overline{x}\overline{y})A_2(y)R_{23}(y)A_2(\overline{y})\hfill \\ & & =A_1(x)R_{12}(x)A_1(y)A_2(y)R_{13}(xy)R_{23}(y)A_1(\overline{x}\overline{y})A_2(\overline{y})\hfill \\ & & =A_1(xy)A_2(y)R_{12}(x)R_{13}(xy)R_{23}(y)A_1(\overline{x}\overline{y})A_2(\overline{y})\hfill \\ & & =A_1(xy)A_2(y)R_{23}(y)R_{13}(xy)R_{12}(x)A_1(\overline{x}\overline{y})A_2(\overline{y})\hfill \\ & & =A_2(y)R_{23}(y)A_1(xy)R_{13}(xy)A_1(\overline{y})A_2(\overline{y})R_{12}(x)A_1(\overline{x})\hfill \\ & & =A_2(y)R_{23}(y)A_2(\overline{y})A_1(xy)R_{13}(xy)A_1(\overline{y})R_{12}(x)A_1(\overline{x})\hfill \\ & & =A_2(y)R_{23}(y)A_2(\overline{y})A_1(xy)R_{13}(xy)A_1(\overline{x}\overline{y})A_1(x)R_{12}(x)A_1(\overline{x})\hfill \\ & & =_{23}(y)_{13}(xy)_{12}(x).\hfill \end{array}`$
The above calculation shows that $`(x)`$ also satisfies (1). The connection between $`R(x)`$ and $`(x)`$ in terms of the gradation of the underlying affine algebraic structure for these solutions is discussed in . For our purposes, we need not describe this in detail, and simply refer to $`R(x)`$ and $`(x)`$ as *gauge equivalent*. The essential point is that the limit of $`(x)`$ as $`x\mathrm{}`$ yields different quantum R matrices depending on the choice of $`A(x)`$.
## 3 Trigonometric R matrix with gauge parameters
Applying the above to a particular example, the following trigonometric R matrix $`\stackrel{ˇ}{}^{r,s}(u)`$ with *gauge parameters* $`r`$ and $`s`$ arises from the one-parameter family of minimal typical representations of $`U_q[gl(2|1)]`$. (Here, we have replaced variable $`x`$ with $`u`$, defined by $`xq^u`$.) The operator $`\stackrel{ˇ}{}^{r,s}(u)`$ has $`36`$ nonzero components, and is scaled such that its first component is $`1`$. It satisfies the TYBE in the additive form:
$`\stackrel{ˇ}{}_{12}(u)\stackrel{ˇ}{}_{23}(u+v)\stackrel{ˇ}{}_{12}(v)=\stackrel{ˇ}{}_{23}(v)\stackrel{ˇ}{}_{12}(u+v)\stackrel{ˇ}{}_{23}(u).`$ (6)
(Confirming this involves manipulating expressions in a total of $`6`$ variables, viz: the representation variables $`q`$ and $`\alpha `$, the spectral variables $`u`$ and $`v`$, and the gauge parameters $`r`$ and $`s`$.) Explicitly, $`\stackrel{ˇ}{}^{r,s}(u)`$ is:
$`\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},{\displaystyle \frac{[\alpha +u]}{[\alpha u]}}\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},{\displaystyle \frac{[\alpha +u][1+\alpha +u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`{\displaystyle \frac{[\alpha ]}{[\alpha u]}}\left\{\begin{array}{c}r^u\overline{q}^ue_{12}^{12},s^u\overline{q}^ue_{13}^{13}\\ \overline{r}^uq^ue_{21}^{21},\overline{s}^uq^ue_{31}^{31}\end{array}\right\},{\displaystyle \frac{[\alpha ][1+\alpha ]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}r^us^u\overline{q}^{2u}e_{14}^{14}\\ \overline{r}^u\overline{s}^uq^{2u}e_{41}^{41}\end{array}\right\},`$
$`{\displaystyle \frac{1}{\mathrm{\Delta }^2[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}\overline{r}^us^uf(\overline{q})e_{23}^{23}\\ r^u\overline{s}^uf(q)e_{32}^{32}\end{array}\right\},{\displaystyle \frac{[1+\alpha ][\alpha +u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}s^u\overline{q}^ue_{24}^{24},r^u\overline{q}^ue_{34}^{34}\\ \overline{s}^uq^ue_{42}^{42},\overline{r}^uq^ue_{43}^{43}\end{array}\right\},`$
$`{\displaystyle \frac{[u]}{[\alpha u]}}\left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},{\displaystyle \frac{[1u][u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},`$
$`{\displaystyle \frac{[u]^2}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\},{\displaystyle \frac{[u][\alpha +u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},`$
$`{\displaystyle \frac{[\alpha ]^{\frac{1}{2}}[1+\alpha ]^{\frac{1}{2}}[u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}+{\displaystyle \frac{r^u}{q^u}}q^{\frac{1}{2}}\left\{\begin{array}{c}e_{32}^{14}\\ e_{14}^{32}\end{array}\right\},{\displaystyle \frac{q^u}{r^u}}\overline{q}^{\frac{1}{2}}\left\{\begin{array}{c}e_{23}^{41}\\ e_{41}^{23}\end{array}\right\},+{\displaystyle \frac{q^u}{s^u}}q^{\frac{1}{2}}\left\{\begin{array}{c}e_{41}^{32}\\ e_{32}^{41}\end{array}\right\},{\displaystyle \frac{s^u}{q^u}}\overline{q}^{\frac{1}{2}}\left\{\begin{array}{c}e_{14}^{23}\\ e_{23}^{14}\end{array}\right\}\end{array}\right\},`$
where $`f(q)2q+q^{2u}(q\overline{q})+q^{1+2\alpha }+\overline{q}^{1+2\alpha }`$ and we have defined $`\mathrm{\Delta }q\overline{q}`$, as well as writing $`[x]`$ for $`(q^x\overline{q}^x)/(q\overline{q})`$.
This solution of the TYBE originates in the following trigonometric R matrix $`\stackrel{ˇ}{R}(u)\stackrel{ˇ}{}^{r=1,s=1}(u)`$, which may be regarded as a *gauge-free* version of $`\stackrel{ˇ}{}^{r,s}(u)`$.
$`\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},{\displaystyle \frac{[\alpha +u]}{[\alpha u]}}\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},{\displaystyle \frac{[\alpha +u][1+\alpha +u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`{\displaystyle \frac{[\alpha ]}{[\alpha u]}}\left\{\begin{array}{c}\overline{q}^ue_{12}^{12},\overline{q}^ue_{13}^{13}\\ q^ue_{21}^{21},q^ue_{31}^{31}\end{array}\right\},{\displaystyle \frac{[\alpha ][1+\alpha ]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}\overline{q}^{2u}e_{14}^{14}\\ q^{2u}e_{41}^{41}\end{array}\right\},`$
$`{\displaystyle \frac{1}{\mathrm{\Delta }^2[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}f(\overline{q})e_{23}^{23}\\ f(q)e_{32}^{32}\end{array}\right\},{\displaystyle \frac{[1+\alpha ][\alpha +u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}\overline{q}^ue_{24}^{24},\overline{q}^ue_{34}^{34}\\ q^ue_{42}^{42},q^ue_{43}^{43}\end{array}\right\},`$
$`{\displaystyle \frac{[u]}{[\alpha u]}}\left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},{\displaystyle \frac{[1u][u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},`$
$`{\displaystyle \frac{[u]^2}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\},{\displaystyle \frac{[u][\alpha +u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},`$
$`{\displaystyle \frac{[\alpha ]^{\frac{1}{2}}[1+\alpha ]^{\frac{1}{2}}[u]}{[\alpha u][1+\alpha u]}}\left\{\begin{array}{c}+\overline{q}^{u\frac{1}{2}}\left\{\begin{array}{c}e_{32}^{14}\\ e_{14}^{32}\end{array}\right\},q^{u\frac{1}{2}}\left\{\begin{array}{c}e_{23}^{41}\\ e_{41}^{23}\end{array}\right\},+q^{u+\frac{1}{2}}\left\{\begin{array}{c}e_{41}^{32}\\ e_{32}^{41}\end{array}\right\},\overline{q}^{u+\frac{1}{2}}\left\{\begin{array}{c}e_{14}^{23}\\ e_{23}^{14}\end{array}\right\}\end{array}\right\}.`$
The above solution of (6) is obtained from the representation theory of the quantum superalgebra $`U_q[gl(2|1)]`$ and the *Baxterization* procedure (e.g. see ). It is related to $`R(u)`$ through a transposition:
$`\stackrel{ˇ}{R}(u)_{bd}^{ac}=R(u)_{bd}^{ca}.`$
From it we may construct a quantum R matrix $`\stackrel{ˇ}{R}=lim_u\mathrm{}\stackrel{ˇ}{R}(u)`$, and build a link invariant .
By the following procedure, $`\stackrel{ˇ}{}^{r,s}(u)`$ may be obtained from $`\stackrel{ˇ}{R}(u)`$. Let us define:
$`(u)_{bd}^{ac}A(u)_e^aR(u)_{fd}^{ec}A(u)_b^f,`$ (21)
where $`A(u)`$ is a $`4\times 4`$ matrix satisfying the properties (4) which ensure that $`(u)`$ also satisfies (1). Specifically, we will use the following diagonal $`A(u)A^{r,s}(u)`$ containing gauge parameters $`r`$ and $`s`$:
$`A^{r,s}(u)\mathrm{diag}\{1,r^u,s^u,r^us^u\},`$ (22)
thus we will write $`^{r,s}(u)`$ rather than $`(u)`$. We define $`\stackrel{ˇ}{}^{r,s}(u)`$ by again transposing:
$`\stackrel{ˇ}{}^{r,s}(u)_{bd}^{ac}^{r,s}(u)_{bd}^{ca}.`$ (23)
This $`\stackrel{ˇ}{}^{r,s}(u)`$ satisfies (6), and indeed is the object which we first introduced.
To investigate the possibility of constructing link invariants, we will select appropriate gauge choices $`rr_i(q)`$, $`ss_i(q)`$ and take the spectral limit $`u\mathrm{}`$ of $`\stackrel{ˇ}{}^{r_i,s_i}(u)`$ to yield a quantum R matrix $`\stackrel{ˇ}{R}^i`$:
$`\stackrel{ˇ}{R}^i\underset{u\mathrm{}}{lim}\stackrel{ˇ}{}^{r_i,s_i}(u),`$ (24)
which satisfies the (non-parametric) QYBE:
$`\stackrel{ˇ}{R}_{12}\stackrel{ˇ}{R}_{23}\stackrel{ˇ}{R}_{12}=\stackrel{ˇ}{R}_{23}\stackrel{ˇ}{R}_{12}\stackrel{ˇ}{R}_{23},`$ (25)
which is graphically depictable in terms of braids. (In the above, $`i`$ is an index to number the gauge choice, and the functions $`r_i`$ and $`s_i`$ may also contain variable $`q`$.)
Substituting (21) and (22) into (23) allows us to write (no sums on $`b`$ and $`c`$):
$`\stackrel{ˇ}{R}^{r,s}(u)_{bd}^{ac}=A^{r,s}(u)_c^c\stackrel{ˇ}{R}(u)_{bd}^{ac}A^{r,s}(u)_b^b,`$
hence we have for grading choice $`i`$ a quantum R matrix $`R^i`$ with components:
$`(\stackrel{ˇ}{R}^i)_{bd}^{ac}=\underset{u\mathrm{}}{lim}A^{r_i,s_i}(u)_c^c\stackrel{ˇ}{R}(u)_{bd}^{ac}A^{r_i,s_i}(u)_b^b.`$
The gauge choices that we shall make are depicted in Table 1.
From Case 1, we recover the ungauged situation.
## 4 Quantum R matrices
In the spectral limit $`u\mathrm{}`$, our trigonometric R matrix $`\stackrel{ˇ}{}^{r_i,s_i}(u)`$ becomes a quantum R matrix $`\stackrel{ˇ}{R}^i`$ in variables $`q`$ and $`\alpha `$, where $`i`$ is the index of the gauge choice (24). Here, we present the $`\stackrel{ˇ}{R}^i`$, in terms of ‘internal’ variables, $`pq^{\alpha +\frac{1}{2}}`$ and $`Qq^{\frac{1}{2}}`$, which simplify computations.
$`\stackrel{ˇ}{R}^1`$ has $`26`$ nonzero components:
$`\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},p^4\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`p\overline{Q}(p\overline{Q}\overline{p}Q)\left\{\begin{array}{c}e_{21}^{21}\\ e_{31}^{31}\end{array}\right\},p^3\overline{Q}(pQ\overline{p}\overline{Q})\left\{\begin{array}{c}e_{42}^{42}\\ e_{43}^{43}\end{array}\right\},`$
$`p^2(p\overline{Q}\overline{p}Q)(pQ\overline{p}\overline{Q})\left\{\begin{array}{c}e_{41}^{41}\end{array}\right\},p^2(Q^2\overline{Q}^2)\left\{\begin{array}{c}e_{32}^{32}\end{array}\right\},`$
$`p\overline{Q}\left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},p^3\overline{Q}\left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},p^2\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\},`$
$`p^2(p\overline{Q}\overline{p}Q)^{\frac{1}{2}}(pQ\overline{p}\overline{Q})^{\frac{1}{2}}\left\{\begin{array}{c}\overline{Q}\left\{\begin{array}{c}e_{23}^{41}\\ e_{41}^{23}\end{array}\right\},Q\left\{\begin{array}{c}e_{41}^{32}\\ e_{32}^{41}\end{array}\right\}\end{array}\right\}.`$
$`\stackrel{ˇ}{R}^2`$ has $`20`$ nonzero components:
$`\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},p^4\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`p\overline{Q}(p\overline{Q}\overline{p}Q)\left\{\begin{array}{c}e_{21}^{21}\end{array}\right\},p^3\overline{Q}(pQ\overline{p}\overline{Q})\left\{\begin{array}{c}e_{43}^{43}\end{array}\right\},`$
$`p\overline{Q}\left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},p^3\overline{Q}\left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},p^2\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\},`$
$`p^2\overline{Q}(p\overline{Q}\overline{p}Q)^{\frac{1}{2}}(pQ\overline{p}\overline{Q})^{\frac{1}{2}}\left\{\begin{array}{c}e_{23}^{41}\\ e_{41}^{23}\end{array}\right\}.`$
$`\stackrel{ˇ}{R}^3`$ has $`17`$ nonzero components:
$`\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},p^4\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`p^2(Q^2\overline{Q}^2)\left\{\begin{array}{c}e_{32}^{32}\end{array}\right\},`$
$`p\overline{Q}\left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},p^3\overline{Q}\left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},p^2\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\}.`$
$`\stackrel{ˇ}{R}^4`$ has $`16`$ nonzero components:
$`\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},p^4\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`p\overline{Q}\left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},p^3\overline{Q}\left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},p^2\overline{Q}^2\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},p^2\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\}.`$
Each $`\stackrel{ˇ}{R}^i`$ has a distinct set of eigenvalues, with some overlap. Immediately, the construction of $`\stackrel{ˇ}{}(u)`$ reminds us that the three distinct (diagonal) components $`(\stackrel{ˇ}{R}^i)_{jj}^{jj}`$ must be (gauge-independent) eigenvalues. In the spectral limit $`u\mathrm{}`$, these become $`1`$, $`q^{2\alpha }`$ (twice), and $`q^{4\alpha +2}`$ respectively. No matter what the gauge, we will have a minimum of these $`3`$ eigenvalues; in fact $`\stackrel{ˇ}{R}^1`$ contains only these $`3`$. Overall, we expect a maximum of $`15`$ distinct eigenvalues; these are described in Table 2.
## 5 Ambient isotopy link invariants
A state model for evaluation of link invariants of ambient isotopy based on a quantum R matrix $`\stackrel{ˇ}{R}`$, is defined by two parameters: a representation of the braid generator $`\sigma `$, and a representation of the ‘left handle’ $`C`$ . In our case, we seek a scaling factor $`\kappa `$ such that $`\sigma ^{\pm 1}=\kappa ^{\pm 1}\stackrel{ˇ}{R}^{\pm 1}`$, and a $`4\times 4`$ matrix $`C`$ such that (Einstein summation convention):
$$C_c^d(\sigma ^{\pm 1})_{db}^{ca}=\delta _b^a,$$
(41)
is satisfied.
This ensures that that the value of the link invariant over a single loop of writhe is unity, viz that the first Reidemeister move is satisfied. This is depicted in Figure 1. In fact, this only determines $`\kappa `$ up to a factor of $`\pm 1`$, but we shall ignore one case.
As $`\stackrel{ˇ}{R}`$ necessarily satisfies the QYBE (25), abstract tensors built from $`\sigma `$ are invariant under the second and third Reidemeister moves, hence we may construct representations of arbitrary braids from $`\sigma `$. As all links may be represented by braids combined with left handles, together these are sufficient parameters.
Our state model then yields a polynomial invariant of oriented of $`(1,1)`$ tangles. Evaluation of such a model for arbitrary links is described in our previous work . Here, we use the same principles and Mathematica code – only the state model parameters are changed.
### Case 1
Suitable $`C`$ that satisfy (41) are:
$`\overline{\kappa }\overline{p}^4Q^2A=\kappa \overline{Q}^2A,\mathrm{where}A=\mathrm{diag}\{+Q^2,Q^2,\overline{Q}^2,+\overline{Q}^2\},`$
hence $`\kappa =\overline{p}^2Q^2`$ suffices to yield: $`C=p^2A`$. The associated braid generator $`\sigma `$ is:
$`\overline{p}^2Q^2\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},1\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},p^2Q^2\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`\overline{p}Q(p\overline{Q}\overline{p}Q)\left\{\begin{array}{c}e_{21}^{21}\\ e_{31}^{31}\end{array}\right\},pQ(pQ\overline{p}\overline{Q})\left\{\begin{array}{c}e_{42}^{42}\\ e_{43}^{43}\end{array}\right\},`$
$`Q^2(p\overline{Q}\overline{p}Q)(pQ\overline{p}\overline{Q})\left\{\begin{array}{c}e_{41}^{41}\end{array}\right\},Q^2(Q^2\overline{Q}^2)\left\{\begin{array}{c}e_{32}^{32}\end{array}\right\},`$
$`\overline{p}Q\left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},pQ\left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},Q^2\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\},`$
$`Q^2(p\overline{Q}\overline{p}Q)^{\frac{1}{2}}(pQ\overline{p}\overline{Q})^{\frac{1}{2}}\left\{\begin{array}{c}\overline{Q}\left\{\begin{array}{c}e_{23}^{41}\\ e_{41}^{23}\end{array}\right\},Q\left\{\begin{array}{c}e_{41}^{32}\\ e_{32}^{41}\end{array}\right\}\end{array}\right\}.`$
These choices of course lead to the Links–Gould invariant $`LG`$ of .
For Cases 2–4, there are problems in finding $`\kappa `$, as the process yields no consistent choices of $`\kappa `$ without reducing the variables $`p`$ and $`Q`$. (Furthermore, in each case, relaxing the condition that $`C`$ be diagonal immediately yields the conclusion that $`C`$ *must* be diagonal if it does exist.)
### Case 2
Appropriate $`C`$ are:
$`\overline{\kappa }\overline{p}^3Q\mathrm{diag}\{+pQ,pQ,\overline{p}\overline{Q},+\overline{p}\overline{Q}\},`$
$`\kappa p\overline{Q}\mathrm{diag}\{+\overline{p}Q,\overline{p}Q,p\overline{Q},+p\overline{Q}\}.`$
The only solution to this system is found by setting $`p=\pm 1`$, whence $`\kappa =Q`$. In this case, we have:
$`C=\mathrm{diag}\{+Q,Q,\overline{Q},+\overline{Q}\}`$
and the following one-variable braid generator $`\sigma `$ (note the imaginary components):
$`Q\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},\overline{Q}\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},Q\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`(Q\overline{Q})\left\{\begin{array}{c}e_{21}^{21},e_{43}^{43}\end{array}\right\},`$
$`\pm \left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},\pm \left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},\overline{Q}\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},Q\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\},`$
$`i(Q\overline{Q})\left\{\begin{array}{c}e_{23}^{41}\\ e_{41}^{23}\end{array}\right\}.`$
After the scaling, the $`7`$ distinct eigenvalues of $`\stackrel{ˇ}{R}^2`$ coalesce to $`4`$ of $`\sigma `$; they are $`\{q^2,\overline{q}^2\pm 1\}\{Q,\overline{Q},\pm 1\}`$. Inspection of the eigenvectors shows that the $`\sigma `$ is not diagonizable, as it has only $`14`$ (out of $`16`$) distinct eigenvectors.
The resulting invariant is nontrivial. It turns out to be the Alexander–Conway invariant in variable $`q`$. (This is an experimental observation only, but has been verified for all prime knots of up to $`10`$ crossings.)
### Case 3
Appropriate $`C`$ are:
$`\overline{\kappa }\overline{p}^2Q^2\mathrm{diag}\{+p^2\overline{Q}^2,Q^4,\overline{Q}^4,+\overline{p}^2\overline{Q}^2\}\text{and}`$
$`\kappa p^2\overline{Q}^2\mathrm{diag}\{+\overline{p}^2Q^2,Q^4,\overline{Q}^4,+p^2Q^2\}.`$
The only solutions to this system leave us with an integer invariant (e.g. $`p=Q=1`$ and $`\kappa =1`$). In any case, setting $`Q=1`$ means that the R matrix degenerates to being a special case of Case 4, so we ignore these solutions.
### Case 4
Appropriate $`C`$ are:
$`\overline{\kappa }\overline{p}^2\mathrm{diag}\{p^2,Q^2,Q^2,+\overline{p}^2\}\text{and}`$
$`\kappa p^2\mathrm{diag}\{\overline{p}^2,\overline{Q}^2,\overline{Q}^2,+p^2\}.`$
Again, the only solutions to this system leave us with an integer invariant. For example, for $`p=Q=\pm 1`$ and $`\kappa =1`$, we obtain $`C=\mathrm{diag}\{+1,1,1,+1\}`$, and the braid generator:
$`\left\{\begin{array}{c}e_{11}^{11}\end{array}\right\},\left\{\begin{array}{c}e_{22}^{22}\\ e_{33}^{33}\end{array}\right\},\left\{\begin{array}{c}e_{44}^{44}\end{array}\right\},`$
$`\pm \left\{\begin{array}{c}e_{21}^{12},e_{31}^{13}\\ e_{12}^{21},e_{13}^{31}\end{array}\right\},\pm \left\{\begin{array}{c}e_{42}^{24},e_{43}^{34}\\ e_{24}^{42},e_{34}^{43}\end{array}\right\},\left\{\begin{array}{c}e_{41}^{14}\\ e_{14}^{41}\end{array}\right\},\left\{\begin{array}{c}e_{32}^{23}\\ e_{23}^{32}\end{array}\right\}.`$
The $`10`$ distinct eigenvalues of $`\stackrel{ˇ}{R}^4`$ coalesce to $`2`$ of $`\sigma `$; they are just $`\left\{\pm 1\right\}`$, and $`\sigma `$ is of course diagonizable.
The resulting invariant is integer, and in fact trivial. This has been confirmed by the application of the ‘Matveev $`\mathrm{\Delta }`$$``$ test’ . This test applies a little theorem which shows that if an invariant cannot distinguish the braids $`\sigma _1\overline{\sigma }_2\sigma _1`$ and $`\sigma _2\overline{\sigma }_1\sigma _2`$, then it cannot distinguish any knot from the unknot, since repeated interchanges of these braids are sufficient to untangle any knot.
In fact, applying this test to $`\stackrel{ˇ}{R}^4`$ shows that, regardless of $`C`$ or any special choice of the variables $`p`$ and $`Q`$, any associated invariant will be trivial. (Applying the test to Cases 2 and 3 shows that we *may* have a nontrivial invariant.)
## 6 Regular isotopy link invariants
For Cases 2 and 3, in general, we fail to build a link invariant of ambient isotopy as we cannot satisfy (41). We now build invariants associated with these cases that are only of regular isotopy.
### Case 2
The choice $`\kappa =\overline{p}^2Q`$ and $`C=\overline{p}\mathrm{diag}\{+Q,Q,\overline{Q},+\overline{Q}\}`$, yields the symmetric results:
$`(\mathrm{tr}I)(CI)\sigma ^{\pm 1}=\mathrm{diag}\{\overline{p},\overline{p},p,p\}^{\pm 1}.`$
Applying our link invariant evaluation engine to these parameters, we obtain an open tangle invariant which is a diagonal $`4\times 4`$ matrix. For a knot $`K`$ (presented as the closure $`\widehat{\beta }`$ of a braid $`\beta `$), we obtain the corresponding *regularly isotopic* $`(1,1)`$-tangle invariant of the following form:
$`\mathrm{diag}\{\underset{\text{(twice)}}{\underset{}{\overline{p}^w\mathrm{\Delta }_K(Q^2\overline{p}^2)}},\underset{\text{(twice)}}{\underset{}{p^w\mathrm{\Delta }_K(Q^2p^2)}}\},`$
where $`w`$ is the writhe of $`\beta `$, and $`\mathrm{\Delta }_K`$ is the Alexander–Conway invariant of $`K`$. Again, this result is an experimental observation only, checked to be valid for all prime knots of up to $`10`$ crossings. Setting $`p=1`$ recovers our previous observation that the only possible invariant of ambient isotopy is $`\mathrm{\Delta }_K`$.
### Case 3
The choices $`\kappa =\overline{p}^2`$ and $`C=\mathrm{diag}\{1,Q^2,\overline{Q}^2,1\}`$, yield the symmetric results:
$`(\mathrm{tr}I)(CI)\sigma ^{\pm 1}=\mathrm{diag}\{+\overline{p}^2,\overline{Q}^4,\overline{Q}^4,+p^2\}^{\pm 1}.`$
Again, we obtain an open-tangle invariant which is a diagonal $`4\times 4`$ matrix:
$`\mathrm{diag}\{\overline{p}^{2w},\underset{\text{twice}}{\underset{}{(\overline{Q}^4)^wV_K(Q^4)}},p^{2w}\},`$
where $`V_K`$ is the Jones polynomial of the link $`K`$. Again, we emphasise that this is an experimental observation, known to be valid for all prime knots of up to $`10`$ crossings.
## 7 Conclusions
Our calculations show that it is possible to access several invariants from a single solution of the TYBE. Specifically, we have shown that it is possible to recover both the Jones and the Alexander–Conway invariants from the solution of the TYBE which was originally employed to define the Links–Gould invariant $`LG`$ .
Repeating this process for other solutions may well uncover hitherto unknown invariants! Particularly, new solutions to the TYBE recently reported in warrant investigation in this context.
## Acknowledgements
Jon Links is grateful to the Australian Research Council for financial support and the cherry blossoms of “Hanami” in Kyoto, April 2000 for inspiration.
David De Wit’s research at Kyoto University is funded by a Postdoctoral Fellowship for Foreign Researchers (# P99703), provided by the Japan Society for the Promotion of Science. Dōmo arigatō gozaimashita!
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# Hidden Variables or Positive Probabilities?
## 1 Introduction
With the introduction of his celebrated inequalities in 1964, John Bell provided the basis for an experimental test to distinguish quantum mechanics from local hidden-variable theories. Since that time the universal interpretation of the results has been that quantum mechanics violates Bell’s inequalities due to its “nonlocal” character, whereas local hidden variable theories satisfy the inequalities because, as their name implies, they are “local.”
The situation is actually not so transparent. Bohr taught us to be aware of ambiguous language. Although derivations of Bell’s inequalities are evidently based on Einstein’s “locality” condition, couched in various phrases such as “principle of separability” and so forth, mathematically all derivations make an identical assumption, specifically: hidden-variable theories introduce a set of a priori positive-definite probabilities P that are not predicted by quantum mechanics. In Bohm’s classic version of the Einstein-Podolsky-Rosen experiment, for example, a particle in a spin-singlet state decays into two daughter particles with zero total angular momentum (see, e.g., Sakurai’s text or Sudarshan and Rothman , henceforth SR). According to local hidden-variable theories there is an a priori positive-definite probability that the daughter particles will be detected with spins “up” along a chosen axis. Quantum mechanics, on the other hand, assumes that the daughter particles are in a superposition of states and so, by definition, there can be no a priori probability P such that their spins will be detected along a given direction.
Contrary to this view, in SR we pointed out that quantum mechanics does predict a set of a priori probabilities, in exactly the same way as do hidden-variable theories, but the quantum probabilities are not positive-definite. They are nevertheless meaningful in that when applied to physical situations they give the standard quantum-mechanical answers, in particular the usual violation of Bell’s inequalities. Given the exact analogy in producing the two sets of probabilities the distinction between “local” hidden-variable theories and “non-local” quantum mechanics is dissolved. From this point of view one merely has two competing theories that give two different sets of probabilities; it is unsurprising that hidden-variables theories fail experimental tests of Bell’s inequalities because they used the wrong set of probabilities for a quantum-mechanical problem.
The notion of “extended” probabilities dates back to Dirac and we have not been the only authors to suggest that they can resolve the EPR paradox (see ) but, needless to say, the SR argument has not found widespread acceptance. Recently, several rather old papers, in particular one by Eberhard entitled “Bell’s Theorem Without Hidden Variables,” have come to our attention. Eberhard’s paper is of interest because it claims to show that a more general version of Bell’s inequalities, known as the CHSH inequality (after Clauser, Horne, Shimony and Holt) , is violated by quantum mechanics, and that the CHSH inequality can be demonstrated solely on the basis of the locality principle, without the introduction of hidden variables. (A slightly later paper by Peres gives an almost identical argument; one by Stapp is in some respects similar.) At first sight these proofs appear to assume little more than $`2<2\sqrt{2}`$. On closer inspection, however, we find that they “play into our hands,” i.e., they may not make an explicit statement about hidden variables but they do assume a set of positive-definite probabilities. We now demonstrate this is so, reinforcing the contention in SR that, despite any words employed, the crucial mathematical assumption in derivations of Bell’s inequalities is not locality but positive probability.
## 2 The Eberhard Argument
Eberhard considers two identical apparata, $`A`$ and $`B`$, at two different locations. On apparatus $`A`$ is a knob $`a`$ that can be turned to two positions, 1 and 2. On apparatus $`B`$ is a knob $`b`$ that can also be turned to two positions, 1 and 2. With its knob at either position apparatus $`A`$ can record a series of events. It is not important exactly what the events are, but we assume that for each event each apparatus can measure only one of two possible outcomes, which for simplicity we take to be $`\pm 1`$. When the knob $`a`$ is in the 1 position, we designate the outcome of the $`jth`$ event as $`\alpha _{1j}`$, with similar notation for position 2 and knob $`b`$. For each event we can thus in principle have: $`\alpha _{1j}=\pm 1,\alpha _{2j}=\pm 1,\beta _{1j}=\pm 1,\beta _{2j}=\pm 1`$. However, for each measurement we will choose only one setting on each apparatus, so a given event will produce a pair of readings, such as $`\alpha _1=1,\beta _2=1`$. (Here and below we suppress the subscript $`j`$ when it will not cause confusion.)
For a series of $`N`$ measurements Eberhard next defines a quantity $`C`$, such that
$$C=\frac{1}{N}\underset{j=1}{\overset{N}{}}\alpha _j\beta _j$$
(2.1)
We see that $`C=<\alpha _j\beta _j>`$, the statistical mean of the $`N`$ products $`\alpha _j\beta _j`$. No restriction is placed on the fraction of the $`N`$ measurements for which the $`\alpha `$’s and $`\beta `$’s come out positive or negative, but note that each product $`\alpha _j\beta _j=1`$ when $`\alpha `$ and $`\beta `$ have the same sign and $`\alpha _j\beta _j=1`$ when they have opposite signs. Thus $`C`$ represents the fraction of events in which $`\alpha `$ and $`\beta `$ have the same sign minus the fraction in which they have opposite sign.
Because each knob has two positions, there are four possible versions of $`C`$. That is, we can define
$`C_{11}`$ $`=`$ $`<\alpha _1\beta _1>`$
$`C_{12}`$ $`=`$ $`<\alpha _1\beta _2>`$
$`C_{21}`$ $`=`$ $`<\alpha _2\beta _1>`$
$`C_{22}`$ $`=`$ $`<\alpha _2\beta _2>`$ (2.2)
(sum on $`j`$ understood). Here, $`C_{11}`$ is just the above statistical mean when knobs $`a`$ and $`b`$ are both in position 1, and so forth.
Now, for each event let
$$\gamma \alpha _1\beta _1+\alpha _1\beta _2+\alpha _2\beta _1\alpha _2\beta _2.$$
(2.3)
Then, the statistical mean of $`\gamma `$ is just
$`<\gamma >`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}\gamma _j`$ (2.4)
$`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}(\alpha _1\beta _1+\alpha _1\beta _2+\alpha _2\beta _1\alpha _2\beta _2)`$
$``$ $`C_{11}+C_{12}+C_{21}C_{22},`$
where in the second line we have again suppressed $`j`$.
The locality condition enters the discussion when we attempt to put bounds on $`<\gamma >`$. Recall that a knob will be set to either position 1 or 2 for each measurement. We assume that a measurement on $`A`$ is independent of a measurement on $`B`$. The $`\alpha `$’s and $`\beta `$’s are thus treated independently. This is the locality condition.
At this point a digression is necessary. Eberhard states that only one setting of each knob (position 1 or 2) will be used for each measurement, and that thus only one $`\alpha `$ or $`\beta `$ is recorded for each event. However, if this were indeed the case, then for each measurement only one term in $`\gamma `$ would survive (one product $`\alpha \beta `$) and the upper bound on $`\gamma `$ would be 1 (cf. Eqs. (2.3) and (2.7)). That the upper bound is 2 shows that mathematically all four possible terms $`\alpha \beta `$ are present in $`\gamma `$. Consequently, not only are the $`\alpha `$’s being taken to be independent of the $`\beta `$’s but $`\alpha _1`$ ($`\beta _1`$) is being treated as independent of $`\alpha _2`$ ($`\beta _2`$). The rationale for including all $`\alpha `$’s and $`\beta `$’s in $`\gamma `$ simultaneously comes from a 1971 suggestion of Stapp . Stapp, Eberhard (and Peres in his nearly identical thought experiment), are actually considering all possible outcomes of the measurements in a hypothetical ensemble space. By doing so they intend to show that any conceivable outcome of the experiment is violated by quantum mechanics.
One can take several attitudes toward such a procedure. A first possible attitude is that it is illegitimate to speculate about the results of unperformed experiments. In other words, if one takes the quantity $`\gamma `$ literally, the knobs must be set to two positions at once, a physical impossibility. A second view is that it is indeed legitimate to think about all possible outcomes of an experiment<sup>1</sup><sup>1</sup>1This concept is often referred to as “counterfactual definiteness,” after Stapp. and that if one does so, one is forced to the conclusion that quantum mechanics is nonlocal. In fact, there is a third possible viewpoint. As we discuss below, the $`\gamma `$’s are derivable from the “master probabilities” employed in a standard derivation of Bell’s inequalities, quantities that are not directly measurable but nevertheless have physical consequences. Hence both the Eberhard procedure and the standard derivation suffer from exactly the same ambiguities. For the moment it is not important which philosophy one adopts; we merely treat $`\gamma `$ as a mathematical quantity, as Eberhard does. At the same time, however, we see that by treating all the $`\alpha `$’s and $`\beta `$’s as independent, mathematically the locality condition becomes indistinguishable from the general assumption of independent variables.
In any case, following Eberhard we assume 16 possible values for each $`\gamma `$. At this stage of the exposition, Eberhard goes through an elaborate argument to show that $`\gamma 2`$ always. However, let us redistribute the terms in Eq. (2.3) and write
$$\gamma =\alpha _1(\beta _1+\beta _2)+\alpha _2(\beta _1\beta _2).$$
(2.5)
Because $`\beta _1`$ and $`\beta _2`$ are equal or of opposite sign, if the first term is nonzero, the second term is zero and vice versa. Thus we can see trivially that $`\gamma =\pm 2`$ always and $`|\gamma |=2`$, period.
But by the triangle inequality we know that
$$|\frac{1}{N}\underset{j=1}{\overset{N}{}}(\alpha _1\beta _1+\alpha _1\beta _2+\alpha _2\beta _1\alpha _2\beta _2)|\frac{1}{N}\underset{j=1}{\overset{N}{}}|(\alpha _1\beta _1+\alpha _1\beta _2+\alpha _2\beta _1\alpha _2\beta _2)|$$
(2.6)
Yet from Eq. (2.4) and Eq. (2.3) this is by definition
$`|C_{11}+C_{12}+C_{21}C_{22}|`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{j=1}{\overset{N}{}}}|\gamma _j|`$ (2.7)
$`=`$ $`{\displaystyle \frac{1}{N}}\times N\times 2`$
The CHSH inequality follows immediately:
$$|C_{11}+C_{12}+C_{21}C_{22}|2,$$
(2.8)
or, in more compact notation,
$$|𝒞|2.$$
(2.9)
Eberhard next considers a quantum-mechanical experiment in which two photons are emitted in the directions of $`A`$ and $`B`$ by an atom between them. The photons are detected by polarizers; each $`\alpha `$ ($`\beta `$) is taken to be +1 when one polarization is detected and -1 when the other is detected. Unfortunately, at this point the paper becomes very unclear. Eberhard merely asserts without calculation that for each of the $`C`$’s in Eq. (2.2), quantum mechanics predicts that “if the number of events $`N`$ is large enough, then $`Ccos(2a2b)`$,” where $`2a2b`$ is twice the angle between the polarizers. Actually, no approximation is necessary. For spin-1/2 particles, the correct prediction is
$$𝒞_{qm}=3cos\theta cos3\theta ,$$
(2.10)
which we derive below, and in which $`\theta `$ is the angle between polarizers. (The result for photons will be the same if $`\theta `$ is taken to be twice the angle between polarizers.) Note that for $`\theta =45^o`$ (2.10) gives $`𝒞_{qm}=2\sqrt{2}2`$. Therefore, quantum mechanics violates the CHSH inequality, just as it does the Bell inequalities.
As mentioned above, the demonstration seems to assume almost nothing: no hidden variables, merely “locality,” which implies that a certain mathematical quantity $`\gamma `$ always equals $`\pm `$ 2. However, on closer inspection we find that more than an assumption of independent $`\alpha ^{}s`$ and $`\beta `$’s is being made. In the first place, the value 2 on the right-hand side of Eq. (2.8) is entirely arbitrary and results merely from the choice of $`\pm 1`$ as the “eigenvalues” for $`\alpha `$ and $`\beta `$. One could have equally well chosen $`\pm 1000`$. In that case, however, one would necessarily have to assume that the corresponding quantum experiment also had eigenvalues of $`\pm 1000`$. This matter is not so serious, but it nevertheless illustrates that the CHSH inequality is not a purely mathematical assertion; a real measurement does lurk in the background.
The central issue lies elsewhere. Eberhard’s version of CHSH inequality is a statement about the statistical mean of $`\gamma `$, and therefore it does deal with a probability distribution over the $`\gamma `$. Moreover, the frequency that a particular $`\gamma `$ occurs is clearly taken to be positive. That probabilities should be positive-definite is usually regarded as self-evident, but because the assumption is the crux of the matter, we spend a moment examining it. (In the Appendix we detail where other authors have made the same assumption.)
As mentioned, there are 16 possible combinations of $`\alpha _1\beta _1+\alpha _1\beta _2+\alpha _2\beta _1\alpha _2\beta _2(=\gamma )`$, of which eight have the value +2 and eight have the value -2. In a sequence of $`N`$ measurements, let us suppose that +2 occurs $`n_1`$ times and -2 occurs $`n_2`$ times, such that $`n_1+n_2=N`$. Then
$$𝒞=\frac{2}{N}[n_1n_2].$$
(2.11)
If all frequencies are equal, i.e. $`n_1=n_2`$, then $`𝒞=0`$. If $`n_2=0`$, then $`𝒞=2`$ and if $`n_1=0`$ then $`𝒞=2`$. But here we have assumed that both $`n_1`$ and $`n_2`$ are positive-definite. If $`n_2<0`$, then $`𝒞>2`$. In other words, the step leading to the second line in Eq. (2.7) is valid only when $`|n|=n`$.
The notion of “extended” (non-positive-definite) probabilities has been considered by a surprising number of prominent investigators, but the majority of physicists continue to regard them with distaste, if not revulsion. Nevertheless, the quantum violation of the bound on $`𝒞`$ is effectively due to the fact that quantum mechanics allows negative probabilities. In the next section we examine this claim in greater detail.
## 3 Quantum Mechanical Probabilities
Before deriving Eq.(2.10), it will be helpful to summarize the procedure for obtaining the standard Bell Inequalities in order to point out similarities to the CHSH-Eberhard experiment. The reader is referred to SR or Sakurai for additional details; see also the Appendix. Like its successor, Bell’s theorem is valid for local hidden-variable theories, which involve only classical probabilities. In a typical derivation such as Sakurai’s one assumes that spin measurements may be made along any of three axes, a, b and c. A system of decaying atoms emits $`N`$ particles of which a certain fraction are taken to be, say, of the type (a+, b+, c+) $`(+++)`$, which designates spin up along all three axes. To ensure zero total angular momentum, each emitted particle of type (+++) must be paired with one of type ($``$). There are eight such spin combinations in all, as listed in Table 1.
The probability that (+++) is emitted (and in the case of hidden variables, detected) is defined simply as $`P(+++)=N(+++)/N`$. One can immediately object that such a probability is unphysical because to determine it requires three simultaneous spin measurements on a system of two particles, which is impossible. To eliminate this difficulty, one forms pairwise probabilities of the type $`P(𝐚+,𝐛+)P(++)`$, which represents the joint probability that the first particle will be found + along a and the second particle + along b. This is easily done. From the table, the total number of particles such that the first particle’s spin is + along a is $`N(++)+N(+)`$, which must be paired with $`N(+)+N(++)`$, the total number of particles for which the second particle’s spin is + along b. This combination is labeled $`N_3+N_5`$. Next one forms triangle-type inequalities such as
$$N_3+N_5(N_2+N_5)+(N_3+N_7),$$
(3.1)
which is obviously true, since we have just added positive numbers to $`N_3+N_5`$. Dividing by $`N`$ gives by definition
$$P(𝐚+,𝐛+)P(𝐚+,𝐜+)+P(𝐜+,𝐛+),$$
(3.2)
one of the Bell inequalities. Eq. (3.2) involves only one measurement on each particle and so represents a physically realizable situation. Note that the “three-probabilities” $`P(+++)`$ were reduced to pairwise probabilities $`P(++)`$ by summing over the spins on the extraneous axis, in the above example c. We emphasize that, just as was the case for the CHSH inequality, the Bell inequality is valid only if the N’s and hence the P’s are taken to be positive-definite. In SR we demonstrated that one can form quantum probabilities $`P(+++)`$, analogous to the classical probabilities, then sum over the third argument exactly as above to get pairwise quantum probabilities $`P(++)`$ that violate (3.2) in the usual way.
By this point the reader will have noticed a similarity between the $`\gamma `$’s in Eberhard’s experiment and the three-probabilities here. Authors who derive the generalized Bell inequalities introduce $`\gamma `$ as a measure of correlations between real and imagined experiments but, as mentioned, if one takes it literally it amounts to having the apparatus knobs set on two positions simultaneously. This would seem to represent the same sort of physical impossibility as that of making three simultaneous spin measurements on two particles. Indeed, we will demonstrate in Section 4 that the two procedures are identical: Introducing an ensemble of hypothetical measurements is exactly equivalent to assuming a “master probability distribution” that requires more than two simultaneous spin measurements on two particles. Before doing so, however, we return to the Eberhard derivation.
Eberhard’s experiment involves four axes, $`𝐚_\mathrm{𝟏},𝐚_\mathrm{𝟐},𝐛_\mathrm{𝟏},𝐛_\mathrm{𝟐}`$, rather than three, but otherwise is almost identical to the standard derivation of Bell’s inequalities and so it is not surprising that the above procedure can be followed to demonstrate a violation of the CHSH inequality. We first need to compute the quantum pairwise probabilities of the type just mentioned, $`P(𝐚+,𝐛+)`$. There are several ways to do this. Following SR, we write the quantum-mechanical projection operator for spin-1/2 particles as
$$\mathrm{\Pi }(𝐚\pm )=\frac{1}{2}(\mathrm{𝟏}\pm 𝝈\mathbf{}𝒂).$$
(3.3)
In this equation we are representing the Pauli spin matrices as a vector, $`𝝈=\widehat{𝐢}\sigma _𝐱+\widehat{𝐣}\sigma _𝐲+\widehat{𝐤}\sigma _𝐳`$. Thus $`𝝈\mathbf{}𝒂=\sigma _xa_x+\sigma _ya_y+\sigma _za_z`$ represents a traceless, $`2\times 2`$ matrix and 1 is the unit matrix. Now, the expectation value of any operator $`𝒪`$ can be written $`<𝒪>=\mathrm{Tr}(\rho 𝒪)`$, where $`\rho `$ is the density matrix $`\mathrm{diag}(1/2,1/2)`$ for an initially unpolarized beam. The probability of finding the first particle in the + state along a is thus $`\mathrm{Tr}(\rho \mathrm{\Pi }(𝐚))=1/2`$. Similarly, the joint probability $`P(𝐚+,𝐛\pm )`$ of finding the first particle in the + state along a and the second particle in the $`\pm `$ state along b is
$`P(𝐚+,𝐛\pm )`$ $`=`$ $`{\displaystyle \frac{1}{2}}Tr\mathrm{\Pi }(𝐚)\mathrm{\Pi }(𝐛\pm )`$ (3.4)
$`=`$ $`{\displaystyle \frac{1}{8}}Tr\{(\mathrm{𝟏}+𝝈\mathbf{}𝒂)(\mathrm{𝟏}\pm 𝝈\mathbf{}𝒃)\}`$
$`=`$ $`{\displaystyle \frac{1}{4}}(1\pm 𝐚𝐛).`$
Here, use has been made of the standard identity (see )
$$(𝝈\mathbf{}𝒂)(𝝈\mathbf{}𝒃)=(𝐚𝐛)\mathrm{𝟏}+i𝝈(𝐚\times 𝐛).$$
(3.5)
Because the Pauli matrix is traceless, taking the trace of (3.5) yields 2$`𝒂\mathbf{}𝒃`$.
Equation (3.4) is simply a sophisticated way of writing Malus’ law. The first factor of $`1/2`$ in (3.4) gives the probability of detecting a particle in the + state along the a axis. The remaining factor $`1/2(1+𝒂\mathbf{}𝒃)=1/2(1+cos\theta )`$, where $`\theta `$ is the angle between polarizers. For photons(where $`\theta `$ is taken to be the double angle) this then represents the usual decrease in intensity with $`cos^2\theta `$. For a Bohm-type experiment,which assumes an (antisymmetric) spin-singlet state, one should choose the $``$ on the right of (3.4)when computing $`P(𝐚+,𝐛+)`$ to conserve angular momentum. With either sign, by inserting (3.4) into (3.2), it is straightforward to show that quantum mechanics violates Bell’s inequalities.
For the Eberhard experiment we take the knob settings $`a_1,a_2,b_1,b_2`$ to represent the position of the polarizers on the measuring devices. Recall that his quantities $`C=<\alpha \beta >`$ represented the fraction of events in which $`\alpha `$ and $`\beta `$ had the same sign minus the fraction in which they had opposite signs, irrespective of whether an individual spin is $`+`$ or $``$. Evidently the equivalent quantum expression is $`1/2(1+𝐚𝐛)1/2(1𝐚𝐛)`$. Then
$$𝒞_{qm}=𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}.$$
(3.6)
If the axes are chosen to be coplanar such that $`𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}=𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}=𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}=cos\theta `$ and $`𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}=cos3\theta `$, then (3.6) gives exactly (2.10), which violates the CHSH inequality for $`\theta =45^o`$.
The derivation of (2.10) just given involved only pairwise probabilities and did not go beyond standard quantum mechanics. With the projection-operator formalism, however, it is not difficult to write down the joint probability for four “simultaneous” spin measurements among four axes. An example would be $`P(++++)`$, in analogy to the classical three-probability mentioned earlier that appears in the derivation of Bell’s inequality. Extending (3.4) to four arguments we take
$$P(\lambda 𝐚_\mathrm{𝟏},\mu 𝐚_\mathrm{𝟐},\nu 𝐛_\mathrm{𝟏},\tau 𝐛_\mathrm{𝟐})=\frac{1}{2}Tr\{\mathrm{\Pi }(\lambda 𝐚_\mathrm{𝟏})\mathrm{\Pi }(\mu 𝐚_\mathrm{𝟐})\mathrm{\Pi }(\nu 𝐛_\mathrm{𝟏})\mathrm{\Pi }(\tau 𝐛_\mathrm{𝟐})\},$$
(3.7)
where $`\lambda ,\mu ,\nu ,\tau `$ are chosen as $`\pm 1`$ to represent up or down. For the symmetric case this is
$$P(\lambda 𝐚_\mathrm{𝟏},\mu 𝐚_\mathrm{𝟐},\nu 𝐛_\mathrm{𝟏},\tau 𝐛_\mathrm{𝟐})=\frac{1}{32}Tr\{(\mathrm{𝟏}+\lambda 𝝈\mathbf{}𝒂_\mathrm{𝟏})(\mathrm{𝟏}+\mu 𝝈\mathbf{}𝒂_\mathrm{𝟐})(\mathrm{𝟏}+\nu 𝝈\mathbf{}𝒃_\mathrm{𝟏})(\mathrm{𝟏}+\tau 𝝈\mathbf{}𝒃_\mathrm{𝟐})\}$$
(3.8)
We will need the antisymmetric expression later to make the subtraction just done above. Assuming that a measurement of + on knob $`a`$ requires $``$ on knob $`b`$, the antisymmetric case will be the same expression as (3.8)with the signs on the $`b`$’s reversed. We calculate only the symmetric case and state the results for the antisymmetric case as needed.
Working out (3.8) and making frequent use of the identity (3.5) yields
$`P(\lambda 𝐚_\mathrm{𝟏},\mu 𝐚_\mathrm{𝟐},\nu 𝐛_\mathrm{𝟏},\tau 𝐛_\mathrm{𝟐})={\displaystyle \frac{1}{16}}\{1+\lambda \mu 𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}+\lambda \nu 𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}+\lambda \tau 𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}`$
$`+\mu \nu 𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}+\mu \tau 𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}+\nu \tau 𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}`$
$`+\mathbf{ı}\lambda \mu \nu (𝐚_\mathrm{𝟏}\times 𝐚_\mathrm{𝟐})𝐛_\mathrm{𝟏}+\mathbf{ı}\lambda \mu \tau (𝐚_\mathrm{𝟏}\times 𝐚_\mathrm{𝟐})𝐛_\mathrm{𝟐}`$
$`+\mathbf{ı}\lambda \nu \tau (𝐛_\mathrm{𝟏}\times 𝐛_\mathrm{𝟐})𝐚_\mathrm{𝟏}+\mathbf{ı}\mu \nu \tau (𝐛_\mathrm{𝟏}\times 𝐛_\mathrm{𝟐})𝐚_\mathrm{𝟐}`$
$`+\lambda \mu \nu \tau [(𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐})(𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐})+\mathrm{𝟐}(𝐚_\mathrm{𝟏}\times 𝐚_\mathrm{𝟐})(𝐛_\mathrm{𝟏}\times 𝐛_\mathrm{𝟐})]\}.`$ (3.9)
Notice that this expression is complex due to the imaginary elements of $`\sigma _y`$. If we desire a real result to eventually make contact with the usual quantum predictions, we can easily eliminate the imaginary terms. Note that $`\mathrm{\Pi }(\lambda 𝐚_\mathrm{𝟏})\mathrm{\Pi }(\mu 𝐚_\mathrm{𝟐})\mathrm{\Pi }(\nu 𝐛_\mathrm{𝟏})\mathrm{\Pi }(\tau 𝐛_\mathrm{𝟐})`$ has been written in an arbitrary order; it is not symmetric in the arguments. There are $`4!`$ permutations of the arguments in this expression, twelve even and twelve odd. In (3.9) each imaginary term is a triple scalar product, which is invariant under even permutations and changes sign under odd permutations. Thus these terms vanish under symmetrization, as does the double cross product in the last line. The symmetrized version of (3.9) is
$`P(\lambda 𝐚_\mathrm{𝟏},\mu 𝐚_\mathrm{𝟐},\nu 𝐛_\mathrm{𝟏},\tau 𝐛_\mathrm{𝟐})={\displaystyle \frac{1}{16}}\{1+\lambda \mu 𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}+\lambda \nu 𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}+\lambda \tau 𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}`$
$`+\mu \nu 𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}+\mu \tau 𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}+\nu \tau 𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}`$
$`+{\displaystyle \frac{1}{3}}\lambda \mu \nu \tau [(𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐})(𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐})+(𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏})(𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐})+(𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐})(𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐})]\},`$ (3.10)
which is entirely real. <sup>2</sup><sup>2</sup>2It is not actually necessary to symmetrize (3.9). One can leave it as a complex expression, but when the sum over the extraneous arguments is performed as in (3.11), the imaginary terms cancel and the result will be entirely real, as before. However, the complex four-probability is not symmetric in the arguments.
It is now easy to read off the various four-probabilities, $`P(++++),P()`$ etc. for each case merely by choosing the required signs of $`\lambda ,\mu ,\nu ,\tau `$. The sixteen possibilities are listed for convenience in Table II. Note that these four-probabilities do sum to one and therefore in that respect behave as ordinary probabilities. However, although it is perhaps not evident from inspection, several of these probabilities can become negative. We plot $`P(+++)`$ and $`P(++)`$ in Figure 1. The antisymmetric $`P`$’s can be obtained from the symmetric ones merely merely by flipping the signs on the two $`b`$’s.
From these four-probabilities one can form the quantity $`𝒞_{qm}`$ in Eq. (3.6) in exact analogy to the procedure used for deriving the Bell inequalities. To compute $`P(𝐚_\mathrm{𝟏}+,𝐛_\mathrm{𝟏}+)`$, for example, we only care that the first particle will be found + along $`𝐚_\mathrm{𝟏}`$ and the second particle will be found $`+`$ along $`𝐛_\mathrm{𝟏}`$. As before, we count all such possibilities by summing over the two extraneous arguments, $`𝐚_\mathrm{𝟐}`$ and $`𝐛_\mathrm{𝟐}`$. Thus, for the symmetric wavefunction,
$$P(𝐚_\mathrm{𝟏}+,𝐛_\mathrm{𝟏}+)=P(+\mathrm{\_}\mathrm{\_}+\mathrm{\_}\mathrm{\_})=P(++++)+P(+++)+P(++)+P(+++)$$
(3.11)
Reading off these $`P`$’s from Table II and performing the sum yields
$$\frac{1}{4}(1+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}),$$
(3.12)
which is exactly Eq. (3.4). For the antisymmetric wave function one obtains $`1/4(1𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏})`$. Similar expressions are obtained for the other three pairwise probabilities. Clearly, subtracting the antisymmetric expressions from the symmetric ones and adding the four terms leads back to Eq. (3.6) for $`𝒞_{qm}`$. This procedure must work because the four-probabilities are symmetric in all the arguments; summing over any of them produces an equal number of terms of opposite sign, which cancel out, leaving the usual quantum pairwise probabilities.
## 4 Discussion and Conclusions
We have shown that, like the Bell inequalities, the CHSH inequality assumes positive-definite probabilities and that quantum mechanics breaks both inequalities effectively because it introduces negative weights to the measurements. These negative four-probabilities enter the derivation in exactly the same way as the classical three-probabilities entered the derivation of the Bell’s inequalities. If they are unphysical, it is not necessarily because they are negative, but because it is impossible to make four simultaneous spin measurements on two particles. By the same token, it is impossible to make three simultaneous spin measurements on two particles. In any case, neither the classical three-probabilities found in Bell’s theorem, nor the four-probabilities that figure here are actually measured. Both merely serve as “master distributions” from which to derive the usual pairwise probabilities, classical and quantum, which are both positive-definite. To reiterate our earlier remarks, from this point of view it is not surprising that the Bell and CHSH inequalities are violated by experimental tests; they merely used the wrong set of probabilities for a quantum-mechanical problem.
Although one might choose to reject negative probabilities as unphysical, one should not reject the notion of master probability distributions in favor of correlations between real and imaginary experiments because the two procedures are identical! Recall again that Eberhard’s quantity $`C_{11}`$ was $`C_{11}=\frac{1}{N}_{j=1}^N\alpha _{1j}\beta _{1j}`$, which represented the fraction of events $`\alpha _1\beta _1`$ that had the same sign minus the fraction that had opposite sign. Thus by definition we can write
$$C_{11}=P(𝐚_\mathrm{𝟏}+,𝐛_\mathrm{𝟏}+)+P(𝐚_\mathrm{𝟏},𝐛_\mathrm{𝟏})[P(𝐚_\mathrm{𝟏}+,𝐛_\mathrm{𝟏})+P(𝐚_\mathrm{𝟏}+,𝐛_\mathrm{𝟏})].$$
(4.1)
Now, in exact analogy with the procedure of Section 3 we imagine that these pairwise probabilities can be derived from a master distribution involving all four axes $`𝐚_\mathrm{𝟏},𝐚_\mathrm{𝟐},𝐛_\mathrm{𝟐},𝐛_\mathrm{𝟑}`$. In that case, as in Eq (3.11), $`P(++)=P(𝐚_\mathrm{𝟏}+,𝐛_\mathrm{𝟏}+)=P(++++)+P(+++)+P(++)+P(+++)`$, with analogous expressions for $`P(),P(+)`$ and $`P(+)`$. There are thus 16 terms that contribute to $`C_{11}`$, similarly for $`C_{12},C_{21}`$ and $`C_{22}`$. Writing out all 64 terms yields for $`𝒞=<\gamma >`$:
$`𝒞=`$ $`2\{P(++++)+P()+P(+++)+P(+)`$ (4.2)
$`+P(+++)+P(+)+P(++)+P(++)`$
$`P(+++)P(+)P(+++)P(+)`$
$`P(++)P(++)P(++)P(++)\}`$
These $`P`$’s are general and may be taken to be either classical or quantum. Notice half enter with positive sign and half with negative. If all the probabilities are equal, then $`𝒞=0`$. If those that enter with negative sign are zero, then $`𝒞=2`$ and if those that enter with positive sign are zero, then $`𝒞=2`$. All this is in complete agreement with the analysis of Section 2. Clearly, if the $`P`$’s are positive-definite then $`𝒞2`$, but if the probabilities are allowed to become negative then this bound is violated. If the $`P`$’s are assumed to be quantum, they take on the values given by Table II. In this case, inserting those values into (4.2) gives exactly (3.6), as before.
This demonstration shows clearly that the $`\gamma `$’s can be derived from a master probability distribution which involves simultaneous spin measurements along four axes. The only difference between the classical and quantum cases is that in the former we assume the probabilities are positive-definite. The master distributions themselves cannot be regarded as any more or less meaningful than the space of hypothetical measurements, because the procedures are exactly equivalent. Indeed, we see that there is no difference between the Eberhard procedure and the usual derivation of Bell’s inequalities.
There remains the problem of interpretation. Most people insist that probability be defined in terms of relative frequency of events, in which case it must be positive-definite. In quantum mechanics, however, although one can define the expectation value in terms of the square of the wave amplitude, which corresponds to a relative-frequency interpretation, an alternate procedure is available. The expectation value may also be taken as a functional of the dynamical variables under consideration, for example position and momentum. Classically, one might consider a Maxwellian distribution of particles in phase space; integrating over position or momentum would give the marginal probability distribution for the conjugate variable. But in quantum mechanics, the uncertainty principle precludes precise simultaneous knowledge of noncommuting variables. If one attempts to associate a functional with a distribution over noncommuting variables, such that an integration over one of them gives the correct marginal distribution for the other, then one finds that the distribution function must in places become negative. This is the well known Wigner Distribution.
In the case of spin, the different components of angular momentum do not commute; hence no ordinary (positive-definite) probability distribution can be defined over the various components simultaneously. Any distribution will share with the Wigner distribution the property that it will become negative in some region of “phase space.” For example, in the spin-1/2 systems we have been considering, the probability of finding $`S_z`$ in the $`+`$ state and $`S_x`$ in the $`+`$ state is given by taking the trace of the product of the projection operators, as we have done earlier. Now, given a state with $`S_x=+`$, the probability is 1/2 for finding $`S_z=+`$, and 1/2 for $`S_z=`$. Suppose, however, that many measurements show $`S_z=+`$, always, but that $`S_x=+`$ appears with probability $`\lambda `$ and $`S_x=`$ appears with probability $`1\lambda `$ ($`0\lambda 1`$). The probability for finding $`S_z=`$ must be then be $`(1/2)\lambda +(1/2)(1\lambda )=1/2`$. On the one hand the probability of $`S_z=`$ must equal zero. On the other hand, no mixture of $`S_x=+`$ and $`S_x=`$ can give a zero probability for $`S_z=`$.
This is quite a general property of noncommuting variables and has little to do with quantum mechanics. In such situations the best that one can ask for is that the probability distribution give the correct marginal distribution for one of the variables, in our case one component of angular momentum. This is what has been found in the present paper. The probability distribution for simultaneous measurements along three or more axes are not positive-definite, but the marginal distributions that give correlations between two spin components are, and are in accord with the standard predictions of quantum mechanics.
The main point of this paper has been that assumptions beyond locality do enter into derivations of Bell’s inequalities. It is worth mentioning yet another tacit assumption: that space is flat. The notion of parallel and antiparallel spins is only well defined for flat space where the measurement axes (the “z” axes) can be taken to be everywhere fixed relative to one another. In curved space there is no universal definition of parallel and one can only compare spins in distant locations by parallel transporting the measurement axes . In the case of nonnegligible gravitational fields, then, the “nonlocal” EPR correlation between two particles, to the extent that they can be said to exist at all, must be the result of parallel transport, a local phenomenon.
Returning to probabilities, we find ourselves in a strange situation. If one insists that probabilities remain positive-definite, we are forced to use vague and imprecise concepts, such as “local” or “nonlocal” to describe the outcome of the EPR experiment. On the other hand, we are able formulate the precise mathematical conditions necessary for the violation of the Bell and CHSH inequalities, although at the cost of introducing negative probabilities. Most investigators would say that a unified, physical interpretation of negative probabilities is, in fact, exactly what is currently lacking. To be sure, Feynman conceded (see and ; also )that all the results of quantum mechanics can be analyzed in terms of negative probabilities but he remained skeptical about the utility of such an approach and that a useful meaning could be attached to it. Nevertheless, many of the interpretational problems associated with negative probabilities stem from an insistence on viewing them within the framework of relative frequencies. This is clearly “no go.” We have shown that a more natural framework for their interpretation arises when one considers the expectation value as a measure of probability over noncommuting variables. One can even go further than we have and consider complex probability measures (), which also involve expectation values. Under such circumstances it is well to bear in mind that imaginary numbers are more similar to rotations than to real numbers. One should also bear in mind the very word “imaginary,” an obsolete relic of their original status.
Acknowledgements We would like to thank Sebastiano Sonego for bringing our attention to the Eberhard and Peres papers and explaining a few details of the former. T.R. would also like to thank Gabe Spalding for helping to check some algebra.
Note added: Since this paper was initially posted, José Cereceda has come to essentially the same conclusions (see quant-ph/0010091).
## Appendix
Many researchers appear unwilling to accept that any assumptions beyond locality are employed in the derivations of Bell’s inequalities. We now list a few of the proofs we have found and point out explicitly where the assumption of positive probabilities enters.
Bell 64. In Bell’s original proof he defines two quantities $`A(\stackrel{}{a},\lambda )=\pm 1`$, $`B(\stackrel{}{b},\lambda )=\pm 1`$. He defines a normalized probability distribution $`\rho (\lambda )`$, such that $`𝑑\lambda \rho (\lambda )=1`$. The expectation value of the spin components $`\stackrel{}{\sigma _1}\stackrel{}{a}`$ and $`\stackrel{}{\sigma _2}\stackrel{}{b}`$ is
$$P(\stackrel{}{a},\stackrel{}{b})=𝑑\lambda \rho (\lambda )A(\stackrel{}{a},\lambda )B(\stackrel{}{b},\lambda ),$$
(A.1)
which he shows can be written (his equation 14) as
$$P(\stackrel{}{a},\stackrel{}{b})=𝑑\lambda \rho (\lambda )A(\stackrel{}{a},\lambda )A(\stackrel{}{b},\lambda ).$$
(A.2)
When another vector $`\stackrel{}{c}`$ is involved, one has
$$P(\stackrel{}{a},\stackrel{}{b})P(\stackrel{}{a},\stackrel{}{c})=𝑑\lambda \rho (\lambda )[A(\stackrel{}{a},\lambda )A(\stackrel{}{b},\lambda )A(\stackrel{}{a},\lambda )A(\stackrel{}{c},\lambda )]$$
(A.3)
Bearing in mind that$`A(\stackrel{}{b},\lambda )=1/A(\stackrel{}{b},\lambda )`$ one can rewrite this as
$$P(\stackrel{}{a},\stackrel{}{b})P(\stackrel{}{a},\stackrel{}{c})=𝑑\lambda \rho (\lambda )A(\stackrel{}{a},\lambda )A(\stackrel{}{b},\lambda )[A(\stackrel{}{b},\lambda )A(\stackrel{}{c},\lambda )1].$$
(A.4)
Bell then asserts
$$|P(\stackrel{}{a},\stackrel{}{b})P(\stackrel{}{a},\stackrel{}{c})|𝑑\lambda \rho (\lambda )[A(\stackrel{}{b},\lambda )A(\stackrel{}{c},\lambda )1],$$
(A.5)
where, of course, $`|A(\stackrel{}{a},\lambda )A(\stackrel{}{b},\lambda )|=1`$. However, stricly speaking the triangle inequality gives
$$|P(\stackrel{}{a},\stackrel{}{b})P(\stackrel{}{a},\stackrel{}{c})|𝑑\lambda |\rho (\lambda )|[A(\stackrel{}{b},\lambda )A(\stackrel{}{c},\lambda )1],$$
(A.6)
which is equal to (A.5) only when $`|\rho |=\rho `$, i.e., when $`\rho 0`$.
CHSH. The CHSH paper makes the same assumption at the identical point in their derivation, in their first (unnumbered) equation.
Peres. Peres’ derivation is almost identical to Eberhard’s and makes the same assumption of positive weights in the same step, i.e. between steps 1 and 2 of Eq. (2.7) of this paper.
Stapp 71. Stapp’s 1971 proof is very similar to Bell’s. He arrives at an expression (below his equation 8)
$$\sqrt{2}\frac{1}{N}\underset{j}{}|n_{2j}^{\prime \prime }n_{2j}^{}1|,$$
(A.7)
where $`n_{2j}^{\prime \prime }=\pm 1`$ and $`n_{2j}^{}=\pm 1`$. He then shows this leads to the contradiction $`\sqrt{2}1`$. However, if the $`n`$’s are $`\pm 1`$, then the summand can only have values 0,2. If $`N_1`$ and $`N_2`$ are the frequencies with which these two values occur, and $`N_1+N_2=N`$, then the right hand side can be written
$$\frac{1}{N}[N_1\times 0+N_2\times 2]=\frac{2N_2}{N}=\frac{2(NN_1)}{N}=2(1\frac{N_1}{N}).$$
(A.8)
As in the Eberhard argument, a contradiction can always be avoided by taking $`N_1`$ negative.
Stapp 85. Stapp establishes a contradiction by demonstrating (his Eq. 8) that
$$\frac{1}{n}\underset{i=1}{\overset{n}{}}\left[\sqrt{2}r_{Ai}(\widehat{\lambda }_a)+r_{Bi}(\widehat{\lambda }_a)+r_{Bi}(\widehat{\lambda }_b)\right]^2>(\sqrt{2}2)^2,$$
(A.9)
where $`r_{Ai}(\widehat{\lambda }_a)=\pm 1`$, $`r_{Bi}(\widehat{\lambda }_a)=\pm 1`$ and $`r_{Bi}(\widehat{\lambda }_b)=\pm 1`$ . However, since the $`r`$’s are $`\pm 1`$, the summand can have only one of three values: $`(\sqrt{2})^2`$, $`(2+\sqrt{2})^2`$ and $`(2\sqrt{2})^2`$. Then the above expression can be written as
$$\frac{1}{n}\left[n_1(\sqrt{2})^2+n_2(\sqrt{2}+2)^2+n_3(2\sqrt{2})^2\right],$$
(A.10)
where $`n_1,n_2,n_3`$ are the frequencies with which the three terms occur and $`n_1+n_2+n_3=n`$. Squaring out and combining terms yields
$$\frac{2(n_1+n_2+n_3)}{n}+\frac{2n_2(2+\sqrt{2})}{n}+\frac{2n_3(2\sqrt{2})}{n}.$$
(A.11)
Assuming $`n`$ and $`n_3`$ positive, this expression can become negative if
$$n_2<\frac{(n+n_3(2\sqrt{2}))}{2+\sqrt{2}},$$
(A.12)
in other words, if $`n_2`$ is sufficiently negative.
Bell 71. A proof that has been cited as qualitatively different than the others is Bell’s 1971 proof . This proof is basically the same as the CHSH proof. In Bell’s 1971 version the probability density is also explicitly taken to be positive definite. The only difference is that now $`|A(\stackrel{}{a},\lambda )|1`$ and $`|B(\stackrel{}{b},\lambda )|1`$. (In our notation this corresponds to $`|\alpha _i|1`$ and $`|\beta _i|1`$.) This change merely strengthens the upper bound on the classical correlations. That is, in our equation (2.5), whereas previously $`|\gamma |=2`$, now $`|\gamma |2`$. The rest of the derivation is consequently unaffected and the CHSH inequality continues to hold. Furthermore, our demonstration of the equivalence of the Eberhard procedure with the ”master probability distribution” procedure is also unaffected, since Eq. (4.2) made no assumption about the values of the $`P`$’s.
TABLES AND FIGURES
TABLE I. Spin combinations for standard Bell inequalities. Hidden-variable models assume that spin-1/2 particles can be emitted with $`\pm `$ spin along each of three axes, a, b and c. The notation ($`+++`$) etc., means spin up along all three axes. The eight possible spin combinations are shown. To ensure conservation of angular momentum, a particle of the type $`(+++)`$ must be paired with one of ($``$) and so on.
| Population | Particle 1 | Particle 2 |
| --- | --- | --- |
| $`N_1`$ | ($`+++`$) | ($``$) |
| $`N_2`$ | ($`++`$) | ($`+`$) |
| $`N_3`$ | ($`++`$) | ($`+`$) |
| $`N_4`$ | ($`++`$) | ($`+`$) |
| $`N_5`$ | ($`+`$) | ($`++`$) |
| $`N_6`$ | ($`+`$) | ($`++`$) |
| $`N_7`$ | ($`+`$) | ($`++`$) |
| $`N_8`$ | ($``$) | ($`+++`$) |
TABLE II. Four probabilities. Shown are the four-probabilities from symmetric wavefunction as computed from Eq. (3.10). The quantity
$`\mathrm{\Delta }\frac{1}{3}\left[(𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐})(𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐})+(𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏})(𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐})+(𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐})(𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐})\right]`$.
Note that these probabilities sum to one. The four-probabilities for the antisymmetric wave function can be obtained by flipping last two signs, i.e., $`P(++++)_{AS}=P(++)_s,P(+++)_{AS}=P(+)_s`$, etc.
| $`P(++++)=P()=\frac{1}{16}\{1+𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}+𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}+\mathrm{\Delta }\}`$ |
| --- |
| $`P(+++)=P(+)=\frac{1}{16}\{1𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}+𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}\mathrm{\Delta }\}`$ |
| $`P(+++)=P(+)=\frac{1}{16}\{1𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}+𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}\mathrm{\Delta }\}`$ |
| $`P(+++)=P(+)=\frac{1}{16}\{1+𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}\mathrm{\Delta }\}`$ |
| $`P(+++)=P(+)=\frac{1}{16}\{1+𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}\mathrm{\Delta }\}`$ |
| $`P(++)=P(++)=\frac{1}{16}\{1+𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}+𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}+\mathrm{\Delta }\}`$ |
| $`P(++)=P(++)=\frac{1}{16}\{1𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}+\mathrm{\Delta }\}`$ |
| $`P(++)=P(++)=\frac{1}{16}\{1𝐚_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}+𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}+𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}𝐛_\mathrm{𝟏}𝐛_\mathrm{𝟐}+\mathrm{\Delta }\}`$ |
TABLE III. Four probabilities as functions of polarizer angles. Shown are the same four-probabilities as on Table II for the configuration $`𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟏}=𝐚_\mathrm{𝟏}𝐛_\mathrm{𝟐}=𝐛_\mathrm{𝟏}𝐚_\mathrm{𝟐}=cos\theta `$ and $`𝐚_\mathrm{𝟐}𝐛_\mathrm{𝟐}=cos3\theta `$. Now $`\mathrm{\Delta }=1/3(cos^2\theta +cos^22\theta +cos\theta \mathrm{cos}3\theta )`$. With the identities $`cos2\theta =2cos^2\theta 1`$ and $`cos3\theta =4cos^3\theta 3cos\theta `$ all the probabilities can be written in terms of one parameter, $`cos\theta C`$. This form makes it more plausible that some of the $`P`$’s can become negative.
| $`P(++++)=P()=\frac{1}{16}\{1+3cos\theta +2cos2\theta +cos3\theta +\mathrm{\Delta }\}=\frac{1}{16}\{4C^3+4C^21+\mathrm{\Delta }\}`$ |
| --- |
| $`P(+++)=P(+)=\frac{1}{16}\{1cos\theta +cos3\theta \mathrm{\Delta }\}=\frac{1}{16}\{4C^34C+1\mathrm{\Delta }\}`$ |
| $`P(+++)=P(+)=\frac{1}{16}\{1+cos\theta cos3\theta \mathrm{\Delta }\}=\frac{1}{16}\{4C^3+4C+1\mathrm{\Delta }\}`$ |
| $`P(+++)=P(+)=\frac{1}{16}\{1cos\theta +cos3\theta \mathrm{\Delta }\}=\frac{1}{16}\{4C^34C+1\mathrm{\Delta }\}`$ |
| $`P(+++)=P(+)=\frac{1}{16}\{1+cos\theta cos3\theta \mathrm{\Delta }\}=\frac{1}{16}\{4C^3+4C+1\mathrm{\Delta }\}`$ |
| $`P(++)=P(++)=\frac{1}{16}\{1+2cos2\theta 3cos\theta cos3\theta +\mathrm{\Delta }\}=\frac{1}{16}\{4C^3+C^21+\mathrm{\Delta }\}`$ |
| $`P(++)=P(++)=\frac{1}{16}\{1cos\theta 2cos2\theta +cos3\theta +\mathrm{\Delta }\}=\frac{1}{16}\{4C^34C^24C+3+\mathrm{\Delta }\}`$ |
| $`P(++)=P(++)=\frac{1}{16}\{1+cos\theta 2cos2\theta cos3\theta +\mathrm{\Delta }\}=\frac{1}{16}\{4C^34C^2+4C+3+\mathrm{\Delta }\}`$ |
FIG. 1. Four-probabilities from Table III. (a) Plot of $`16P(+++)`$. (b) Plot of $`16P(++)`$. Note that these quantities become negative.
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# Untitled Document
On the Frobenius integrability of certain holomorphic $`p`$-forms
Jean-Pierre Demailly
Dedicated to Professor Hans Grauert, on the occasion of his 70th birthday
Abstract. The goal of this note is to exhibit the integrability properties (in the sense of the Frobenius theorem) of holomorphic $`p`$-forms with values in certain line bundles with semi-negative curvature on a compact Kähler manifold. There are in fact very strong restrictions, both on the holomorphic form and on the curvature of the semi-negative line bundle. In particular, these observations provide interesting information on the structure of projective manifolds which admit a contact structure: either they are Fano manifolds or, thanks to results of Kebekus-Peternell-Sommese-Wisniewski, they are biholomorphic to the projectivization of the cotangent bundle of another suitable projective manifold.
1. Main results
Recall that a holomorphic line bundle $`L`$ on a compact complex manifold is said to be pseudo-effective if $`c_1(L)`$ contains a closed positive $`(1,1)`$-current $`T`$, or equivalently, if $`L`$ possesses a (possibly singular) hermitian metric $`h`$ such that the curvature current $`T=\mathrm{\Theta }_h(L)=i\overline{}\mathrm{log}h`$ is nonnegative. If $`X`$ is projective, $`L`$ is pseudo-effective if and only if $`c_1(L)`$ belongs to the closed cone of $`H_{\mathrm{}}^{1,1}(X)`$ generated by classes of effective divisors (see \[Dem90, 92\]). Our main result is
Main Theorem. Let $`X`$ be a compact Kähler manifold. Assume that there exists a pseudo-effective line bundle $`L`$ on $`X`$ and a nonzero holomorphic section $`\theta H^0(X,\mathrm{\Omega }_X^pL^1)`$, where $`0pn=dimX`$. Let $`𝒮_\theta `$ be the coherent subsheaf of germs of vector fields $`\xi `$ in the tangent sheaf $`T_X`$, such that the contraction $`i_\xi \theta `$ vanishes. Then $`𝒮_\theta `$ is integrable, namely $`[𝒮_\theta ,𝒮_\theta ]𝒮_\theta `$, and $`L`$ has flat curvature along the leaves of the $`(`$possibly singular$`)`$ foliation defined by $`𝒮_\theta `$.
Before entering into the proof, we discuss several consequences. If $`p=0`$ or $`p=n`$, the result is trivial (with $`𝒮_\theta =T_X`$ and $`𝒮_\theta =0`$, respectively). The most interesting case is $`p=1`$.
Corollary 1. In the above situation, if the line bundle $`LX`$ is pseudo-effective and $`\theta H^0(X,\mathrm{\Omega }_X^1L^1)`$ is a nonzero section, the subsheaf $`𝒮_\theta `$ defines a holomorphic foliation of codimension $`1`$ in $`X`$, that is, $`\theta d\theta =0`$.
We now concentrate ourselves on the case when $`X`$ is a contact manifold, i.e. $`dimX=n=2m+1`$, $`m1`$, and there exists a form $`\theta H^0(X,\mathrm{\Omega }_X^1L^1)`$, called the contact form, such that $`\theta (d\theta )^mH^0(X,K_XL^{m1})`$ has no zeroes. Then $`𝒮_\theta `$ is a codimension $`1`$ locally free subsheaf of $`T_X`$ and there are dual exact sequences
$$0L\mathrm{\Omega }_X^1𝒮_\theta ^{}0,0𝒮_\theta T_XL^{}0.$$
The subsheaf $`𝒮_\theta T_X`$ is said to be the contact structure of $`X`$. The assumption that $`\theta (d\theta )^m`$ does not vanish implies that $`K_XL^{m+1}`$. In that case, the subsheaf is not integrable, hence $`L`$ and $`K_X`$ cannot be pseudo-effective.
Corollary 2. If $`X`$ is a compact Kähler manifold admitting a contact structure, then $`K_X`$ is not pseudo-effective, in particular the Kodaira dimension $`\kappa (X)`$ is equal to $`\mathrm{}`$.
The fact that $`\kappa (X)=\mathrm{}`$ had been observed previously by Stéphane Druel \[Dru98\]. In the projective context, the minimal model conjecture would imply (among many other things) that the conditions $`\kappa (X)=\mathrm{}`$ and “$`K_X`$ non pseudo-effective” are equivalent, but a priori the latter property is much stronger (and in large dimensions, the minimal model conjecture still seems far beyond reach!)
Corollary 3. If $`X`$ is a compact Kähler manifold with a contact structure and with second Betti number $`b_2=1`$, then $`K_X`$ is negative, i.e., $`X`$ is a Fano manifold.
Actually the Kodaira embedding theorem shows that the Kähler manifold $`X`$ is projective if $`b_2=1`$, and then every line bundle is either positive, flat or negative. As $`K_X`$ is not pseudo-effective it must therefore be negative. In that direction, Boothby \[Boo61\], Wolf \[Wol65\] and Beauville \[Bea98\] have exhibited a natural construction of contact Fano manifolds. Each of the known examples is obtained as a homogeneous variety which is the unique closed orbit in the projectivized (co)adjoint representation of a simple algebraic Lie group. Beauville’s work (\[Bea98\], \[Bea99\]) provides strong evidence that this is the complete classification in the case $`b_2=1`$.
We now come to the case $`b_22`$. If $`Y`$ is an arbitrary compact Kähler manifold, the bundle $`X=P(T_Y^{})`$ of hyperplanes of $`T_Y`$ has a contact structure associated with the line bundle $`L=𝒪_X(1)`$. Actually, if $`\pi :XY`$ is the canonical projection, one can define a contact form $`\theta H^0(X,\mathrm{\Omega }_X^1L^1)`$ by setting
$$\theta (x)=\theta (y,[\xi ])=\xi ^1\pi ^{}\xi =\xi ^1\underset{1jp}{}\xi _jdy_j,p=dimY,$$
at every point $`x=(y,[\xi ])X`$, $`\xi T_{Y,y}^{}\{0\}`$ (observe that $`\xi L_x=𝒪_X(1)_x`$). Morever $`b_2(X)=1+b_2(Y)2`$. Conversely, Kebekus, Peternell, Sommese and Wiśniewski \[KPSW\] have recently shown that every projective algebraic manifold $`X`$ such that
(i)$`X`$ has a contact structure,
(ii)$`b_22`$,
(iii)$`K_X`$ is not nef (numerically effective) is of the form $`X=P(T_Y^{})`$ for some projective algebraic manifold $`Y`$. However, the condition that $`K_X`$ is not nef is implied by the fact that $`K_X`$ is not pseudo-effective. Hence we get
Corollary 4. If $`X`$ is a contact projective manifold with $`b_22`$, then $`X`$ is a projectivized hyperplane bundle $`X=P(T_Y^{})`$ associated with some projective manifold $`Y`$.
The Kähler case of corollary 4 is still unsolved, as the proof of \[KPSW\] heavily relies on Mori theory (and, unfortunately, the extension of Mori theory to compact Kähler manifolds remains to be settled $`\mathrm{}`$).
I would like to thank Arnaud Beauville, Frédéric Campana, Stefan Kebekus and Thomas Peternell for illuminating discussions on these subjects. The present work was written during a visit at Göttingen University, on the occasion of a colloquium in honor of Professor Hans Grauert for his 70th birthday.
2. Proof of the Main Theorem
In some sense, the proof is just a straightforward integration by parts, but there are slight technical difficulties due to the fact that we have to work with singular metrics.
Let $`X`$ be a compact Kähler manifold, $`\omega `$ the Kähler metric, and let $`L`$ be a pseudo-effective line bundle on $`X`$. We select a hermitian metric $`h`$ on $`L`$ with nonnegative curvature current $`\mathrm{\Theta }_h(L)0`$, and let $`\phi `$ be the plurisubharmonic weight of the metric $`h`$ in any local trivialisation $`L_{|U}U\times \mathrm{}`$. In other words, we have
$$\xi _h^2=|\xi |^2e^{\phi (x)},\xi ^{}_h^{}^2=|\xi ^{}|^2e^{\phi (x)}$$
for all $`xU`$ and $`\xi L_x`$, $`\xi ^{}L^1`$. We then have a Chern connection $`=_h^{}+\overline{}`$ acting on all $`(p,q)`$-forms $`f`$ with values in $`L^1`$, given locally by
$$_\phi f=e^\phi (e^\phi f)=f+\phi f$$
in every trivialization $`L_{|U}`$. Now, assume that there is a holomorphic section $`\theta H^0(X,\mathrm{\Omega }_X^pL^1)`$, i.e., a $`\overline{}`$-closed $`(p,0)`$ form $`\theta `$ with values in $`L^1`$. We compute the global $`L^2`$ norm
$$_X\{_h^{}\theta ,_h^{}\theta \}_h^{}\omega ^{np1}=_Xe^\phi _\phi \theta \overline{_\phi \theta }\omega ^{np1}$$
where $`\{,\}_h^{}`$ is the natural sesquilinear pairing sending pairs of $`L^1`$-valued forms of type $`(p,q)`$, $`(r,s)`$ into $`(p+s,q+r)`$ complex valued forms. The right hand side is of course only locally defined, but it explains better how the forms are calculated, and also all local representatives glue together into a well defined global form; we will therefore use the latter notation as if it were global. As
$$d\left(e^\phi \theta \overline{_\phi \theta }\omega ^{np1}\right)=e^\phi _\phi \theta \overline{_\phi \theta }\omega ^{np1}+(1)^pe^\phi \theta \overline{\overline{}_\phi \theta }\omega ^{np1}$$
and $`\overline{}_\phi \theta =\overline{}\phi \theta `$, an integration by parts via Stokes theorem yields
$$_Xe^\phi _\phi \theta \overline{_\phi \theta }\omega ^{np1}=(1)^p_Xe^\phi \overline{}\phi \theta \overline{\theta }\omega ^{np1}.$$
These calculations need a word of explanation, since $`\phi `$ is in general singular. However, it is well known that the $`i\overline{}`$ of a plurisubharmonic function is a closed positive current, in particular
$$i\overline{}(e^\phi )=e^\phi (i\phi \overline{}\phi +i\overline{}\phi )$$
is positive and has measure coefficients. This shows that $`\phi `$ is $`L^2`$ with respect to the weight $`e^\phi `$, and similarly that $`e^\phi \overline{}\phi `$ has locally finite measure coefficients. Moreover, the results of \[Dem92\] imply that there is a decreasing sequence of metrics $`h_\nu ^{}`$ and corresponding weights $`\phi _\nu \phi `$, such that $`\mathrm{\Theta }_{h_\nu }C\omega `$ with a fixed constant $`C>0`$ (this claim is in fact much weaker than the results of \[Dem92\], and easy to prove e.g. by using convolutions in suitable coordinate patches and a standard gluing technique). Now, the results of Bedford-Taylor \[BT76, BT82\] applied to the uniformly bounded functions $`e^{c\phi _\nu }`$, $`c>0`$, imply that we have local weak convergence
$$e^{\phi _\nu }\overline{}\phi _\nu e^\phi \overline{}\phi ,e^{\phi _\nu }\phi _\nu e^\phi \phi ,e^{\phi _\nu }\phi _\nu \overline{}\phi _\nu e^\phi \phi \overline{}\phi ,$$
possibly afting adding $`C^{}|z|^2`$ to the $`\phi _\nu `$’s to make them plurisubharmonic. This is enough to justify the calculations. Now, we take care of signs, using the fact that $`i^{p^2}\theta \overline{\theta }0`$ whenever $`\theta `$ is a $`(p,0)`$-form. Our previous equality can be rewritten
$$_Xe^\phi i^{(p+1)^2}_\phi \theta \overline{_\phi \theta }\omega ^{np1}=_Xe^\phi i\overline{}\phi i^{p^2}\theta \overline{\theta }\omega ^{np1}.$$
Since the left hand side is nonnegative and the right hand side is nonpositive, we conclude that $`_\phi \theta =0`$ almost everywhere, i.e. $`\theta =\phi \theta `$ almost everywhere. The formula for the exterior derivative of a $`p`$-form reads
$$\begin{array}{ccc}\hfill d\theta (\xi _0,\mathrm{},\xi _p)& =\underset{0jp}{}(1)^j\xi _j\theta (\xi _0,\mathrm{},\widehat{\xi _j},\mathrm{},\xi _p)\hfill & \\ & +\underset{0j<kp}{}(1)^{j+k}\theta ([\xi _j,\xi _k],\xi _0,\mathrm{},\widehat{\xi _j},\mathrm{},\widehat{\xi _k},\mathrm{},\xi _p).\hfill & ()\hfill \end{array}$$
If two of the vector fields – say $`\xi _0`$ and $`\xi _1`$ – lie in $`𝒮_\theta `$, then
$$d\theta (\xi _0,\mathrm{},\xi _p)=(\phi \theta )(\xi _0,\mathrm{},\xi _p)=0$$
and all terms in the right hand side of $`()`$ are also zero, except perhaps the term $`\theta ([\xi _0,\xi _1],\xi _2,\mathrm{},\xi _p)`$. We infer that this term must vanish. Since this is true for arbitrary vector fields $`\xi _2,\mathrm{},\xi _p`$, we conclude that $`[\xi _0,\xi _1]𝒮_\theta `$ and that $`𝒮_\theta `$ is integrable.
The above arguments also yield strong restrictions on the hermitian metric $`h`$. In fact the equality $`\theta =\phi \theta `$ implies $`\overline{}\phi \theta =0`$ by taking the $`\overline{}`$. Fix a smooth point in a leaf of the foliation, and local coordinates $`(z_1,\mathrm{},z_n)`$ such that the leaves are given by $`z_1=c_1,\mathrm{},z_r=c_r`$ ($`c_i=`$constant), in a neighborhood of that point. Then $`𝒮_\theta `$ is generated by $`/z_{r+1},\mathrm{},/z_n`$, and $`\theta `$ depends only on $`dz_1,\mathrm{},dz_r`$. This implies that $`^2\phi /z_j\overline{z}_k=0`$ for $`j,k>r`$, in other words $`(L,h)`$ has flat curvature along the leaves of the foliation. The main theorem is proved.
References
\[Bea98\]Beauville, A.: Fano contact manifolds and nilpotent orbits; Comm. Math. Helv., 73 (4) (1998), 566–583.
\[Bea99\]Beauville, A.: Riemannian holonomy and Algebraic Geometry; Duke/alg-geom preprint 9902110, (1999).
\[Boo61\]Boothby, W.: Homogeneous complex contact manifolds; Proc. Symp. Pure Math. 3 (Differential Geometry) (1961) 144–154.
\[BT76\]Bedford, E., Taylor, B.A.: The Dirichlet problem for a complex Monge-Ampère equation; Invent. Math. 37 (1976) 1–44.
\[BT82\]Bedford, E., Taylor, B.A.: A new capacity for plurisubharmonic functions; Acta Math. 149 (1982) 1–41.
\[Dem90\]Demailly, J.-P.: Singular hermitian metrics on positive line bundles; Proceedings of the Bayreuth conference “Complex algebraic varieties”, April 2-6, 1990, edited by K. Hulek, T. Peternell, M. Schneider, F. Schreyer, Lecture Notes in Math. n$`^{}`$1507, Springer-Verlag, 1992.
\[Dem92\]Demailly, J.-P.: Regularization of closed positive currents and Intersection Theory; J. Alg. Geom. 1 (1992), 361–409.
\[Dru98\]Druel, S.: Contact structures on algebraic $`5`$-dimensional manifolds; C.R. Acad. Sci. Paris, 327 (1998), 365–368.
\[KPSW\]Kebekus, S., Peternell, Th., Sommese, A.J. and Wiśniewski, J.A.: Projective Contact Manifolds; preprint 1999, to appear in Inventiones Math.
\[Wol65\]Wolf, J.: Complex homogeneous contact manifolds and quaternionic symmetric spaces; J. Math. Mech. 14 (1965) 1033–1047.
(version of April 13, 2000, printed on )
Jean-Pierre Demailly Université de Grenoble I, Département de Mathématiques, Institut Fourier 38402 Saint-Martin d’Hères, France e-mail: demailly@ujf-grenoble.fr
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# Radiation Pressure as a Source of Decoherence
## I Introduction
Superposition states have an important role in the formalism of quantum mechanics. However, they are in flagrant contradiction with our classical world when the components correspond to macroscopically distinguishable states. The reason why these states are not encountered round the corner is decoherence, a process by which the interaction between the degrees of freedom of the system in question with any other degrees of freedom, either internal or external (the so-called environment), leads to a suppression of the coherence between the components of the superposition . Even if this coupling is very weak, the decoherence rate may be huge, resulting in a very fast decay of these “weird” states and in the emergence of the classical world. Recent developments in technology now allow one to study in real-time the process of decoherence in the lab. For example, over the past several years techniques have been developed to generate mesoscopic superpositions of motional states of trapped ions , and of photon states in cavity quantum electrodynamics . In these cases decoherence due to the coupling with the ambient reservoirs was observed, confirming the expectation that the decoherence rate is faster, the larger and more separated the state components are . Recently another experiment has succeeded in “engineering” the environment in the context of trapped ions, studying scaling laws of decoherence theory for a variety of reservoirs in a wide range of parameters .
Usually, decoherence is analyzed in the framework of heuristic models that describe phenomenological dissipation (the reservoir is taken to be a collection of harmonic oscillators, coupled linearly to the position operator of the system ). In this paper, we consider instead an ab initio model for decoherence of a particle in a harmonic potential, scattering the radiation field (at temperature $`T`$), which then plays the role of the reservoir. Starting from first principles, we show that the resulting radiation pressure coupling with the field transforms an initial pure superposition state of the particle into a statistical mixture.
Of special relevance is the limit $`T=0.`$ In this case, the reservoir is the quantum vacuum field, which dissipates the mechanical energy of the oscillating particle (or ‘mirror’). This effect is associated to the emission of pairs of photons, the so-called dynamical Casimir effect. Much work has been done on quantum radiation from moving mirrors . Important properties like the spectrum of the emitted radiation , the time evolution of the energy-momentum tensor , the total radiated intensity and the dissipative radiation pressure on the particle (radiation reaction force corresponding to the photon emission effect) have been considered. Here we focus our attention on the particle as the system of interest, and show that decoherence is a consequence of the entanglement between particle and field two-photon states. This result has fundamental implications, for it shows that any particle not transparent to the radiation is unavoidably under the action of decoherence through the radiation pressure coupling with vacuum fluctuations.
The zero temperature limit was briefly discussed in our previous letter . This article presents results for finite values of temperature, as well as a detailed discussion of the case $`T=0.`$ The formalism relies on the 1D scalar model for the field, but extensions to 3D electromagnetic field are also discussed, allowing us to address the question of orders of magnitude. The paper is organized as follows. In section II we start from the Hamiltonian model for the radiation pressure coupling, and then derive a master equation for the particle. In section III we discuss how the environment selects a prefered basis in the particle’s Hilbert space, the pointer basis. In Section IV we derive a general relation between decoherence and damping rates at arbitrary temperature by means of the fluctuation-dissipation theorem. The zero and high temperature limits, including extensions to the 3D electromagnetic model, are discussed in Sections V and VI. Section VII contains our conclusions. Finally, in the appendix an alternative, simpler derivation of the decoherence rate is given, which is based on the entanglement between the particle and two-photon states.
## II The model
Most treatments of the dynamical Casimir effect are based on the assumption that the mirror follows a prescribed trajectory, thus neglecting the recoil effect. However, in this paper we want to focus on the mirror as a dynamical quantum system, hence the need to tackle the full mirror-plus-field dynamics. This has already been addressed in the framework of linear response theory in order to calculate the fluctuations of the position of a dispersive mirror driven by vacuum radiation pressure , and related calculations have been performed in Refs. to derive mass corrections caused by the interaction with the field.
We consider a nonrelativistic partially reflecting mirror of mass $`M`$ (with position $`q`$ and momentum $`p`$) in a harmonic potential of frequency $`\omega _0`$, and under the action of radiation pressure. We take a scalar field in 1+1 dimensions, which mimics the electromagnetic field modes that propagate along the direction perpendicular to the plane of the mirror. Extensions to the real 3+1 case are analyzed in Secs. V and VI. We neglect third and higher order terms in $`v/c`$, where $`v`$ is the mirror’s velocity (we set $`c=1`$ hereafter, except when an explicit evaluation of orders of magnitude is required). Our starting point is the Hamiltonian formalism developed in Refs. and (although these references consider a free mirror, the extension to the harmonic oscillator is straightforward). The total Hamiltonian is
$$H=H_M+H_F+H_{\mathrm{int}},$$
(1)
where
$$H_M=\frac{p^2}{2M}+\frac{M\omega _0^2}{2}q^2,$$
(2)
is the harmonic oscillator Hamiltonian for the mirror, and
$$H_F=\frac{dx}{2}\left[\mathrm{\Pi }^2+(_x\varphi )^2\right]+\mathrm{\Omega }\varphi ^2(x=0)$$
(3)
is the free Hamiltonian for the field $`\varphi `$ and its momentum canonically conjugated $`\mathrm{\Pi }=_t\varphi `$. The second term in the r.-h.-s. of Eq. (3) is associated to the boundary condition of a partially-reflecting mirror at rest at $`x=0.`$ In the context of the plasma sheet model of Ref. , it corresponds to the kinetic energy of the plasma charged particles. The coupling constant $`\mathrm{\Omega }`$ plays the role of a transparency frequency, since from Eq. (3) one derives the boundary condition
$`_x\varphi (0^+)_x\varphi (0^{})=2\mathrm{\Omega }\varphi (0)`$
($`\varphi `$ is continuous at $`x=0`$), which yields a frequency-dependent reflection amplitude :
$$R(\omega )=i\frac{\mathrm{\Omega }}{\omega +i\mathrm{\Omega }}.$$
(4)
Finally, the interaction Hamiltonian is given by
$$H_{\mathrm{int}}=\frac{p𝒫}{M}+\frac{𝒫^2}{2M}\frac{1}{2}\mathrm{\Omega }\varphi ^2(x=0)\frac{p^2}{M^2},$$
(5)
where $`𝒫=𝑑x_x\varphi _t\varphi `$ is the field momentum operator. $`H_{\mathrm{int}}`$ describes, to second order in $`v/c,`$ the modification of the boundary condition for the field due to the motion of the mirror, which in its turn is affected by the field radiation pressure. Thus, it provides a coupling between the harmonic oscillator and the field, to be treated within perturbation theory. The small perturbation parameter is $`v/c,`$ and not the transparency frequency $`\mathrm{\Omega },`$ which may be arbitrarily large. The first term in Eq. (5) is responsible for the effect of decoherence to be discussed here. It also accounts for the effects of emission of photons, dissipation of the mirror’s energy, and part of the mass correction.
We calculate the density matrix $`\stackrel{~}{\rho }(t)`$ of the combined mirror-plus-field system using second order perturbation theory, and trace over the field operators to derive the master equation for the mirror’s density matrix $`\rho (t)`$ . We assume that at $`t=0`$ the mirror and field are not correlated: $`\stackrel{~}{\rho }(0)=\rho (0)\rho _F,`$ where $`\rho _F`$ is the density matrix of the field (assumed to be in some steady state; later in this section we take a thermal equilibrium state). We find
$$i\mathrm{}\dot{\rho }(t)=[H_M,\rho (t)]\mathrm{\Omega }\frac{\varphi ^2(0)}{2M^2}[p^2,\rho (t)]$$
(6)
$`{\displaystyle \frac{i}{2\mathrm{}M^2}}{\displaystyle _0^t}𝑑t^{}\left([p,[p^I(t^{}),\rho (t)]]\sigma (t^{})+[p,\{p^I(t^{}),\rho (t)\}]\xi (t^{})\right),`$
where the superscript $`I`$ indicates the operators to be taken in the interaction picture. The second term in the r.-h.-s. of Eq. (6) is the contribution in first-order of perturbation theory of the $`p^2`$ term in the interaction Hamiltonian (see Eq. (5)). It corresponds to a (cut–off dependent) mass correction given by
$$\mathrm{\Delta }M_1=\mathrm{\Omega }\varphi ^2(0),$$
(7)
as already found in Refs. and . The (anti-)symmetric second order correlation function ($`\xi `$) $`\sigma `$ is defined as
$$\sigma (t)=C(t)+C(t),$$
(8)
$$\xi (t)=C(t)C(t),$$
(9)
with
$$C(t)=𝒫^I(t)𝒫^I(0)𝒫^2.$$
(10)
When computing the correlation functions, we take the unperturbed field, which corresponds to the static boundary condition (eigenfunctions of $`H_F`$).
Replacing the free evolution for $`p^I(t^{})`$ in (6) yields
$$i\mathrm{}\dot{\rho }=[H_M\frac{\mathrm{\Delta }M(t)}{M}\frac{p^2}{2M},\rho ]\mathrm{\Gamma }(t)[p,\{q,\rho \}]\frac{i}{\mathrm{}}D_1(t)[p,[p,\rho ]]\frac{i}{\mathrm{}}D_2(t)[p,[q,\rho ]].$$
(11)
The total mass correction is $`\mathrm{\Delta }M=\mathrm{\Delta }M_1+\mathrm{\Delta }M_2,`$ where $`\mathrm{\Delta }M_2,`$ as well as the remaining coefficients in (11), originate from the first term in the r.-h.-s. of Eq. (5), taken in second-order perturbation theory. Their meanings are best understood when writing the Fokker-Planck equation for the Wigner function $`W(x,p,t):`$
$$_tW=(1\mathrm{\Delta }M/M)\frac{p}{M}_xW+M\omega _0^2x_pW+2\mathrm{\Gamma }_x(xW)+D_1\frac{^2}{x^2}WD_2\frac{^2}{xp}W.$$
(12)
$`\mathrm{\Delta }M_2`$ and the damping coefficient $`\mathrm{\Gamma }`$ are calculated from the anti-symmetric correlation function:
$$\mathrm{\Delta }M_2(t)=\frac{i}{\mathrm{}}_0^t𝑑t^{}\mathrm{cos}(\omega _0t^{})\xi (t^{}),$$
(13)
$$\mathrm{\Gamma }(t)=\frac{i\omega _0}{2M\mathrm{}}_0^t𝑑t^{}\mathrm{sin}(\omega _0t^{})\xi (t^{});$$
(14)
whereas the diffusion coefficients are associated to the symmetric correlation function:
$$D_1(t)=\frac{1}{2M^2}_0^t𝑑t^{}\mathrm{cos}(\omega _0t^{})\sigma (t^{}),$$
(15)
$$D_2(t)=\frac{\omega _0}{2M}_0^t𝑑t^{}\mathrm{sin}(\omega _0t^{})\sigma (t^{}).$$
(16)
We assume the field to be in a thermal state (temperature $`T`$), and take the following strategy to calculate the momentum correlation functions. The time derivative of the field momentum is minus the radiation pressure force on the mirror :
$$\frac{d𝒫^I}{dt}=2\mathrm{\Omega }\varphi (0,t)\overline{}_x\varphi (0,t),$$
(17)
where $`\overline{}_x\varphi (0,t)=[_x\varphi (0^+)+_x\varphi (0^{})]/2.`$ Using Eq. (17), we calculate $`C(t)`$ by integrating the correlation function of the field calculated at $`x=0:`$
$`C(t)`$ $`=`$ $`(2\mathrm{\Omega })^2{\displaystyle _{\mathrm{}}^t}dt_1{\displaystyle _{\mathrm{}}^0}dt_2[\varphi (0,t_1)\overline{}_x\varphi (0,t_1)\varphi (0,t_2)\overline{}_x\varphi (0,t_2)`$ (19)
$`\varphi (0,t_1)\overline{}_x\varphi (0,t_1)\varphi (0,t_2)\overline{}_x\varphi (0,t_2)]`$
The equal-time second-order correlation function in (19) corresponds to the force on the static (single) mirror. It vanishes since the vacuum radiation pressures exerted on each side of the mirror are in perfect equilibrium. On the other hand, the fourth-order correlation function may be expressed as a sum of second order correlation functions (with the fields taken at different times), which are calculated with the help of the normal mode expansion for the field. They are directly connected to the average number of photons $`n_\omega =1/[\mathrm{exp}(\mathrm{}\omega /T)1]`$ in the mode of frequency $`\omega `$ at temperature $`T`$ (we take the Boltzmann constant $`k_\mathrm{B}=1`$). It is useful to write the result in the Fourier domain, the Fourier transform of the anti-symmetric correlation function $`\xi (t)`$ being defined as
$$\xi [\omega ]=𝑑t\mathrm{exp}(i\omega t)\xi (t).$$
(20)
Eqs. (9) and (19) yield
$$\xi [\omega ]=\xi ^0[\omega ]+\xi ^T[\omega ],$$
(21)
where
$$\xi ^0[\omega ]=(2/\pi )\mathrm{}^2\mathrm{\Omega }\zeta (\omega /\mathrm{\Omega }),$$
(22)
with
$$\zeta (u)=\mathrm{ln}(1+u^2)/(2u)+(\mathrm{arctan}u)/u^21/u,$$
(23)
represents the correlation function at $`T=0`$ (vacuum fluctuations), whereas
$$\xi ^T[\omega ]=\frac{2\mathrm{}^2\mathrm{\Omega }^2}{\pi \omega ^2}_0^{\mathrm{}}𝑑\omega ^{}\frac{\omega ^{}}{\omega ^{}{}_{}{}^{2}+\mathrm{\Omega }^2}\left[G(\omega ,\omega ^{})G(\omega ,\omega ^{})\right],$$
(24)
with
$$G(\omega ,\omega ^{})=|\omega ^{}\omega |\left(n_{|\omega ^{}\omega |}ϵ(\omega ^{}\omega )n_\omega ^{}\right),$$
(25)
represents the thermal fluctuations ($`ϵ`$ is the sign function).
Symmetric and anti-symmetric correlation functions for a system in thermal equilibrium are related in a very general way :
$$\sigma [\omega ]=\frac{\xi [\omega ]}{\mathrm{tanh}(\frac{\mathrm{}\omega }{2T})}.$$
(26)
According to Eqs. (14) and (15), this result provides a general relation between diffusion and damping, in the spirit of the fluctuation-dissipation theorem. This relation is particularly simple for the asymptotic values of the coefficients $`\mathrm{\Gamma }(t)`$ and $`D_1(t)`$ at $`t\mathrm{}.`$ Since the integrands in Eqs. (14) and (15) are even functions of time, we may extend the integration range to $`\mathrm{},`$ yielding
$$\mathrm{\Gamma }=\frac{\omega _0}{4M\mathrm{}}\xi [\omega _0]$$
(27)
and
$$D_1=\frac{1}{4M^2}\sigma [\omega _0].$$
(28)
Thus, the asymptotic values of $`\mathrm{\Gamma }`$ and $`D_1`$ are directly connected to the fluctuations at the mechanical frequency $`\omega _0,`$ allowing us to derive, from Eq. (26), a simple and general relation between these two coefficients. On the other hand, no such simple connection exists for the remaining time-dependent coefficients, $`\mathrm{\Delta }M_2`$ and $`D_2,`$ whose asymptotic values result from the joint contribution of the whole spectrum of fluctuations .
Combining Eqs. (26)–(28), we find
$$D_1=\frac{\mathrm{}}{M\omega _0}\frac{\mathrm{\Gamma }}{\mathrm{tanh}(\frac{\mathrm{}\omega _0}{2T})},$$
(29)
a clear manifestation of the fluctuation-dissipation theorem. According to Eq. (29), the temperature dependence of the diffusion coefficient is determined, apart from the $`T`$ dependence of the damping coefficient $`\mathrm{\Gamma }`$ (to be discussed later), by the relative importance of the thermal fluctuations (and their corresponding energy $`T`$) with respect to quantum fluctuations (and their corresponding zero-point energy $`\mathrm{}\omega _0/2`$). In the high-temperature limit, $`T\mathrm{}\omega _0/2,`$ Eq. (29) yields $`D_1=2T\mathrm{\Gamma }/(M\omega _0^2).`$ In the theory of Brownian motion, $`\mathrm{\Gamma }`$ is a $`T`$ independent phenomenological constant, and hence the diffusion coefficient is a linear function of temperature in this limit. Here, however, $`\mathrm{\Gamma }`$ has an explicit temperature dependence, to be analyzed in Sec. V.
From Eq. (29), we shall derive a relation between decoherence and damping time scales, valid for any temperature $`T.`$ Before considering a specific superposition state, however, we discuss, in the next section, the degree of sensitivity of different states in the Hilbert space to the action of decoherence. We also analyze in more detail the precise meaning of $`t\mathrm{}`$ (in the particular case of $`T=0`$), in order to know how fast the coefficients approach their asymptotic values. From Eqs. (14) and (15) alone it may be shown, in a general way, that a sufficient condition is $`t1/\omega _0,`$ but in some cases the convergence may be much faster.
## III Pointer states
Different criteria have been introduced in the literature in order to find out the states in the Hilbert space that are most robust under the interaction with the environment and behave more classically . Here we shall follow the one introduced by Zurek, the so-called “predictability sieve”. The idea is to take every possible state of the Hilbert space, calculate its entropy, and order the states in a tower according to increasing entropy. The most classical states are those that lie at the bottom of that tower, and correspond to the most predictable ones. For these ‘pointer’ states, information loss due to the interaction with the environment is minimal. This philosophy is put in quantitative terms by minimizing the linear entropy of the system,
$$s[\rho ]1\mathrm{tr}\rho ^2,$$
(30)
which is zero for a pure state and greater than zero for a statistical mixture. In general this is a difficult problem because complicated entanglement between system and environment develops on account of their mutual interaction. So far, results have been successfully derived assuming that the initial state of the system is pure. Here we follow the same approach, and calculate the rate of entropy increase starting from the master equation (11). We assume that the state is nearly pure at time $`t`$ to find
$$\dot{s}(t)=2\mathrm{\Gamma }(t)(s(t)1)+\frac{4D_1(t)}{\mathrm{}^2}\mathrm{\Delta }p^2+\frac{2D_2(t)}{\mathrm{}^2}\sigma _{q,p}$$
(31)
where $`\mathrm{\Delta }p^2p^2p^2`$ is the momentum dispersion and $`\sigma _{q,p}\{q,p\}2pq`$. Here $`\mathrm{}=\mathrm{tr}(\mathrm{}\rho )`$, and all operators are evaluated at the same time $`t`$. The first term in (31) leads to a decrease of entropy $`s(t)=1\mathrm{exp}(2_0^t\mathrm{\Gamma }(t^{})𝑑t^{})`$, hence damping tries to localize the state competing against diffusion. This decrease is independent of the initial state of the system, and therefore is irrelevant for determining the pointer states.
We assume that the typical decoherence time scale is much larger than the period of free oscillation $`2\pi /\omega _0,`$ so that we may integrate Eq. (31) to find the entropy at an intermediate time $`\tau =n\mathrm{\hspace{0.17em}2}\pi /\omega _0.`$ We take $`n1,`$ allowing us to replace the time dependent coefficients by their constant asymptotic values, but assume that $`\tau `$ is much shorter than the decoherence time scale, in order to be consistent with the small-entropy approximation underlying Eq. (31). Moreover, in this weak coupling limit, we may take the free evolution (corresponding to the harmonic oscillator Hamiltonian $`H_M`$) for the mirror’s operators in Eq. (31). The correlation function $`\sigma _{q,p}`$ oscillates around zero, and then does not contribute to $`s(\tau ),`$ whereas the free evolution of $`\mathrm{\Delta }p^2(t)`$ mixes up position and momentum fluctuations, yielding
$$s(\tau )=2\tau \frac{D_1}{\mathrm{}^2}\left[(\mathrm{\Delta }p)_0^2+(M\omega _0)^2(\mathrm{\Delta }q)_0^2M\mathrm{}\omega _0\right],$$
(32)
where $`(\mathrm{\Delta }p)_0^2`$ and $`(\mathrm{\Delta }q)_0^2`$ represent the dispersions for the initial state. From Eq. (32), we find that the minimum entropy given the constraint $`\mathrm{\Delta }q\mathrm{\Delta }p\mathrm{}/2`$ is for $`\mathrm{\Delta }q{}_{}{}^{2}=\mathrm{}/(2M\omega _0),`$ $`\mathrm{\Delta }p{}_{}{}^{2}=M\mathrm{}\omega _0/2.`$ Hence, as in the Caldeira-Leggett model, and for any temperature of the field, the pointer basis consists of coherent states . In this weak coupling approximation, the minimum value corresponds to $`s(\tau )=0,`$ hence the increase of entropy of a coherent state is a higher order effect.
The crucial approximation in the derivation of (32) from (31) is the replacement of the time dependent coefficients by their finite, constant asymptotic values. It is instructive to analyze in more detail the behavior of the coefficients and its connection with entropy production. As an example, we take $`T=0,`$ and consider first the perfectly-reflecting limit, which corresponds to $`\omega _0\mathrm{\Omega },`$ for in this case the relevant field modes have frequencies much smaller than the mirror’s transparency frequency. In Fig. 1 we plot the diffusion and damping coefficients as functions of $`\omega _0t`$ for $`\omega _0/\mathrm{\Omega }=10^4`$ and $`T=0.`$ The damping coefficient $`\mathrm{\Gamma }`$ approaches its asymptotic value very fast, for $`t1/\mathrm{\Omega },`$ whereas $`D_1(t)`$ develops an initial jolt for times of the order of $`\mathrm{\Omega }^1`$ and then decreases to the asymptotic value $`(D_1)_{\mathrm{perf}}=\mathrm{}^2\omega _0/(12\pi M^2)`$ for $`t1/\omega _0.`$ When we integrate Eq. (31) over many periods of oscillation, the contribution to the entropy of the initial jolt is negligible, allowing us to replace $`D_1`$ by its asymptotic value.
In Ref. , it was shown that no net entropy is produced for the Caldeira–Leggett model with an adiabatic environment, since all the time-dependent coefficients are oscillating functions around a zero mean. At first sight, the same would happen in our model when $`\omega _0\mathrm{\Omega },`$ for in this case the dominant field frequencies are slow with respect to the mirror’s translational time scale. However, as discussed in Section V, the spectral density $`\xi ^0(\omega )`$ decays too slowly for $`\omega \mathrm{\Omega },`$ and as a consequence field frequencies of the order of $`\omega _0`$ provide a significant contribution even in this limit. Thus, one cannot ascribe a frequency cut-off to the environment such that the typical frequency of the system $`\omega _0`$ is much greater than the maximum frequency of the environment. Therefore, the vacuum field does not behave adiabatically in the sense of . In our case instead, the diffusion coefficients oscillate around a non-zero value, leading to a net entropy increase. In Fig. 2, we plot the diffusion and damping coefficients as functions of $`\omega _0t`$ for $`\omega _0/\mathrm{\Omega }=10^4`$ and $`T=0.`$ They oscillate around their asymptotic values with (angular) frequency $`\omega _0`$ and with an amplitude of oscillation that decays in a time $`t1/\mathrm{\Omega }1/\omega _0.`$ The oscillatory terms do not contribute to the entropy increase when we average over many oscillations. Hence Eq. (32) also holds in this case, although the rate of entropy increase is much smaller than in the perfectly-reflecting limit.
## IV Decoherence versus damping
In this section, we derive a general relation between damping and decoherence time scales, starting from the fluctuation-dissipation result given by Eq. (29). As an extreme case of decoherent dynamics, we consider a superposition of two coherent states, since they correspond to the pointer states, according to the results of Sec. III. Specifically, we take at $`t=0`$ the even superposition state $`|\psi _\mathrm{e}=(|\alpha +|\alpha )/\sqrt{2},`$ with $`\alpha =iP_0/\sqrt{2M\mathrm{}\omega _0},`$ so that the coherent states are initially along the momentum axis in phase space, and $`\pm P_0`$ are the average values of momentum of the components at $`t=0.`$ We also assume that $`|\alpha |1,`$ hence the average energy of the state components is much larger than the zero-point energy. The corresponding Wigner function is
$$W=W_m+\frac{1}{\pi \mathrm{}}\mathrm{exp}\left[\frac{q^2}{2(\mathrm{\Delta }q_0)^2}\frac{p^2}{2(\mathrm{\Delta }p_0)^2}\right]\mathrm{cos}(\frac{2P_0q}{\mathrm{}}),$$
(33)
where $`\mathrm{\Delta }q_0=\sqrt{\mathrm{}/(2M\omega _0)}`$ and $`\mathrm{\Delta }p_0=\mathrm{}/(2\mathrm{\Delta }q_0)`$ are the position and momentum uncertainties of the ground state. $`W_m`$ corresponds to the statistical mixture
$$\rho _m=(|\alpha \alpha |+|\alpha \alpha |)/2.$$
(34)
In phase space, $`W_m`$ has two Gaussian peaks along the momentum axis at $`\pm P_0`$. $`\rho _m`$ is a classical state in the sense that $`W_m,`$ being positive defined, may be interpreted as a probability distribution in phase space. On the other hand, the nonclassical nature of the superposition state is featured by the remaining term in Eq. (33), representing the coherent interference between the two state components, and which oscillates into negative values along the position axis.
Diffusion along position, associated to the coefficient $`D_1,`$ averages out the oscillations of the interference term at a rate $`1/t_\mathrm{d},`$ to be calculated from the Fokker-Planck equation (12. According to Eq. (33), the oscillations are faster the higher the value of $`P_0,`$ so that $`t_\mathrm{d}`$ is a decreasing function of $`|\alpha |.`$ As in Section III, we assume that decoherence is very slow, $`1/t_\mathrm{d}\omega _0,`$ so that several free oscillations take place before coherence is lost. In this limit, the particle has enough time to probe the harmonic potential before diffusion takes place, and as a consequence decoherence is governed by the asymptotic value of $`D_1,`$ which is directly connected to the field flucutuations at the frequency of oscillation $`\omega _0,`$ according to Eq. (28). This condition holds for most experiments, where mesoscopic superpositions are employed so as to render decoherence slow enough to be measured . Moreover, it always applies in the case of vacuum radiation pressure ($`T=0`$), as shown in Sec. V. Diffusion is maximum when the state components are along the momentum axis: from (33), we find $`_q^2W(2P_0/\mathrm{})^2W;`$ and vanishes when the two wavepackets reach the turning points in the harmonic potential. The average over many oscillation yields
$$\frac{1}{t_\mathrm{d}}=\frac{1}{2}D_1\left(\frac{_q^2W}{W}\right)_{\mathrm{max}}=\frac{2P_0^2D_1}{\mathrm{}^2},$$
(35)
that combined with the fluctuation-dissipation theorem, Eq. (29), yields the following result for the decoherence time $`t_\mathrm{d}:`$
$$t_\mathrm{d}=\frac{1}{4|\alpha |^2}\mathrm{tanh}(\frac{\mathrm{}\omega _0}{2T})\frac{1}{\mathrm{\Gamma }}.$$
(36)
A $`T=0`$ (or more generally, for $`T\mathrm{}\omega _0`$), Eq. (36) yields $`t_d=1/(4|\alpha |^2\mathrm{\Gamma }).`$ This result may be written in terms of the distance $`\mathrm{\Delta }P=2P_0`$ between the two components in phase space at $`t=0,`$ or in terms of the distance $`\mathrm{\Delta }Q=\mathrm{\Delta }P/(M\omega _0)`$ at $`t=\pi /2\omega _0:`$
$$t_d=4\left(\frac{\mathrm{\Delta }p_0}{\mathrm{\Delta }P}\right)^2\frac{1}{\mathrm{\Gamma }}=4\left(\frac{\mathrm{\Delta }q_0}{\mathrm{\Delta }Q}\right)^2\frac{1}{\mathrm{\Gamma }}.$$
(37)
The interpretation of (37) is clear: decoherence is faster, the more separated the state components in phase space are. Here the zero-point fluctuations define the reference of distance in phase space. At high temperatures, on the other hand, this reference is provided by the thermal de Broglie wavelength $`\lambda _T=\mathrm{}/\sqrt{2MT}.`$ In fact, (36) yields, for $`T\mathrm{}\omega _0,`$
$$t_d=\frac{\mathrm{}\omega _0}{2T}\frac{1}{4|\alpha |^2}\frac{1}{\mathrm{\Gamma }}=2\left(\frac{\lambda _T}{\mathrm{\Delta }Q}\right)^2\frac{1}{\mathrm{\Gamma }}.$$
(38)
Eq. (38) also shows that the ratio between decoherence and damping rates is larger at high temperatures by the factor $`2T/(\mathrm{}\omega _0).`$
When written in terms of distances in phase space, the results above are also valid for more general superposition states, like $`(|0+|\alpha )/\sqrt{2}.`$ Moreover, their range of validity is not limited to the radiation pressure coupling considered here. In fact, Eqs. (37) and (38) are in perfect agreement with the results obtained in the framework of the Caldeira-Legget phenomenological model for quantum dissipation . This is hardly surprising, since they rely on general properties of the correlation functions associated to the fluctuation-dissipation theorem. Eq. (36), which interpolates the low and high temperature limits, is also discussed in Ref. , in the context of a two-level system. The temperature dependence for the ratio between decoherence and damping times has a simple interpretation: at $`T>0,`$ the time scale for the relaxation of the populations is shorter than $`1/\mathrm{\Gamma }`$ exactly by the factor $`\mathrm{tanh}(\mathrm{}\omega _0/2T),`$ on account of the contribution of absorption and stimulated emission. Here this factor originates from the general relation between symmetric and anti-symmetric correlation functions, Eq. (26), which is at the heart of the fluctuation-dissipation theorem.
The peculiarities of the radiation pressure model considered here are contained in the damping rate $`\mathrm{\Gamma }.`$ Rather than a phenomenological input parameter, it is computed from first principles, first for $`T=0`$ in Sec. V, and then for $`T\mathrm{}\omega _0`$ in Sec. VI.
## V Vacuum field
At $`T=0,`$ the spectral density is given by Eqs. (22) and (23). This result is more easily obtained from the following argument (a similar method, applied for the force correlation function, may be found in Refs. and ). Since $`𝒫`$ is quadratic in the field operators, the correlation function $`C(t)`$ may be calculated from the two-photon matrix elements of the momentum operator as follows:
$$C(t)=\frac{1}{2}_0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _20|𝒫(t)|\omega _1,\omega _2\omega _1,\omega _2|𝒫|0.$$
(39)
We have $`0|𝒫(t)|\omega _1,\omega _2=\mathrm{exp}[i(\omega _1+\omega _2)t]0|𝒫(0)|\omega _1,\omega _2`$ since only the annihilation operators contribute, and hence
$$C[\omega ]=\pi _0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _2|0|𝒫(0)|\omega _1,\omega _2|^2\delta (\omega \omega _1\omega _2).$$
(40)
Thus, at $`T=0`$ the fluctuations at the (positive) frequency $`\omega `$ originate from two-photon states $`|\omega _1,\omega _2`$ such that $`\omega _1+\omega _2=\omega .`$ In the dynamical Casimir effect, the oscillation at the mechanical frequency $`\omega _0`$ gives rise to the emission of pairs of photons of frequencies $`\omega _1`$ and $`\omega _2,`$ such that $`\omega _0=\omega _1+\omega _2.`$ On the other hand, according to (27), $`\mathrm{\Gamma }`$ originates from the fluctuations at frequency $`\omega _0,`$ and hence
$$\mathrm{\Gamma }=\frac{\pi }{4}\frac{\omega _0}{M\mathrm{}}_0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _2|\omega _1,\omega _2|𝒫|0|^2\delta (\omega _0\omega _1\omega _2),$$
(41)
rendering explicit the connection between damping and the photon emission effect. In the appendix, we present an alternative derivation of (41), starting from the two-photon emission probabilities and making use of energy conservation.
Eq. (40) also shows that $`C[\omega ]`$ vanishes for negative frequencies, and as a consequence, $`\sigma [\omega ]=ϵ(\omega )\xi [\omega ]`$ in agreement with Eq. (26). Finally, the result of Eq. (22) follows from (40) by using again Eq. (17. In Fig. 3 , we plot $`\zeta (\omega /\mathrm{\Omega })`$ as a function of its argument. According to Eq. (22), the transparency frequency $`\mathrm{\Omega }`$ defines a scale for the behavior of the spectrum of fluctuations in vacuum. For $`\omega \mathrm{\Omega },`$ the spectrum is linear: $`\zeta (\omega /\mathrm{\Omega })\omega /(6\mathrm{\Omega }),`$ and goes to zero slowly, as $`\mathrm{ln}(\omega /\mathrm{\Omega })/(\omega /\mathrm{\Omega }),`$ for $`\omega \mathrm{\Omega },`$ due to the high-frequency transparency of the mirror.
The damping coefficient at zero temperature is obtained from Eqs. (22) and (27), or alternatively from Eq. (41):
$$\mathrm{\Gamma }=\frac{\mathrm{}\mathrm{\Omega }\omega _0}{2\pi M}\zeta (\frac{\omega _0}{\mathrm{\Omega }}).$$
(42)
In the perfectly-reflecting limit, $`\omega _0\mathrm{\Omega },`$ Eq. (42) yields
$$\mathrm{\Gamma }=\frac{\mathrm{}\omega _0}{12\pi M}\omega _0.$$
(43)
Thus, the damping induced by the Casimir effect is a small perturbation of the free harmonic oscillations. The ratio between the zero-point energy and the rest energy appearing in (43) is also of the order of the recoil velocity of the mirror divided by $`c,`$ which is, as explained in Sec. II, the small parameter of the perturbation approach leading to the master equation (11). For larger values of $`\omega _0/\mathrm{\Omega },`$ the damping as given by (42) is still smaller, since vacuum frequencies of the order of $`\omega _0`$ are not well reflected by the mirror in this case.
Eq. (43) is directly connected to the well-known formula for the dissipative Casimir force on a single perfect moving mirror , $`F=\mathrm{}x^{^{\prime \prime \prime }}/(6\pi ),`$ for the equation of motion then reads
$$x^{^{\prime \prime }}=\omega _0^2x+\frac{\mathrm{}x^{^{\prime \prime \prime }}}{6\pi M},$$
(44)
whose solution when $`\mathrm{}\omega _0/M1`$ is
$`x=x_0e^{i\omega _0t}\mathrm{exp}\left({\displaystyle \frac{\mathrm{}\omega _0^2}{12\pi M}}t\right).`$
The decoherence time scale at $`T=0`$ in the perfectly-reflecting limit is derived from Eqs. (37) and (43):
$$t_\mathrm{d}=\frac{3}{v^2}\frac{2\pi }{\omega _0},$$
(45)
where $`v=P_0/M`$ is the initial velocity of the wavepackets. Being of the order of $`(v/c)^2,`$ the decoherence rate is very small at $`T=0`$ (or, at any rate, in the nonrelativistic limit considered here). Since $`\omega _0t_\mathrm{d}1,`$ decoherence is the cumulative effect of several free oscillations in the harmonic well, which justifies the approach employed in the derivation of (35) and the use of the asymptotic value for $`D_1(t).`$
In order to further understand how the dynamical Casimir effect engenders decoherence, we present, in the appendix, an alternative approach, where we follow the evolution of the complete oscillator-plus-field quantum state. It shows that the superposition state decoheres because the two wavepacket components oscillating out-of-phase yield amplitudes for emission of photon pairs with opposite signs. As a consequence, an entangled mirror-plus-field state is developed, given by
$$|\mathrm{\Psi }_{\mathrm{\Delta }t}=B(\mathrm{\Delta }t)|\psi _\mathrm{e}|0+\frac{1}{2}_0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _2b(\omega _1,\omega _2;\mathrm{\Delta }t)|\psi _\mathrm{o}|\omega _1,\omega _2,$$
(46)
where $`|\psi _\mathrm{o}=(|\alpha |\alpha )/\sqrt{2}`$ is the odd superposition state, $`b(\omega _1,\omega _2;\mathrm{\Delta }t)`$ is the amplitude for emission of a photon pair with frequencies $`\omega _1`$ and $`\omega _2`$ during $`\mathrm{\Delta }t`$ (the explicit expressions are given in the appendix), and $`B`$ is the amplitude for persistence in the vacuum state:
$$|B(\mathrm{\Delta }t)|^2=1\frac{1}{2}_0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _2|b(\omega _1,\omega _2,\mathrm{\Delta }t)|^2.$$
(47)
The density operators of the odd and even superposition states differ by the sign of the interference term, $`\rho _{\mathrm{int}}=\rho \rho _m`$ ($`\rho _m`$ is defined in Eq. (34)). Accordingly, when computing the reduced density matrix of the mirror, $`\rho (\mathrm{\Delta }t)=\mathrm{tr}_{\mathrm{field}}(|\mathrm{\Psi }_{\mathrm{\Delta }t}\mathrm{\Psi }|),`$ the contribution of the two-photon states in Eq. (46) reduces the coherence of the state. With the help of Eq. (47), we find
$$\mathrm{\Delta }\rho _{\mathrm{int}}\rho _{\mathrm{int}}(\mathrm{\Delta }t)\rho _{\mathrm{int}}(0)=\frac{1}{2}_0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _2|b(\omega _1,\omega _2;\mathrm{\Delta }t)|^2\rho _{\mathrm{int}}(0).$$
(48)
The two-photon probabilities are proportional to $`\mathrm{\Delta }t`$ and connected to the damping rate $`\mathrm{\Gamma }`$ as discussed in the appendix. Hence Eq. (48) yields
$$\frac{d\rho _{\mathrm{int}}}{dt}\frac{\mathrm{\Delta }\rho _{\mathrm{int}}}{\mathrm{\Delta }t}=\frac{\rho _{\mathrm{int}}}{t_\mathrm{d}},$$
(49)
with $`t_\mathrm{d}`$ given by Eq. (37).
In this derivation, the expression for the emission amplitudes per se are not of any relevance — only its connection with the damping rate $`\mathrm{\Gamma }`$ is important. This connection is based on the principle of energy conservation: the energy of the oscillator is damped at the rate at which energy is radiated. Since this argument also holds for the real 3+1 electromagnetic field, we may extend our results by replacing the 3D result for $`\mathrm{\Gamma }`$ into Eq. (37). The dissipative dynamical Casimir force on an oscillating (frequency $`\omega _0`$) perfectly-reflecting sphere was obtained in Ref. . Usually, the sphere is very small when compared with the wavelength of the relevant vacuum fluctuations, which are of the order of $`2\pi /\omega _0.`$ When $`\omega _0R1,`$ the force on the sphere of radius $`R`$ is given by
$$F=\frac{\mathrm{}R^6}{648\pi }x^{(9)},$$
(50)
where $`x^{(9)}`$ is the ninth time derivative of the position of the sphere. Following again the method of Eq. (44), we calculate the damping rate $`\mathrm{\Gamma }`$ from the equation of motion. We find
$$\mathrm{\Gamma }=\frac{1}{1296\pi }\frac{\mathrm{}\omega _0^8R^6}{M}$$
(51)
showing that the coupling with the vacuum field is reduced, as compared with the 1D case, by the (very small) factor $`(\omega _0R)^6.`$ This reduction factor accounts for the inefficient coupling of the small particle, which scatters field modes of very long wavelengths. Using Eq. (37), we find that the decoherence time increases by the same factor:
$$t_\mathrm{d}=\frac{324}{v^2}\frac{1}{(\omega _0R)^6}\frac{2\pi }{\omega _0},$$
(52)
hence decoherence through radiation pressure is a tiny effect at $`T=0.`$ At finite temperatures, however, the effect may be significant, as discussed in the next section.
## VI High-temperature limit
In this section, we compute the damping and decoherence rates when $`T\mathrm{}\omega _0.`$ In this limit, vacuum fluctuations are negligible when compared with thermal fluctuations, and the dominant contribution in Eq. (21) comes from $`\xi ^T,`$ which is given by Eq. (24). When the temperature is also higher than the cut-off energy $`\mathrm{}\mathrm{\Omega },`$ all relevant frequencies in (24), which are smaller or of the order of $`\mathrm{\Omega },`$ are much smaller than $`T/\mathrm{}.`$ Then, we may take the approximation $`n_\omega ^{}T/(\mathrm{}\omega ^{}),`$ yielding
$$\xi ^T[\omega _0]=2\frac{\mathrm{}\mathrm{\Omega }T}{\omega _0}.$$
(53)
Replacing (53) into (27) yields
$$\mathrm{\Gamma }=\frac{\mathrm{\Omega }T}{2M},$$
(54)
in agreement with the result for the viscous radiation pressure force obtained in Ref. : $`F=\mathrm{\Omega }T\dot{q}(t).`$
From a practical point-of-view, the opposite limit, $`\mathrm{}\omega _0T\mathrm{}\mathrm{\Omega }`$ is more interesting for particles that scatter visible light ($`\mathrm{\Omega }10^{16}\mathrm{Hz}`$). In this case, the corresponding reflectivity amplitude $`R(\omega )`$ is approximately constant for the field modes whose frequencies are smaller or of the order of $`T/\mathrm{}.`$ As a consequence, we may neglect the Lorentzian fall-off in (24). Moreover, we replace the thermal photon number $`n_{\omega ^{}\omega _0}`$ in (25) by
$$n_{\omega ^{}\omega _0}\left[\mathrm{exp}(\mathrm{}\omega ^{}/T)(1\mathrm{}\omega _0/T)1\right]^1.$$
(55)
Neglecting second and higher order terms in $`\mathrm{}\omega _0/T,`$ we find
$$G(\omega _0,\omega ^{})=\frac{\omega ^{}e^{\frac{\mathrm{}\omega ^{}}{T}}}{(e^{\frac{\mathrm{}\omega ^{}}{T}}1)^2}\frac{\mathrm{}\omega _0}{T}.$$
(56)
From (24) and (56) we find
$$\xi ^T[\omega _0]=\frac{4\pi }{3}\frac{T^2}{\omega _0},$$
(57)
and then
$$\mathrm{\Gamma }=\frac{\pi }{3}\frac{T^2}{M\mathrm{}},$$
(58)
which is also in agreement with Ref. . It corresponds to the high-temperature, perfectly-reflecting limit. Here $`T`$ plays the role of frequency cut-off instead of $`\mathrm{\Omega },`$ so that the resulting damping rate is independent of the latter.
The dissipative force in the high temperature limit may be interpreted as the effect of Doppler shift of the reflected thermal photons . For a photon of frequency $`\omega ,`$ the frequency shift is $`\mathrm{\Delta }\omega =\pm 2\omega \dot{q},`$ the plus and minus signs applying for counter and co-propagating cases. Hence the motion gives rise to an unbalance between the radiation pressure exerted on each side of the mirror, corresponding to a momentum transfer $`\mathrm{\Delta }P=2\mathrm{\Delta }E\dot{q},`$ where $`\mathrm{\Delta }E`$ is the reflected energy during a time interval $`\mathrm{\Delta }t.`$ In terms of the density of modes $`g(\omega ),`$ we have
$$\mathrm{\Delta }E=_0^{\mathrm{}}𝑑\omega |R(\omega )|^2g(\omega )n_\omega \mathrm{}\omega ,$$
(59)
where $`|R(\omega )|^2`$ represents the mirror reflectivity (the reflection amplitude $`R`$ is given by Eq. (4)). From $`\mathrm{\Delta }E,`$ the friction force is obtained through
$$F=2\frac{\mathrm{\Delta }E}{\mathrm{\Delta }t}\dot{q}.$$
(60)
In the 1D case, the density of modes is frequency independent: $`g(\omega )d\omega =(L/\pi )d\omega `$, where $`L=\mathrm{\Delta }t`$ is the length of the quantization box. When replaced into Eq. (59), this result leads, with the help of (60), to expressions for the force in agreement with our results for $`\mathrm{\Gamma },`$ except for a factor of $`2`$ when $`\mathrm{}\omega _0T\mathrm{}\mathrm{\Omega }`$ . In the 3D case, on the other hand, we have
$$g(\omega )d\omega =\frac{V}{\pi ^2}\omega ^2d\omega ,$$
(61)
where $`V=A\mathrm{\Delta }t`$ is the quantization volume, $`A`$ being the surface of the mirror (in this case, for simplicity, we assume a flat rather than spherical mirror). In the limit $`\mathrm{}\omega _0T\mathrm{}\mathrm{\Omega },`$ Eqs. (59) and (61) yield
$$\frac{\mathrm{\Delta }E}{\mathrm{\Delta }t}=\frac{\mathrm{}A}{\pi ^2}_0^{\mathrm{}}𝑑\omega \frac{\omega ^3}{\mathrm{exp}\left(\mathrm{}\omega /T\right)1}=\frac{\pi ^2}{15}\frac{AT^4}{\mathrm{}^3}.$$
(62)
As expected, the reflected power features the $`T^4`$ dependence of the Stefan-Boltzmann law, since it is proportional to the total blackbody radiation energy in this limit. The friction force is found by replacing Eq. (62) into (60). The resulting damping coefficient is given by
$$\mathrm{\Gamma }=\frac{\pi ^2}{15}\frac{A}{\mathrm{}^3}\frac{T^4}{M}.$$
(63)
Since we have neglected diffraction at the borders of the mirror, this result only applies when the mirror is much larger than the thermal photon wavelength $`\lambda _{\mathrm{th}}=2\pi \mathrm{}/T.`$
The decoherence time is then found by replacing (63) into (38), which connects damping and decoherence in the high-temperature limit (we re-introduce the speed of light $`c`$ in order to allow an evaluation of the orders of magnitude):
$$t_\mathrm{d}=\frac{15}{32\pi ^7}\frac{\lambda _{\mathrm{th}}^5}{cA\mathrm{\Delta }Q^2}.$$
(64)
As a numerical example, we take $`T=50K,`$ which gives $`\lambda _{\mathrm{th}}=2.9\times 10^4\mathrm{m},`$ and $`A=1\mathrm{m}\mathrm{m}^2.`$ In this case, diffraction is negligible, and Eq. (64) yields $`t_\mathrm{d}[\mathrm{s}]=1.0\times 10^{24}/(\mathrm{\Delta }Q^2[\mathrm{m}^2]),`$ showing that decoherence is very fast even when the distance between the wavepackets is, for instance, in the nanometer range — in this case the decoherence time is of the order of a micro-second. Since $`t_\mathrm{d}`$ scales as $`1/T^5,`$ it is still shorter, by a factor $`8\times 10^3,`$ at room temperature. Note, however, that Eq. (64) only applies when $`\omega _0t_\mathrm{d}1,`$ the basic assumption that allowed us to relate decoherence and damping time scales with the help of the fluctuation-dissipation theorem .
## VII Conclusions
As in the Caldeira-Legget model , coherent states are the most robust when the radiation pressure coupling with the quantum field is considered. This is amazingly in line with their well-known status of ‘quasi-classical’ states, i.e., the closest possible representation of a classical oscillation in a harmonic potential well. Superpositions of coherent states decay into a mixture at a rate proportional to the damping rate and to the squared distance in phase space. The ratio between decoherence and damping rates is a simple hyperbolic increasing function of temperature, which interpolates the zero and high temperature limits. It originates from the general relation between symmetric and anti-symmetric correlation functions, associated to the fluctuation-dissipation theorem. Thus, the particular nature of the model for the coupling with the reservoir seems to be immaterial, as far as the connection between damping and decoherence is concerned. Note, however, that the validity of this result is limited by the assumption that decoherence is slow compared to the free oscillations.
We have shown that the radiation pressure exerted by thermal photons is a very efficient source of decoherence, although the corresponding energy damping effect, associated to the Doppler frequency shift of the reflected photons, is usually negligible. At $`T=0,`$ the energy damping is associated to the emission of photon pairs (dynamical Casimir effect). The dominant contribution comes from vacuum fluctuations corresponding to wavelengths of the order of $`2\pi c/\omega _0,`$ which is usually much greater than the size of the oscillator. As a consequence, the radiation pressure coupling is inefficient, and both damping and decoherence rates become very small. It is however remarkable, from a theoretical point-of-view, that the mere inclusion of an unavoidable, intrinsically quantum effect, is sufficient (in principle) to engender decoherence, and by that means restoring, although in a very long time scale, the classical world.
We are grateful to A. Calogeracos and G. Barton for correspondence, and to J. Dziarmaga, A. Lambrecht, M.-T. Jaekel and S. Reynaud for discussions. P. A. M. N. thanks CNPq, PRONEX and FAPERJ for partial financial support.
## A Entanglement with two-photon states
In this appendix, we present an alternative, simpler derivation of the decoherence time scale at $`T=0,`$ which shows more clearly how the dynamical Casimir effect modifies the quantum phase of a superposition state and engenders decoherence. Instead of tracing over the field, we follow its evolution during many periods of free oscillation. We first take, at $`t=0,`$ the mirror-plus-field state $`|\alpha |0`$ ($`|0`$ represents the vacuum field state), where $`|\alpha `$ is a coherent state of large amplitude: $`|\alpha |1.`$ We take $`\alpha =i\dot{q}(0)\sqrt{M/2\mathrm{}\omega _0},`$ so that $`|\alpha `$ is a ‘semiclassical’ state associated to a minimum uncertainty wave-packet whose initial velocity is $`\dot{q}(0).`$ We have shown in Sec. IV that the action of the vacuum radiation pressure on the motion of the mirror is a very small perturbation (weak coupling limit). Thus the time evolution may be computed from a simple ‘semi-classical’ model, in which the field evolution is obtained assuming the classical prescribed motion
$$\dot{q}(t)=\dot{q}(0)\mathrm{cos}(\omega _0t),$$
(A1)
where $`q(t)`$ is the position of the mirror at time $`t.`$ The dynamical Casimir effect is described by the interaction Hamiltonian (see Ref. , and compare with the first term in Eq. (5))
$$H_{\mathrm{int}}=\dot{q}(t)𝒫.$$
(A2)
The amplitude $`b(\mathrm{\Delta }t)`$ for the creation of photon pairs corresponding to frequencies $`\omega _1`$ and $`\omega _2`$ at time $`\mathrm{\Delta }t`$ is given by
$$b(\omega _1,\omega _2;\mathrm{\Delta }t)=\frac{i}{\mathrm{}}\omega _1,\omega _2|𝒫|0_0^{\mathrm{\Delta }t}𝑑t^{}e^{i(\omega _1+\omega _2)t^{}}\dot{q}(t^{}).$$
(A3)
According to Eq. (A3), the amplitude depends on the sign of $`\dot{q},`$ which is very important to the discussion of decoherence.
Replacing Eq. (A1) into (A3), we find for the two-photon probabilities
$$|b(\omega _1,\omega _2;\mathrm{\Delta }t)|^2\frac{1}{\mathrm{}^2}|0|𝒫|\omega _1,\omega _2|^2\dot{q}(0)^2$$
(A4)
$`\times {\displaystyle \frac{\mathrm{sin}^2\left[(\omega _1+\omega _2\omega _0)\mathrm{\Delta }t/2\right]}{(\omega _1+\omega _2\omega _0)^2}}.`$
For $`\omega _0\mathrm{\Delta }t1,`$ the r.-h.-s. of Eq. (A4) is sharply peaked around $`\omega _1+\omega _2=\omega _0.`$ Thus, for large times energy is well defined, in agreement with the time-energy uncertainty relation. In this limit, Eq. (A4) yields
$$|b(\omega _1,\omega _2;\mathrm{\Delta }t)|^2\frac{\pi }{2\mathrm{}^2}|0|𝒫|\omega _1,\omega _2|^2\dot{q}(0)^2\mathrm{\Delta }t$$
(A5)
$`\times \delta (\omega _1+\omega _2\omega _0)`$
Since the source of the radiated energy is the motion of mirror, one may expect that the two-photon probabilities are related to the amplitude decay rate $`\mathrm{\Gamma }.`$ The radiated energy during $`\mathrm{\Delta }t`$ is
$$\mathrm{\Delta }E=\frac{1}{2}_0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _2|b(\omega _1,\omega _2;\mathrm{\Delta }t)|^2\mathrm{}(\omega _1+\omega _2),$$
(A6)
which according to Eq. (A5) is proportional to the time interval $`\mathrm{\Delta }t.`$ The energy of the mirror decays as $`dE_M/dt=2\mathrm{\Gamma }E_M,`$ where $`E_M=M\dot{q}(0)^2/2.`$ Hence, from energy conservation we have
$$\mathrm{\Gamma }=\frac{1}{M\dot{q}(0)^2}\frac{\mathrm{\Delta }E}{\mathrm{\Delta }t},$$
leading, with the help of (A5) and (A6), to the representation given by (41).
To analyze the effect of decoherence, we take the field to be initially in the ‘even’ superposition state $`|\psi _\mathrm{e}=(|\alpha +|\alpha )/\sqrt{2},`$ so that the mirror-plus-field state at $`t=0`$ is
$`|\mathrm{\Psi }_0=|\psi _\mathrm{e}|0.`$
By linearity, its time evolution is obtained from the two–photon amplitudes given by (A3):
$$|\mathrm{\Psi }_{\mathrm{\Delta }t}=(|\alpha |\varphi ^+_{\mathrm{\Delta }t}+|\alpha |\varphi ^{}_{\mathrm{\Delta }t})/\sqrt{2},$$
(A7)
where
$$|\varphi ^\pm _{\mathrm{\Delta }t}=B(\mathrm{\Delta }t)|0\pm \frac{1}{2}_0^{\mathrm{}}𝑑\omega _1_0^{\mathrm{}}𝑑\omega _2b(\omega _1,\omega _2;\mathrm{\Delta }t)|\omega _1,\omega _2.$$
(A8)
The already noted sensitivity of the two-photon amplitudes to the phase of the motion of the mirror, which is explicit through the ‘minus’ sign for $`|\varphi ^{}`$ in (A8), generates entanglement between mirror and field. This is discussed in Sec. V, whose starting point is Eq. (46), which is derived by replacing (A8) into (A7).
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# Fluctuation, Dissipation, and Entanglement: the Classical and Quantum Theory of Thermal Magnetic Noise
## 1. Introduction
The engineering control of quantum decoherence is central to several emerging scientific goals. One such goal is the direct observation of molecular structure—in three dimensions and with single-atom resolution—by magnetic resonance force microscopy . Another is the solution of otherwise intractable mathematical problems—like factoring large numbers—by quantum computing .
At present, the main practical obstacle to achieving these goals is that the various mechanisms that cause quantum decoherence are not yet fully cataloged and understood. In consequence, quantum-coherent experiments often reveal unexpectedly fast decoherence rates, due to unanticipated—or even previously unknown—relaxation mechanisms. Similar struggles with decoherence have occupied previous scientific generations; Orbach for example reviews the forty-year struggle to achieve a reasonably comprehensive understanding of relaxation in electron spin resonance.
This article is concerned with a decoherence mechanism which is ubiquitous in quantum-coherent technologies: *thermal magnetic noise*. Physically speaking, thermal magnetic noise is created by the same thermal current fluctuations that create Johnson noise. Thus, thermal magnetic noise is present in any device that contains electric conductors.
We will be mainly concerned with practical engineering aspects of thermal magnetic noise, but we will also give an explicit example—the first in the literature to our knowledge—of a *fluctuation-dissipation-entanglement* theorem. Such theorems may be broadly defined as invertible functional relations between the dissipative kernel of a system and the system’s quantum entanglement with a thermal reservoir, such that the entanglement determines the dissipation and *vice versa*.
This article is organized as follows. Prior work relating to thermal magnetic noise is reviewed in Section 1.1. The main new results are summarized in Section 2, these results consist of closed-form expressions for the magnetic noise spectral density (Section 2.1), quantum decoherence (Section 2.2), and a fluctuation-dissipation-entanglement theorem (Section 2.3). The practical implications of these results for quantum-coherent engineering are discussed in Section 3 via worked examples for trapped atom experiments (Section 3.1), magnetic resonance force microscopy (Section 3.2), and quantum computing (Section 3.3). Some topics for further research are suggested in Section 4.
Algebraically tedious derivations are relegated to three appendices. Appendix A derives a closed-form expression for the magnetic noise spectral density, Appendix B derives quantum Langevin and Bloch equations that describe decoherence in the presence of thermal magnetic noise, and Appendix C derives fluctuation-dissipation-entanglement theorems for both spin-$`\frac{1}{2}`$ particles and harmonic oscillators, as well as relations between dissipation and renormalization.
### 1.1. Prior work relating to thermal magnetic noise
The quantum field theory of electromagnetic fluctuations in lossy linear media is presented in textbooks by Landau, Lifshitz, and Piaevskii and by Rytov *et al.* . Far from a warm body, this theory describes fluctuations which are simply the familiar phenomenon of black-body radiation. Near to a warm body, the situation is considerably more complex. Non-radiating terms in the field theory become important, and create phenomena such as the Casimir force (identical to the Van der Waals force), which is familiar from chemistry and force microscopy. Recently, Dorofeyev *et al.* have observed the attractive, fluctuating and dissipative components of the Casimir force, which were in good agreement with field-theoretic predictions.
The Casimir force arises mainly from fluctuating electric fields, which do not directly couple to the spin magnetic moments. Our investigation instead focuses on thermal fluctuations in magnetic fields.
To the best of our knowledge, the theory and experiments of Varpula and Poutanen , as reviewed and extended by Nenonen, Montonen, and Katila , were the first experimental observation and theoretical analysis of thermal magnetic fluctuations. Their work focussed on biomagnetism experiments, in which magnetic fluctuations originate in the copper walls of shielded rooms; such noise can be easily large enough to obscure biomagnetic signals. They presented a phenomenological model of thermal magnetic noise in which metallic conductors were modeled as collections of independent thermally excited resistive elements. The resulting predictions were in excellent accord with experiment.
Quantum-coherent technologies operate in a vastly different regime from biomagnetism experiments. Ambient temperatures $`T`$ are as low as is experimentally feasible—a few Kelvins down to millikelvins—and observation frequencies $`\omega `$ are megahertz to gigahertz, such that $`\mathrm{}\omega k_\mathrm{B}T`$.
Quantum-coherent phenomena dominate this physical regime. The applicability of the Varpula-Poutanen noise model then becomes uncertain, because the model is derived from a purely phenomenological description of thermal noise in room-temperature conductors. We have therefore sought to improve and extend the Varpula-Poutanen model in five respects:
1. Rigorous quantum mechanical and thermodynamic foundations have been provided via the fluctuation-dissipation theorem.
2. The results now encompass superconducting materials and/or materials with non-vanishing magnetic susceptibility.
3. A closed-form expression for the thermal magnetic noise spectral density has been obtained.
4. Bloch equations describing quantum decoherence have been derived via a quantum Langevin formalism and the independent oscillator heat bath model of Ford, Lewis, and O’Connell.
5. Fluctuation-dissipation-entanglement theorems have been constructed, and the link between dissipation and renormalization has been clarified.
In carrying through this program, we have been able to prove that the Varpula-Poutanen model is *exact* in the physical regime relevant to biomagnetism: real conductivity, macroscopic length scales, and large temperature. This explains why the Varpula-Poutanen model yields predictions which are in excellent accord with experimental measurements of biomagnetism , and is a tribute to their pioneering physical insight.
## 2. Physical motivation and main results
The necessary existence of thermal magnetic noise can be deduced by considering a loop antenna held near a conducting slab, as shown in Fig. 2. If a voltmeter is placed across the terminals of such an antenna, a fluctuating voltage $`V(t)`$ will be observed. The spectral density<sup>1</sup><sup>1</sup>1Our normalization for spectral densities is $`S_V(\omega )=_{\mathrm{}}^{\mathrm{}}𝑑\tau e^{i\omega \tau }V(0)V(\tau )`$, where $`V(0)V(\tau )`$ is the voltage autocorrelation. Thus $`V^2=1/(2\pi )_{\mathrm{}}^{\mathrm{}}𝑑\omega S_V(\omega )`$, so that spectral densities are two-sided (positive and negative frequencies), with bandwidths given in Hertz. $`S_V(\omega )`$ of this *Johnson noise* is related to the complex antenna impedance $`Z(\omega )`$ by the well-known equation
(2.1) $`S_V(\omega )`$ $`=\text{Re}\left(Z(\omega )\right)\mathrm{}\omega \mathrm{coth}\left({\displaystyle \frac{\mathrm{}\omega }{2k_\text{B}T}}\right),`$ in general;
$`\text{Re}\left(Z(\omega )\right)2k_\text{B}T,`$ for $`k_\text{B}T\mathrm{}\omega `$.
Here $`\text{Re}\left(Z(\omega )\right)`$ is the resistive impedance. This result is just the fluctuation-dissipation theorem as it applies to fluctuating voltages in dissipative impedances.
Now we ask, what is the physical origin of these voltage fluctuations? We suppose that the coil has negligible intrinsic resistance, in which case the voltage can only have been induced by a magnetic flux threading the antenna loop. In turn, this magnetic flux can only have been generated by thermally excited currents in the nearby conducting slab.
So wherever Johnson noise is present, thermal magnetic noise will be present also, because the same underlying current fluctuations generate both phenomena.
### 2.1. Thermal magnetic spectral densities
We will begin by presenting closed-form expressions for the spectral density of thermal magnetic noise, and discussing their physical significance. Deriving these expressions is tedious, and is deferred until Appendix A.
Since the thermal magnetic field $`𝑩(t)`$ is a vector, its autocorrelation $`𝑪_𝑩(\tau )`$ and spectral density $`𝑺_𝑩(\omega )`$ are matrix-valued functions defined by
(2.2) $`𝑪_𝑩(\tau )`$ $`=𝑩(0)\mathbf{}𝑩(\tau ),`$
(2.3) $`𝑺_𝑩(\omega )`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{i\omega \tau }𝑪_𝑩(\tau )𝑑\tau `$
where $`\mathbf{}`$ denotes the outer product of two vectors: $`(𝒂\mathbf{}𝒃)_{ij}a_ib_j`$. As shown in Section A.3, $`𝑺_𝑩(\omega )`$ is of the general form:
(2.4) $`𝑺_𝑩(\omega )`$ $`=𝚪(\omega )\mathrm{}\omega \mathrm{coth}\left({\displaystyle \frac{\mathrm{}\omega }{2k_\text{B}T}}\right)`$
$`=\left[𝑰+\widehat{𝒏}\mathbf{}\widehat{𝒏}\right]\mathrm{\Gamma }(\omega )\mathrm{}\omega \mathrm{coth}\left({\displaystyle \frac{\mathrm{}\omega }{2k_\text{B}T}}\right)`$
where $`\widehat{𝒏}`$ is the unit vector normal to the slab surface and $`𝑰`$ is the 3$`\times `$3 identity matrix.
The *magnetic dissipation tensor* $`𝚪(\omega )`$ in (2.4) plays the same fundamental role that resistance plays in Johnson noise, in the sense that—as we will see—knowledge of $`𝚪(\omega )`$ quantitatively determines a broad spectrum of fluctuation, dissipation, and quantum entanglement phenomena.
#### 2.1.1. Device parameters
In general, we will concerned with thermal magnetic noise induced at a distance $`d`$ from a conductive slab of thickness $`t`$, as illustrated in Fig. 1. The slab is assumed to have a frequency-dependent complex conductivity $`\sigma (\omega )`$ and permeability $`\mu (\omega )`$; for compactness we suppress the frequency-dependence of $`\sigma `$ and $`\mu `$. All results are presented in S.I. units, and all time-dependence is assumed to be $`e^{i\omega t}`$. See Section A.3 for a statement of Maxwell’s equations in this convention.
The four device parameters $`\{d,t,\sigma ,\mu \}`$ may be regarded as fundamental; all other parameters are derived from them. The single most important derived parameter is the *skin depth* $`\lambda |\omega \mu \sigma |^{1/2}`$. Another important derived parameter is the phase $`\varphi `$ of the complex conductivity, as defined by $`𝒋=\sigma 𝑬=|\sigma |e^{i\varphi }`$, where $`𝑬`$ is the E-field and $`𝒋`$ is the current.
#### 2.1.2. Superconductors
For the special case of an ideal superconductor, the London equation can be stated in the form $`i\mu _0\omega 𝒋=\lambda _\text{L}^2𝑬`$, where $`\lambda _\text{L}`$ is the *London penetration depth*. Comparing this with the definition of the conductivity $`𝒋=\sigma 𝑬`$, we see that for an ideal superconductor $`\sigma =i\lambda _\text{L}^2/(\mu _0\omega )`$, $`\lambda =\lambda _\text{L}`$, and $`\varphi =\pi /2`$. Thus ideal superconductors do not have infinite conductivity, but rather have a wholly imaginary conductivity.
Real superconductors approximate ideal superconductors only at zero temperature. At finite temperatures a mixture of superconducting and normal phase charge carriers is present, such that the conductivity has an intermediate phase $`\varphi (0,\pi /2)`$. See for a review of the literature, including measurements of $`\sigma `$ at finite temperature.
#### 2.1.3. Single paramagnetic or diamagnetic conducting slabs
Most nonferrous metals and insulators have relative magnetic permeability $`K\mu /\mu _01`$, where $`\mu `$ is the permeability of the material (the slab in our case), and $`\mu _0`$ is the permeability of the vacuum. In Appendix A we derive the following asymptotic dependance of the magnetic dissipation for $`K=1`$
(2.5)
$$\mathrm{\Gamma }(\omega )=\frac{\mu _0^2\text{Re}(\sigma )}{64\pi }\times \{\begin{array}{cc}\frac{t}{d(t+d)},\hfill & \lambda \text{min}\{d,\sqrt{dt}\};\hfill \\ \frac{3\lambda ^3}{d^4\mathrm{cos}(\pi /4\varphi /2)},\hfill & \lambda \text{min}\{d,t\};\hfill \\ \frac{2\lambda ^4(dt8\lambda ^2\mathrm{sin}\varphi )}{d^5t^2}.\hfill & t\lambda \sqrt{dt}.\hfill \end{array}$$
Collectively, these three limits span the entire $`\{d,t,\sigma ,\omega \}`$ design space. The following expression smoothly interpolates all three limits:
(2.6)
$$\mathrm{\Gamma }(\omega )\frac{3\mu _0^2\text{Re}(\sigma )\lambda ^3t\mathrm{cos}\varphi }{64\pi \left(3\lambda ^3d(d+t)+(1e^{\alpha _\text{c}})td^2(d+2\lambda )^2\mathrm{cos}(\pi /4\varphi /2)\right)}.$$
It is constructed by adding the first and second asymptotic limits of (2.5) inversely, then introducing a transition parameter
(2.7)
$$\alpha _\text{c}=(dt+8\lambda ^2\mathrm{sin}\varphi )/(2d\lambda \mathrm{cos}(\pi /4\varphi /2))$$
such that the third asymptotic limit also is smoothly interpolated.
A numerical survey shows that the resulting closed-form expression (2.6) predicts noise levels that are no more than $`1.75\text{ dB}`$ too large and $`0.5\text{ dB}`$ too small everywhere in $`\{d,t,\sigma ,\omega \}`$ design space—an accuracy which suffices for almost all practical design work.
#### 2.1.4. The case of two conducting slabs
The case of magnetic fluctuations measured at a point midway between two identical conducting slabs is also solved in Appendix A, and the resulting asymptotic expressions are precisely double those given in (2.5). A numerical survey shows that doubling the one-slab noise yields a slight overestimate of the two-slab midpoint noise, but the error is no more than $`0.6\text{ dB}`$ for all $`\{d,t,\sigma ,\omega \}`$. Thus, to a very good approximation two adjacent slabs each generate independent noise at the midpoint.
#### 2.1.5. The general case
In the general case in which both lossy magnetic permeability and electric conductivity are combined with a skin depth $`\lambda `$ comparable to $`d`$, no analytic results are presently available. Appendix A presents a one-dimensional integral (A.18) which must be numerically evaluated to determine the magnetic spectral densities.
#### 2.1.6. Comparison with experiment
Although our general integral expression (A.18) is superficially of a different form than the Varpula-Poutanen model , and is derived from a very different underlying model, a numerical comparison shows that these two models yield identical predictions in the physical regime relevant to biomagnetic experiments: real conductivity, audio frequencies, vanishing magnetic susceptibility, and large temperature. Since in this regime the Varpula-Poutanen model has been shown to agree very well with experiment, we may regard the model presented in this article as having been at least partially validated.
More extensive comparisons of our model with the Varpula-Poutanen model are not feasible, because the V-P model does not contain parameters corresponding to the magnetic susceptibility $`\mu `$, conductivity phase $`\varphi `$, or the quantum of action $`\mathrm{}`$.
We note also that *no* model of thermal magnetic noise has yet been experimentally tested at microscopic scales, low temperatures, and high frequencies, in materials that are wholly or partially superconducting, or are magnetic—which is precisely the regime of greatest importance to quantum-coherent engineering.
### 2.2. Quantum decoherence
Now we turn our attention to the quantum decoherence induced by thermal magnetic fields. We assume a two-state system coupled to the magnetic thermal field $`𝑩(t)`$ via the Hamiltonian
$$H=\gamma \left(𝑩_0+𝑩(t)\right)\mathbf{}𝒔$$
where $`\gamma `$ is the gyromagnetic ratio of the spin, $`𝑩_0`$ is a constant polarizing field, and the spin angular momentum $`𝒔`$ satisfies the usual commutation relations $`[s_i,s_j]=i\mathrm{}ϵ_{ijk}s_k`$.
Our results will be expressed in terms of the rotating-frame magnetic spectral density $`𝑺_{𝑩_{\text{rot}}}(\omega )`$, which is given in terms of the laboratory-frame density $`𝑺_𝑩(\omega )`$ (2.4) via
(2.8)
$$\begin{array}{c}𝑺_{𝑩_{\text{rot}}}(\omega )=\left(\widehat{𝒃}\mathbf{}\widehat{𝒃}\right)\text{tr}\left[\widehat{𝒃}\mathbf{}\widehat{𝒃}\mathbf{}𝑺_𝑩(\omega )\right]+\hfill \\ \hfill \left(𝑰\widehat{𝒃}\mathbf{}\widehat{𝒃}\right)\frac{1}{2}\text{tr}\left[(𝑰\widehat{𝒃}\mathbf{}\widehat{𝒃})\mathbf{}𝑺_𝑩(\omega _0)\right],\end{array}$$
where the precession frequency $`\omega _0`$ about the unit axis of spin precession $`\widehat{𝒃}`$ is given by $`\gamma 𝑩_0\omega _0\widehat{𝒃}`$, and $`\text{tr}[\mathrm{}]`$ is the matrix trace (the sum of diagonal elements).
In Appendix B we derive the quantum Langevin equations for a spin-$`\frac{1}{2}`$ particle coupled to a thermal magnetic field, and from them we show that the time evolution of the spin density matrix is described by Bloch-type equations with relaxation times $`T_1`$, $`T_2`$, and $`T_{1\rho }`$ given by
(2.9a) $`{\displaystyle \frac{1}{T_1}}`$ $`={\displaystyle \frac{1}{2}}\gamma ^2\text{tr}\left[(𝑰\widehat{𝒃}\mathbf{}\widehat{𝒃})\mathbf{}𝑺_{𝑩_{\text{rot}}}(0)\right]`$
(2.9b) $`{\displaystyle \frac{1}{T_2}}`$ $`={\displaystyle \frac{1}{2T_1}}+{\displaystyle \frac{1}{2}}\gamma ^2\text{tr}\left[\widehat{𝒃}\mathbf{}\widehat{𝒃}\mathbf{}𝑺_{𝑩_{\text{rot}}}(0)\right]`$
(2.9c) $`{\displaystyle \frac{1}{T_{1\rho }}}`$ $`={\displaystyle \frac{1}{2}}\gamma ^2\text{tr}\left[(𝑰\widehat{𝒃}_1\mathbf{}\widehat{𝒃}_1)\mathbf{}𝑺_{𝑩_{\text{rot}}}(\omega _1)\right]`$
Here $`T_{1\rho }`$ is the relaxation time in the presence of a radio-frequency (RF) applied field—commonly called a “spin-locking” field. In the frame co-rotating with the RF field, it appears as a constant B-field $`𝑩_1=B_1\widehat{𝒃}_1`$ whose characteristic precession frequency is $`\omega _1=\gamma B_1`$. Substituting (2.42.8) in (2.9a2.9b), we obtain $`\{T_1,T_2,T_{1\rho }\}`$ explicitly in terms of the magnetic damping coefficient $`\mathrm{\Gamma }(\omega )`$ given in (2.6)
(2.10a) $`{\displaystyle \frac{1}{T_1}}`$ $`={\displaystyle \frac{1}{2}}\gamma ^2(3\mathrm{cos}^2\theta )\mathrm{\Gamma }(\omega _0)\mathrm{}\omega _0\mathrm{coth}\left({\displaystyle \frac{\mathrm{}\omega _0}{2k_\text{B}T}}\right)`$
(2.10b) $`{\displaystyle \frac{1}{T_2}}`$ $`={\displaystyle \frac{1}{2T_1}}+{\displaystyle \frac{1}{2}}\gamma ^2(1+\mathrm{cos}^2\theta )\mathrm{\Gamma }(0)\mathrm{\hspace{0.17em}2}k_\text{B}T`$
(2.10c) $`{\displaystyle \frac{1}{T_{1\rho }}}`$ $`={\displaystyle \frac{1+\mathrm{cos}^2\beta }{2T_1}}+`$
$`{\displaystyle \frac{1}{2}}\gamma ^2\mathrm{sin}^2\beta (1+\mathrm{cos}^2\theta )\mathrm{\Gamma }(\omega _1)\mathrm{}\omega _1\mathrm{coth}\left({\displaystyle \frac{\mathrm{}\omega _1}{2k_\text{B}T}}\right)`$
Here $`\mathrm{cos}\theta =\widehat{𝒃}\mathbf{}\widehat{𝒏}`$ and $`\mathrm{cos}\beta =\widehat{𝒃}\mathbf{}\widehat{𝒃}_1`$, where $`\widehat{𝒃}_1`$ is the RF field axis in the rotating frame, $`\widehat{𝒃}`$ is the polarization axis, and $`\widehat{𝒏}`$ is the unit vector normal to the slab surface.
Note that $`\mathrm{\Gamma }(\omega )`$ appears with varying frequency arguments $`\omega \{0,\omega _1,\omega _0\}`$; this has important design consequences because in most cases $`\mathrm{\Gamma }(\omega )`$ exhibits a strong frequency dependence.
As far as practical quantum engineering is concerned, expressions (2.10a2.10c) are the main “deliverable” of this article.
### 2.3. A fluctuation-dissipation-entanglement theorem
From a purely formal point of view, fluctuation-dissipation-entanglement theorems exist for a simple reason: the same Hamiltonian matrix elements that control fluctuation and dissipation also control certain measures of quantum entanglement; Appendix C discusses this point of view.
Physically speaking, we reason as follows. We consider a two-state quantum system, which—as in the real world—interacts weakly with a thermal reservoir. We adjust the temperature of the reservoir to zero, and ask: what is the probability that the two-state system is *not* in its ground state? To the extent that this probability is non-zero, it describes an irreducible quantum entanglement with the thermal reservoir.
Now, at zero temperature the classical entanglement probability is zero, and even the Bloch equations (B.19), which have a firmer quantum justification, relax spin systems to their ground state at zero temperature, and thus predict zero entanglement.
However, a higher-order calculation reveals that the entanglement probability is finite even at zero temperature. We calculate this probability as follows. As before, we consider a spin-$`\frac{1}{2}`$ particle magnetically coupled to a thermal reservoir. We polarize the spin via an external field $`𝑩_0`$, which lifts the degeneracy of the spin system along a polarization axis $`\omega _0\widehat{𝒑}=\gamma 𝑩_0`$, with $`\gamma `$ the gyromagnetic ratio of the spin and $`\omega _0`$ the (by convention positive) precession frequency. We let $`|\psi _0`$ be the ground state of the isolated spin, and we define a projection operator $`P_{}=𝕀|\psi _0\mathbf{}\psi _0|`$, with $`𝕀`$ the identity operator. By construction, the expectation value of $`P_{}`$ vanishes for an isolated ground state: $`\psi _0|P_{}|\psi _0=0`$.
Of course, no real-world spin system is perfectly isolated.<sup>2</sup><sup>2</sup>2 If nothing else, the spin can exchange quanta with the vacuum, which may be regarded as a zero-temperature thermal reservoir. It follows that fluctuation-dissipation-entanglement theorems can be constructed in quantum field theories, where their significance remains to be elucidated. We therefore enlarge our Hilbert space to encompass a thermal reservoir with a potentially infinite set of basis states $`\{|\varphi _i;i0,1,\mathrm{},\mathrm{}\}`$. We generalize $`P_{}`$ to include the thermal reservoir by defining a new projection operator $``$
(2.11)
$$\underset{i}{}|\varphi _iP_{}\varphi _i|$$
which projects onto the subspace in which the spin is *not* in its ground state.
It is a well-defined problem to calculate, perturbatively, the expectation $`\mathrm{\Psi }_0||\mathrm{\Psi }_0`$, where $`|\mathrm{\Psi }_0`$ the ground state of the combined spin-plus-thermal-reservoir system. Keeping in mind that $``$ and $`|\mathrm{\Psi }_0`$ depend on $`\widehat{𝒑}`$ and $`\omega _0`$, we can define a scalar entanglement function $`E(\widehat{𝒑},\omega _0)\mathrm{\Psi }_0||\mathrm{\Psi }_0`$. In Appendix C we show that $`E(\widehat{𝒑},\omega _0)`$ is related to the magnetic dissipation tensor $`𝚪(\omega )`$ by the following *fluctuation-dissipation-entanglement theorem:*
(2.12) $`E(\widehat{𝒑},\omega _0)`$ $`={\displaystyle \frac{\gamma ^2\mathrm{}}{4\pi }}𝔊_2\{\omega \text{tr}\left[(𝑰\widehat{𝒑}\mathbf{}\widehat{𝒑})\mathbf{}𝚪(\omega )\right];\omega _0\}`$
Here $`𝔊_2`$ is a *Stieltj i e s transform*, in Bateman’s notation , as defined by
(2.13)
$$𝔊_\rho \{f(x);y\}=_0^{\mathrm{}}𝑑xf(x)(x+y)^\rho $$
We note that Stieltj i e s transforms are invertible, such that a measurement of the entanglement function $`E(\widehat{𝒑},\omega _0)`$ suffices in principle to determine the magnetic dissipation tensor $`𝚪(\omega )`$ and *vice versa*.
For practical quantum engineering purposes, the main utility of this theorem is that it allows us to determine the hard-to-measure entanglement from the easier-to-measure (or predict) dissipation. For the particular case of thermal magnetic noise, the frequency dependence (2.6) of $`𝚪(\omega )`$ is such that the Stieltj i e s integral is absolutely convergent, so that the entanglement of real-world devices can be readily be calculated.
#### 2.3.1. Approximate expression for the entanglement
In general, the Stieltj i e s integral (2.12) cannot be evaluated in closed form. However, a simple and physically meaningful approximate evaluation is possible. We begin by remarking that dissipative kernels always have a cut-off frequency $`\omega _\text{c}`$, because otherwise their associated noise spectral density would carry infinite power. For the particular case of thermal magnetic noise originating in a thick plate, the cut-off frequency is such that $`\lambda (\omega _\text{c})d`$. We assume the frequency of interest $`\omega _0`$ is small compared to the cut-off frequency, so that $`\omega _\text{c}\omega _0`$, and we brutally approximate $`𝚪(\omega )`$ as constant over the range $`\omega (\omega _0,\omega _\text{c})`$, and zero elsewhere. Then the Stieltj i e s integral (2.12) can be evaluated in closed form as
(2.14) $`E(\widehat{𝒑},\omega _0)`$ $`{\displaystyle \frac{\gamma ^2\mathrm{}}{4\pi }}\text{tr}\left[(𝑰\widehat{𝒑}\mathbf{}\widehat{𝒑})\mathbf{}𝚪(\omega _0)\right]\mathrm{ln}(\omega _\text{c}/\omega _0)`$ from (2.12)
$`={\displaystyle \frac{\gamma ^2\mathrm{}}{4\pi }}\text{tr}\left[(𝑰\widehat{𝒑}\mathbf{}\widehat{𝒑})\mathbf{}{\displaystyle \frac{𝑺_𝑩(\omega _0)}{\mathrm{}\omega _0}}\right]\mathrm{ln}(\omega _\text{c}/\omega _0)`$ by (2.4) for $`T0`$
$`={\displaystyle \frac{\mathrm{ln}(\omega _\text{c}/\omega _0)}{2\pi \omega _0T_1}}`$ by (2.8) and (2.9a)
where $`T_1`$ is evaluated at zero temperature.
## 3. Quantum engineering implications
We will now apply (2.42.10) in the design analysis of representative quantum technologies. Our goal is partly to illustrate practical applications of our results, and partly to identify areas where further research is needed.
A small but important point: when quoting numerical values $`1/\sqrt{\text{Hz}}`$ we shall embrace the usual engineering convention that bandwidths encompass only positive frequencies; this requires the insertion of an additional factor of two when evaluating two-sided spectral densities, *e.g.* (2.4).
### 3.1. Atom and ion traps
We begin by considering thermal magnetic noise at centimeter length scales and audio frequencies, at room temperature. This is the same regime considered by Varpula and Poutanen , and it is also a regime typical of at least some atomic physics experiments .
Specifically, we will calculate the spectral density of the thermal magnetic fields between room-temperature two copper slabs, each $`t=1\text{cm}`$ thick, spaced $`2d=2\text{cm}`$ apart. We are particularly interested in the zero-frequency spectral density in the $`\widehat{𝒏}`$ direction normal to the slab surface, because this describes thermally induced fluctuations in the background polarizing field. At room temperature $`T=300\text{K}`$ the conductivity of high-purity copper is of order $`5.9\times 10^7(\mathrm{\Omega }\text{ m})^1`$. From (2.4) and (2.6), with $`d=t=1\text{cm}`$ and $`\varphi =0`$, and taking into account that we have two plates whose noise is additive, we find at zero frequency $`(\widehat{𝒏}\mathbf{}𝑺_𝑩(0\text{Hz})\mathbf{}\widehat{𝒏})^{1/2}1.2\text{pT}/\sqrt{\text{Hz}}`$, rolling off to $`(\widehat{𝒏}\mathbf{}𝑺_𝑩(100\text{Hz})\mathbf{}\widehat{𝒏})^{1/2}0.6\text{pT}/\sqrt{\text{Hz}}`$ at 100 Hertz.
Until recently, such picoTesla fluctuations would have been viewed as being of no practical consequence. However, magnetic fields change sign under time reversal, and hence magnetic fluctuations locally violate time-reversal invariance, and so these fluctuations must be understood and controlled in high-precision tests of fundamental physics.
For example, the most stringent experimental limit on the electric dipole moment $`d_\text{E}`$ of a fundamental particle is $`d_\text{E}8.7\times 10^{28}e\text{ cm}`$ for the mercury isotope $`{}_{}{}^{199}\text{Hg}`$ (here $`e`$ is the electron charge). Typically, dipole moments are measured by observing the change in precession frequency induced by an applied electric field of magnitude $`E`$. It follows that $`𝑺_𝑩`$ can be expressed as an equivalent dipole noise spectral density $`S_d=(\mathrm{}\gamma /(2E))^2\widehat{𝒏}\mathbf{}𝑺_𝑩(0)\mathbf{}\widehat{𝒏}`$. For $`{}_{}{}^{199}\text{Hg}`$, the gyromagnetic ratio $`\gamma /(2\pi )=7.59\text{MHz/Tesla}`$. This yields a zero-frequency equivalent dipole noise for our example of room-temperature copper plates of
$$ES_d^{1/2}1.94\times 10^{20}\text{eV}/\sqrt{\text{Hz}}.$$
For a typical electric field $`E10^4\text{V/cm}`$, the equivalent dipole noise would therefore be $`S_d^{1/2}1.94\times 10^{24}e\text{ cm}/\sqrt{\text{Hz}}`$.
This is a substantial noise level: it would naively require an averaging time on the order of $`10^7`$ seconds to better the published electric dipole moment limit . And it would not be entirely straightforward to reduce the thermal magnetic noise by making the copper shields thinner or less conductive, because this would defeat their purpose of shielding the experiment from external fields.
Fortunately, another effect provides mitigation: trapped atom experiments typically measure the net signal from atoms in many different regions of a cell. To the extent that $`n`$ independent regions are averaged, the equivalent dipole noise power will be reduced by a factor of $`1/n`$. Some of the formalism necessary for calculating $`n`$ is set forth in Appendix A.1. Detailed calculations would be strongly dependent on the particular design chosen.
### 3.2. Magnetic resonance force microscopy
Magnetic resonance force microscopy (MRFM) is a quantum-coherent technology whose ultimate objective is to produce magnetic resonance images of individual molecules *in situ*, nondestructively, in three dimensions, with Angstrom resolution. Such a technology would allow much of molecular biology to be conducted as an observational science—along the lines of, *e.g.*, astronomy—rather than an experimental science.
We remark that a single-spin MRFM imaging device can be alternatively regarded as a first-generation solid-state quantum computer, in which the individual qubits carry binary information about the presence or absence of a spin spin at a specified atomic coordinate.
There are two main experimental challenges in MRFM. The first is that single-spin signal forces are exceedingly small, of order $`10^{18}\text{N}`$ for electron magnetic moments. The MRFM community’s design options for achieving the required sensitivity are reasonably well understood (but challenging to implement in practice): reduce the mass of the cantilever, increase its damping time, reduce the ambient temperature, and employ a sharper magnetic tip with a higher field gradient.
The second—emerging—challenge in MRFM is that the spin state must maintain its quantum coherence for a time long enough to detect the signal force. Here the design issues are *not* yet well understood. We will illustrate some of these issues by calculating the effects of thermal magnetic noise on spin relaxation in a typical MRFM environment.
We model the magnetic tip as a sphere with radius $`r=1\mu \text{m}`$ and uniform magnetization $`\mu _0M=1\text{T}`$. The electron magnetic moment is located at a separation $`d=50\text{nm}`$ from the tip surface. At this distance the polarizing field is $`B_0=2\mu _0Mr^3/(3(r+d)^3)=0.58\text{T}`$, and the spin precession frequency is $`\gamma B_0/(2\pi )=16.1\text{GHz}`$, where $`\gamma /(2\pi )=2.8\text{MHz}/T`$ is the gyromagnetic ration of the electron. We take the ambient temperature to be $`T=4\text{K}`$. We note that $`dr`$ and we therefore—roughly—model the tip as a slab of thickness $`t=d`$. Finally, the direction of the polarizing field is taken to be parallel to the vector normal to the tip suface ($`\theta =0`$).
With the parameters $`\{d,t,T,\gamma ,\omega _0,\theta \}`$ now specified (see Section 2.1.1 for further discussion), the electron spin relaxation rates $`1/T_1(\sigma )`$ and $`1/T_2(\sigma )`$ are predicted as functions of the tip conductivity $`\sigma `$ by expressions (2.10a2.10b), with results as shown in Fig. 3. For completeness, relaxation times for proton magnetic moments at the same distance from the tip are also shown.
#### 3.2.1. Design lessons for MRFM
It is clear that tip conductivity is a key engineering parameter in MRFM. For example, if the tip conductivity is as great as that of pure elemental iron ($`4\times 10^9\mathrm{\Omega }\text{-m}`$ at $`4\text{K}`$) or nickel ($`1.7\times 10^{10}\mathrm{\Omega }\text{-m}`$ at $`4\text{K}`$), the predicted relaxation rate is $`1/T_110^2\text{s}^1`$. Such rates are much too rapid to permit coherent spin imaging.
Real-world MRFM tips typically are composed not of pure elements, but of magnetically “hard” alloys such as samarium-cobalt or neodymium-iron-boron. These alloys are commerically prepared as powders, without much regard for minimizing lattice defects: the resulting conductivity is very likely to be substantially reduced relative to the values quoted above for pure elements. There is not much data in the literature relating to the electrical conductivity of ferromagnetic particles at cryogenic temperature. Misra *et al.* have measured the conductivity of sputtered Cu/Cr multilayered films to be $`2.5\times 10^7\mathrm{\Omega }\text{-m}`$ in the zero-temperature limit. If we take this value to be representative of micron-scale ferromagnetic tips at cryogenic temperatures—in the absence of better data—the predicted relaxation rate is $`1/T_110\text{s}^1`$. This is still an uncomfortably rapid relaxation rate.
If these predictions are correct, the MRFM community has at least five design options: (1) reduce the temperature, (2) reduce the tip conductivity (*e.g.*, by employing an insulating magnetic garnet tip), (3) employ superconducting tips (*e.g.*, generating gradients via trapped vortices), (4) detect electron moments incoherently, or (5) focus on coherent proton detection, for which the predicted relaxation times are much longer.
We see how vital it is to achieve a thorough understanding of all the mechanisms that influence spin relaxation. It is sobering, therefore, to reflect upon some of the thermal reservoirs that we have *not* mentioned, which surely are present in ferromagnetic tips: spin waves, thermally excited domain wall motions, and unpaired electron spins in passivating oxide layers, to name three. As discussed in Section 4, much work remains to be done.
### 3.3. Quantum computing
Kane has described a design for a solid-state quantum computer in which the qubits are the electron and nuclear spin quantum numbers of $`{}_{}{}^{31}\text{P}`$ donor sites in a silicon lattice. Now we will consider the effects of thermal magnetic noise within Kane’s device.
In Kane’s design, the $`{}_{}{}^{31}\text{P}`$ donor sites carry both nuclear and electron quantum numbers, which are coupled by a Hamiltonian $`H_{\text{int}}`$ of the form
(3.1)
$$H_{\text{int}}=\gamma _\text{e}𝑩_0\mathbf{}𝒔_\text{e}+\gamma _\text{n}𝑩_0\mathbf{}𝒔_\text{n}+4\mathrm{}^2A𝒔_\text{e}\mathbf{}𝒔_\text{n}$$
Here $`𝒔_\text{e}`$ and $`𝒔_\text{n}`$ are the electron and $`{}_{}{}^{31}\text{P}`$ spin operators, $`\gamma _\text{e}`$ and $`\gamma _\text{n}`$ are the gyromagnetic ratios (in a convention where both ratios are positive), and the hyperfine coupling $`A`$ has the value $`\mathrm{}A/(2\pi )=29\text{MHz}`$.
The computational action occurs between the two lowest energy levels of this system, and the presence of the hyperfine coupling modifies these lowest-energy levels in two respects. First, the energy separation of the two levels depends on $`A`$ via an equation given by Kane as
(3.2)
$$\mathrm{}\omega _0=\gamma _\text{n}B_0+2A+𝒪(A^2)$$
Second, the two lowest energy levels of the system couple to external time-dependent B-fields $`𝑩_{\text{ext}}(t)`$ via an effective Hamiltonian $`H_{\text{eff}}`$ which a straightforward perturbative calculation shows to be
(3.3)
$$H_{\text{eff}}=\frac{\mathrm{}\gamma _\text{n}}{2}𝑩_{\text{ext}}(t)\mathbf{}𝑲\mathbf{}𝝈$$
where $`𝝈\{\sigma _x,\sigma _y,\sigma _z\}`$ is a vector of Pauli matrices acting on the two lowest-energy states, and $`𝑲`$ is a dimensionless matrix given by
(3.4)
$$𝑲=\widehat{𝒃}\mathbf{}\widehat{𝒃}+\left(1+\frac{2A}{\gamma _\text{n}B_0}\right)\left(𝑰\widehat{𝒃}\mathbf{}\widehat{𝒃}\right)+𝒪(A^2)$$
where $`\widehat{𝒃}`$ is a unit vector in the direction of the applied B-field. We observe that the hyperfine coupling $`A`$ acts to increase the coupling of the two lowest-energy states to thermal magnetic noise. This increased coupling is most conveniently taken into account by the replacement in (2.8) of the laboratory-frame spectral density $`𝑺_𝑩`$ by an effective spectral density $`𝑺_𝑩^{\text{eff}}`$
(3.5)
$$𝑺_𝑩𝑺_𝑩^{\text{eff}}𝑲\mathbf{}𝑺_𝑩\mathbf{}𝑲$$
Subsequent calculations of relaxation rates are unaltered. In Kane’s design, the renormalized coupling $`𝑲`$ increases the predicted relaxation rates by a factor of order $`\left(1+(2A)/(\gamma _\text{n}B_0)\right)^27.2`$.
Per Kane’s design , we specify a polarizing field $`B_0=2\text{T}`$ and a temperature $`T=100\text{mK}`$. We model Kane’s A-gates and J-gates as metallic pads of thickness $`t=5\text{nm}`$ at a distance $`d=10\text{nm}`$ from the $`{}_{}{}^{31}\text{P}`$ centers. The direction of the polarizing field is taken to be parallel to the vector normal to the pad surface ($`\theta =0`$). The electron-spin and nuclear-spin relaxation rates $`1/T_1`$ and $`1/T_2`$ are then predicted as functions of the pad conductivity $`\sigma `$ by expressions (2.10a2.10b), with results as shown in Fig. 3.
#### 3.3.1. Design lessons for solid-state quantum computing
Kane’s article optimistically quotes the magnetic resonance literature to the effect
> at temperatures $`T1.5\text{K}`$ the electron spin relaxation time is thousands of seconds, and \[…\] at millikelvin temperatures the phonon limited $`{}_{}{}^{31}\text{P}`$ relaxation time is of the order of $`10^{18}`$ seconds, making this system ideal for quantum computation.
We see that thermal magnetic noise originating in the A- and J-gates is likely to substantially increase these relaxation rates. Fortunately, to the extent that calculation occurs only between the two lowest-energy states, the predicted rapid relaxation of the higher-energy electron states may not pose a problem. And it is gratifying that the predicted $`{}_{}{}^{31}\text{P}`$ relaxation rates remain reasonably small even when the amplifying effects of the hyperfine coupling are taken into account.
These predictions should be regarded very cautiously, however. As we discuss at greater length in the following section, even the notion of conductivity becomes suspect at these length scales and temperatures. In our view, the main design lesson is that all potential noise sources must be treated seriously.
## 4. Conclusions
We have developed a general theory of thermal magnetic fluctuations near the surface of conductive and/or magnetically permeable slabs. We have applied the resulting closed-form expression for the magnetic spectral density (Section 2.1) to practical problems in atomic physics experiments, magnetic resonance force microscopy, and solid-state quantum computing (Sections 3.13.3). The theory predicts magnetic noise levels that are large enough to require substantial care to minimize their effects.
More generally, our formalism indicates that if the dissipative kernel of a linearly damped system is known, then the consequences for fluctuation, dissipation, entanglement, and renormalization are fully determined, via fluctuation-dissipation theorems (eqs. (2.4) and (C.9)), fluctuation-dissipa- tion-entanglement theorems (eqs. (2.12) and (C.10)), and renormalization-dissipation relations (eqs. (C.14) and (C.15)).
We postulate that these fluctuation-dissipation-entanglement-renormaliz-ation relations are *universal*, in the sense that they apply to *all* thermal reservoirs, without regard for the internal structure of the reservoir, with the sole restriction that the coupling to the reservoir must be weak enough to be linearly dissipative.
A major limitation of our article is that we derived our results in the context of a specific model of thermal reservoirs—the independent oscillator (IO) model—and hence we have *not* proved universality. Ford, Lewis, and O’Connell’s pioneering article on independent oscillator models proves universality in the restricted sense that these models are shown to describe all forms of fluctuation and dissipation that are causal and linearly dissipative. An important milestone for further research would be to prove universality in an expanded context which included entanglement and renormalization in addition to fluctuation and dissipation.
There are also substantial reasons to expect that measurement processes might also be subsumed under a suitably expanded notion of universality. We reason as follows. Quantum-coherent technologies necessarily include some means for reading-out binary data; in quantum computing the data are qubits, in MRFM experiments the data are the presence or absense of a spin at a given atomic coordinate. In most cases measurement is a continuous interferometric process, achieved *e.g.* via weak interactions of a cantilever with laser-supplied photons.<sup>3</sup><sup>3</sup>3We remark that the detailed quantum dynamics of the interferometer-cantilever-spin system remains to be worked out. From an engineering point of view, the measurement photons comprise a “good” reservoir—effectively at zero temperature—whose finely tuned interactions compete with many different “bad” thermal reservoirs that contribute only noise. The quantum system seeks to achieve equilibrium with the “good” and “bad” reservoirs simultaneously; the resulting quantum dynamics are traditionally described by master equations of the Fokker-Planck type. A unified formalism of this sort, treating measurement and noise processes on a common footing, would offer many practical advantages in quantum-coherent engineering.
In carrying through these calculations, we have come to appreciate an emerging common ground—the technical challenge of preserving and manipulating quantum coherence—that is shared by magnetic resonance force microscopy and quantum computation. Single-spin MRFM imaging devices (if they can be demonstrated) would constitute first-generation solid-state quantum computers, in which the individual qubits carried binary information about the presence or absence of a spin spin at a specified atomic coordinate.
Arguably the most daunting challenge in developing quantum-coherent technologies is achieving simultaneous engineering control of a panoply of thermal reservoirs:
* phonons and surface vibrational modes
* ferro- and ferrimagnetic spin waves and domain walls
* conduction band excitations in metals and semiconductors
* paramagnetic and nuclear spins
* mobile molecules and charges on surfaces
* lattice defects and impurities
* super-conducting condensates admixed with normal-phase conductors
This list includes most of the excitations that are commonly studied by condensed matter physicists. We remark it is not enough to achieve good engineering control of some or even most of these noise sources; *all* must be controlled simultaneously, as even one poorly-controlled noise source is enough to destroy quantum coherence.
The dissipative kernels of these excitations are poorly understood at the ultra-low temperatures and mesoscopic length scales that are characteristic of quantum-coherent technologies. For example, even Maxwell’s equations in conducting media—an exceedingly well-studied area of physics—will require re-interpretation due to our still-limited understanding of mesoscopic conductivity .
In view of these difficulties, Landauer has modestly proposed that the following disclaimer be appended to all quantum computing proposals:
> This proposal, like all proposals for quantum computation, relies upon speculative technology, does not in its current form take into account all possible sources of noise, unreliability and manufacturing error, and probably will not work.
Our article may be read as a response to Landauer’s concerns: we have analyzed thermal magnetic noise as a “possible source of noise and unreliability” in MRFM imaging and in Kane’s proposal for a quantum computer. More ambitiously, we have done so in the context of a general formalism which might be applied across a broad range of thermal reservoirs. In doing so, we hope to have contributed to the emerging task of providing solid and well-organized foundations for the development of quantum-coherent technologies.
## Appendix A Fluctuation and dissipation
Many readers will be acquainted with the fluctuation-dissipation theorem as it applies to voltage noise—known as Johnson noise—in resistive circuit elements . The theorem can be briefly stated as follows: if a net charge $`q(t)=qe^{i\omega t}`$ is passed through a frequency-dependent impedance $`Z(\omega )`$, such that the current $`i(t)=q(t)/t=i\omega qe^{i\omega t}`$, and the power $`P(\omega )`$ dissipated in the impedance is $`P(\omega )=\text{Re}(Z(\omega ))q^2\omega ^2`$, then the spectrum of thermal voltage fluctuations is given in terms of the dissipative impedance $`\text{Re}(Z)`$ by (2.1).
The general fluctuation-dissipation theorem allows us to deduce, by an exact analogy, that if a time-dependent magnetic moment $`𝒎(t)=𝒎e^{i\omega t}`$ is adjacent to a conductive slab, such that the power dissipated in the slab is $`P=𝒎\mathbf{}𝚪(\omega )\mathbf{}𝒎\omega ^2`$, then the spectrum of the magnetic fluctuations is given in terms of the magnetic dissipation matrix $`𝚪(\omega )`$ via (2.4).
Discussions which justify the above line of reasoning are presented in many graduate-level textbooks , but these discussions tend to be rather abstract and lengthy. Some readers will prefer the following physical derivation, which is completely general and rigorous.
We will deduce the spectrum of magnetic fluctuations from the well-known spectrum of thermal voltage fluctuations (Johnson noise).
Let $`A`$ be the area of the pickup loop in Fig. 2, and let $`\widehat{𝒎}`$ be the unit vector normal to the loop. We let $`I`$ and $`𝒎`$ be the real coefficients of the loop current $`I(t)=Ie^{i\omega t}`$ and loop magnetic moment $`𝒎(t)=𝒎e^{i\omega t}`$, such that $`𝒎=AI\widehat{𝒎}`$. We wish to compute the spectral density $`\widehat{𝒎}\mathbf{}𝑺_𝑩\mathbf{}\widehat{𝒎}`$ of the magnetic field traversing the loop. We reason as follows:<sup>4</sup><sup>4</sup>4We write $`2k_\text{B}T`$ in place of $`\mathrm{}\omega \mathrm{coth}(\mathrm{}\omega /(2k_\text{B}T))`$ purely for economy of notation.
$$\widehat{𝒎}\mathbf{}𝑺_𝑩\mathbf{}\widehat{𝒎}\begin{array}{cc}=S_\mathrm{\Phi }/A^2,\hfill & \text{magnetic flux }\mathrm{\Phi }=A\widehat{𝒎}\mathbf{}𝑩;\hfill \\ =S_V/(\omega A)^2,\hfill & \text{Faraday’s law of Induction;}\hfill \\ =\text{Re}(Z)/(\omega A)^22k_\text{B}T,\hfill & \text{thermal voltage noise per (}\text{2.1}\text{);}\hfill \\ =P/(\omega IA)^22k_\text{B}T,\hfill & \text{because }P=\text{Re}(Z)I^2\text{;}\hfill \\ =𝒎\mathbf{}𝚪\mathbf{}𝒎/(IA)^22k_\text{B}T,\hfill & \text{definition of }𝚪\text{ in terms of }P\text{;}\hfill \\ =\widehat{𝒎}\mathbf{}𝚪\mathbf{}\widehat{𝒎}2k_\text{B}T,\hfill & \text{substitute }𝒎=IA\widehat{𝒎}\text{.}\hfill \end{array}$$
Since the result holds for arbitrary loop orientation $`\widehat{𝒎}`$, and both $`𝑺_𝑩`$ and $`𝚪`$ are symmetric matrices, it must be that $`𝑺_𝑩(\omega )=𝚪(\omega )\mathrm{\hspace{0.17em}2}k_\text{B}T`$, *QED.*
### A.1. Spatial correlations and gradients
This reasoning is readily extended to encompass spatial correlations and gradients. Let $`n`$ magnetic moments $`[𝒎]\{𝒎_1,𝒎_2,\mathrm{},𝒎_n\}`$ be located at coordinates $`\{𝒙_1,𝒙_2,\mathrm{},𝒙_n\}`$. Assuming a time dependence $`e^{i\omega t}`$, the dissipated power is calculable from the geometry and conductivity of the system—at least in principle—and will be of the general form $`P=[𝒎]\mathbf{}\left[𝚪\right]\mathbf{}[𝒎]\omega ^2`$, where $`\left[𝚪\right]`$ is the obvious $`n`$–point generalization of $`𝚪`$. Then the $`n`$–point generalization of (2.4) is simply
(A.1)
$$𝑺_{[𝑩]}=\left[𝚪\right]\mathrm{}\omega \mathrm{coth}\left(\frac{\mathrm{}\omega }{2k_\text{B}T}\right)$$
In this fashion all questions involving $`n`$–point spectral densities of magnetic noise can be reduced to the calculation of $`\left[𝚪\right]`$. Although we will not do so in this article, the spectral density of magnetic gradients can be determined by considering the limiting case in which two coordinates approach each other.
### A.2. Classical and quantum calculations
Students in particular may be troubled that our calculation of the dissipation tensor $`𝚪(\omega )`$ will be entirely classical, via Maxwell’s equations. After all, doesn’t the title of this article promise “the classical *and quantum* theory of thermal magnetic noise”? The quantum part of the promise will be fullfilled in Appendix B, to follow, in which we present a rigorous quantum theory of a spin coupled to a thermal reservoir, and show that the classical dissipative kernel determines the quantum dissipation also.
### A.3. Solving the Maxwell equations
To calculate power dissipation within the slab, and thus determine $`𝚪(\omega )`$, we seek solutions of Maxwell’s equations in their standard form for time dependence $`e^{i\omega t}`$ :
(A.2) $`\mathbf{}\mathbf{\times }𝑬`$ $`=i\omega 𝑩`$ $`\mathbf{}\mathbf{}𝑫`$ $`=\rho _\text{c}`$
$`\mathbf{}\mathbf{\times }𝑯`$ $`=𝒋+i\omega 𝑫`$ $`\mathbf{}\mathbf{}𝑩`$ $`=0`$
together with the constitutive and divergence equations for the current
(A.3) $`𝒋`$ $`=(\sigma 𝑬+𝒋^\text{s})`$ $`\mathbf{}\mathbf{}𝒋`$ $`=i\omega \rho _\text{c}`$
Here $`𝑩`$ is the magnetic field, $`𝑯=𝑩/\mu `$ is the magnetic intensity, $`\mu `$ is the magnetic permeability, $`𝑬`$ is the electric field, $`𝑫=ϵ𝑬`$ is the electric displacement, $`𝒋`$ is the total current, $`ϵ`$ is the electric permittivity, $`\rho _\text{c}`$ is the charge density, $`\sigma `$ is the conductivity, and $`𝒋^\text{s}`$ is the externally-applied dipole source current.
#### A.3.1. The scalar Maxwell equations
We begin our solution of the Maxwell equations with the help of a simplifying *ansatz*: that the irrotational part of $`𝑬`$ and $`𝑫`$ vanishes, *i.e.*, the charge density $`\rho _\text{c}`$ is zero within the slabs and at their interfaces. Furthermore, we will assume that the displacement current $`i\omega 𝑫=i\omega ϵ𝑬`$ is negligibly small compared to the conduction current $`\sigma 𝑬`$. This approximation is quite accurate for all the cases we will consider, *e.g.*, for room temperature copper at one GHz we have $`\omega ϵ/\sigma 10^9`$.<sup>5</sup><sup>5</sup>5We remark that the irrotational *ansatz* and the neglect of displacement currents must be embraced simultaneously, because individually the *ansatz* and the approximation yield a mathematically inconsistent set of Maxwell equations. See, *e.g.*, for a discussion.
It is natural to write a divergence-free $`𝑬`$ field as a curl; we write $`𝑬`$ as
(A.4a)
$$𝑬=i\omega \mathbf{}\mathbf{\times }(\psi \widehat{𝒏})$$
where $`\psi `$ is a scalar vorticity, and we recall that $`\widehat{𝒏}`$ is the (constant) unit vector normal to the slab surfaces. Per our *ansatz*, (A.4a) guarantees that $`\mathbf{}\mathbf{}𝑫=0`$ everywhere within a slab (boundary conditions will be considered shortly). In terms of $`\psi `$, the Maxwell equation $`\mathbf{}\mathbf{\times }𝑬=i\omega 𝑩`$ becomes
(A.4b)
$$𝑩=\widehat{𝒏}(^2\psi )\mathbf{}(\widehat{𝒏}\mathbf{}\mathbf{}\psi )$$
which, furthermore, guarantees $`\mathbf{}\mathbf{}𝑩=0`$. The sole remaining Maxwell equation $`\mathbf{}\mathbf{\times }𝑯=\sigma 𝑬+𝒋^\text{s}`$ is then equivalent to the scalar wave equation
(A.4c)
$$^2\psi =i\mu \sigma \psi +j^\text{s}$$
where $`j^\text{s}`$ is a scalar dipole source whose nature will be specificed shortly.
It remains only to enforce the usual boundary conditions of electrodynamics: continuity of the normal components of $`𝑩`$ and $`𝑫`$ and the transverse components of $`𝑬`$ and $`𝑯`$. These are equivalent to the following boundary conditions on $`\psi `$ at the juncture of two slabs $`i`$ and $`j`$
(A.4d) $`\psi _i`$ $`=\psi _j`$
(A.4e) $`(\widehat{𝒏}\mathbf{}\mathbf{}\psi _i)/\mu _i`$ $`=(\widehat{𝒏}\mathbf{}\mathbf{}\psi _j)/\mu _j`$
The problem is now reduced to the solution of the scalar wave equation (A.4c), subject to the between-slab boundary conditions (A.4d). The $`𝑬`$ and $`𝑩`$ fields are given by (A.4a) and (A.4b), which now are purely constitutive equations.
#### A.3.2. The scalar dipole source
For a (vector) magnetic dipole $`𝒎`$ located within a region of zero conductivity, it is convenient to specify the scalar source $`j^\text{s}`$ implicitly by $`j^\text{s}=^2\psi ^\text{s}`$, with the source field $`\psi ^\text{s}`$ given by
(A.5)
$$\psi ^\text{s}(𝒙)=\frac{\mu }{4\pi }\left(\frac{\widehat{𝒏}}{|𝒙|}+\frac{\widehat{𝒙}(𝑰\widehat{𝒏}\widehat{𝒏})}{|𝒙|+𝒙\widehat{𝒏}}\right)𝒎$$
Inserting $`\psi ^\text{s}`$ in the constitutive equation (A.4b) yields for $`𝑩`$
(A.6)
$$𝑩(𝒙)=\frac{\mu }{4\pi |𝒙|^3}(3\widehat{𝒙}\widehat{𝒙}𝑰)𝒎$$
which (by construction) is the dipole B-field generated by a magnetic moment $`𝒎`$. This justifies our specification of $`\psi ^\text{s}`$ as the source field. The possibility of writing a general dipole field in this form is the key step leading to the success of our scalar *ansatz*.
As a technical point, the vector potential $`\psi ^\text{s}(𝒙)\widehat{𝒏}`$ describes an E-field which is singular along a line extending from $`𝒙=0`$ to infinity in the $`\widehat{𝒏}`$ direction (away from the slab). As we will show in Section A.3.6, the singular portion of the E-field is irrotational, and it can be written as the gradient of a harmonic potential which disappears when we introduce a finite charge density on the surface of the slab. The dissipated power within the slab is not appreciably affected by the currents required to create this surface charge, so for the time being the singularity can simply be ignored.
#### A.3.3. Solving the scalar Maxwell equations
The scalar equation (A.4c) for $`\psi `$ is most easily solved in cylindrical coordinates $`\{r,\phi ,z\}`$. It is further convenient to work with the Bessel transform $`\stackrel{~}{\psi }_m(\rho ,z)`$ defined by<sup>6</sup><sup>6</sup>6The Bessel functions $`J_m`$ satisfy an orthogonality relation
$$_0^{\mathrm{}}r𝑑rJ_m(\rho r)J_m^{}(\rho ^{}r)=\delta _{mm^{}}\delta (\rho \rho ^{})/\rho $$
(A.7)
$$\stackrel{~}{\psi }_m(\rho ,z)=\frac{1}{2\pi }_0^{2\pi }𝑑\phi _0^{\mathrm{}}r𝑑re^{im\phi }J_m(\rho r)\psi (r,\phi ,z)$$
whose inverse is
(A.8)
$$\psi (r,\phi ,z)=\underset{m=\mathrm{}}{\overset{m=\mathrm{}}{}}e^{im\phi }_0^{\mathrm{}}𝑑\rho \rho J_m(\rho r)\stackrel{~}{\psi }_m(\rho ,z)$$
Here $`J_m`$ is a Bessel function of integer argument. The scalar wave equation for $`\stackrel{~}{\psi }_m(\rho ,z)`$ is
(A.9)
$$\frac{^2}{z^2}\stackrel{~}{\psi }_m(\rho ,z)=i\omega \mu \sigma \stackrel{~}{\psi }_m(\rho ,z)+\stackrel{~}{j}_m^\text{s}(\rho ,z)$$
whose general homogenous solution is
(A.10)
$$\stackrel{~}{\psi }_m(\rho ,z)=\stackrel{~}{\psi }_m^+(\rho )e^{kz}+\stackrel{~}{\psi }_m^{}(\rho )e^{kz}$$
where we have introduced the complex wavenumber $`k\sqrt{i\omega \mu \sigma +\rho ^2}`$.
Our next task is to explicitly solve for the Bessel coefficients $`\stackrel{~}{\psi }_m(\rho )`$. We divide space into three regions {I,II,III}, where Region I is the nonconductive region that contains the dipole source, Region II is the conductive slab of thickness $`t`$ beginning at a distance $`d`$ from the source, and Region III is the nonconducting region on the opposite side of the slab (see Fig. 1).
Within Regions I and III we assume vanishing conductivity $`\sigma =0`$ and a purely real permeability $`\mu _0`$. In contrast, within Region II we allow both $`\sigma `$ and $`\mu `$ to assume complex values. Thus energy dissipation can occur only within Region II.
Requiring that $`\stackrel{~}{\psi }_m(\rho ,z)0`$ as $`|z|\mathrm{}`$ leads to the general solution:
(A.11)
$$\stackrel{~}{\psi }_m(\rho ,z)=\{\begin{array}{ccc}\stackrel{~}{\psi }_m^{\text{I}+}(\rho )e^{\rho z}+\stackrel{~}{\psi }_m^\text{s}(\rho ,z)& & \text{Region I (source);}\hfill \\ & |& z=d\hfill \\ \stackrel{~}{\psi }_m^{\text{II}+}(\rho )e^{kz}+\stackrel{~}{\psi }_m^{\text{II}}(\rho )e^{kz}& +z& \text{Region II (slab);}\hfill \\ & & z=d+t\hfill \\ \stackrel{~}{\psi }_m^{\text{III}}(\rho )e^{\rho z}& & \text{Region III (far side).}\hfill \end{array}$$
For illustrations of Regions I, II, and III, see Figs. 12. The source term $`\stackrel{~}{\psi }_m^\text{s}(\rho ,z)`$ in Region I is found by substituting (A.5) in (A.7):
(A.12) $`\stackrel{~}{\psi }_0^\text{s}(\rho ,z)`$ $`={\displaystyle \frac{\mu _0}{4\pi \rho }}e^{\rho |z|}\widehat{𝒏}\mathbf{}𝒎`$
(A.13) $`\stackrel{~}{\psi }_{\pm 1}^\text{s}(\rho ,z)`$ $`={\displaystyle \frac{\mu _0}{8\pi \rho }}e^{\rho |z|}(\widehat{𝒙}\pm i\widehat{𝒚})\mathbf{}𝒎`$
where we have introduced Cartesian basis vectors $`\widehat{𝒙}\mathbf{\times }\widehat{𝒚}=\widehat{𝒏}`$. Note the singularity at $`z=0`$; this reflects our assumption of a point dipole source. Section A.3.6 derives explicitly nonsingular expressions for $`𝑩`$ and $`𝑬`$ originating from finite-sized current sources, however the present dipole result is adequate for our main purpose of estimating power dissipation in the slab.
With the source term now specified, the boundary equations (A.4d) for the Bessel coefficients $`\{\stackrel{~}{\psi }_0^{\text{I}+},\stackrel{~}{\psi }_0^{\text{II}+},\stackrel{~}{\psi }_0^{\text{II}},\stackrel{~}{\psi }_0^{\text{III}}\}`$ can be readily solved. We obtain for the $`m=0`$ coefficients:
(A.14)
$$\begin{array}{c}\left[\begin{array}{c}\stackrel{~}{\psi }_0^{\text{I}+}\\ \stackrel{~}{\psi }_0^{\text{II}+}\\ \stackrel{~}{\psi }_0^{\text{II}}\\ \stackrel{~}{\psi }_0^{\text{III}}\end{array}\right]=\frac{\mu _0}{4\pi }\left[\begin{array}{c}e^{2\rho d}(e^{2kt}1)(K^2\rho ^2k^2)\\ 2e^{(\rho +k)d}(K\rho k)K\rho \\ 2e^{2kt(\rho k)d}(K\rho +k)K\rho \\ 4e^{(\rho +k)}kK\rho \end{array}\right]\times \hfill \\ \hfill \frac{\widehat{𝒏}\mathbf{}𝒎}{\rho \left((K^2\rho ^2+k^2)(e^{2kt}1)+2K\rho k(e^{2kt}+1)\right)}\end{array}$$
where $`K\mu /\mu _0`$ is the relative permeability of the slab. The $`m=\pm 1`$ coefficients are directly proportional to the $`m=0`$ coefficients
(A.15) $`\left[\begin{array}{c}\stackrel{~}{\psi }_{\pm 1}^{\text{I}+}\\ \mathrm{}\end{array}\right]`$ $`={\displaystyle \frac{(\widehat{𝒙}\pm i\widehat{𝒚})\mathbf{}𝒎}{2(\widehat{𝒏}\mathbf{}𝒎)}}\left[\begin{array}{c}\stackrel{~}{\psi }_0^{\text{I}+}\\ \mathrm{}\end{array}\right]`$
When we calculate magnetic spectral densities, this simple proportionality will ensure that $`𝚪(𝑰+\widehat{𝒏}\mathbf{}\widehat{𝒏})`$ at all frequencies and length scales.
#### A.3.4. Energy dissipation
Our next task is to calculate the power $`P`$ dissipated within Region II. From classical electrodynamics we have
(A.16) $`P`$ $`={\displaystyle _V}𝑑V\left(𝒋\mathbf{}𝑬+𝑯\mathbf{}{\displaystyle \frac{𝑩}{t}}\right)`$
which with the help of the Bessel orthogonality relations becomes
(A.17) $`P`$ $`=2\pi {\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}\rho 𝑑\rho {\displaystyle _d^{d+t}}𝑑z\left(\rho ^2\omega ^2\text{Re}(\sigma )+\rho ^4\omega {\displaystyle \frac{\text{Im}(\mu )}{|\mu |^2}}\right)|\psi _m(\rho ,z)|^2`$
Carrying through the $`z`$-integration, we find the scalar magnetic dissipation coefficient $`\mathrm{\Gamma }(\omega )`$, as implicitly defined by $`P\omega ^2\mathrm{\Gamma }(\omega )\left(𝒎\mathbf{}(𝑰+\widehat{𝒏}\mathbf{}\widehat{𝒏})\mathbf{}𝒎\right)`$ :
(A.18)
$$\begin{array}{c}\mathrm{\Gamma }(\omega )=\frac{\mu _0^2}{4\pi }_0^{\mathrm{}}d\rho (\rho ^3\text{Re}(\sigma )+\frac{\rho ^5\text{Im}(\mu )}{\omega |\mu |^2})|K|^2e^{2\rho d}\times \hfill \\ \hfill \frac{\text{Re}\left[ik^{}(e^{2kt}1)(e^{2k^{}t}+1)(K^2\rho ^2k^2)/(\omega \mu \sigma )+\rho K|e^{2kt}1|^2\right]}{\left|(K^2\rho ^2+k^2)(e^{2kt}1)+2K\rho k(e^{2kt}+1)\right|^2}\end{array}$$
where we recall that $`k\sqrt{i\omega \mu \sigma +\rho ^2}`$ and $`K\mu /\mu _0`$. In simplifying this result we have kept only the leading and next-to-leading terms in $`\text{Im}(\mu )`$, *i.e.*, we have assumed $`K`$ is predominantly real.
The case of thermal magnetic fluctuations measured at a point midway between two identical conducting slabs can be solved by exactly similar methods. It is only necessary to replace the Region I homogenous term $`\stackrel{~}{\psi }_m^{\text{I}+}(\rho )e^{\rho z}`$ in (A.11) with the symmetrized term $`\stackrel{~}{\psi }_m^{\text{I}+}(\rho )\mathrm{cosh}(\rho z)`$. The resulting magnetic dissipation coefficient is
(A.19)
$$\begin{array}{c}\mathrm{\Gamma }^{}(\omega )=\frac{\mu _0^2}{4\pi }_0^{\mathrm{}}d\rho (\rho ^3\text{Re}(\sigma )+\frac{\rho ^5\text{Im}(\mu )}{\omega |\mu |^2})|K|^2e^{2\rho d}\times \hfill \\ \hfill \frac{2\text{Re}\left[ik^{}(e^{2kt}1)(e^{2k^{}t}+1)(K^2\rho ^2k^2)/(\omega \mu \sigma )+\rho K|e^{2kt}1|^2\right]}{\left|(k^2+K^2\rho ^2+e^{2\rho d}(k^2K^2\rho ^2))(e^{2kt}1)+2K\rho k(e^{2kt}+1)\right|^2}\end{array}$$
Note that the double-slab integrand is just twice the single-slab integrand, slightly modified by an additional term $`e^{2\rho d}`$ in the denominator. For the special case $`K=1`$ (the case of primary interest) this additional term can be shown to be non-leading for all values of $`\{\rho ,d,t,\sigma ,\omega \}`$. A numerical survey confirms that $`\mathrm{\Gamma }^{}(\omega )2\mathrm{\Gamma }(\omega )`$ to within very good accuracy, as described in Section 2.1.
#### A.3.5. Asymptotic limits for paramagnetic and diamagnetic slabs
The three asymptotic expressions in (2.5) are obtained by evaluating (A.18) for $`K=1`$ and $`\text{Im}(\mu )=0`$ in the following limits:
(A.20)
$$\mathrm{\Gamma }(\omega )\text{Re}(\sigma )\frac{\mu _0^2}{4\pi }_0^{\mathrm{}}𝑑\rho \times \{\begin{array}{cc}\frac{(1e^{2\rho t})}{8e^{2\rho d}},\hfill & |\sigma |0\text{;}\hfill \\ \frac{\rho ^3\lambda ^3e^{2\rho d}}{2\mathrm{cos}(\pi /4\varphi /2)},\hfill & |\sigma |\mathrm{}\text{;}\hfill \\ \frac{\rho ^3\lambda ^4(t4\rho \lambda ^2\mathrm{sin}\varphi )}{t^2e^{2\rho d}},\hfill & kt1\rho d.\hfill \end{array}$$
The integrations are readily carried out, with the results as in (2.5).
#### A.3.6. The irrotational electric field
Now we have only one chore remaining: repairing the irrotational singularity in the E-field mentioned in Section A.3.2, and verifying that the repair does not alter the energy dissipation just calculated. In the end, we will obtain expressions for the E- and B-fields which are closed-form and explicitly finite.
In outline, the repair is accomplished by adding a purely irrotational E-field to Region I—thereby relaxing the *ansatz* of Section A.3.1 that $`𝑬`$ is purely rotational.<sup>7</sup><sup>7</sup>7We recall *Helmholtz’s theorem*: that any finite vector field that vanishes at infinity can be written uniquely as the sum of an irrotational part (the divergence of a scalar potential) and a rotational part (the curl of a vector potential) . With the added E-field, Faraday’s law $`\mathbf{}\mathbf{\times }𝑬=i\omega 𝑩`$ is still satisfied, because the added E-field is irrotational. Because the added field is confined to Region I, which is nonconductive, and because displacement currents $`i\omega 𝑫`$ are negligibly small, Ampere’s law $`\mathbf{}\mathbf{\times }𝑯=\sigma 𝑬+𝒋^\text{s}+i\omega 𝑫𝒋^\text{s}`$ is still satisfied. Gauss’s law $`\mathbf{}\mathbf{}𝑫=0`$ is still satisfied because the added E-field is the gradient of a harmonic potential. Physically, the added field is generated by surface charge on the slab, and the within-slab currents necessary to sustain the surface charge are negligibly small compared to the eddy currents already computed. Thus, the added E-field does not alter the magnetic dissipation coefficient $`𝚪(\omega )`$ already computed. The remainder of this section is devoted to proving these assertions.
We begin by noting that the E-field singularity in Region I physically corresponds to the field induced by an end-to-end line of electric dipoles along the $`\widehat{𝒏}`$ axis. It is easy to verify that such a singularity can be cancelled by the addition of an (irrotational) E-field which is the gradient of the harmonic potential $`\varphi _1^\text{I}(𝒙)`$,
(A.21)
$$\varphi _1^\text{I}(𝒙)=i\omega \frac{\mu _0}{4\pi }\frac{\widehat{𝒙}\mathbf{}(\widehat{𝒏}\mathbf{\times }𝒎)}{|𝒙|+𝒙\widehat{𝒏}}$$
To satisfy the boundary conditions between Regions I and II, we must also add the E-field of the image potential
(A.22)
$$\varphi _2^\text{I}(𝒙)\varphi _1^\text{I}(2d\widehat{𝒏}𝒙)$$
in order that $`𝑬^\text{I}=\mathbf{}(\varphi _1^\text{I}+\varphi _2^\text{I})`$ be normal to the slab surface, thus ensuring that the transverse component of $`𝑬`$ remains continuous at the boundary.
To sustain the $`𝑬^\text{I}`$, a charge density $`\rho ^\text{I}=ϵ\widehat{𝒏}\mathbf{}𝑬^\text{I}`$ must exist on the surface of the slab, and this surface charge must in turn be created by a normal-to-surface current within the slab $`𝒋^\text{I}=i\omega \rho ^\text{I}\widehat{𝒏}`$. Relative to the intraslab eddy currents current $`\sigma 𝑬`$ which are our main concern, it is easy to show that $`𝒋^\text{I}`$ is of order $`\omega ϵ/\sigma 10^9`$ (see the discussion following (A.2)) and thus is negligible.
Putting all these pieces together, we are now ready to specify the fields in Region I in a manifestly finite manner, assuming only that the spatial size of the current source is reasonably small compared to the distance $`d`$ separating the source from the slab (like the current loop in Fig. 2).<sup>8</sup><sup>8</sup>8Also, we continue to ignore the displacement current $`i\omega 𝑫`$ in Ampere’s law $`\mathbf{}\mathbf{\times }𝑯=𝒋+i\omega 𝑫`$; this is equivalent to the near-field assumption $`\omega dc`$, where $`c=(\mu _0ϵ_0)^{1/2}`$ is the speed of light. This near-field approximation is very well satisfied for most practical problems in quantum-coherent engineering.
From the localized current source, we compute a vector potential $`𝑨(𝒙)`$ in the usual manner
(A.23)
$$𝑨(𝒙)=\frac{\mu _0}{4\pi }_Vd^3𝒙^{\mathbf{}}\frac{𝒋(𝒙^{\mathbf{}})}{|𝒙𝒙^{\mathbf{}}|}\frac{\mu _0}{4\pi }\frac{𝒎\mathbf{\times }\widehat{𝒙}}{|𝒙|^2}$$
where $`𝒋(𝒙^{\mathbf{}})`$ is the (finite) current density within the (finite) source. From the far-field limit of $`𝑨(𝒙)`$ we determine the dipole moment $`𝒎`$ of the source, and under the assumption that the size of the source is small compared to $`d`$, we use $`𝒎`$ to compute the back-reaction vorticity field $`\psi ^\text{I}(𝒙)`$ via equations (A.8A.14) and the back-reaction irrotational field $`\varphi _2^\text{I}`$ via (A.22). With the help of the identity $`i\omega 𝑨(𝒙)=i\omega \mathbf{}\mathbf{\times }(\psi ^\text{s}(𝒙))+\mathbf{}\varphi _1^\text{I}(𝒙)`$, which is exact for a point dipole current source, the net B-field and E-field in Region I can be written as follows
(A.24a) $`𝑩(𝒙)`$ $`=\mathbf{}\mathbf{\times }𝑨(𝒙)\mathbf{}(\widehat{𝒏}\mathbf{}\mathbf{}\psi ^\text{I}(𝒙))`$
(A.24b) $`𝑬(𝒙)`$ $`=i\omega 𝑨(𝒙)i\omega \mathbf{}\mathbf{\times }(\psi ^\text{I}(𝒙)\widehat{𝒏})+\mathbf{}\varphi _2^\text{I}(𝒙)`$
where $`𝑨(𝒙)`$, $`\psi ^\text{I}(𝒙)`$, and $`\varphi _2^\text{I}(𝒙)`$ are individually nonsingular.
Physically speaking, $`𝑨(𝒙)`$ describes the E- and B-fields that are generated directly by the source current, $`\psi ^\text{I}(𝒙)`$ describes the E- and B-fields generated by the induced currents within the slab (treating the source current as a dipole), and $`\varphi _2^\text{I}(𝒙)`$ describes the E-field generated by surface charges on the slab (also treating the source current as a dipole). These expressions were used to generate the field geometry in Fig. 2.
## Appendix B Quantum decoherence
In this section we present an exactly solvable microscopic model of a spin-$`\frac{1}{2}`$ particle magnetically coupled to a heat bath. Our goal is derive the expressions for $`T_1`$, $`T_2`$, and $`T_{1\rho }`$ quoted in (2.9a)-(2.9c).
We follow Ford, Lewis, and O’Connell in modeling the heat bath as a collection of independent harmonic oscillators ; the reader is referred to this article for a thermodynamic justification of this model. We will adjust the bath parameters so as to reproduce the known spectral density of magnetic fluctuations, which we computed in the previous section. The quantum Langevin and optical Bloch equations—which yield closed-form expressions for $`T_1`$, $`T_2`$, and $`T_{1\rho }`$—will then be uniquely determined.
### B.1. The thermal reservoir Hamiltonian
The Hamiltonian $`H`$ of the spin/heat bath system is taken to be
(B.1)
$$H=\gamma 𝑩_0\mathbf{}𝒔+\underset{j}{}\frac{1}{2}\omega _j\left(p_j^2+(q_j\beta _j\widehat{𝒏}_j\mathbf{}𝒔)^2\right)$$
Here $`𝑩_0`$ is the polarizing field, $`\gamma `$ is the gyromagnetic ratio of the spin, and $`𝒔`$ is the angular momentum operator of the spin, satisfying commutation relations $`[s_i,s_j]=i\mathrm{}ϵ_{ijk}s_k`$. The heat bath is described in terms of oscillators with resonant frequency $`\omega _i`$, whose dynamical coodinates $`\{p_i,q_i\}`$ satisfy $`[q_i,p_j]=i\mathrm{}\delta _{ij}`$, which are coupled to the spin with strength $`\beta _j`$.
The physical nature of the heat bath variables need not be otherwise specified. Reader may optionally conceive them as, *e.g.*, thermal excitations of conduction band electrons, magnons in a ferromagnet, excited states of the vacuum in field theory, or in general as excitations of whatever heat bath is furnished by a given quantum technology.
We begin by solving the equation of motion of the heat bath variables. In the Heisenberg picture we have<sup>9</sup><sup>9</sup>9We recall that the Heisenberg equation of motion for an arbitrary operator $`O(t)`$ is $`\dot{O}(t)=[O,H]/(i\mathrm{})`$. The results of this section all follow directly from this equation, using only the commutation relations for $`\{q_i,p_i,𝒔\}`$, plus the fact that commutators in the Heisenberg representation are the same as those in the Schroedinger representation.
(B.2)
$$\ddot{q}_j+\omega _j^2q_j=\beta _j\omega _j^2\widehat{𝒏}_j\mathbf{}𝒔$$
Following Ford *et al.*, we write the formal solution to these equations as
(B.3)
$$q_j(t)=q_j^\text{h}(t)+\beta _j\widehat{𝒏}_j\mathbf{}𝒔(t)_{\mathrm{}}^t𝑑t^{}\mathrm{cos}(\omega _j(tt^{}))\beta _j\widehat{𝒏}_j\mathbf{}\dot{𝒔}(t^{})$$
Here $`q_j^\text{h}(t)`$ is a homogenous solution to the equations of motion (corresponding physically to the evolution of the heat bath in the absence of the back-action of the spin). With the help of this result, the spin operator equation of motion $`\dot{𝒔}=[𝒔,H]/(i\mathrm{})`$ can be readily cast into the form of the following *quantum Langevin equation*<sup>10</sup><sup>10</sup>10As a help to students, we note that the quantum Langevin equation can be written in several equivalent forms, which arise because the commutation relation $`[q_j(t),𝒔(t)]=0`$ allows the heat bath interaction to be written in any of the following fully equivalent forms
$$q_j(t)𝒔(t)=𝒔(t)q_j(t)=\frac{1}{2}(q_j(t)𝒔(t)+𝒔(t)q_j(t))$$
Upon substituting (B.3) for $`q_j(t)`$, these various forms yield fully equivalent (but superficially quite different) Langevin equations. We have chosen to work with the fully symmetrized Langevin equation; none of our main results depend on this choice.
(B.4) $`\dot{𝒔}(t)`$ $`=\gamma 𝑩_0\mathbf{\times }𝒔(t)`$
$`{\displaystyle \frac{1}{2}}\gamma \left(𝑩(t)\mathbf{\times }𝒔(t)𝒔(t)\mathbf{\times }𝑩(t)\right)+`$
$`{\displaystyle \frac{1}{2}}\left[\left(𝑪\mathbf{}𝒔(t)\right)\mathbf{\times }𝒔(t)𝒔(t)\mathbf{\times }\left(𝑪\mathbf{}𝒔(t)\right)\right]+`$
$`{\displaystyle _{\mathrm{}}^t}𝑑t^{}{\displaystyle \frac{1}{2}}\left[\left(𝑮(tt^{})\mathbf{}\dot{𝒔}(t^{})\right)\mathbf{\times }𝒔(t)𝒔(t)\mathbf{\times }\left(𝑮(tt^{})\mathbf{}\dot{𝒔}(t^{})\right)\right]`$
Here $`𝑩(t)`$ is a fluctuating thermal magnetic field which depends only on the homogenous heat bath operators $`q_j^\text{h}(t)`$
(B.5) $`𝑩(t)`$ $`={\displaystyle \underset{j}{}}\gamma ^1\omega _j\beta _jq_j^\text{h}(t)\widehat{𝒏}_j`$
and $`𝑮(tt^{})`$ is a dissipative kernel
(B.6) $`𝑮(tt^{})`$ $`=\{\begin{array}{cc}_j\omega _j\beta _j^2\mathrm{cos}(\omega _j(tt^{}))\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j\hfill & t>t^{}\hfill \\ 0\hfill & t<t^{}\hfill \end{array}`$
Note that $`𝑮(tt^{})`$ is a pure c-number matrix which is independent of the temperature of the heat bath. It is the matrix generalization of what Ford *et al.* call the “memory function” of the heat bath.
The *anisotropy tensor* $`𝑪`$ is a symmetric c-number tensor given by
(B.7)
$$𝑪=\underset{j}{}\omega _j\beta _j^2\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j$$
Physically speaking, $`𝑪`$ describes the interaction of the particle with image currents in the nearby slab. This reflects a very general and physically realistic property of the independent oscillator model: the dynamical equations of a particle are *renormalized* by its interaction with the thermal reservoir.
For the special case of spin-$`\frac{1}{2}`$ particles $`𝑪`$ has no dynamical effects, due to an identity satisfied by spin-$`\frac{1}{2}`$ operators (and not by higher-spin particles)
(B.8)
$$s_is_j+s_js_i=\frac{\mathrm{}^2}{2}𝕀$$
where $`𝕀`$ is the identity operator. This identify ensures that the $`𝑪`$-dependent terms in (B.4) vanish identically; we emphasize that this occurs only for spin-$`\frac{1}{2}`$ particles.
So far, we have not used the fact that the heat bath variables are in thermal equilibrium; we now use this fact to compute the spectral density $`𝑺_𝑩`$ of $`𝑩(t)`$. Not surprisingly, we will discover that $`𝑺_𝑩`$ has a functional form that is closely related to the dissipative kernel $`𝑮`$.
The calculation is straightforward. Denoting a thermal ensemble average by $`\mathrm{}_t`$, we follow Ford *et al.* in noting that the correlation of the homogenous heat bath operators $`\{q_i^\text{h}\}`$ is of the simple form
(B.9)
$$q_i^\text{h}(t)q_j^\text{h}(t^{})_t=\frac{1}{2}\delta _{ij}\mathrm{}\mathrm{coth}\left(\frac{\mathrm{}\omega _j}{2k_\text{B}T}\right)\mathrm{cos}(\omega _j(tt^{}))$$
Combining this with (B.5) yields an explicit expression for $`𝑺_𝑩(\omega )`$
(B.10) $`𝑺_𝑩(\omega )`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau e^{i\omega \tau }{\displaystyle \frac{1}{2}}𝑩(t)\mathbf{}𝑩(t+\tau )+𝑩(t+\tau )\mathbf{}𝑩(t)_t`$
$`={\displaystyle \underset{j}{}}{\displaystyle \frac{\omega _j^2\beta _j^2}{2\gamma ^2}}\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j\mathrm{}\mathrm{coth}\left({\displaystyle \frac{\mathrm{}\omega }{2k_\text{B}T}}\right){\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}e^{i\omega \tau }\mathrm{cos}(\omega _j\tau )`$
$`={\displaystyle \underset{j}{}}{\displaystyle \frac{\pi \omega _j\beta _j^2}{2\gamma ^2}}\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j\mathrm{}\omega \mathrm{coth}\left({\displaystyle \frac{\mathrm{}\omega }{2k_\text{B}T}}\right)\left(\delta (\omega \omega _j)+\delta (\omega +\omega _j)\right)`$
If we compare this result with the real part of the Fourier transform of the dissipative kernel
(B.11) $`\text{Re}\stackrel{\mathbf{~}}{𝑮}(\omega )`$ $`=\text{Re}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑t^{}e^{i\omega (tt^{})}𝑮(tt^{})`$
$`={\displaystyle \frac{\pi }{2}}{\displaystyle \underset{j}{}}\omega _j\beta _j^2\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j\left(\delta (\omega \omega _j)+\delta (\omega +\omega _j)\right)`$
it is apparent that
(B.12)
$$𝑺_𝑩(\omega )=\gamma ^2\text{Re}\stackrel{\mathbf{~}}{𝑮}(\omega )\mathrm{}\omega \mathrm{coth}\left(\frac{\mathrm{}\omega }{2k_\text{B}T}\right)$$
In terms of the magnetic dissipation matrix $`𝚪(\omega )`$ computed in the previous section, this takes an even simpler form
(B.13)
$$\gamma ^2𝚪(\omega )=\text{Re}\stackrel{\mathbf{~}}{𝑮}(\omega )$$
This is the fluctuation-dissipation theorem as it applies to particles of arbitrary spin. For our purposes an alternative relation is even more useful: if $`𝚪(\omega )`$ is known, then the dissipative kernel $`𝑮(tt^{})`$ is given explicitly by
(B.14)
$$𝑮(tt^{})=\{\begin{array}{cc}\frac{\gamma ^2}{\pi }_{\mathrm{}}^{\mathrm{}}𝑑\omega e^{i\omega (tt^{})}𝚪(\omega ),\hfill & t>t^{};\hfill \\ 0,\hfill & t<t^{};\hfill \end{array}$$
as may be verified by writing $`𝚪(\omega )`$ in terms of $`𝑺_𝑩(\omega )`$ using (2.4), then writing $`𝑺_𝑩(\omega )`$ in terms of heat bath variables using (B.10), then carrying out the integration and comparing with (B.6).<sup>11</sup><sup>11</sup>11 It may occur to the reader that by virtue of (B.13) and (B.14), knowledge of $`\text{Re}\stackrel{~}{𝑮}(\omega )`$ suffices to determine $`\text{Im}\stackrel{~}{𝑮}(\omega )`$. This insight leads directly to a *Kramers-Kroenig* relation: $`\text{Im}\stackrel{~}{𝑮}(\omega )`$ $`=\text{Im}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\tau e^{i\omega \tau }𝑮(\tau )`$ $`=\text{Im}{\displaystyle _0^{\mathrm{}}}𝑑\tau {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega ^{}e^{i(\omega ^{}\omega )\tau }{\displaystyle \frac{\gamma ^2}{\pi }}𝚪(\omega ^{})`$ $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\tau {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega ^{}\mathrm{sin}\left((\omega ^{}\omega )\tau \right)\text{Re}\stackrel{~}{𝑮}(\omega ^{})`$ which is the Kramers-Kroenig relation for particles of arbitrary spin interacting with thermal magnetic noise. Most textbooks present such relations in a more compact but less obviously finite form, which can be obtained by rewriting the above result as $`={\displaystyle \frac{1}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\tau {\displaystyle \frac{d}{d\tau }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega ^{}\mathrm{cos}\left((\omega ^{}\omega )\tau \right){\displaystyle \frac{\text{Re}\left(\stackrel{~}{𝑮}(\omega ^{})\stackrel{~}{𝑮}(\omega )\right)}{\omega \omega ^{}}}`$ where the $`\stackrel{~}{𝑮}(\omega )`$ term is introduced to regulate the singularity at $`\omega ^{}=\omega `$, and the $`\omega ^{}`$ integration contour is adjusted in the complex plane so that the net contribution of this added term is zero. The $`\tau `$ integration then becomes trivial and the remaining $`\omega ^{}`$ integration is the Kramers-Kroenig relation in its traditional form (B.15) $`\text{Im}\stackrel{~}{𝑮}(\omega )`$ $`={\displaystyle \frac{1}{\pi }}P{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega ^{}{\displaystyle \frac{\text{Re}\stackrel{~}{𝑮}(\omega ^{})}{\omega \omega ^{}}}`$ where $`P`$ is a principal value.
### B.2. The master equation
We now direct our attention to solving the quantum Langevin equation (B.4). As it stands, this equation is too complicated to solve directly.<sup>12</sup><sup>12</sup>12Students (and engineers) may wish to reflect that the quantum Langevin equation (B.4) is easy to solve in principle—it can be numerically integrated quite readily (the integration is quite straightforward to program in languages like *Mathematica*; this is a good exercise for a student). Heisenberg operators are be represented numerically by matrices of complex numbers, with operator addition and multiplication achieved via ordinary matrix addition and multiplication. As time goes on, the quantum Langevin equation increaingly entangles the spin operator $`𝒔(t)`$ with the thermal reservoir basis states—thus leading, as expected, to fluctuation, dissipation, and entanglement of the spin operators. The only practical difficulty with this program is that a reservoir of $`n`$ particles, each having $`m`$ quantum states, requires $`m^n`$ quantum basis states, with each Heisenberg operator stored as a time-dependent $`m^n\times m^n`$ Hermitian matrix. For $`n`$ as small as two and $`m`$ as small as twenty—too small to describe a realistic thermal reservoir—the storage and multiplication of these exponentially large matrices is enough to overwhelm *any* classical computer. Quantum computers were invented in part to overcome this fundamental *quantum simulation problem*. In particular, $`B(t)`$ is an operator-valued function of heat bath variables whose detailed dynamical behavior—describing every minute fluctuation within the thermal reservoir—we neither know nor wish to know. To make progress, we will average over a thermodynamic ensemble of heat bath variables, and study only a coarse-grained time derivative of $`𝒔(t)`$. Simplified equations obtained by this general strategy are called *master equations*—our goal is to derive a master equation from the quantum Langevin equation.
We begin by defining a rotation matrix $`𝑹(t)`$, satisfying $`\dot{𝑹}(t)=\gamma 𝑩_0\mathbf{\times }𝑹(t)`$, in terms of which we define rotating-frame quantities
$`𝒔_{\text{rot}}(t)`$ $`𝑹(t)\mathbf{}𝒔(t)`$
$`𝑩_{\text{rot}}(t)`$ $`𝑹(t)\mathbf{}𝑩(t)`$
$`𝑮_{\text{rot}}(tt^{})`$ $`𝑹(t)\mathbf{}𝑮(tt^{})\mathbf{}𝑹^{}(t)`$
It follows that $`\dot{𝒔}(t)=𝑹^{}(t)\mathbf{}(\dot{𝒔}_{\text{rot}}(t)\gamma 𝑩_0\mathbf{\times }𝒔_{\text{rot}}(t))`$.
Neglecting $`𝒪(\dot{𝒔}_{\text{rot}}^2)`$ terms—this is the rotating-frame approximation—and specializing to spin-$`\frac{1}{2}`$ particles so that the anisotropy tensor $`𝑪`$ does not enter, the quantum Langevin equation (B.4) becomes
(B.16)
$$\begin{array}{c}\dot{𝒔}_{\text{rot}}(t)=\frac{1}{2}\gamma \left(𝑩_{\text{rot}}(t)\mathbf{\times }𝒔_{\text{rot}}(t)𝒔_{\text{rot}}(t)\mathbf{\times }𝑩_{\text{rot}}(t)\right)+\hfill \\ \hfill \frac{1}{2}_{\mathrm{}}^tdt^{}((𝑮_{\text{rot}}(tt^{})\mathbf{}(\gamma 𝑩_0\mathbf{\times }𝒔_{\text{rot}}(t^{})))\mathbf{\times }𝒔_{\text{rot}}(t^{})\\ \hfill 𝒔_{\text{rot}}(t^{})\mathbf{\times }(𝑮_{\text{rot}}(tt^{})\mathbf{}(\gamma 𝑩_0\mathbf{\times }𝒔_{\text{rot}}(t^{}))))\end{array}$$
Recalling that $`𝑮_{\text{rot}}(\tau )`$ is the temperature-independent “memory function” of the heat bath, we now assume it decorrelates rapidly compared to the rate at which $`𝒔_{\text{rot}}(t)`$ varies. We can then ignore the time-dependence of $`𝒔_{\text{rot}}(t)`$ inside of integrals, and make the following substitution
(B.17)
$$_{\mathrm{}}^t𝑑t^{}𝑮_{\text{rot}}(tt^{})=\gamma ^2𝚪_{\text{rot}}(0)$$
Here the rotating-frame magnetic dissipation matrix $`𝚪_{\text{rot}}(\omega )`$ is related to the laboratory-frame matrix $`𝚪(\omega )`$ via
(B.18)
$$\begin{array}{c}𝚪_{\text{rot}}(\omega )=\left[\widehat{𝒃}\mathbf{}\widehat{𝒃}\right]\text{tr}\left[\widehat{𝒃}\mathbf{}\widehat{𝒃}\mathbf{}𝚪(\omega )\right]+\hfill \\ \hfill \left[𝑰\widehat{𝒃}\mathbf{}\widehat{𝒃}\right]\frac{1}{2}\text{tr}\left[(𝑰\widehat{𝒃}\mathbf{}\widehat{𝒃})\mathbf{}𝚪(\omega _0)\right]\end{array}$$
with $`\gamma 𝑩_0\omega _0\widehat{𝒃}`$. Similarly, $`𝑺_{𝑩_{\text{rot}}}(\omega )`$ is related to $`𝑺_𝑩(\omega )`$ via (2.8). Then, with the help of some straightforward algebraic manipulations,<sup>13</sup><sup>13</sup>13The derivation involves integrating (B.16) to first order in $`𝑮_{\text{rot}}`$ and second order in $`𝑩_{\text{rot}}`$, then substituting (B.17). The rotating-frame commutation relations are then invoked to simplify the cross-terms:
$$[(𝒔_{\text{rot}}(t))_i,(𝒔_{\text{rot}}(t))_j]=i\mathrm{}ϵ_{ijk}(𝒔_{\text{rot}}(t))_k,$$
These are valid because spin commutators are form-invariant under transformation to the Heisenberg picture and to the rotating frame. Further specializing to spin-$`\frac{1}{2}`$ implies
$$\left(𝒔_{\text{rot}}(t)\right)_x^2=\left(𝒔_{\text{rot}}(t)\right)_y^2=\left(𝒔_{\text{rot}}(t)\right)_z^2=\left(\frac{\mathrm{}}{2}\right)^2𝕀$$
where $`𝕀`$ is the identity operator; these identities also are form-invariant. Finally, when we average over a thermal ensemble, we have
$$_{\mathrm{}}^t^{}dt^{}\frac{1}{2}𝑩_{\text{rot}}(t)\mathbf{}𝑩_{\text{rot}}(t^{})+tt^{}=𝑺_{𝑩_{\text{rot}}}(0)$$
as well as (trivially) $`𝕀=1`$. the Langevin equation for spin-$`\frac{1}{2}`$ particles can be cast into the form of a master equation, which turns out to be the following *Bloch equation*
(B.19) $`{\displaystyle \frac{d}{dt}}𝒔_{\text{rot}}(t)=`$ $`{\displaystyle \frac{1}{T_1}}\widehat{𝒃}\mathbf{}\widehat{𝒃}\mathbf{}\left(𝒔_{\text{rot}}(t)𝒔_0\right)`$
$`{\displaystyle \frac{1}{T_2}}(𝑰\widehat{𝒃}\mathbf{}\widehat{𝒃})\mathbf{}𝒔_{\text{rot}}(t)`$
Here $`T_1`$ and $`T_2`$ are as given in (2.9), and we have wrapped the spin operator in an ensemble average $`\mathrm{}`$ in order to express the Bloch equations in terms of the (c-number) vector spin polarization $`𝒔_{\text{rot}}`$, as is traditional. The equilibrium spin polarization $`𝒔_0`$ is found to be
(B.20)
$$𝒔_0=\frac{\mathrm{}}{2}\mathrm{tanh}\left(\frac{\mathrm{}\gamma B_0}{2k_\text{B}T}\right)\widehat{𝒃}$$
which agrees with the well-known thermodynamic expression for spin-$`\frac{1}{2}`$ polarization; this provides a consistency check of the Langevin formalism.
The calculation of the spin-locked relaxation time $`T_{1\rho }`$ is similar, differing only in that transformation to a doubly rotating frame is required. The result is as given in (2.9); details of the derivation will not be given.
We close by noting that our spin relaxation rates differ from the oscillator relaxation rates of Ford *et al.* in one physically important respect: $`T_1`$, $`T_2`$, and $`T_{1\rho }`$ are strongly temperature-dependent, while Ford *et al.* predict an oscillator quality $`Q`$ that is independent of temperature. Formally, the reason for this difference can be traced to the Langevin equation (B.16), in which the dissipative kernel $`𝑮_{\text{rot}}`$ is temperature-independent (this is true for both spins and oscillators). However, when this Langevin equation is integrated to second order to obtain the Bloch equations, the temperature-dependent fluctuations $`𝑩(t)`$ *also* contribute to spin relaxation via the commutation relation $`[s_i,s_j]=i\mathrm{}ϵ_{ijk}s_k`$. In contrast, the oscillator operators $`p`$ and $`q`$ have a pure c-number commutator $`[q,p]=i\mathrm{}`$, and in consequence the fluctuating force exerted by the reservoir does *not* directly contribute to oscillator relaxation.
The difference between spin relaxation and oscillator relaxation can also be understood physically. If we imagine a classical spin that is subject to a randomly fluctuating magnetic field, and no other physical influence, the average spin orientation $`𝒔`$ necessarily will relax toward $`𝒔=0`$. So to the extent that a higher-temperature reservoir creates stronger magnetic fluctuations, spins will relax more rapidly at higher temperature. In contrast, a randomly fluctuating force does not change the average trajectory of an oscillator at all—provided the average force is zero—because oscillator equations of motion are linear, unlike spin equations. So to the extent that a thermal reservoir has a temperature-independent dissipative kernel—like the independent oscillator model—the oscillator quality $`Q`$ will be independent of temperature.
## Appendix C Fluctuation-dissipation-entanglement theorems
Now we will prove fluctuation-dissipation-entanglement theorems for both spin-$`\frac{1}{2}`$ particles and harmonic oscillators in contact with a thermal reservoir.
### C.1. A spin-$`\frac{1}{2}`$ fluctuation-dissipation-entanglement theorem
Section 2.3 has already provided the definitions we need to prove the fluctuation-dissipation-entanglement theorem, and Appendix B has carried through many of the needed calculations. We need only organize our reasoning as follows.
Specializing to spin-$`\frac{1}{2}`$ particles, we write the total Hamiltonian (B.1) in the usual form $`H=H_0+V`$, with the perturbing Hamiltonian $`V`$ given by
(C.1)
$$V=\underset{j}{}\omega _j\beta _jq_j\widehat{𝒏}_j\mathbf{}𝒔$$
The exact ground state $`|\mathrm{\Psi }_0`$ of $`H`$ can be calculated order-by-order in $`V`$ as
$$|\mathrm{\Psi }_0=|\mathrm{\Psi }_0^{(0)}+|\mathrm{\Psi }_0^{(1)}+|\mathrm{\Psi }_0^{(2)}+\mathrm{}$$
in the notation of Landau and Lifschitz . The formal expression for $`|\mathrm{\Psi }_0^{(2)}`$ is quite lengthy, but fortunately we have $`|\mathrm{\Psi }_0^{(0)}=0`$ by construction (2.11), so that cross terms like $`\mathrm{\Psi }_0^{(2)}||\mathrm{\Psi }_0^{(0)}`$ vanish. In consequence, the leading contribution to the entanglement comes entirely from $`\mathrm{\Psi }_0^{(1)}||\mathrm{\Psi }_0^{(1)}`$, for which we have the well-known expression
(C.2)
$$|\mathrm{\Psi }_0^{(1)}=\underset{m0}{}\frac{\psi _m^{(0)}|V|\psi _0^{(0)}}{E_0E_m}|\psi _m^{(0)}$$
The fluctuation-dissipation-entanglement theorem (2.12) asserted in Section 2.3 can now be readily derived
(C.3a) $`E(\widehat{𝒑},\omega _0)\mathrm{\Psi }_0||\mathrm{\Psi }_0\mathrm{\Psi }_0^{(1)}||\mathrm{\Psi }_0^{(1)}`$
(C.3b) $`={\displaystyle \underset{j}{}}{\displaystyle \frac{\mathrm{}\omega _j^2\beta _j^2}{8(\omega _0+\omega _j)^2}}\text{tr}\left[(𝑰\widehat{𝒑}\mathbf{}\widehat{𝒑})\mathbf{}(\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j)\right]`$ by (C.1)
(C.3e) $`={\displaystyle \frac{\mathrm{}}{4\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle \frac{\omega }{(\omega _0+\omega )^2}}\begin{array}{c}({\displaystyle \underset{j}{}}{\displaystyle \frac{\pi \omega _j\beta _j^2}{2}}\delta (\omega \omega _j)\hfill \\ \text{tr}[(𝑰\widehat{𝒑}\mathbf{}\widehat{𝒑})\mathbf{}(\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j)])\hfill \end{array}`$
(C.3f) $`={\displaystyle \frac{\mathrm{}}{4\pi }}𝔊_2\{\omega \text{tr}\left[(𝑰\widehat{𝒑}\mathbf{}\widehat{𝒑})\mathbf{}\text{Re}(\stackrel{\mathbf{~}}{𝑮}(\omega ))\right];\omega _0\}`$ by (B.11)
(C.3g) $`={\displaystyle \frac{\mathrm{}\gamma ^2}{4\pi }}𝔊_2\{\omega \text{tr}\left[(𝑰\widehat{𝒑}\mathbf{}\widehat{𝒑})\mathbf{}𝚪(\omega )\right];\omega _0\}`$ by (B.13)
where (C.3e) was obtained by introducing $`\omega `$ as a dummy variable of integration over $`\delta (\omega \omega _j)`$.
### C.2. An oscillator fluctuation-dissipation-entanglement theorem
A similar fluctuation-dissipation-entanglement theorem exists for harmonic oscillators. We specify the Hamiltonian $`H`$ of the oscillator/thermal reservoir as
(C.4)
$$H=\frac{1}{2}\omega _0\left(p^2+q^2\right)+\underset{j}{}\frac{1}{2}\omega _j\beta _j\left(p_j^2+(q_jq)^2\right)$$
Here $`q`$ and $`p`$ are generalized oscillator coordinates satisfying $`[q,p]=i\mathrm{}`$, whose physical nature is not otherwise specified. As shown by Ford, Lewis, and O’Connell , the Hamiltonian (C.4) implies the following quantum Langevin equation in the Heisenberg picture
(C.5)
$$\ddot{q}(t)+\omega _0^2q(t)+\omega _0_{\mathrm{}}^t𝑑t^{}\mu (tt^{})\dot{q}(t^{})=f(t)$$
where $`f(t)`$ is a fluctuating thermal force, and $`\mu (\tau )`$ is a dissipative “memory function” given by
(C.6)
$$\mu (\tau )=\{\begin{array}{cc}_j\omega _j\beta _j^2\mathrm{cos}(\omega _j\tau )\hfill & \text{for}\tau 0;\hfill \\ 0\hfill & \text{for}\tau <0.\hfill \end{array}$$
The Fourier transform of $`\mu (\tau )`$
(C.7)
$$\stackrel{~}{\mu }(\omega )_{\mathrm{}}^{\mathrm{}}𝑑\tau e^{i\omega \tau }\mu (\tau )$$
has a real part given explicitly by
(C.8)
$$\text{Re}[\stackrel{~}{\mu }(\omega )]=\frac{\pi }{2}\underset{j}{}\omega _j\beta _j^2\left(\delta (\omega \omega _j)+\delta (\omega +\omega _j)\right)$$
which is related to the spectral density $`S_f(\omega )`$ of the thermal force $`f(t)`$ by
(C.9)
$$S_f(\omega )=\omega _0^2\text{Re}[\stackrel{~}{\mu }(\omega )]\mathrm{}\omega \mathrm{coth}\left(\frac{\mathrm{}\omega }{2k_\text{B}T}\right)$$
This is, of course, the celebrated fluctuation-dissipation theorem for harmonic oscillators.
Now we have all the pieces we need to calculate the zero-temperature entanglement from the dissipative kernel $`\stackrel{~}{\mu }(\omega )`$. By reasoning exactly analogous to the spin-$`\frac{1}{2}`$ case (C.3aC.3f) we find
(C.10a) $`\mathrm{\Psi }_0||\mathrm{\Psi }_0`$ $`\mathrm{\Psi }_0^{(1)}||\mathrm{\Psi }_0^{(1)}`$
(C.10b) $`={\displaystyle \underset{j}{}}{\displaystyle \frac{\omega _j^2\beta _j^2}{4(\omega _0+\omega _j)^2}}`$
(C.10c) $`={\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle \frac{\omega }{(\omega _0+\omega )^2}}\text{Re}[\stackrel{~}{\mu }(\omega )]`$
(C.10d) $`={\displaystyle \frac{1}{2\pi }}𝔊_2\{\omega \text{Re}[\stackrel{~}{\mu }(\omega )];\omega _0\}`$
This is the general form of the fluctuation-dissipation-entanglement theorem for harmonic oscillators.
As in Section 2.3.1 for spin-$`\frac{1}{2}`$ particles, oscillator entanglement can be evaluated approximately. We stipulate that $`\stackrel{~}{\mu }(\omega )1/Q`$ for $`\omega (0,\omega _c)`$, with $`Q`$ the oscillator quality and $`\omega _\text{c}`$ the thermal reservoir’s cutoff frequency $`\omega _\text{c}\omega _0`$. Then we find to leading order in $`\omega _0/\omega _\text{c}`$
(C.11)
$$\mathrm{\Psi }_0||\mathrm{\Psi }_0\frac{\mathrm{ln}(\omega _\text{c}/\omega _0)}{2\pi Q}$$
This approximate result is not new. Li, Ford, and O’Connell \[16, eq. 14\] have calculated a zero-temperature energy shift that directly implies (C.11), in the context of their reply to a critique by Senitzky of a still-earlier analysis of energy balance in dissipative systems . However, they did not interpret the energy shift as an entanglement relation, and the generality and rigor of the link between dissipation and entanglement via Stieltj i e s transforms was not pointed out.
### C.3. Oscillator renormalization
For spin-$`\frac{1}{2}`$ particles, we have seen that the thermal reservoir interactions *renormalize* the dynamical equations of the field operators, via the anisotropy tensor $`𝑪`$ (B.4,B.7). We will now show that a similar renormalization occurs for oscillators.
To begin, we remark that in quantum-coherent engineering, renormalization is more than an abstract concept. At least in the context of magnetic resonance force microscopy, we will see that renormalization effects are large—they play a key role in experimental protocols—and are of considerable practical consequence.
Next, we note that oscillator renormalization is much easier to analyze if we recognize at the outset that the operators $`q`$ and $`p`$ and the frequency $`\omega _0`$ which appear in the independent oscillator Hamiltonian (C.4) are *already* renormalized. We know this from the quantum Langevin equation (C.5), in which $`\omega _0`$ appears as the resonant frequency of the oscillator *after* it has been renormalized by interactions with the thermal reservoir. Our task, therefore, is to deduce an unrenormalized “bare” Hamiltonian from the renormalized “dressed” Hamiltonian (C.4). A great virtue of the independent oscillator model is that this calculation can be carried through quite easily.
We define an unrenormalized frequency $`\omega _0^{}`$ by
(C.12a) $`\omega _0^{}`$ $`\left[\omega _0\left(\omega _0+_j\omega _j\beta _j^2\right)\right]^{1/2}`$
and unrenormalized operators $`q^{}`$ and $`p^{}`$ by
(C.12b) $`q^{}`$ $`q(\omega _0^{}/\omega _0)^{1/2}`$
(C.12c) $`p^{}`$ $`p(\omega _0/\omega _0^{})^{1/2}`$
and unrenormalized reservoir couplings $`\beta _j^{}`$ by
(C.12d) $`\beta _j^{}`$ $`\beta _j(\omega _0/\omega _0^{})^{1/2}`$
By construction, the unrenormalized operators $`q^{}`$ and $`p^{}`$ satisfy the canonical commutation relation $`[q^{},p^{}]=i\mathrm{}`$. Furthermore, the unrenormalized parameters $`\omega _0^{}`$ and $`\{\beta _j^{}\}`$ are such that the Hamiltonian—when written in terms of unrenormalized quantities—takes the desired “bare” form $`H=H_0+V`$, with $`H_0`$ and $`V`$ given by
(C.13a) $`H_0`$ $`={\displaystyle \frac{1}{2}}\omega _0^{}\left(p_{}^{}{}_{}{}^{2}+q_{}^{}{}_{}{}^{2}\right)+{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{2}}\omega _j\left(p_j^2+q_j^2\right)`$
(C.13b) $`V`$ $`={\displaystyle \underset{j}{}}\omega _jq_j\beta _j^{}q^{}`$
The virtue of the “bare” Hamiltonian $`H_0`$ is that it clearly shows what happens when the oscillator-reservoir couplings $`\{\beta _j^{}\}`$ are turned off: the resulting “bare” oscillator dynamics are described in terms of unrenormalized canonical operators $`q^{}`$ and $`p^{}`$ and resonant frequency $`\omega _0^{}`$.
The physical meaning of—and the necessity for—renormalization is made clear by the following real-world example. In magnetic resonance force microscopy (MRFM) experiments, the cantilever resonant frequency is routinely monitored as the tip of the cantilever is brought near to a sample. As the tip nears the sample, the Van der Waals potential (equivalent to the Casimir effect) increasingly acts to reduce the cantilever’s spring constant, which renormalizes the resonant frequency $`\omega _0`$ to a lower value than the “bare” value $`\omega _0^{}`$, consistent with (C.12a). Renormalization effects in MRFM are so strong that is commonplace for the renormalized frequency to pass through zero and become imaginary, in which case the cantilever becomes unstable and the tip “snaps in” to the sample. Observation of the renormalized frequency provides a vital MRFM technique for monitoring the tip-sample separation without ever touching the sample. In our experiments, we frequently operate at frequency shifts in excess of 100 Hertz, so that renormalization is by no means a small effect.
Concomitantly, the cantilever’s quantum ground state is also renormalized: far from the sample the ground state satisfies $`(q^{}+ip^{})|\psi _0^{}=0`$, while near to the sample the ground state satisfies $`(q+ip)|\psi _0=0`$. Physically speaking, the renormalized ground state $`|\psi _0`$ and the unrenormalized ground state $`|\psi _0^{}`$ are squeezed relative to one another. Although deliberately squeezing cantilever states is not yet routine practice, it is not too far-fetched to imagine that someday it may be.
We close this section by remarking that—as is typically the case in renormalization theory—the inverse problem of computing the renormalized parameters from the bare parameters is not solvable in closed form.<sup>14</sup><sup>14</sup>14To see this, students should try to invert (C.12aC.12d) to express the renormalized variables $`\{q,p,\omega _0,\beta _j\}`$ as closed-form functions of the bare variables $`\{q^{},p^{},\omega _0^{},\beta _j^{}\}`$. Instead, the renormalized parameters must be calculated perturbatively, order by order in $`\{\beta _j^{}\}`$, and convergence is not guaranteed. In fact, perturbative renormalization is *guaranteed* to fail when the bare Hamiltonian has a spectrum that is unbounded from below, because we are trying to convert the bare Hamiltonian into a manifestly positive-definite renormalized form. Such divergences are perfectly physical: they are realized experimentally whenever a cantilever tip snaps in. Thus, all aspects of renormalization theory—even perturbative divergences—acquire practical significance in quantum-coherent engineering.
### C.4. Renormalization-dissipation relations
Now we show that renormalization parameters can be calculated from the same dissipative kernels that control fluctuation, dissipation, and entanglement. In our nomenclature, these relations do not qualify as “theorems” because the relationship is not invertible—the form of the dissipative kernel cannot be inferred from measurements of renormalization.
For oscillators, renormalization is entirely specified by the ratio of the unrenormalized frequency $`\omega _0^{}`$ to the renormalized frequency $`\omega _0`$, per (C.12aC.12d). From (C.8) and (C.12a) we express the frequency shift in terms of the dissipative kernel $`\text{Re}[\stackrel{~}{\mu }(\omega )]`$
(C.14) $`{\displaystyle \frac{\omega _{0}^{}{}_{}{}^{2}}{\omega _0^2}}`$ $`=1+{\displaystyle \underset{j}{}}\omega _j\beta _j^2`$
$`=1+{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\omega \text{Re}[\stackrel{~}{\mu }(\omega )]`$
For particles of arbitrary spin, renormalization is entirely specified by the anisotropy tensor $`𝑪`$, per the quantum Langevin equation (B.4). From (B.7) amd (B.11) we find the anisotropy tensor in terms of the dissipative kernel
(C.15) $`𝑪`$ $`={\displaystyle \underset{j}{}}\omega _j\beta _j^2\widehat{𝒏}_j\mathbf{}\widehat{𝒏}_j`$
$`={\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\omega \text{Re}[\stackrel{\mathbf{~}}{𝑮}(\omega )]`$
We see that renormalization effects are temperature-independent; this is consistent with a careful reading of the independent-oscillator literature .
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# 1 Introduction
## 1 Introduction
The present work is the first of two papers devoted to the construction of the symplecto–orthogonal quantum supergroup $`\text{SPO}_q(2n|2m)`$ and of some of its comodule superalgebras. In this work, I am going to calculate the $`R`$–matrix of the symplecto–orthogonal quantum superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the vector representation. Once this has been done, we can use the techniques of Ref. (generalized to the super case) to define the corresponding quantum supergroup $`\text{SPO}_q(2n|2m)`$ and to introduce its basic comodule superalgebras. This will be carried out in the subsequent paper .
As the reader will immediately notice, my approach is slightly different from what he/she presumably might expect. Hence a few words of explanation are in order. The starting point, and one of the main goals of the present investigation, was to construct a deformed Weyl superalgebra (i.e., a deformed oscillator algebra) $`W_q(n|m)`$, describing $`n`$ deformed bosons and $`m`$ deformed fermions, and covariant under deformed orthosymplectic transformations (I am grateful to V. Rittenberg for insisting that this problem should be solved). Classically, the bosonic/fermionic commutation relations are invariant under symplectic/orthogonal transformations. Since, in supersymmetry, bosons/fermions are regarded to be even/odd, the natural supersymmetric generalization of the above is that a combined system consisting of $`n`$ bosons and $`m`$ fermions is invariant under the action of the symplecto–orthogonal Lie superalgebra $`𝔰𝔭𝔬(2n|2m)`$, rather than under the action of the orthosymplectic Lie superalgebra $`𝔬𝔰𝔭(2m|2n)`$. From a practical point of view, this distinction is not really important. It is well–known that the Lie superalgebras $`𝔬𝔰𝔭(2m|2n)`$ and $`𝔰𝔭𝔬(2n|2m)`$ are naturally isomorphic: Basically, the transition from $`𝔬𝔰𝔭(2m|2n)`$ to $`𝔰𝔭𝔬(2n|2m)`$ amounts to a shift of the gradation of the vector representation . Nevertheless, I prefer to work from the outset with the natural gradations, and to avoid any shift of gradations.
The second point where I am going to depart from the more familiar formulation is more serious. Since Kac’s basic papers on Lie superalgebras , it has become customary to split the family of Lie superalgebras $`𝔬𝔰𝔭(2m|2n)`$ into two subfamilies, the $`C`$–type algebras, which are those with $`m=1`$, and the $`D`$–type algebras, which are those with $`m2`$. Accordingly, the so–called distinguished basis of the root system is chosen differently for these two subfamilies.
Needless to say, there are good reasons for considering the $`C`$– and $`D`$–type Lie superalgebras separately. In the standard terminology, the former are of type I , while the latter are of type II . This has serious consequences for the general representation theory of these algebras. On the other hand, one must not forget that the root systems of all of the $`𝔬𝔰𝔭(2m|2n)`$ algebras have some bases which resemble those of the $`C`$–type Lie algebras, and others which are similar to those of the $`D`$–type Lie algebras (so that I would prefer to say that these algebras are of CD–type). In particular, for each of the $`𝔬𝔰𝔭(2m|2n)`$ algebras, the root system has a basis, which is of $`C`$–type and contains only one odd simple root (see Section 2). This is the basis I am going to choose (but, of course, for the $`𝔰𝔭𝔬(2n|2m)`$ algebras).
Since the quantum superalgebra associated to a basic classical Lie superalgebra depends on the choice of the basis of the root system, any such choice has non–trivial consequences. The advantage of my choice is that it allows of a simultaneous treatment of all cases, resulting in a unified construction of the corresponding quantum supergroups $`\text{SPO}_q(2n|2m)`$ and of the deformed Weyl superalgebras $`W_q(n|m)`$. The reader might wonder whether the differences between the $`C`$–type and $`D`$–type Lie superalgebras will not show up at some stage of our investigations. But since in the following we only have to consider the vector module $`V`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$ and its tensorial square $`VV`$, such is not the case.
In principle, the $`R`$–matrix in question could be calculated by specializing the formula for the universal $`R`$–matrix given in Ref. (see also Ref. ), or by using the results of Ref. (I am grateful to M. Jimbo and M. Okado for drawing my attention to the latter reference). However, I prefer to proceed differently and to determine the corresponding braid generator $`\widehat{R}`$ by investigating the module structure of $`VV`$. This procedure has the advantage of yielding the spectral decomposition of $`\widehat{R}`$ as well, moreover, at several places it can serve to check the general theory.
The present work is set up as follows. In Section 2 we introduce the Lie superalgebra $`𝔰𝔭𝔬(2n|2m)`$ and fix some notation. In particular, we specify the basis of the root system that we are going to use, and we introduce the corresponding Chevalley–Serre generators of the algebra. Using these data, we define in Section 3 the quantum superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the sense of Drinfeld and Jimbo (generalized to the super case). Basically, we follow Ref. , but some details are different. In Section 4 we introduce the vector module $`V`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$. This is almost trivial, since (in the usual sloppy terminology) this module is undeformed. We also show that, as in the undeformed case, there exists on $`V`$ a $`U_q(𝔰𝔭𝔬(2n|2m))`$–invariant bilinear form, which is unique up to scalar multiples.
In Section 5 we investigate the structure of the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$, in particular, we determine its module endomorphisms. This section is central to the present work. Using the results obtained therein, we can calculate the $`R`$–matrix $`R`$ (equivalently, the braid generator $`\widehat{R}`$ ) of $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the vector representation. This will be carried out in Section 6. In Section 7 we collect some of the basic properties of $`R`$ and $`\widehat{R}`$. Section 8 contains a comparison of our results with known special cases. A brief discussion in Section 9 closes the main body of the paper. There are two appendices: In Appendix A we comment on invariant bilinear forms, in Appendix B we introduce what we have called the partial (super)transposition.
We close this introduction by explaining some of our conventions. The base field will be the field $``$ of complex numbers (in the appendices, we allow for an arbitrary field $`𝕂`$ of characteristic zero). If $`A`$ is an algebra, and if $`V`$ is an arbitrary (left) $`A`$–module, the representative of an element $`aA`$ under the corresponding representation will be denoted by $`a_V`$, and the image of an element $`xV`$ under the module action of $`a`$ will be written in the form $`a_V(x)=ax`$. The multiplication in a Lie superalgebra will be denoted by pointed brackets $`,`$. All algebraic notions and constructions are to be understood in the super sense, i.e., they are assumed to be consistent with the $`_2`$ – gradations and to include the appropriate sign factors.
## 2 Notation and a few comments on the Lie superalgebra $`𝔰𝔭𝔬(2n|2m)`$
Essentially, we use the same type of notation as in Ref. (see also Ref. ), but adapted to the present setting.
We choose two integers $`m,n1`$ and set
$$r=m+n.$$
Let $`V=V_{\overline{0}}V_{\overline{1}}`$ be a $`_2`$ – graded vector space such that
$$dimV_{\overline{0}}=2n,dimV_{\overline{1}}=2m,$$
let $`b`$ be a non–degenerate, even, super–skew–symmetric, bilinear form on $`V`$, and let $`𝔰𝔭𝔬(b)`$ be the Lie superalgebra consisting of all vector space endomorphisms of $`V`$ that leave the form $`b`$ invariant. Then $`𝔰𝔭𝔬(b)`$ is isomorphic to $`𝔰𝔭𝔬(2n|2m)`$.
According to Ref. , the Lie superalgebra $`𝔰𝔭𝔬(b)`$ can be described as follows (note that in the cited reference we have written $`𝔬𝔰𝔭(b)`$ instead of $`𝔰𝔭𝔬(b)`$). Let $`𝔤𝔩(V_{\overline{0}}V_{\overline{1}})`$ be the general linear Lie superalgebra of the $`_2`$ – graded vector space $`V`$, and let
$$\theta :VV𝔤𝔩(V_{\overline{0}}V_{\overline{1}})$$
be the linear map defined by
$$\theta (xy)z=b(y,z)x+(1)^{\xi \eta }b(x,z)y,$$
for all $`xV_\xi `$, $`yV_\eta `$, $`zV`$, with $`\xi ,\eta _2`$. Then the kernel of $`\theta `$ is equal to the subspace of all super–skew–symmetric tensors in $`VV`$, its image is equal to the subalgebra $`𝔰𝔭𝔬(b)`$ of $`𝔤𝔩(V_{\overline{0}}V_{\overline{1}})`$, and $`\theta `$ is an $`𝔰𝔭𝔬(b)`$ – module homomorphism. In particular, $`\theta `$ induces an $`𝔰𝔭𝔬(b)`$ – module isomorphism of the submodule of all super–symmetric tensors in $`VV`$ onto the adjoint $`𝔰𝔭𝔬(b)`$ – module.
Let us make all this more explicit by introducing a suitable basis of $`V`$. In order to do that we need some more notation. Define the index sets
$`I`$ $`=`$ $`\{r,r+1,\mathrm{},2,1,1,2,\mathrm{},r1,r\}`$
$`I_{\overline{0}}`$ $`=`$ $`\{n,n+1,\mathrm{},2,1,1,2,\mathrm{},n1,n\}`$
$`I_{\overline{1}}`$ $`=`$ $`\{r,r+1,\mathrm{},n2,n1,n+1,n+2,\mathrm{},r1,r\}`$
and also
$$\begin{array}{ccc}\hfill J& =& \{r,r+1,\mathrm{},2,1\}\hfill \\ \hfill J_{\overline{0}}& =& \{n,n+1,\mathrm{},2,1\}=JI_{\overline{0}}\hfill \\ \hfill J_{\overline{1}}& =& \{r,r+1,\mathrm{},n2,n1\}=JI_{\overline{1}}.\hfill \end{array}$$
Moreover, define the elements $`\eta _i_2`$; $`iI`$, by
$$\eta _i=\alpha \text{if }iI_\alpha \text{}\alpha _2,$$
and the sign factors
$$\sigma _i=(1)^{\eta _i},\sigma _{ij}=(1)^{\eta _i\eta _j}\text{for all }i,jI$$
$$\tau _j=1\text{and}\tau _j=\sigma _j\text{for all }jJ.$$
Note that
$$\tau _i\tau _i=\sigma _i\text{for all }iI$$
(note also that the mapping $`\pi :II`$ used in Ref. is given by $`\pi (i)=i`$ for all $`iI`$).
Then there exists a homogeneous basis $`(e_i)_{iI}`$ of $`V`$ such that $`e_i`$ is homogeneous of degree $`\eta _i`$, for all $`iI`$, and such that
$$b(e_i,e_j)=\tau _j\delta _{i,j}\text{for all }i,jI.$$
We shall also use the notation
$$C_{ij}=b(e_i,e_j),i,jI.$$
(2.1)
If $`C`$ is the $`I\times I`$– matrix with elements $`C_{ij}`$, and if $`G`$ is the $`I\times I`$– matrix defined by
$$G_{ij}=\sigma _i\delta _{ij}\text{for all }i,jI,$$
(2.2)
we have
$$C^2=G.$$
(2.3)
Besides the basis $`(e_i)_{iI}`$ of $`V`$, we also use the basis $`(f_i)_{iI}`$, which is dual to $`(e_i)`$ with respect to $`b`$ and is defined by
$$b(f_j,e_i)=\delta _{ij}\text{for all }i,jI.$$
Obviously, $`f_i`$ is homogeneous of degree $`\eta _i`$. Explicitly, we have
$$f_i=\underset{jI}{}(C^1)_{ij}e_j=\tau _ie_i\text{for all }iI.$$
Using the two bases $`(e_i)`$ and $`(f_i)`$ of $`V`$, we define the following elements of $`𝔰𝔭𝔬(b)`$:
$$X_{ij}=\theta (e_if_j)\text{for all }i,jI.$$
Let $`(E_{ij})_{i,jI}`$ be the basis of $`𝔤𝔩(V_{\overline{0}}V_{\overline{1}})`$ that canonically corresponds to the basis $`(e_i)_{iI}`$ of $`V`$, i.e.,
$$E_{ij}(e_k)=\delta _{jk}e_i\text{for all }i,j,kI.$$
Then we obtain
$$X_{ij}=E_{ij}+\sigma _{ij}\underset{k,\mathrm{}}{}C_{ik}(C^1)_j\mathrm{}E_\mathrm{}k=E_{ij}+\sigma _{ij}\tau _i\tau _jE_{j,i}\text{for all }i,jI.$$
In particular, we have
$$X_{ii}=E_{ii}E_{i,i}\text{for all }iI.$$
According to the properties of the map $`\theta `$, the elements $`X_{ij}`$ generate the vector space $`𝔰𝔭𝔬(b)`$, moreover, the super–symmetry of $`\theta `$ implies that
$$\tau _jX_{i,j}=\sigma _{ij}\tau _iX_{j,i}\text{for all }i,jI.$$
Thus we have
$$X_{i,i}=0\text{for all }iI_{\overline{1}}.$$
Let $`𝔥`$ be the subspace of $`𝔰𝔭𝔬(b)`$ that is spanned by the elements $`X_{ii}`$, $`iI`$. Obviously, $`𝔥`$ consists of those elements of $`𝔰𝔭𝔬(b)`$ whose matrices with respect to the basis $`(e_i)`$ are diagonal, and the $`X_{jj}`$ with $`jJ`$ form a basis of $`𝔥`$.
Define, for every $`iI`$, the linear form $`\epsilon _i`$ on $`𝔥`$ by
$$H(e_i)=\epsilon _i(H)e_i\text{for all }H𝔥.$$
Then it is easy to see that
$$\epsilon _i=\epsilon _i\text{for all }iI,$$
and that
$$\epsilon _i(X_{jj})=\delta _{ij}\text{for all }i,jJ.$$
Thus $`(\epsilon _j)_{jJ}`$ is the basis of $`𝔥^{}`$ that is dual to the basis $`(X_{jj})_{jJ}`$ of $`𝔥`$.
Since $`\theta `$ is an $`𝔰𝔭𝔬(b)`$ – module homomorphism, it follows that
$$H,X_{ij}=(\epsilon _i\epsilon _j)(H)X_{ij}$$
for all $`H𝔥`$ and all $`i,jI`$ (recall that the multiplication in a Lie superalgebra is denoted by pointed brackets). We conclude that $`𝔥`$ is a Cartan subalgebra of $`𝔰𝔭𝔬(b)`$, that
$$\mathrm{\Delta }=\{\epsilon _i\epsilon _j|i,jI;ji\text{ and }j<i,\text{ or }j=iI_{\overline{0}}\}$$
(2.4)
is the root system of $`𝔰𝔭𝔬(b)`$ with respect to $`𝔥`$, and that $`X_{ij}`$ is a (non–zero) root vector corresponding to the root $`\epsilon _i\epsilon _j`$ (with $`i,j`$ as specified on the right hand side of Eqn. (2.4)). The root $`\epsilon _i\epsilon _j`$ is even/odd depending on whether $`\sigma _i\sigma _j=\pm 1`$.
In order to introduce an adequate bilinear form on $`𝔥^{}`$, we recall that the invariant bilinear form
$$(X,Y)\frac{1}{2}\text{Str}(XY)$$
on $`𝔰𝔭𝔬(b)`$ is non–degenerate and super–symmetric; consequently, its restriction to $`𝔥`$ is likewise. Let $`(|)`$ denote the bilinear form on $`𝔥^{}`$ that is inverse to this restriction. By definition, we have
$$(\lambda |\mu )=\frac{1}{2}\text{Str}(H_\lambda H_\mu )$$
for all $`\lambda ,\mu 𝔥^{}`$, where, for example, the element $`H_\lambda 𝔥`$ is uniquely determined through the equation
$$\lambda (H)=\frac{1}{2}\text{Str}(H_\lambda H)\text{for all }H𝔥.$$
(2.5)
It is easy to check that
$$H_\lambda =\underset{jJ}{}\sigma _j\lambda (X_{jj})X_{jj}\text{for all }\lambda 𝔥^{},$$
and that
$$(\epsilon _i|\epsilon _j)=\sigma _i\delta _{ij}\text{for all }i,jJ.$$
Let us now specify the basis of the root system $`\mathrm{\Delta }`$ that we are going to use in the following. It is equal to $`(\alpha _j)_{jJ}`$, where the simple roots $`\alpha _j`$ are defined by
$$\alpha _j=\{\begin{array}{cc}\epsilon _j\epsilon _{j+1}\hfill & \text{for }rj2\hfill \\ 2\epsilon _1\hfill & \text{for }j=1.\hfill \end{array}$$
Note that $`\alpha _{n1}`$ is the sole odd simple root. The corresponding Chevalley–Serre generators of $`𝔰𝔭𝔬(b)`$ are denoted by $`E_j^v`$, $`F_j^v`$, $`H_j^v`$; $`jJ`$, and are introduced as follows. First of all, we choose
$$\begin{array}{ccc}\hfill E_j^v& =& \{\begin{array}{cccc}X_{j,j+1}& =& E_{j,j+1}\sigma _j\sigma _{j,j+1}E_{j1,j}\hfill & \text{for }rj2\hfill \\ \frac{1}{2}X_{1,1}& =& E_{1,1}\hfill & \text{for }j=1\hfill \end{array}\hfill \\ \hfill F_j^v& =& \{\begin{array}{cccc}X_{j+1,j}& =& E_{j+1,j}\sigma _{j+1}\sigma _{j+1,j}E_{j,j1}\hfill & \text{for }rj2\hfill \\ \frac{1}{2}X_{1,1}& =& E_{1,1}\hfill & \text{for }j=1.\hfill \end{array}\hfill \end{array}$$
Remark 2.1. It is easy to check that
$$\sigma _{j,j+1}=\sigma _{j+1}\text{for }rj2,$$
(2.6)
or, equivalently, that
$$\sigma _{j,j+1}=\sigma _j\text{for }1jr1.$$
(2.7)
This implies that in the equation for $`E_j^v`$, $`rj2`$, the factor $`\sigma _{j,j+1}`$ might be replaced by $`\sigma _{j+1}`$, and in the equation for $`F_j^v`$, $`rj2`$, the $`\sigma `$–factors might be dropped. I prefer to keep the $`\sigma `$–factors as they stand: They have an immediate meaning in terms of the sign rules of supersymmetry, and by modifying the equations for the $`E_j^v`$ and $`F_j^v`$ as mentioned above, we might well end up in the unpleasant situation where we would have to check (possibly implicitly) the equations (2.6) or (2.7) again and again.
Using the elements $`E_j^v`$ and $`F_j^v`$, we define the generators $`H_j^v𝔥`$ as usual by
$$H_j^v=E_j^v,F_j^v\text{for all }jJ.$$
(2.8)
More explicitly, we find
$$H_j^v=\{\begin{array}{cc}X_{jj}\sigma _j\sigma _{j+1}X_{j+1,j+1}\hfill & \text{for }rj2\hfill \\ X_{1,1}\hfill & \text{for }j=1.\hfill \end{array}$$
Then the generators $`E_j^v`$, $`F_j^v`$, $`H_j^v`$ satisfy the following familiar relations, which hold for all $`i,jJ`$:
$$H_i^v,H_j^v=0$$
(2.9)
$$E_i^v,F_j^v=\delta _{ij}H_j^v$$
(2.10)
$$H_i^v,E_j^v=a_{ij}E_j^v,H_i^v,F_j^v=a_{ij}F_j^v.$$
(2.11)
Here, $`A=(a_{ij})_{i,jJ}`$ is the Cartan matrix, whose elements are defined by
$$a_{ij}=\alpha _j(H_i^v)\text{for all }i,jJ.$$
The Cartan matrix is tridiagonal. For $`n2`$, it takes the following form:
$$A=\left(\begin{array}{cccccccccccc}\hfill 2& \hfill 1& & \multicolumn{9}{c}{}\\ \hfill 1& \hfill 2& \hfill 1& \multicolumn{9}{c}{}\\ \multicolumn{12}{c}{}\\ \multicolumn{12}{c}{}\\ \multicolumn{3}{c}{}& \hfill 1& \hfill 2& \hfill 1& \multicolumn{6}{c}{}\\ \multicolumn{4}{c}{}& \hfill 1& \hfill 0& \hfill \mathrm{\hspace{0.33em}\hspace{0.33em}1}& \multicolumn{5}{c}{}\\ \multicolumn{5}{c}{}& \hfill 1& \hfill \mathrm{\hspace{0.33em}\hspace{0.33em}2}& \hfill 1& \multicolumn{4}{c}{}\\ \multicolumn{12}{c}{}\\ \multicolumn{12}{c}{}\\ \multicolumn{8}{c}{}& \hfill 1& \hfill 2& \hfill 1& \\ \multicolumn{8}{c}{}& & \hfill 1& \hfill 2& \hfill 2\\ \multicolumn{8}{c}{}& & & \hfill 1& \hfill 2\end{array}\right),$$
where the zero on the diagonal has the row and column number $`n1`$. For $`n=1`$, the zero is in position $`(2,2)`$, and the lower right corner of $`A`$ is equal to
$$\left(\begin{array}{ccc}\hfill 2& \hfill 1& \hfill \mathrm{\hspace{0.33em}\hspace{0.33em}0}\\ \hfill 1& \hfill 0& \hfill \mathrm{\hspace{0.33em}\hspace{0.33em}2}\\ \hfill 0& \hfill 1& \hfill \mathrm{\hspace{0.33em}\hspace{0.33em}2}\end{array}\right).$$
Finally, for $`m=n=1`$ the Cartan matrix is given by
$$A=\left(\begin{array}{cc}\hfill 0& \hfill 2\\ \hfill 1& \hfill 2\end{array}\right).$$
Thus the Dynkin diagram of $`𝔰𝔭𝔬(b)𝔰𝔭𝔬(2n|2m)`$ with respect to our basis of the root system takes the form
Dynkin diagram of the Lie superalgebra $`𝔰𝔭𝔬(2n|2m)`$
Remark 2.2. It may be helpful to comment on the rules according to which the generators $`H_i^v`$ (and hence $`E_i^v`$ and $`F_i^v`$) have been chosen. If the simple root $`\alpha _i`$ is even, we choose $`H_i^v`$ such that $`\alpha _i(H_i^v)=2`$. For odd simple roots, the situation is more complicated. In the present case, the sole odd simple root $`\alpha _i`$, $`i=n1`$, is such that $`(\alpha _i|\alpha _i)=0`$. Then it follows that $`a_{ii}=0`$, and the element $`H_i^v`$ is usually chosen such that (for this index $`i`$)
$$a_{ij}\text{for all }jJ,$$
and such that these $`a_{ij}`$ don’t have a common divisor. This fixes the $`a_{ij}`$ up to a common sign factor, which (according to Kac) is chosen such that $`a_{i,i+1}>0`$ (assuming that $`i+1J`$ and that $`a_{i,i+1}0`$). These conventions are introduced simply for convenience, and they are of little (if any) importance. Note that, for $`m=n=1`$, we haven’t followed these conventions: The first row of the Cartan matrix could be divided by $`2`$, and for $`𝔰𝔩(2|1)𝔰𝔭𝔬(2|2)`$ this is usually done. Our choice is motivated by the wish for a unified treatment of all cases.
The relations (2.9) – (2.11) given above are not sufficient to characterize the Lie superalgebra $`𝔰𝔭𝔬(2n|2m)`$ completely, there are certain Serre–type and supplementary relations which must also be satisfied. We don’t give these relations here, but only mention that they can be read off from the relations (3.5) – (3.13) by setting $`q=1`$.
## 3 Definition of the quantum superalgebra <br>$`U_q(𝔰𝔭𝔬(2n|2m))`$
The notation introduced in the preceding section will now be used to define the quantum superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$. Basically, we are going to follow Ref. , however, there will be differences in detail.
Define the diagonal $`J\times J`$– matrix $`D`$ by
$$D=(d_i\delta _{ij})_{i,jJ}=\text{diag}(\underset{m}{\underset{}{1,1,\mathrm{},1}},\underset{n1}{\underset{}{1,1,\mathrm{},1}},2).$$
It is chosen such that
$$(DA)_{ij}=(\alpha _i|\alpha _j)\text{for all }i,jJ.$$
In particular, the matrix $`DA`$ is symmetric.
Let $`q`$ be a non–zero complex number, and assume that $`q`$ is not a root of unity. We use the abbreviation
$$q_i=q^{d_i}.$$
Then we have
$$q_i^{a_{ij}}=q^{(\alpha _i|\alpha _j)}\text{for all }i,jJ.$$
Now the quantum superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ is defined to be the universal associative superalgebra (with unit element) with generators $`K_i`$ , $`K_i^1`$, $`E_i`$ , $`F_i`$ ; $`iJ`$, and certain relations to be specified below. The $`_2`$ – gradation is fixed by requiring that $`E_{n1}`$ and $`F_{n1}`$ be odd, while all the other generators are even. (Needless to say, one has to check that the relations are compatible with this requirement.) The relations are the following, they are assumed to hold for all $`i,jJ`$ :
$$K_iK_i^1=K_i^1K_i=1$$
(3.1)
$$K_iK_j=K_jK_i$$
(3.2)
$$K_iE_jK_i^1=q_i^{a_{ij}}E_j,K_iF_jK_i^1=q_i^{a_{ij}}F_j$$
(3.3)
$$E_i,F_j=\delta _{ij}\frac{K_iK_i^1}{q_iq_i^1}.$$
(3.4)
In addition, the generators $`E_i`$ satisfy certain Serre–type and supplementary relations among themselves, as do the generators $`F_i`$. We only write the relations for the $`E_i`$, those for the $`F_i`$ are obtained from these by simply replacing $`E`$ by $`F`$.
In the subsequent relations, it is always assumed that $`i,jJ`$. Suppose first that the root $`\alpha _i`$ is even, i.e., that $`in1`$ . Then we have
$$E_i,E_j=0\text{for }a_{ij}=0$$
(3.5)
$$E_i^2E_j(q+q^1)E_iE_jE_i+E_jE_i^2=0\text{for }i3\text{ , }|ij|=1.$$
(3.6)
If $`\alpha _2`$ is even, i.e., if $`n2`$ , we have (as in the case of symplectic Lie algebras)
$$E_2^2E_3(q+q^1)E_2E_3E_2+E_3E_2^2=0$$
(3.7)
$$E_2^3E_1(q^2+1+q^2)E_2^2E_1E_2+(q^2+1+q^2)E_2E_1E_2^2E_1E_2^3=0.$$
(3.8)
In all cases, the generators $`E_2`$ and $`E_1`$ satisfy
$$E_1^2E_2(q^2+q^2)E_1E_2E_1+E_2E_1^2=0.$$
(3.9)
Next we recall that $`\alpha _{n1}`$ is the sole odd simple root, and that this root is isotropic. Correspondingly, we have
$$E_{n1},E_j=0\text{for }a_{n1,j}=0,$$
(3.10)
in particular,
$$E_{n1}^2=0.$$
(3.11)
Finally, there are the following supplementary relations. If $`m,n2`$ , we have
$$E_{n1},E_n,E_{n1},E_{n2}_q_{q^1}=0,$$
(3.12)
and for $`n=1`$, $`m3`$ we have
$$E_2,E_3,E_2,E_1,E_2,E_3,E_4_q_q_{q^2}_{q^1}_q=0.$$
(3.13)
The last two relations are expressed in terms of so–called $`q`$–supercommutators. We recall the definition: If $`A`$ is any associative superalgebra, if $`p`$ is any non–zero complex number, and if $`XA_\xi `$ and $`YA_\eta `$ , with $`\xi ,\eta _2`$, the $`p`$–supercommutator of $`X`$ and $`Y`$ is defined by
$$X,Y_p=XYp(1)^{\xi \eta }YX.$$
Obviously, we have
$$X,Y_1=X,Y.$$
As shown below, the Serre–type relations can also be expressed in terms of $`q`$–supercommutators.
The superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ is converted into a Hopf superalgebra by means of structure maps, which are fixed by the following equations:
coproduct
$$\begin{array}{ccc}\hfill \mathrm{\Delta }(K_i^{\pm 1})& =& K_i^{\pm 1}K_i^{\pm 1}\hfill \\ \hfill \mathrm{\Delta }(E_i)& =& E_i1+K_iE_i\hfill \\ \hfill \mathrm{\Delta }(F_i)& =& F_iK_i^1+1F_i\hfill \end{array}$$
counit
$$\epsilon (K_i^{\pm 1})=1,\epsilon (E_i)=\epsilon (F_i)=0$$
antipode
$$S(K_i^{\pm 1})=K_i^1,S(E_i)=K_i^1E_i,S(F_i)=F_iK_i.$$
In the subsequent series of remarks, we collect some elementary properties of the Hopf superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$.
Remark 3.1. Let $`𝐐`$ be the root lattice of $`𝔰𝔭𝔬(2n|2m)`$, i.e.,
$$𝐐=\underset{iJ}{}\alpha _i.$$
Then the algebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ admits a unique $`𝐐`$ – gradation such that, for all $`iJ`$, the element $`E_i`$ is homogeneous of degree $`\alpha _i`$, $`F_i`$ is homogeneous of degree $`\alpha _i`$, and $`K_i`$ is homogeneous of degree $`0`$. In view of a more general definition to be given later, the $`𝐐`$ – degree of a $`𝐐`$ – homogeneous element is called its weight. If an element $`XU_q(𝔰𝔭𝔬(2n|2m))`$ has the weight $`\lambda 𝐐`$, it satisfies
$$K_jXK_j^1=q^{(\alpha _j|\lambda )}X\text{for all }jJ.$$
Conversely, if an element $`XU_q(𝔰𝔭𝔬(2n|2m))`$ satisfies this condition, it is $`𝐐`$ – homogeneous of weight $`\lambda `$ (since $`q`$ is not a root of unity). Note that the structure maps $`\mathrm{\Delta }`$, $`\epsilon `$, and $`S`$ are $`𝐐`$ – homogeneous of degree zero.
Remark 3.2. The antipode $`S`$ is bijective. To prove this, we show that $`S^2`$ is bijective. Indeed, it is easy to check that, for all $`iJ`$,
$$S^2(K_i)=K_i,S^2(E_i)=q^{(\alpha _i|\alpha _i)}E_i,S^2(F_i)=q^{(\alpha _i|\alpha _i)}F_i.$$
(3.14)
Since $`S^2`$ is an algebra endomorphism of $`U_q(𝔰𝔭𝔬(2n|2m))`$, and since a suitable set of monomials in the generators $`K_i^{\pm 1}`$, $`E_i`$, and $`F_i`$ forms a basis of $`U_q(𝔰𝔭𝔬(2n|2m))`$, this implies our claim.
Actually, $`S^2`$ is an inner automorphism of the algebra $`U_q(𝔰𝔭𝔬(2n|2m))`$. Let $`2\rho `$ denote the sum of the even positive roots minus the sum of the odd positive roots of $`𝔰𝔭𝔬(2n|2m)`$. Explicitly, we have
$$2\rho =2\underset{i=r}{\overset{n1}{}}(i+2n+1)\epsilon _i2\underset{i=n}{\overset{1}{}}i\epsilon _i,$$
and it is easy to check that
$$(2\rho |\alpha _i)=(\alpha _i|\alpha _i)\text{for all }iJ.$$
Given an arbitrary linear combination of the simple roots $`\alpha _i`$ with coefficients $`r_i`$ ,
$$\lambda =\underset{iJ}{}r_i\alpha _i𝐐,$$
we define
$$K_\lambda =\underset{iJ}{}K_i^{r_i}.$$
In particular, we have
$$K_{\alpha _i}=K_i\text{for all }iJ.$$
Then the Eqns. (3.14) immediately imply that
$$S^2(X)=K_{2\rho }XK_{2\rho }^1\text{for all }XU_q(𝔰𝔭𝔬(2n|2m)).$$
Remark 3.3. Obviously, there is a certain symmetry between the $`E`$ and $`F`$ generators of $`U_q(𝔰𝔭𝔬(2n|2m))`$. To make this more explicit, we note that there is a unique algebra endomorphismus
$$\phi :U_q(𝔰𝔭𝔬(2n|2m))U_q(𝔰𝔭𝔬(2n|2m)),$$
such that for all $`iJ`$
$$\phi (E_i)=F_i,\phi (F_i)=(1)^{\gamma _i}E_i,\phi (K_i^{\pm 1})=K_i^1,$$
where $`\gamma _i_2`$ is the degree of $`E_i`$. This endomorphism is homogeneous of $`_2`$ – degree zero, and we have $`\phi ^4=\text{id}`$ . Consequently, $`\phi `$ is an automorphism of the (associative) superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$.
Now let $`U_q(𝔰𝔭𝔬(2n|2m))^{\text{cop}}`$ be the bi–superalgebra which, regarded as a $`_2`$ – graded algebra, coincides with $`U_q(𝔰𝔭𝔬(2n|2m))`$, but whose coalgebra structure is opposite (in the super sense) to that of $`U_q(𝔰𝔭𝔬(2n|2m))`$. Then it is easy to check that
$$\phi :U_q(𝔰𝔭𝔬(2n|2m))U_q(𝔰𝔭𝔬(2n|2m))^{\text{cop}}$$
(3.15)
is a homomorphism of bi–superalgebras. Since $`\phi `$ is bijective, this implies that $`U_q(𝔰𝔭𝔬(2n|2m))^{\text{cop}}`$ is a Hopf superalgebra, and that $`\phi `$ is a Hopf superalgebra isomorphism. As is well-known, it follows (once again) that $`S`$ is bijective, and that $`S^1`$ is the antipode of $`U_q(𝔰𝔭𝔬(2n|2m))^{\text{cop}}`$.
The Serre–type and the supplementary relations can be written in various ways. Before we do that, we remind the reader of the definition of the adjoint representation of a Hopf superalgebra $`H`$. This is a (graded) representation of the superalgebra $`H`$ on the graded vector space $`H`$, it is denoted by ad , and is defined as follows. Let $`X`$ be an arbitrary element of $`H`$, and set
$$\mathrm{\Delta }(X)=\underset{^a}{}X_a^1X_a^2,$$
with homogeneous elements $`X_a^1,X_a^2H`$, of degree $`\xi _a^1`$ and $`\xi _a^2`$, respectively. Then $`\text{ad}X`$ (the representative of $`X`$) is given by
$$(\text{ad}X)(Y)=\underset{^a}{}(1)^{\xi _a^2\eta }X_a^1YS(X_a^2),$$
for all homogeneous elements $`YH_\eta `$, where $`\eta _2`$. We note that $`\text{ad}X`$ is a generalized derivation in the sense that, if $`Y^{}`$ is another element of $`H`$,
$$(\text{ad}X)(YY^{})=\underset{^a}{}(1)^{\xi _a^2\eta }(\text{ad}X_a^1)(Y)(\text{ad}X_a^2)(Y^{}).$$
Now suppose that $`S`$ bijective. This implies that $`H^{\text{cop}}`$ (see the analogous definition of $`U_q(𝔰𝔭𝔬(2n|2m))^{\text{cop}}`$ given above) is a Hopf superalgebra with antipode $`S^1`$. Let $`\overline{\text{ad}}`$ be the adjoint representation of $`H^{\text{cop}}`$. Then $`\overline{\text{ad}}`$ is a graded representation of $`H`$ in $`H`$, it is given by
$$(\overline{\text{ad}}X)(Y)=\underset{^a}{}(1)^{\xi _a^1(\xi _a^2+\eta )}X_a^2YS^1(X_a^1),$$
and it satisfies
$$(\overline{\text{ad}}X)(YY^{})=\underset{^a}{}(1)^{\xi _a^1(\xi _a^2+\eta )}(\overline{\text{ad}}X_a^2)(Y)(\overline{\text{ad}}X_a^1)(Y^{}).$$
Let us now choose $`H=U_q(𝔰𝔭𝔬(2n|2m))`$. Then the isomorphism $`\phi `$ given in (3.15) shows that
$$\overline{\text{ad}}(\phi (X))=\phi \stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(\text{ad}X)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}\phi ^1\text{for all }XU_q(𝔰𝔭𝔬(2n|2m)).$$
(3.16)
Moreover, for every element $`XU_q(𝔰𝔭𝔬(2n|2m))`$ of weight $`\lambda `$ we have
$`(\text{ad}E_i)(X)`$ $`=`$ $`E_i,X_{q^{(\alpha _i|\lambda )}}`$ (3.17)
$`(\overline{\text{ad}}F_i)(X)`$ $`=`$ $`F_i,X_{q^{(\alpha _i|\lambda )}}.`$ (3.18)
We note that in the proof of these equations we only have to use the first resp. second of the relations (3.3) but none of the other defining relations.
Now Eqn. (3.17) implies that the left hand side of Eqn. (3.5) is equal to
$$E_i,E_j=(\text{ad}E_i)(E_j),$$
the left hand side of Eqn. (3.6) is equal to
$$E_i,E_i,E_j_{q^1}_{q^{\pm 1}}=(\text{ad}E_i)^2(E_j),$$
the left hand side of Eqn. (3.7) is equal to
$$E_2,E_2,E_3_{q^1}_{q^{\pm 1}}=(\text{ad}E_2)^2(E_3),$$
the left hand side of Eqn. (3.8) is equal to
$$E_2,E_2,E_2,E_1_{q^2}_{q^{\pm 2}}=(\text{ad}E_2)^3(E_1),$$
the left hand side of Eqn. (3.9) is equal to
$$E_1,E_1,E_2_{q^2}_{q^{\pm 2}}=(\text{ad}E_1)^2(E_2),$$
the left hand side of Eqn. (3.10) is equal to
$$E_{n1},E_j=(\text{ad}E_{n1})(E_j),$$
the left hand side of Eqn. (3.12) is equal to
$`E_{n1},E_n,E_{n1},E_{n2}_q_{q^1}`$ (3.19)
$`=`$ $`(\text{ad}E_{n1})(\text{ad}E_n)(\text{ad}E_{n1})(E_{n2}),`$
and the left hand side of Eqn. (3.13) is equal to
$`E_2,E_3,E_2,E_1,E_2,E_3,E_4_q_q_{q^2}_{q^1}_q`$ (3.20)
$`=`$ $`(\text{ad}E_2)(\text{ad}E_3)(\text{ad}E_2)(\text{ad}E_1)(\text{ad}E_2)(\text{ad}E_3)(E_4).\text{}`$
Remark 3.4. Using Eqn. (3.11) and the fact that $`E_{n2}`$ and $`E_n`$ commute, it is easy to see that the left hand side of Eqn. (3.19) is invariant under the substitution $`qq^1`$. Somewhat unexpectedly, it seems that this is not the case for the left hand side of Eqn. (3.20), even if one assumes that all the relations for the $`E`$–generators except Eqn. (3.13) are satisfied.
Remark 3.5. Taking the defining relations for granted except (3.13), one can show that the expressions in Eqn. (3.20) are annihilated by all $`\text{ad}F_i`$. Since, quite generally, we have
$$(\text{ad}F_i)(X)=F_i,XK_i\text{for all }iJ\text{ and all }XU_q(𝔰𝔭𝔬(2n|2m)),$$
the same is true when acting with $`F_i,`$. This shows that by “acting” on the relation (3.13) with the generators $`F_i`$, we cannot derive new relations for the $`E`$–generators.
Up to now we have only discussed the Serre–type and supplementary relations for the $`E`$–generators. Of course, similar comments can be made for the $`F`$–generators as well, but with ad replaced by $`\overline{\text{ad}}`$. In fact, all we have to do is to apply the isomorphism $`\phi `$ given in (3.15) and to recall Eqn. (3.16).
We close this section by a remark on the weights of a $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`W`$. In the present work, all $`U_q(𝔰𝔭𝔬(2n|2m))`$–modules will be weight modules, in the sense that the representatives $`(K_j)_W`$, $`jJ`$, are simultaneously diagonalizable, and such that, for any common eigenvector $`x`$ of these operators, we have
$$K_jx=q^{(\alpha _j|\lambda )}x\text{for all }jJ,$$
with a linear form
$$\lambda \underset{iJ}{}\epsilon _i.$$
Since $`q`$ is not a root of unity, the linear form $`\lambda `$ is uniquely fixed by these conditions and is called the weight of $`x`$. This definition generalizes the definition of the weight of an element of $`U_q(𝔰𝔭𝔬(2n|2m))`$: In that case the representation considered is the adjoint representation ad or its modified version $`\overline{\text{ad}}`$.
## 4 The vector module $`V`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$
Let us now discuss the vector module $`V`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$. The definition of $`V`$ is easy, since (in the usual sloppy terminology) the vector module of $`U_q(𝔰𝔭𝔬(2n|2m))`$ is undeformed. More precisely, let $`V`$ be the graded vector space introduced in Section 2, and let $`E_i^v`$, $`F_i^v`$, and $`H_i^v`$; $`iJ`$, be the linear operators on $`V`$ defined there. Define the linear operators $`K_i^v`$, $`iJ`$, by
$$K_i^v=q_i^{H_i^v}\text{for all }iJ.$$
Since the operators $`H_i^v`$ are diagonalizable, with eigenvalues $`0,\pm 1`$, the operators $`K_i^v`$ are well–defined and, obviously, invertible. It is easy to see that the operators $`E_i^v`$, $`F_i^v`$, and $`(K_i^v)^{\pm 1}`$ satisfy the defining relations of the generators $`E_i`$ , $`F_i`$ , and $`K_i^{\pm 1}`$; $`iJ`$. Hence there exists a unique graded representation $`\pi `$ of the algebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ in $`V`$ such that
$$\pi (E_i)=E_i^v,\pi (F_i)=F_i^v,\pi (K_i^{\pm 1})=(K_i^v)^{\pm 1}\text{for all }iJ.$$
The graded vector space $`V`$, endowed with this representation, will be called the vector module of $`U_q(𝔰𝔭𝔬(2n|2m))`$.
Remark 4.1. The reader might suspect that checking the seventh order relation (3.13) might be quite tedious. Actually, this is not the case. Let $`\text{Lgr}(V)`$ be the superalgebra of all linear operators in $`V`$. It is well–known that $`\text{Lgr}(V)`$ is an $`𝔰𝔭𝔬(b)`$–module in a canonical way, and its weights (with respect to the Cartan subalgebra $`𝔥`$) are the linear forms $`\epsilon _i\epsilon _j`$ ; $`i,jI`$. Since $`𝔰𝔭𝔬(b)`$ acts on $`\text{Lgr}(V)`$ by superderivations, any product with one factor $`E_1^v`$, three factors $`E_2^v`$, two factors $`E_3^v`$, and one factor $`E_4^v`$ has the weight $`\epsilon _4+\epsilon _3+\epsilon _2\epsilon _1`$ . Since this is not a weight of $`\text{Lgr}(V)`$, every such product is equal to zero, and this implies the relation to be proved.
Let $`(e_i)_{iI}`$ be the basis of $`V`$ used in Section 2. Then we have
$$K_je_i=q^{(\alpha _j|\epsilon _i)}e_i\text{for all }jJ\text{ and all }iI.$$
Stated differently, $`e_i`$ is a weight vector with weight $`\epsilon _i`$, just as in the undeformed case.
Our next goal is to show that there exists a unique (up to scalar multiples) $`U_q(𝔰𝔭𝔬(2n|2m))`$–invariant bilinear form on $`V`$. (For a few comments on invariant bilinear forms, see Appendix A.) Let $`b`$ be a bilinear form on $`V`$, and let $`\stackrel{~}{b}`$ be the linear form on $`VV`$ canonically corresponding to $`b`$. Then $`b`$ is $`U_q(𝔰𝔭𝔬(2n|2m))`$–invariant if and only if
$$\stackrel{~}{b}(X(xy))=\epsilon (X)\stackrel{~}{b}(xy)$$
for all $`XU_q(𝔰𝔭𝔬(2n|2m))`$ and all $`x,yV`$ (see Eqn. (A.2)). The condition that
$$\stackrel{~}{b}(K_j(e_ie_k))=\stackrel{~}{b}(e_ie_k)\text{for all }jJ\text{ and all }i,kI$$
is satisfied if and only if
$$\stackrel{~}{b}(e_ie_k)=0\text{for all }i,kI\text{ with }i+k0.$$
(4.1)
In particular, this implies that $`b`$ must be homogeneous of degree zero.
Taking Eqn. (4.1) for granted, the conditions
$$\stackrel{~}{b}(E_j(e_ie_k))=0\text{for all }jJ\text{ and all }i,kI$$
and
$$\stackrel{~}{b}(F_j(e_ie_k))=0\text{for all }jJ\text{ and all }i,kI$$
both yield the same system of linear equations for the elements $`b(e_i,e_k)`$. This system has a unique (up to scalar multiples) solution. Choosing a suitable normalization, the invariant bilinear form $`b^q`$ we are looking for is given by
$$b^q(e_i,e_k)=C_{ik}^q\text{for all }i,kI,$$
where
$$C_{i,k}^q=C_{i,i}^q\delta _{i,k}\text{for all }i,kI,$$
and where the coefficients $`C_{i,i}^q`$ are given by
$$C_{i,i}^q=\{\begin{array}{cc}q^i& \text{for }1in\hfill \\ q^{i2n2}& \text{for }n1ir\hfill \\ q^i& \text{for }1in\hfill \\ q^{2ni}& \text{for }n+1ir.\hfill \end{array}$$
Obviously, the matrix $`C^q=(C_{ij}^q)_{i,jI}`$ is invertible, i.e., the bilinear form $`b^q`$ is non–degenerate. We note that $`C^{q=1}=C`$ (see Eqn. (2.1)), moreover, we have
$$C_{i,i}^qC_{i,i}^q=\{\begin{array}{cc}\hfill 1& \text{for }iI_{\overline{0}}\hfill \\ \hfill q^2& \text{for }iI_{\overline{1}}.\hfill \end{array}$$
(4.2)
Thus the matrix $`(C^q)^2`$ is not equal to $`G`$ (recall the Eqns. (2.2), (2.3)).
## 5 The structure of the module $`VV`$
We now are ready to tackle a crucial intermediate problem, namely, to determine the structure of the tensorial square of the vector module $`V`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$. In the undeformed case, this structure is known. It turns out that in the deformed case, the structure is completely analogous. In particular, for $`n=m`$, the module $`VV`$ is not completely reducible. (Actually, if adequately interpreted, the investigations of the present section apply also to the case $`q=1`$.)
To begin with, we stress that the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$ has the same weights (with the same multiplicities) as in the undeformed case: For all $`i,jI`$, the tensor $`e_ie_j`$ has the weight $`\epsilon _i+\epsilon _j`$.
As expected, $`VV`$ contains a unique (up to scalar multiples) $`U_q(𝔰𝔭𝔬(2n|2m))`$–invariant element, i.e., a non–zero element $`a`$ such that
$$Xa=\epsilon (X)a\text{for all }XU_q(𝔰𝔭𝔬(2n|2m)).$$
The invariance of $`a`$ under the action of the generators $`K_j`$ ; $`jJ`$, is equivalent to the fact that $`a`$ has the weight zero, i.e., that $`a`$ is a linear combination of the following form
$$a=\underset{iI}{}c_ie_ie_i,$$
with some coefficients $`c_i`$ , $`iI`$.
For an element $`a`$ of this form, the conditions
$$E_ja=0\text{for all }jJ$$
and
$$F_ja=0\text{for all }jJ$$
both yield the same system of linear equations for the coefficients $`c_i`$ . This system has a unique (up to scalar multiples) solution. Choosing a suitable normalization, the element $`a`$ is given by
$$a=\underset{i,kI}{}((C^q)^1)_{ik}e_ie_k,$$
where $`C^q`$ is the matrix found in Section 4. Of course, this result might have been anticipated. More explicitly, we have
$$a=\underset{i=1}{\overset{n}{}}(q^ie_ie_iq^ie_ie_i)+q\underset{i=n+1}{\overset{r}{}}(q^{i2n1}e_ie_i+q^{i+2n+1}e_ie_i).$$
It is useful to calculate $`\stackrel{~}{b}^q(a)`$, where $`\stackrel{~}{b}^q`$ is the linear form on $`VV`$ defined in Section 4. Setting
$$d=nm,$$
(5.1)
we obtain
$$\stackrel{~}{b}^q(a)=\frac{1}{q^21}((q^{2d}1)q^2(q^{2d}1))=\frac{q^dq^d}{qq^1}(q^{d+1}+q^{d1}).$$
(5.2)
Note that $`\stackrel{~}{b}^q(a)=0`$ if $`n=m`$. This is a first indication that there will be problems in the case $`n=m`$.
The rest of the present section will now be devoted to prove the following statements.
a) The $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$ is the direct sum of two submodules $`(VV)_s`$ and $`(VV)_a`$:
$$VV=(VV)_s(VV)_a,$$
where in the undeformed case $`(VV)_s`$ corresponds to the subspace of super–symmetric and $`(VV)_a`$ to the subspace of super–skew–symmetric tensors in $`VV`$.
b) The submodule $`(VV)_s`$ is irreducible.
c) The submodule $`(VV)_a`$ contains a submodule $`(VV)_a^0`$ of codimension one, and $`(VV)_s(VV)_a^0`$ is the kernel of the linear form $`\stackrel{~}{b}^q`$ found in Section 4.
d) If $`nm`$, the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_a^0`$ is irreducible, and $`(VV)_a`$ is the direct sum of the submodules $`(VV)_a^0`$ and $`a`$ :
$$(VV)_a=(VV)_a^0a\text{if }nm.$$
e) If $`n=m`$, we have
$$a(VV)_a^0\text{if }n=m,$$
and $`(VV)_a^0`$ does not have a module complement in $`(VV)_a`$.
f) If $`n=m2`$, the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_a^0/a`$ is irreducible.
g) For $`n=m=1`$, there exist two submodules $`V_4`$ and $`\overline{V}_4`$ of $`(VV)_a^0`$ such that
$$V_4+\overline{V}_4=(VV)_a^0,V_4\overline{V}_4=a,$$
and the modules $`V_4/a`$ and $`\overline{V}_4/a`$ are irreducible.
h) Let $`P_s`$ be the projector of $`VV`$ onto $`(VV)_s`$ with kernel $`(VV)_a`$, and let
$$K:VVVV$$
be the linear map defined by
$$K(u)=\stackrel{~}{b}^q(u)a\text{for all }uVV.$$
(5.3)
Then $`\text{id}_{VV}`$, $`P_s`$, and $`K`$ form a basis of the space of all $`U_q(𝔰𝔭𝔬(2n|2m))`$–module endomorphisms of $`VV`$.
In the proof of these claims, we shall obtain more detailed information on the submodules mentioned above. In particular, we shall construct bases of the vector spaces $`(VV)_s`$, $`(VV)_a^0`$, and $`(VV)_a`$.
### 5.1 The module $`(VV)_s`$
As already mentioned above, in the undeformed case the module $`(VV)_s`$ corresponds to the subspace of all super–symmetric tensors in $`VV`$. This subspace is an irreducible $`𝔰𝔭𝔬(b)`$–submodule of $`VV`$ and has the highest weight $`\epsilon _r+\epsilon _{r+1}`$. If there exists a corresponding primitive vector in the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$, it must be a linear combination of $`e_re_{r+1}`$ and $`e_{r+1}e_r`$ . Indeed, there is a unique (up to scalar multiples) linear combination of these tensors that is annihilated by all $`E_j`$, $`jJ`$. Choosing a suitable normalization, it is equal to
$$s_{r,r+1}=e_re_{r+1}+\sigma _{r,r+1}q^1e_{r+1}e_r.$$
By definition, $`(VV)_s`$ is the submodule of $`VV`$ generated by $`s_{r,r+1}`$.
Let us define the following elements of $`VV`$ :
$$\begin{array}{cccc}\hfill s_{i,j}& =& e_ie_j+\sigma _{i,j}q^1e_je_i\hfill & \text{for }i,jI\text{}i<j\text{ but }ij\hfill \\ \hfill s_{i,i}& =& e_ie_i\hfill & \text{for }iI_{\overline{0}},\hfill \end{array}$$
(5.4)
furthermore,
$$s_1=e_1e_1+q^2e_1e_1,$$
(5.5)
and for $`2jr`$
$$\begin{array}{ccc}\hfill s_j& =& q^{\sigma _{j1}}e_{j+1}e_{j1}+\sigma _{j1}q^1e_{j1}e_{j+1}\hfill \\ & & \sigma _{j1}\sigma _je_je_j\sigma _{j1}q^1q^{\sigma _j}e_je_j.\hfill \end{array}$$
(5.6)
It turns out that the tensors (5.4), (5.5), and (5.6) form a basis of $`(VV)_s`$.
First of all, one shows that the tensors (5.4), (5.5), and (5.6) can be obtained by iterated action of the $`F`$–generators on $`s_{r,r+1}`$. Next one proves that the vector space spanned by these tensors is a $`U_q(𝔰𝔭𝔬(2n|2m))`$–submodule of $`VV`$. Since these tensors obviously are linearly independent, this implies our claim.
Next we have to show that the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_s`$ is irreducible. This can be done as follows. First, we prove the following statement:
If $`x`$ is a non–zero element of $`(VV)_s`$, there exists a monomial $`P`$ in the $`F`$–generators such that $`Px`$ is a non–zero scalar multiple of $`s_{r1,r}`$.
A moment’s thought shows that this is a consequence of the following fact:
If $`x`$ is a (non–zero) weight vector of $`(VV)_s`$ whose weight is different from $`\epsilon _{r1}+\epsilon _r`$ (i.e., if $`x`$ is not a scalar multiple of $`s_{r1,r}`$ ), there exists an index $`jJ`$ such that $`F_jx0`$ .
Since $`s_{r,r+1}`$ is a cyclic vector of the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_s`$, the irreducibility of this module will follow if we can show that there exists a monomial in the $`E`$–generators that maps $`s_{r1,r}`$ onto a non–zero scalar multiple of $`s_{r,r+1}`$. Similar as above, this is a consequence of the following fact:
If $`x`$ is a (non–zero) weight vector of $`(VV)_s`$ whose weight is different from $`\epsilon _r+\epsilon _{r+1}`$ (i.e., if $`x`$ is not a scalar multiple $`s_{r,r+1}`$ ), there exists an index $`jJ`$ such that $`E_jx0`$ .
The proof of the foregoing statements amounts to easy but lengthy calculations. Let us mention that one needs the intermediate result that
$$a(VV)_s.$$
Summarizing, we have proved that the statement b) above is correct.
Obviously, we have
$$\stackrel{~}{b}^q(s_{r,r+1})=0.$$
Since the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_s`$ is irreducible, we conclude that it is contained in the kernel of $`\stackrel{~}{b}^q`$. This proves part of statement c).
### 5.2 The module $`(VV)_a^0`$
Basically, our treatment of the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_a^0`$ follows similar lines to that of $`(VV)_s`$, however, in the cases $`n=m`$ there are several complications.
In the undeformed case, the module $`(VV)_a^0`$ corresponds to the subspace of all super–skew–symmetric tensors in $`VV`$ with a vanishing “symplectic trace” (i.e., which belong to the kernel of $`\stackrel{~}{b}`$). As an $`𝔰𝔭𝔬(b)`$–submodule of $`VV`$, it is generated by the tensors $`e_re_r`$ and $`e_re_r`$, and for $`(n,m)(1,1)`$, each of these tensors alone is already sufficient.
In the present deformed case, it is easy to see that $`e_re_r`$ is annihilated by the $`E`$–generators, and that $`e_re_r`$ is annihilated by the $`F`$–generators. Accordingly, we define $`(VV)_a^0`$ to be the $`U_q(𝔰𝔭𝔬(2n|2m))`$–submodule of $`VV`$ generated by $`e_re_r`$ and $`e_re_r`$.
Let us define the following elements of $`VV`$ :
$$\begin{array}{cccc}\hfill a_{i,j}& =& e_ie_j\sigma _{i,j}qe_je_i\hfill & \text{for }i,jI\text{}i<j\text{ but }ij\hfill \\ \hfill a_{i,i}& =& e_ie_i\hfill & \text{for }iI_{\overline{1}},\hfill \end{array}$$
(5.7)
and for $`2jr`$
$$\begin{array}{ccc}\hfill a_j& =& q^{\sigma _{j1}}e_{j+1}e_{j1}\sigma _{j1}qe_{j1}e_{j+1}\hfill \\ & & \sigma _{j1}\sigma _je_je_j+\sigma _{j1}qq^{\sigma _j}e_je_j.\hfill \end{array}$$
(5.8)
(I hope there is no risk to confound the tensors $`a_{i,j}`$ with the elements of the Cartan matrix.) It turns out that the tensors (5.7) and (5.8) form a basis of $`(VV)_a^0`$.
First one proves that the vector space $`U`$ spanned by the tensors (5.7) and (5.8) is a $`U_q(𝔰𝔭𝔬(2n|2m))`$–submodule of $`VV`$. Since these tensors obviously are linearly independent, they form a basis of $`U`$.
Next one shows that $`a_{r,r}`$ generates this module, provided that $`r3`$ (i.e., provided that $`(n,m)(1,1)`$ ). In the case $`n=m=1`$, we have
$$a_2=a\text{if }m=n=1,$$
and $`a_{2,2}`$ generates the $`U_q(𝔰𝔭𝔬(2|2))`$–submodule
$$V_4=a_{2,2}a_{2,1}a_{2,1}a,$$
while $`a_{2,2}`$ generates the $`U_q(𝔰𝔭𝔬(2|2))`$–submodule
$$\overline{V}_4=a_{2,2}a_{1,2}a_{1,2}a.$$
Obviously, we have
$$V_4+\overline{V}_4=U,V_4\overline{V}_4=a,$$
and it is easy to see that the $`U_q(𝔰𝔭𝔬(2|2))`$–modules $`V_4/a`$ and $`\overline{V}_4/a`$ are irreducible. This proves statement g), and we also have shown that in all cases
$$U=(VV)_a^0.$$
Let us next show that the sum of the subspaces $`(VV)_s`$ and $`(VV)_a^0`$ of $`VV`$ is direct. Obviously, it is sufficient to prove the analogous statement for the corresponding weight spaces. For non–zero weights, this is trivial. To prove the claim for the weight zero, it is sufficient to show that the tensors $`s_i`$ , $`1ir`$, and $`a_j`$ , $`2jr`$, are linearly independent. To show this, we observe that
$$s_j=u_j+q^1v_j,a_j=u_jqv_j\text{for }2jr,$$
(5.9)
where the tensors $`u_j`$ and $`v_j`$; $`2jr`$, are defined by
$$\begin{array}{ccc}\hfill u_j& =& q^{\sigma _{j1}}e_{j+1}e_{j1}\sigma _{j1}\sigma _je_je_j\hfill \\ \hfill v_j& =& \sigma _{j1}e_{j1}e_{j+1}\sigma _{j1}q^{\sigma _j}e_je_j.\hfill \end{array}$$
Consequently, we have to prove that the tensors $`s_1`$ and $`u_j`$ , $`v_j`$ , $`2jr`$, are linearly independent. This follows from the obvious fact that the $`2r`$ tensors $`e_1e_1`$, $`s_1`$ , $`u_j`$ , $`v_j`$ span the same subspace of $`VV`$ as the $`2r`$ tensors $`e_ie_i`$, $`iI`$, namely, the weight space of $`VV`$ corresponding to the weight zero.
The proof above shows that the codimension of $`(VV)_s(VV)_a^0`$ in $`VV`$ is equal to one. Obviously, $`\stackrel{~}{b}^q`$ vanishes on the tensors $`e_re_r`$ and $`e_re_r`$, hence also on the submodule $`(VV)_a^0`$ generated by them. As noted earlier, $`\stackrel{~}{b}^q`$ also vanishes on $`(VV)_s`$. Since $`\stackrel{~}{b}^q`$ is a non–zero linear form on $`VV`$, it follows that $`(VV)_s(VV)_a^0`$ is the kernel of $`\stackrel{~}{b}^q`$. This proves the last claim of statement c).
Using Eqn. (5.2), our last result implies that
$$a(VV)_s(VV)_a^0\text{if and only if }n=m.$$
(5.10)
Since the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_s`$ is irreducible and not one–dimensional, it follows that
$$a(VV)_a^0\text{for }n=m.$$
(5.11)
Indeed, it can be shown that
$$a=\underset{j=2}{\overset{n+1}{}}[j1]a_j\underset{j=n+2}{\overset{2n}{}}[2n+1j]a_j\text{for }n=m,$$
(5.12)
where, for all integers $`s`$, the $`q`$–number $`[s]`$ is defined by
$$[s]=[s]_q=\frac{q^sq^s}{qq^1}.$$
For $`n=1`$, Eqn. (5.12) is just the equation $`a=a_2`$ mentioned earlier.
Using Eqn. (5.10) and the fact that $`(VV)_s(VV)_a^0`$ is equal to the kernel of $`\stackrel{~}{b}^q`$, it follows that, for $`n=m`$, this submodule does not have a module complement in $`VV`$. In fact, any such complement would have to be a trivial one–dimensional submodule of $`VV`$, and hence would be spanned by an invariant element of $`VV`$. But the invariant elements of $`VV`$ are the scalar multiples of $`a`$. Combined with Eqn. (5.11), this yields statement e).
Finally, to answer questions of irreducibility, we prove the following technical results.
Suppose that $`r3`$, and that $`x(VV)_a^0`$ is a (non–zero) weight vector which is neither proportional to $`a_{r,r}`$ nor to $`a`$. Then there exists an index $`jJ`$ such that $`F_jxa`$ (in particular, we have $`F_jx0`$ ).
Suppose that $`r3`$, and that $`x(VV)_a^0`$ is a (non–zero) weight vector which is neither proportional to $`a_{r,r}`$ nor to $`a`$. Then there exists an index $`jJ`$ such that $`E_jxa`$ (in particular, we have $`E_jx0`$ ).
As in the case of $`(VV)_s`$, these results follow from easy but lengthy calculations. Once they are established, it is easy to draw the following conclusions.
If $`nm`$, the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_a^0`$ is irreducible.
If $`n=m2`$, the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_a^0/a`$ is irreducible.
These results prove statement f) and the first claim of statement d).
### 5.3 The module $`(VV)_a`$
Our next task is to construct the submodule $`(VV)_a`$ of $`VV`$. In the case $`nm`$, this is easy. Recalling Eqn. (5.10) and the fact that $`(VV)_s(VV)_a^0`$ has the codimension one in $`VV`$, it follows that
$$VV=(VV)_s(VV)_a^0a\text{if }nm.$$
(5.13)
As we know, the submodules on the right hand side of this equation are irreducible, moreover, they are obviously non–isomorphic. This implies that Eqn. (5.13) is the unique decomposition of the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$ into irreducible submodules. Setting
$$(VV)_a=(VV)_a^0a,$$
we have proved the statements a)–d) in the case $`nm`$.
Unfortunately, this type of reasoning is not possible in the case $`n=m`$. Since we want to obtain a unified treatment of the problem, we start all over again and present an approach which is applicable in all cases.
To begin with, we note that every module complement of $`(VV)_s`$ in $`VV`$ must take the form $`(VV)_a^0g`$, where $`g`$ is a weight vector of $`VV`$ of weight zero. Indeed, since the tensors $`e_re_r`$ and $`e_re_r`$ do not belong to $`(VV)_s`$, and since the corresponding weights have multiplicity one, these two tensors and hence the submodule generated by them must be contained in every module complement. Obviously, a subspace of this type is a module complement of $`(VV)_s`$ if and only if $`g(VV)_s(VV)_a^0`$ and if $`E_jg`$ and $`F_jg`$ belong to $`(VV)_a^0`$, for all $`jJ`$.
Since $`g`$ is of zero weight, it takes the form
$$g=\underset{iI}{}g_ie_ie_i,$$
(5.14)
with some coefficients $`g_i`$ . If $`E_jg`$ is non–zero, it is a weight vector with non–zero weight and hence belongs to $`(VV)_s(VV)_a^0`$. Consequently, $`E_jg`$ lies in $`(VV)_a^0`$ if and only if its component in $`(VV)_s`$ is equal to zero. It follows that we have $`E_jg(VV)_a^0`$ if and only if the coefficients $`g_i`$ satisfy the following system of linear equations:
$$g_1+q^2g_1=0$$
(5.15)
$$qg_{j+1}qq^{\sigma _j}g_j=\sigma _jg_j\sigma _{j+1}q^{\sigma _{j+1}}g_{j1}\text{for }rj2.$$
(5.16)
The condition that $`F_jg(VV)_a^0`$ for all $`jJ`$ is equivalent to the same system of equations.
The general solution of this system can easily be described: We can choose $`g_1`$ , $`g_2`$ , …, $`g_r`$ arbitrarily, and then the coefficients $`g_1`$ , $`g_2`$ , …, $`g_r`$ are uniquely fixed.
Let $`X_a`$ be the subspace of $`VV`$ consisting of all tensors of the form (5.14) such that the coefficients $`g_i`$ satisfy the system (5.15), (5.16). According to the foregoing result, this subspace is $`r`$–dimensional. Obviously, $`X_a`$ contains the $`(r1)`$–dimensional weight space of $`(VV)_a^0`$ corresponding to the weight zero. On the other hand, $`X_a`$ does not contain any non–zero elements of $`(VV)_s`$ (indeed, any such element would be invariant and hence proportional to $`a`$ ). It follows that
$$(VV)_a=(VV)_a^0+X_a$$
is a module complement of $`(VV)_s`$ in $`VV`$.
By this latter property, $`(VV)_a`$ is uniquely fixed. In fact, let $`(VV)_a^{}`$ be an arbitrary module complement of $`(VV)_s`$ in $`VV`$. As noted at the beginning of this discussion, $`(VV)_a^{}`$ contains the submodule $`(VV)_a^0`$. Let $`X_a^{}`$ be the weight space of $`(VV)_a^{}`$ corresponding to the weight zero. Of course, $`X_a^{}`$ is $`r`$–dimensional. Consider an arbitrary element $`gX_a^{}`$. For every $`jJ`$, the elements $`E_jg`$ and $`F_jg`$ lie in $`(VV)_a^{}`$ and have a non–zero weight, which implies that they are elements of $`(VV)_a^0`$. By the definition of $`X_a`$, this proves that $`gX_a`$. Thus we have shown that $`X_a^{}X_a`$, and for reasons of dimension, this implies that $`X_a^{}=X_a`$. It follows that $`X_a(VV)_a^{}`$ , hence that
$$(VV)_a(VV)_a^{},$$
and finally
$$(VV)_a=(VV)_a^{},$$
as claimed.
We proceed by choosing, in a unified way, an element $`tX_a`$ that does not belong to $`(VV)_a^0`$. Let $`t`$ be the tensor of the form (5.14) whose coefficients $`g_i^t`$ are fixed by the requirement that
$$g_1^t=1,g_j^t=0\text{for }rj2.$$
Recall that the coefficients $`g_i^t`$ with $`1ir`$ can then be calculated by means of the system (5.15), (5.16). We obtain
$$t=e_1e_1q^2e_1e_1+(qq^1)\underset{i=2}{\overset{n+m}{}}((C^q)^1)_{i,i}e_ie_i.$$
(5.17)
It is easy to check that
$$\stackrel{~}{b}^q(t)=q^1(q^{2d+2}+1)$$
(where $`d=nm`$; see Eqn. (5.1)). Thus $`t`$ does not belong to the kernel of $`\stackrel{~}{b}^q`$ (since $`q`$ is not a root of unity).
According to the results of the present subsection, the tensor $`t`$ and the tensors (5.7) and (5.8) form a basis of the vector space $`(VV)_a`$.
Remark 5.1. We note that some other, simpler looking choices for the tensor $`t`$ are possible. For example, the tensor $`e_re_r+e_re_r`$ is a candidate. However, it is not at all obvious that such a choice would simplify the subsequent calculations.
The reader will easily convice himself/herself that, at this stage, we have proved the statements a)–g).
### 5.4 The module endomorphisms of $`VV`$
In the present subsection we are going to prove statement h), i.e., that the linear operators $`\text{id}_{VV}`$, $`P_s`$, and $`K`$ form a basis of the space of all endomorphisms of the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$ (recall that $`P_s`$ is the projector of $`VV`$ onto $`(VV)_s`$ with kernel $`(VV)_a`$, and that the map $`K`$ has been defined in Eqn. (5.3)). Obviously, the maps $`\text{id}_{VV}`$ and $`P_s`$ are module endomorphisms, and since $`\stackrel{~}{b}^q`$ and $`a`$ are invariant, the same is true of $`K`$.
Now let $`Q`$ be a module endomorphism of $`VV`$, i.e., an even linear map of $`VV`$ into itself that commutes with the action of $`U_q(𝔰𝔭𝔬(2n|2m))`$. Since every weight vector of $`VV`$ of weight $`\epsilon _r+\epsilon _{r+1}`$, that is annihilated by all $`E_j`$, $`jJ`$, is proportional to $`s_{r,r+1}`$, there exists a constant $`c_s`$ such that
$$Q(s_{r,r+1})=c_ss_{r,r+1}.$$
Similarly, since the multiplicity of the weight $`2\epsilon _r`$ is equal to one, we have
$$Q(a_{r,r})=c_aa_{r,r},$$
with some constant $`c_a`$. It follows that
$$\begin{array}{cccc}\hfill Q(x)& =& c_sx\hfill & \text{for all }x(VV)_s\hfill \\ \hfill Q(y)& =& c_ay\hfill & \text{for all }y(VV)_a^0.\hfill \end{array}$$
(5.18)
Indeed, since the tensor $`s_{r,r+1}`$ generates the submodule $`(VV)_s`$ , the first of these equations follows immediately. Similarly, for $`r3`$ the tensor $`a_{r,r}`$ generates the submodule $`(VV)_a^0`$ , which implies the second equation in this case.
In the case $`r=2`$ , i.e., for $`m=n=1`$ , we argue as follows. Quite generally, we also have
$$Q(a_{r,r})=\overline{c}_aa_{r,r},$$
with some constant $`\overline{c}_a`$. Then, for $`m=n=1`$, we conclude as above that
$$\begin{array}{cccc}\hfill Q(y)& =& c_ay\hfill & \text{for all }yV_4\hfill \\ \hfill Q(\overline{y})& =& \overline{c}_a\overline{y}\hfill & \text{for all }\overline{y}\overline{V}_4.\hfill \end{array}$$
Since $`a`$ lies in $`V_4`$ and in $`\overline{V}_4`$, this implies that
$$c_a=\overline{c}_a,$$
and since
$$V_4+\overline{V}_4=(VV)_a^0,$$
it follows that the second of the Eqns. (5.18) holds for $`n=m=1`$ as well.
Now let
$$P_a=\text{id}_{VV}P_s$$
be the projector of $`VV`$ onto $`(VV)_a`$ with kernel $`(VV)_s`$. Then the equations derived above show that
$$Q^{}=Qc_sP_sc_aP_a$$
is a module endomorphism of $`VV`$, that vanishes on
$$\text{ker}\stackrel{~}{b}^q=(VV)_s(VV)_a^0.$$
Consequently, it induces a module homomorphism
$$(VV)/\text{ker}\stackrel{~}{b}^qVV.$$
The module on the left hand side is one–dimensional and trivial, and all invariants in $`VV`$ are proportional to $`a`$. This implies that
$$Q^{}=c_0K,$$
with some constant $`c_0`$, which proves our claim.
We close this subsection by the remark that the maps $`\text{id}_{VV}`$, $`P_s`$, and $`K`$ commute one with another.
## 6 Calculation of the $`R`$–matrix
At last, we are prepared to calculate the $`R`$–matrix $`R`$ or, equivalently, the braid generator $`\widehat{R}`$ of $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the vector representation. By definition, $`R`$ is the representative of the universal $`R`$–matrix $``$ in the vector representation,
$$R=_{VV},$$
and $`\widehat{R}`$ is given by
$$\widehat{R}=PR,$$
where
$$P:VVVV$$
denotes the twist operator (in the super sense), which is given by
$$P(xy)=(1)^{\xi \eta }yx,$$
for all $`xV_\xi `$, $`yV_\eta `$, with $`\xi ,\eta _2`$. To calculate the $`R`$–matrix (or the braid generator) means to calculate its matrix elements with respect to the basis
$`(e_ie_j)_{i,jI}`$ of $`VV`$.
Remark 6.1. Due to the fact that $``$ is given in terms of a formal power series, the foregoing remarks need to be amended. See below for further details.
In order to perform the calculation we observe that $`\widehat{R}`$ is an endomorphism of the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`VV`$. According to Section 5.4, this implies that $`\widehat{R}`$ is a linear combination of $`\text{id}_{VV}`$, $`K`$, and $`P_s`$. Since the matrix elements of $`\text{id}_{VV}`$ and $`K`$ are known, our task consists of two pieces, namely, to calculate the projector $`P_s`$ (or any module endomorphism of $`VV`$ “containing” $`P_s`$ in a non–trivial way), and to find the aforementioned linear combination. As we are going to see, the second problem can easily be dealt with once the first problem has been solved.
Obviously, every module endomorphism of $`VV`$ maps each of the weight spaces into itself. Consequently, the first problem splits into a number of subproblems, one for each of the weight spaces of $`VV`$. Since the weight spaces corresponding to the non–zero weights are at most two–dimensional, the corresponding subproblems are trivial, and we are left with the subproblem corresponding to the zero weight. Basically, this latter problem amounts to writing the tensors $`e_ie_i`$, $`iI`$, as linear combinations of the tensors (5.4) – (5.8) and $`t`$, i.e., we have to invert a certain $`(2r\times 2r)`$–matrix, whose elements are rational functions of $`q`$.
Unfortunately, the corresponding calculations turn out to be rather tedious. Accordingly, I don’t present the details of this calculation but only mention two points. First, in the course of the calculations I have taken advantage of the tensors $`u_j`$ and $`v_j`$ introduced in Eqn. (5.9) and of the resulting equations
$$P_s(u_j)=q(q+q^1)^1s_j,P_s(v_j)=(q+q^1)^1s_j.$$
Secondly, I have applied the following simple trick. In Ref. , the formulae for $`\widehat{R}`$ are slightly simpler than those for $`P_s`$. On the other hand, again according to Ref. (see also Ref. ), it is tempting to conjecture that
$$\widehat{R}=\widehat{R}^{},$$
(6.1)
where the endomorphism $`\widehat{R}^{}`$ of $`VV`$ is defined by
$$\widehat{R}^{}=(q+q^1)P_sq^1𝕀𝕀(qq^1)(1+q^{2d+2})^1K.$$
(6.2)
Here and in the following, $`𝕀`$ denotes the unit operator of $`V`$:
$$𝕀=\text{id}_V.$$
Accordingly, I haven’t calculated $`P_s`$ but rather the operator $`\widehat{R}^{}`$. The Eqn. (6.1) can then be proved at a later stage (which solves the second problem mentioned at the beginning of this section).
A long calculation shows that for $`1jn`$
$$\begin{array}{ccc}\hfill \widehat{R}^{}(e_je_j)& =& (qq^1)\underset{i=1}{\overset{n}{}}q^{ji}e_ie_i\hfill \\ & & +(qq^1)\underset{i=n+1}{\overset{n+m}{}}q^{2n+ij}e_ie_i\hfill \\ & & (qq^1)\underset{i=1}{\overset{j1}{}}q^{ij}e_ie_i\hfill \\ & & +(qq^1)e_je_j+q^1e_je_j\hfill \end{array}$$
$$\begin{array}{ccc}\hfill \widehat{R}^{}(e_je_j)& =& (qq^1)\underset{i=j+1}{\overset{n}{}}q^{ji}e_ie_i\hfill \\ & & (qq^1)\underset{i=n+1}{\overset{n+m}{}}q^{2n+i+j}e_ie_i\hfill \\ & & +q^1e_je_j,\hfill \end{array}$$
and for $`n+1jn+m`$
$$\begin{array}{ccc}\hfill \widehat{R}^{}(e_je_j)& =& (qq^1)\underset{i=1}{\overset{n}{}}q^{2ni+j2}e_ie_i\hfill \\ & & (qq^1)\underset{i=n+1}{\overset{n+m}{}}q^{4n+i+j2}e_ie_i\hfill \\ & & +(qq^1)\underset{i=1}{\overset{n}{}}q^{2n+i+j2}e_ie_i\hfill \\ & & (qq^1)\underset{i=n+1}{\overset{j1}{}}q^{ji}e_ie_i\hfill \\ & & +(qq^1)e_je_jqe_je_j\hfill \end{array}$$
$$\begin{array}{ccc}\hfill \widehat{R}^{}(e_je_j)& =& (qq^1)\underset{i=j+1}{\overset{n+m}{}}q^{ij}e_ie_i\hfill \\ & & qe_je_j.\hfill \end{array}$$
It might have been difficult to unify these equations in a concise formula. Fortunately, Ref. suggests that, for all $`iI`$, we have
$$\begin{array}{ccc}\hfill \widehat{R}^{}(e_ie_i)& =& \sigma _iq^{\sigma _i}e_ie_i+(qq^1)\theta (i>i)e_ie_i\hfill \\ & & (qq^1)C_{i,i}^q\underset{k>i}{}((C^q)^1)_{k,k}e_ke_k,\hfill \end{array}$$
where, for all $`i,jI`$, the symbol $`\theta (j>i)`$ is defined by
$$\theta (j>i)=\{\begin{array}{cc}\hfill 1& \text{if }j>i\hfill \\ \hfill 0& \text{otherwise\hspace{0.17em}.}\hfill \end{array}$$
It is not difficult to see that this is indeed the case.
The remaining tensors $`\widehat{R}^{}(e_ie_j)`$, with $`i,jI`$, are easily determined:
If $`i<j`$, but $`ij`$, we obtain
$$\begin{array}{ccc}\hfill \widehat{R}^{}(e_ie_j)& =& \sigma _{i,j}e_je_i+(qq^1)e_ie_j\hfill \\ \hfill \widehat{R}^{}(e_je_i)& =& \sigma _{j,i}e_ie_j.\hfill \end{array}$$
On the other hand, we find for all $`iI`$
$$\widehat{R}^{}(e_ie_i)=\sigma _iq^{\sigma _i}e_ie_i.$$
Let us next prove Eqn. (6.1). Needless to say, at this point we have to make contact with the theory of the universal $`R`$–matrix. Fortunately, only very little of that theory is needed. Hence it should be sufficient to recall a few simple facts. For more details, we refer the reader to Ref. .
Usually, the theory is formulated in the framework of formal power series in one indeterminate $`h`$. In particular, the complex parameter $`q`$ is replaced by
$$q=e^h,$$
and the corresponding quantum superalgebra $`U_h`$ is a topological Hopf superalgebra over the ring $`[[h]]`$ of formal power series in $`h`$. The Cartan subalgebra $`𝔥`$ of $`𝔰𝔭𝔬(2n|2m)`$ is regarded as a subspace of $`U_h`$, and in this sense, the elements $`K_i`$ are given by
$$K_i=\mathrm{exp}(hH_{\alpha _i}),H_{\alpha _i}=d_iH_i,$$
where the elements $`H_\lambda `$ have been defined in Eqn. (2.5), and the elements $`H_i=H_i^v`$ in Eqn. (2.8).
In this setting, instead of $`V`$ we have to consider the $`[[h]]`$–module $`V[[h]]`$ that is obtained from $`V`$ by an extension of the domain of scalars from $``$ to $`[[h]]`$:
$$V[[h]]=V\begin{array}{c}\\ \end{array}[[h]].$$
It is well–known that $`V[[h]]`$ is a graded $`U_h`$–module in a natural way, and it is this module which in the present setting takes the role of the vector module. More explicitly, the elements $`e_i1V[[h]]`$, $`iI`$, form a basis of the $`[[h]]`$–module $`V[[h]]`$, and it is customary to identify $`e_i1`$ with $`e_i`$. With this convention, the action of the generators $`E_j`$, $`F_j`$, $`K_j`$; $`jJ`$, on the basis elements is given by the same formulae as in Section 4 (of course, with a different meaning of $`q`$ ).
The tensor product (over $`[[h]]`$) of $`V[[h]]`$ with itself is also a graded $`U_h`$–module. On the other hand, it is known that
$$V[[h]]\begin{array}{c}\\ \left[\left[h\right]\right]\end{array}V[[h]](V\begin{array}{c}\\ \end{array}V)\begin{array}{c}\\ \end{array}[[h]]$$
(as graded $`[[h]]`$–modules), and the elements $`e_ie_j1`$; $`i,jI`$, form a basis of this $`[[h]]`$–module. Once again, we identify $`e_ie_j1`$ with $`e_ie_j`$, and then the action of the generators $`E_j`$, $`F_j`$, $`K_j`$ is given by the same formulae as in Section 5.
Using these conventions, the arguments of Section 5 can be adopted almost verbatim. Of course, we have to keep in mind that $`[[h]]`$ is not a field but only a ring. Correspondingly, at various instances we have to observe that a scalar is not only different from zero but even invertible. Moreover, the usual concept of an irreducible module is not useful here: If $`W`$ is a graded $`U_h`$–module, then $`hW`$ is a graded submodule of $`W`$ and, in general, different from $`W`$.
In particular, the $`U_h`$–module endomorphisms of $`V[[h]]_{[[h]]}V[[h]]`$ are linear combinations of the identity map, $`K`$, and $`\widehat{R}^{}`$ (interpreted as $`[[h]]`$–linear maps of $`(VV)[[h]]`$ into itself). Now we can prove Eqn. (6.1) (in the present setting), but let us first complete our survey. Eqn. (6.1) shows that $`\widehat{R}`$ depends on $`h`$ only through $`e^h`$ (a result which immediately follows by inspection of the formula for the universal $`R`$–matrix). In fact, its matrix elements are Laurent polynomials in $`e^h`$. Substituting for $`e^h`$ the complex number $`q`$ we started with, we obtain the braid generator that we want to calculate.
Before we can proceed to the calculation proper, we remind the reader of the general form of the universal $`R`$–matrix. Using the fact that the vector spaces $`(𝔥𝔥)^{}`$ and $`𝔥^{}𝔥^{}`$ are canonically isomorphic, it is obvious that there exists a unique tensor $`B𝔥𝔥`$ such that
$$(\lambda \mu )(B)=(\lambda |\mu )\text{for all }\lambda ,\mu 𝔥^{}.$$
In terms of this tensor, we have
$$=e^{hB}(11+\mathrm{}),$$
(6.3)
where the dots stand for an infinite sum of terms of the form $`XX^{}`$, in which $`X`$ and $`X^{}`$ are weight vectors of the quantum superalgebra with non–zero (and opposite) weights (see Refs. , ).
Now we are ready to prove Eqn. (6.1). According to the preceding discussion, $`\widehat{R}`$ is a $`[[h]]`$–linear combination of $`\widehat{R}^{}`$, $`𝕀𝕀`$, and $`K`$ (regarded as $`[[h]]`$–linear maps of $`(VV)[[h]]`$ into itself). Equivalently, this means that
$$_{(VV)[[h]]}=aP\widehat{R}^{}+bP+cPK,$$
(6.4)
with some coefficients $`a,b,c[[h]]`$. In order to determine these coefficients, we apply Eqn. (6.4) to the tensors $`e_ie_j`$ and keep only the diagonal terms, i.e., the terms proportional to $`e_ie_j`$. On the left hand side, the terms indicated by the dots in Eqn. (6.3) do not contribute, and we are left with
$$\begin{array}{ccc}\hfill (e^{hB})(e_ie_j)& =& q^{(\epsilon _i|\epsilon _j)}e_ie_j\hfill \\ & =& \{\begin{array}{cc}e_ie_j\hfill & \text{if }ij,j\hfill \\ q^{\sigma _i}e_ie_i\hfill & \text{if }j=i\hfill \\ q^{\sigma _i}e_ie_i\hfill & \text{if }j=i.\hfill \end{array}\hfill \end{array}$$
Using the formulae for $`\widehat{R}^{}(e_ie_j)`$ obtained above, the analogous terms on the right hand side can easily be calculated. Comparing both sides, we obtain the following equations:
For $`ij,j`$
$$1=a,$$
for $`i=j`$
$$q^{\sigma _i}=aq^{\sigma _i}+b\sigma _i,$$
for $`i=j`$
$$q^{\sigma _i}=aq^{\sigma _i}+c\sigma _i.$$
The unique solution of these equations is
$$a=1,b=c=0.$$
This implies that
$$_{(VV)[[h]]}=P\widehat{R}^{},$$
which proves our claim.
Remark 6.1. It is well–known that if $``$ is a universal $`R`$–matrix for a Hopf superalgebra $`H`$, then so is $`_{21}^1`$ (we are using the standard notation). Moreover, if $`V`$ is any graded $`H`$–module, and if $`\widehat{R}`$ is the braid generator in $`VV`$ with respect to $``$, then $`\widehat{R}^1`$ is the braid generator in $`VV`$ with respect to $`_{21}^1`$. Thus we can apply the preceding discussion to $`_{21}^1`$ and $`\widehat{R}^1`$. First of all, we conclude that
$$P\widehat{R}^1=(_{21}^1)_{(VV)[[h]]}=a^{}P\widehat{R}^{}+b^{}P+c^{}PK,$$
with some coefficients $`a^{},b^{},c^{}[[h]]`$. Since the tensor $`B`$ obviously is symmetric, we conclude from Eqn. (6.3) that
$$_{21}^1=(11+\mathrm{})e^{hB},$$
where the dots stand for terms similar to those in Eqn. (6.3). Proceeding as above, we can show that
$$a^{}=1,c^{}=b^{}=qq^1,$$
which implies that
$$\begin{array}{ccc}\hfill \widehat{R}^1& =& \widehat{R}^{}(qq^1)𝕀𝕀+(qq^1)K\hfill \\ & =& (q+q^1)P_sq𝕀𝕀+(qq^1)(1+q^{2d2})^1K.\hfill \end{array}$$
It is easy to show directly that the operator on the right hand side really is the inverse of $`\widehat{R}=\widehat{R}^{}`$, which is a first check that our calculations are correct.
Summarizing part of the results of the present section, we have shown that
$$\begin{array}{ccc}\hfill \widehat{R}& =& \underset{i}{}\sigma _iq^{\sigma _i}E_{i,i}E_{i,i}+\underset{i}{}\sigma _iq^{\sigma _i}E_{i,i}E_{i,i}\hfill \\ & & +\underset{ij,j}{}\sigma _{i,j}E_{j,i}E_{i,j}\hfill \\ & & +(qq^1)\underset{i<j}{}E_{i,i}E_{j,j}\hfill \\ & & (qq^1)\underset{i<j}{}((C^q)^1)_{j,j}C_{i,i}^qE_{j,i}E_{j,i},\hfill \end{array}$$
or equivalently, that
$$\begin{array}{ccc}\hfill R& =& \underset{i}{}q^{\sigma _i}E_{i,i}E_{i,i}+\underset{i}{}q^{\sigma _i}E_{i,i}E_{i,i}\hfill \\ & & +\underset{ij,j}{}E_{i,i}E_{j,j}\hfill \\ & & +(qq^1)\underset{i<j}{}\sigma _{i,j}E_{j,i}E_{i,j}\hfill \\ & & (qq^1)\underset{i<j}{}\sigma _j((C^q)^1)_{j,j}C_{i,i}^qE_{j,i}E_{j,i}.\hfill \end{array}$$
In these equations, the indices $`i`$ and $`j`$ run through the index set $`I`$, subject to conditions as specified. The equations hold in both settings, the one in terms of formal power series, and the one where $`q`$ is a complex number. It should also be noted that $`E_{i,j}E_{k,\mathrm{}}`$ denotes the normal (non–super) tensor product of linear mappings (the graded tensor product of two linear maps $`f`$ and $`g`$ is denoted by $`f\overline{}g`$).
## 7 Properties of $`R`$ and $`\widehat{R}`$
In the present section, we want to collect some of the basic relations satisfied by $`R`$ and $`\widehat{R}`$. First of all, we recall the following equations, which have been derived in the preceding section:
$$\widehat{R}=(q+q^1)P_sq^1𝕀𝕀(qq^1)(1+q^{2d+2})^1K$$
(7.1)
$$\widehat{R}^1=(q+q^1)P_sq𝕀𝕀+(qq^1)(1+q^{2d2})^1K.$$
(7.2)
Using the results of Section 5 (in particular, Eqn. (5.2)), the first of these equations implies that
$$\begin{array}{cccc}\hfill \widehat{R}(x)& =& qx\hfill & \text{for all }x(VV)_s\hfill \\ \hfill \widehat{R}(y)& =& q^1y\hfill & \text{for all }y(VV)_a^0\hfill \\ \hfill \widehat{R}(a)& =& q^{2d1}a\hfill & .\hfill \end{array}$$
Moreover, the linear map induced by $`\widehat{R}`$ in the one–dimensional $`U_q(𝔰𝔭𝔬(2n|2m))`$–module $`(VV)_a/(VV)_a^0`$ is equal to the multiplication by $`q^{2d1}`$. For $`nm`$, this follows from the last equation.
Next we recall that the $`U_q(𝔰𝔭𝔬(2n|2m))`$–module endomorphisms of $`VV`$ commute one with another. In particular, we have
$$P_sK=KP_s=0$$
$$\widehat{R}K=K\widehat{R}=q^{2d1}K.$$
On the other hand, the Eqns. (7.1) and (7.2) imply that
$$\widehat{R}\widehat{R}^1=(qq^1)(𝕀𝕀K).$$
Obviously, this equation is equivalent with
$$\widehat{R}^2(qq^1)\widehat{R}𝕀𝕀=(qq^1)\widehat{R}K,$$
and hence also with
$$(\widehat{R}q𝕀𝕀)(\widehat{R}+q^1𝕀𝕀)=(qq^1)q^{2d1}K.$$
(7.3)
Since the image of the operator $`K`$ is contained in $`a`$ , it follows that
$$(\widehat{R}q𝕀𝕀)(\widehat{R}+q^1𝕀𝕀)(\widehat{R}+q^{2d1}𝕀𝕀)=\mathrm{\hspace{0.17em}0}.$$
(7.4)
It is easy to see that the polynomial $`(Xq)(X+q^1)(X+q^{2d1})`$ involved in Eqn. (7.4) is the minimal polynomial of the operator $`\widehat{R}`$.
The preceding equations can be used to write the spectral projectors of $`\widehat{R}`$ as polynomials in $`\widehat{R}`$ (to the extent in which these projectors exist). For example, we find
$$P_s=\frac{(\widehat{R}+q^1)(\widehat{R}+q^{2d1})}{(q+q^1)(q+q^{2d1})}.$$
Moreover, we stress that according to Eqn. (7.3), the operator $`K`$ can be written as a polynomial in $`\widehat{R}`$. This fact (which is not true in the undeformed case $`q=1`$) will turn out to be crucial in the construction of the quantum supergroup $`\text{SPO}_q(2n|2m)`$.
Let us now derive two relations which are related to the fact that on $`V`$ there exists an invariant bilinear form, namely, the form $`b^q`$ found in Section 4. We use the results of Appendix A and argue as in Ref. , for the original setting in which $`q`$ is a complex number. The reader who is not satisfied by this sloppy procedure may either reformulate everything in terms of formal power series, or else regard the final result Eqn. (7.5) as a conjecture which has to be (and has been) checked independently.
Let
$$f_{\mathrm{}}:VV^{\mathrm{gr}},f_r:VV^{\mathrm{gr}}$$
be the linear maps associated to $`b^q`$ (see Appendix A). Like $`b^q`$ they are homogeneous of degree zero.
We write the universal $`R`$–matrix of $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the form
$$=\underset{^s}{}R_s^1R_s^2,$$
where $`R_s^1,R_s^2U_q(𝔰𝔭𝔬(2n|2m))`$. It is well–known that
$$^1=(S\text{id})().$$
This implies that
$$R^1=_{VV}^1=(S\text{id})()_{VV}=\underset{^s}{}S(R_s^1)_V\overline{}(R_s^2)_V,$$
where $`\overline{}`$ denotes the tensor product of linear maps in the graded sense. Using Eqn. (A.5), we conclude that
$$\begin{array}{ccc}\hfill R^1& =& \underset{^s}{}f_r^1\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}((R_s^1)_V)^{\mathrm{st}}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}f_r\overline{}(R_s^2)_V\hfill \\ & =& (f_r^1\overline{}𝕀)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(\underset{^s}{}((R_s^1)_V)^{\mathrm{st}}\overline{}(R_s^2)_V)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(f_r\overline{}𝕀)\hfill \\ & =& (f_r^1𝕀)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}R^{\mathrm{st}_1}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(f_r𝕀),\hfill \end{array}$$
where $`^{\mathrm{st}_1}`$ denotes the super–transposition of the first tensorial factor (for more details, see Appendix B).
Similarly, we can start from the equation
$$^1=(\text{id}S^1)()$$
and derive that
$$R^1=_{VV}^1=\underset{^s}{}(R_s^1)_V\overline{}S^1(R_s^2)_V.$$
According to Eqn. (A.3), this implies that
$$\begin{array}{ccc}\hfill R^1& =& \underset{^s}{}(R_s^1)_V\overline{}f_{\mathrm{}}^1\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}((R_s^2)_V)^{\mathrm{st}}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}f_{\mathrm{}}\hfill \\ & =& (𝕀\overline{}f_{\mathrm{}}^1)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(\underset{^s}{}(R_s^1)_V\overline{}((R_s^2)_V)^{\mathrm{st}})\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(𝕀\overline{}f_{\mathrm{}})\hfill \\ & =& (𝕀f_{\mathrm{}}^1)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}R^{\mathrm{st}_2}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(𝕀f_{\mathrm{}}),\hfill \end{array}$$
where $`^{\mathrm{st}_2}`$ denotes the super–transposition of the second tensorial factor. Summarizing, we have shown that
$$R^1=(f_r^1𝕀)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}R^{\mathrm{st}_1}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(f_r𝕀)=(𝕀f_{\mathrm{}}^1)\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}R^{\mathrm{st}_2}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}(𝕀f_{\mathrm{}}).$$
(7.5)
Note that these equations imply that $`R^{\mathrm{st}_1}`$ and $`R^{\mathrm{st}_2}`$ are invertible.
The equations (7.5) can be checked directly. To do that we need the matrices of the linear maps $`f_{\mathrm{}}`$ and $`f_r`$. If $`(e_i^{})_{iI}`$ is the basis of $`V^{\mathrm{gr}}`$ dual to $`(e_i)_{iI}`$, we find for all $`jI`$
$$f_{\mathrm{}}(e_j)=\underset{iI}{}C_{j,i}^qe_i^{},f_r(e_j)=\underset{iI}{}\sigma _{j,i}C_{i,j}^qe_i^{}.$$
We also need a formula for $`R^1`$. Using the equation
$$R^1=\widehat{R}^1P=PRP(qq^1)(PKP),$$
we derive that
$$\begin{array}{ccc}\hfill R^1& =& \underset{i}{}q^{\sigma _i}E_{i,i}E_{i,i}+\underset{i}{}q^{\sigma _i}E_{i,i}E_{i,i}\hfill \\ & & +\underset{ij,j}{}E_{i,i}E_{j,j}\hfill \\ & & (qq^1)\underset{i<j}{}\sigma _{i,j}E_{j,i}E_{i,j}\hfill \\ & & +(qq^1)\underset{i<j}{}\sigma _i((C^q)^1)_{j,j}C_{i,i}^qE_{j,i}E_{j,i}.\hfill \end{array}$$
Recalling the formulae for the partial super–transpose given in Appendix B, it is now not difficult to show that the equations (7.5) are indeed satisfied.
A closer look at the formula for $`R^1`$ reveals that, somewhat unexpectedly, $`R_q^1`$ is not equal to $`R_{q^1}`$ (we are using the obvious notation). This fact is closely related to Eqn. (4.2).
In the purely symplectic case considered in Ref. it is known that $`\widehat{R}^{\mathrm{t}_1\mathrm{t}_2}`$ is equal to $`\widehat{R}`$ (where $`^{\mathrm{t}_1}`$ and $`^{\mathrm{t}_2}`$ denote the usual transposition of the first resp. second tensorial factor). For reasons similar to those above, I have not been able to derive an analogous equation in the present setting.
We proceed by recalling that the general theory of quasitriangular Hopf superalgebras implies that $`R`$ satisfies the graded Yang–Baxter equation. Equivalently, this means that $`\widehat{R}`$ satisfies the braid relation
$$(\widehat{R}𝕀)(𝕀\widehat{R})(\widehat{R}𝕀)=(𝕀\widehat{R})(\widehat{R}𝕀)(𝕀\widehat{R}).$$
It would be worth–while to check this relation directly, but I haven’t done that.
Finally, I have shown by explicit calculation that $`\widehat{R}`$ and $`K`$ satisfy the following relations:
$$(𝕀K)(\widehat{R}^{\pm 1}𝕀)(𝕀K)=q^{\pm (2d+1)}(𝕀K)$$
$$(K𝕀)(𝕀\widehat{R}^{\pm 1})(K𝕀)=q^{\pm (2d+1)}(K𝕀).$$
Summarizing part of the results of the present section, we conclude that $`\widehat{R}`$ and $`K`$ generate representations of the Birman, Wenzl, Murakami algebras , as defined in Ref. (with $`z=q^{2d+1}`$).
## 8 Comparison with known special cases
In a few special cases, the $`R`$–matrix calculated in this work has already been known.
1. The case $`m=0`$
It should be obvious to the reader that our results apply in the case $`m=0`$ as well, and this case has been settled in Ref. . Actually, I have used this fact throughout the whole investigation: It enabled me to check my calculations and to guess a concise expression for the $`R`$–matrix. For greater clarity, let us mark the entries of the present work by the subscript “here” and those of Ref. by the subscript “RTF”, moreover, let us indicate the dependence on the parameter $`q`$ by the superscript $`q`$. Then we have
$$R_{\mathrm{here}}^q=R_{\mathrm{RTF}}^q,\widehat{R}_{\mathrm{here}}^q=\widehat{R}_{\mathrm{RTF}}^q,$$
and also
$$C_{\mathrm{here}}^q=C_{\mathrm{RTF}}^q,K_{\mathrm{here}}^q=K_{\mathrm{RTF}}^q.$$
2. The case $`n=0`$
This case is more interesting. Once again, this case has been treated in Ref. . On the other hand, the calculations of the present work don’t make sense in this case from the outset, since the root system of the Lie algebra $`𝔬(2m)`$ does not have a basis of the type used here. Nevertheless, we find
$$R_{\mathrm{here}}^q=R_{\mathrm{RTF}}^{(q^1)},\widehat{R}_{\mathrm{here}}^q=\widehat{R}_{\mathrm{RTF}}^{(q^1)},$$
furthermore, we have
$$C_{\mathrm{here}}^q=q^1C_{\mathrm{RTF}}^{(q^1)},K_{\mathrm{here}}^q=K_{\mathrm{RTF}}^{(q^1)}.$$
The change from $`q`$ to $`q^1`$ under the transition from the even to the odd case is a known phenomenon. On the other hand, the sign factor in the formula for $`\widehat{R}`$ is easily understood: In the purely odd case, the supersymmetric twist is equal to minus the normal (non–graded) twist.
3. The case $`m=n=1`$
In this case, the universal $`R`$–matrix and the $`R`$–matrix in the vector representation have been given in Ref. . However, these authers have worked with a basis of the root system which consists of two odd roots (see Ref. ). Actually, I have applied the approach of the present paper also to this case, and (after some obvious adjustments) indeed have obtained the $`R`$–matrix of Ref. .
## 9 Discussion
We have calculated the $`R`$–matrix of the symplecto–orthogonal quantum superalgebra $`U_q(𝔰𝔭𝔬(2n|2m))`$ in the vector representation, and we have derived its most important properties. In a subsequent work , we shall use this $`R`$–matrix to construct the corresponding quantum supergroup $`\text{SPO}_q(2n|2m)`$ and its basic comodule superalgebras.
A special feature of the present work is that we have used a somewhat unusual basis of the root system of $`𝔰𝔭𝔬(2n|2m)`$. This was dictated by the wish for a unified treatment of all cases, and by the assumption that the basis of the root system should contain only one odd root. If one drops this last requirement, there is another possibility: For $`m2`$, one chooses Kac’s distinguished basis, whereas for $`m=1`$ (i.e., for the $`C`$–type Lie superalgebras) one chooses a basis containing two odd simple roots. In the latter case, the Dynkin diagram looks as follows:
Dynkin diagram of the Lie superalgebra $`𝔰𝔭𝔬(2n|2)`$
For $`n=m=1`$, this is just the choice mentioned in the preceding section. It should be interesting to calculate the $`R`$–matrix also under these assumptions.
In Section 6 we had to use a formulation of the theory in terms of formal power series. This made our arguments somewhat clumsy. Of course, we could exclusively use the language of formal power series. The present formulation has been chosen in view of possible applications.
Appendix
## Appendix A Invariant bilinear forms
In the following, the base field will be an arbitrary field $`𝕂`$ of characteristic zero. Let $`\mathrm{\Gamma }`$ be an abelian group, and let $`\sigma `$ be a commutation factor on $`\mathrm{\Gamma }`$ with values in $`𝕂`$. It is well-known that the class of $`\mathrm{\Gamma }`$–graded vector spaces, endowed with the usual tensor product of graded vector spaces and with the twist maps defined by means of $`\sigma `$, forms a tensor category (see Ref. for details). A (generalized) Hopf algebra $`H`$ living in this category is called a $`\sigma `$–Hopf algebra. More explicitly, $`H`$ is a $`\mathrm{\Gamma }`$–graded associative algebra with a unit element, and it is endowed with a coproduct $`\mathrm{\Delta }`$, a counit $`\epsilon `$, and an antipode $`S`$, which satisfy the obvious axioms (in the category). In particular, this implies that the structure maps $`\mathrm{\Delta }`$, $`\epsilon `$, and $`S`$ are homogeneous of degree zero. In the following, we shall freely use the notation and results of Ref. . (For generalized Hopf algebras living in more general categories, see Ref. .)
Let $`V`$ and $`W`$ be two graded (left) $`H`$–modules. Then $`VW`$ and $`\text{Lgr}(V,W)`$ have a canonical structure of a graded $`H`$–module as well. (Recall that $`\text{Lgr}(V,W)`$ denotes the space of all linear maps of $`V`$ into $`W`$ which can be written as a sum of homogeneous linear maps of $`V`$ into $`W`$.) If $`U`$ is a third graded $`H`$–module, there exists a canonical isomorphism of graded $`H`$–modules,
$$\lambda :\text{Lgr}(V,\text{Lgr}(W,U))\text{Lgr}(VW,U),$$
(A.1)
which is defined by
$$\lambda (f)(xy)=(f(x))(y),$$
for all $`f\text{Lgr}(V,\text{Lgr}(W,U))`$, $`xV`$, and $`yW`$.
The next thing to be mentioned is that an element $`x`$ of a graded $`H`$–module $`V`$ is said to be $`H`$–invariant (or simply invariant) if
$$hx=\epsilon (h)x\text{for all }hH.$$
(Recall that, quite generally, the dot denotes a module action.) Let $`g\text{Lgr}(V,W)`$ be homogeneous of degree $`\gamma `$. Then $`g`$ is invariant if and only if it is $`H`$–linear in the graded sense, i.e., if and only if
$$g(hx)=\sigma (\gamma ,\eta )hg(x),$$
for all elements $`hH_\eta `$, $`\eta \mathrm{\Gamma }`$, and all $`xV`$.
In the following, we choose $`U=𝕂`$, the trivial $`H`$–module. Then $`\text{Lgr}(VW,𝕂)=(VW)^{\mathrm{gr}}`$ is the graded dual of $`VW`$. It is well–known that, regarded as a graded vector space, this space is canonically isomorphic to $`\text{Lgr}_2(V,W;𝕂)`$, the space of all bilinear forms on $`V\times W`$ that can be written as a sum of homogeneous bilinear forms on $`V\times W`$. The canonical isomorphism is used to transfer the $`H`$–module structure from $`\text{Lgr}(VW,𝕂)`$ to $`\text{Lgr}_2(V,W;𝕂)`$. For every bilinear form $`b\text{Lgr}_2(V,W;𝕂)`$, the corresponding linear form on $`VW`$ will be denoted by $`\stackrel{~}{b}`$.
Now let
$$b:V\times W𝕂$$
be a bilinear form on $`V\times W`$ that is homogeneous of degree $`\beta `$ , let
$$\stackrel{~}{b}:VW𝕂$$
be the associated linear form on $`VW`$, and let
$$f_{\mathrm{}}:VW^{\mathrm{gr}}$$
be the linear map canonically corresponding to $`\stackrel{~}{b}`$, i.e.,
$$(f_{\mathrm{}}(x))(y)=b(x,y)$$
for all $`xV`$ and $`yW`$. (Choosing $`U=𝕂`$ in Eqn. (A.1), this is to say that $`\lambda (f_{\mathrm{}})=\stackrel{~}{b}`$.) Then $`\stackrel{~}{b}`$ and $`f_{\mathrm{}}`$ are homogeneous of degree $`\beta `$ . According to the preceding discussion, the following statements are equivalent:
1) The bilinear form $`b`$ , or equivalently, the linear form $`\stackrel{~}{b}`$ , is $`H`$–invariant, i.e., we have
$$\sigma (\eta ,\beta )\stackrel{~}{b}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}S(h)_{VW}=\epsilon (h)\stackrel{~}{b}$$
for all elements $`hH_\eta `$, $`\eta \mathrm{\Gamma }`$.
2) The linear form $`\stackrel{~}{b}`$ is $`H`$–linear in the graded sense, i.e., we have
$$\stackrel{~}{b}(h(xy))=\sigma (\beta ,\eta )\epsilon (h)\stackrel{~}{b}(xy)$$
for all elements $`hH_\eta `$, $`\eta \mathrm{\Gamma }`$, and all $`xV`$, $`yW`$. Since $`\epsilon (h)=0`$ if $`\eta 0`$, this is equivalent with
$$\stackrel{~}{b}(h(xy))=\epsilon (h)\stackrel{~}{b}(xy)$$
(A.2)
for all elements $`hH`$, $`xV`$, and $`yW`$.
3) The linear map $`f_{\mathrm{}}`$ is $`H`$–invariant.
4) The linear map $`f_{\mathrm{}}`$ is $`H`$–linear in the graded sense, i.e., we have
$$f_{\mathrm{}}(hx)=\sigma (\beta ,\eta )hf_{\mathrm{}}(x)$$
for all elements $`hH_\eta `$, $`\eta \mathrm{\Gamma }`$, and all $`xV`$. This is equivalent with
$$f_{\mathrm{}}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}h_V=\sigma (\beta ,\eta )(S(h)_W)^\mathrm{T}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}f_{\mathrm{}}$$
(A.3)
for all elements $`hH_\eta `$, $`\eta \mathrm{\Gamma }`$, and also with
$$b(hx,y)=\sigma (\eta ,\xi )b(x,S(h)y)$$
(A.4)
for all elements $`hH_\eta `$, $`\eta \mathrm{\Gamma }`$, all $`xV_\xi `$, $`\xi \mathrm{\Gamma }`$, and all $`yW`$. Recall that <sup>T</sup> denotes the $`\sigma `$–transposition. (In the super case, the super–transposition will be denoted by <sup>st</sup>.)
To proceed, we note that, apart from $`f_{\mathrm{}}`$, the bilinear form $`b`$ yields a second linear map, namely, the map
$$f_r:WV^{\mathrm{gr}}$$
which, for all elements $`xV_\xi `$, $`\xi \mathrm{\Gamma }`$, and $`yW_\eta `$, $`\eta \mathrm{\Gamma }`$, is given by
$$(f_r(y))(x)=\sigma (\eta ,\xi )b(x,y).$$
If
$$\nu _V:V(V^{\mathrm{gr}})^{\mathrm{gr}},\nu _W:W(W^{\mathrm{gr}})^{\mathrm{gr}}$$
are the canonical injections (in the graded sense, see Ref. ), we have
$$f_r=f_{\mathrm{}}^\mathrm{T}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}\nu _W,f_{\mathrm{}}=f_r^\mathrm{T}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}\nu _V.$$
Then the condition (A.3) is equivalent with
$$f_r\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}S(h)_W=\sigma (\beta ,\eta )(h_V)^\mathrm{T}\stackrel{{\scriptscriptstyle \genfrac{}{}{0pt}{}{\text{}}{}}}{\text{}}f_r,$$
(A.5)
for all elements $`hH_\eta `$, $`\eta \mathrm{\Gamma }`$. Indeed, the Eqns. (A.3) and (A.5) can be derived from each other by composing their $`\sigma `$–transposes with $`\nu _W`$ and $`\nu _V`$, respectively.
## Appendix B Partial transposition
In the following, the base field will be an arbitrary field $`𝕂`$ of characteristic zero. Let $`\mathrm{\Gamma }`$ be an abelian group, and let $`\sigma `$ be a commutation factor on $`\mathrm{\Gamma }`$ with values in $`𝕂`$. All gradations considered in this appendix will be $`\mathrm{\Gamma }`$–gradations. We shall freely use the notation and results of Ref. .
Let $`V`$ and $`W`$ be two finite–dimensional graded vector spaces. Then there exists a unique linear map
$${}_{}{}^{\mathrm{T}_1}:\text{Lgr}(VW,VW)\text{Lgr}(V^{\mathrm{gr}}W,V^{\mathrm{gr}}W),$$
such that
$$(f\overline{}g)^{\mathrm{T}_1}=f^\mathrm{T}\overline{}g,$$
and a unique linear map
$${}_{}{}^{\mathrm{T}_2}:\text{Lgr}(VW,VW)\text{Lgr}(VW^{\mathrm{gr}},VW^{\mathrm{gr}}),$$
such that
$$(f\overline{}g)^{\mathrm{T}_2}=f\overline{}g^\mathrm{T},$$
for all $`f\text{Lgr}(V,V)`$ and all $`g\text{Lgr}(W,W)`$. (Recall that $`\overline{}`$ denotes the graded tensor product of linear maps, and that <sup>T</sup> denotes the $`\sigma `$–transposition.) We call $`^{\mathrm{T}_1}`$ and $`^{\mathrm{T}_2}`$ the partial $`\sigma `$–transposition of the first and second tensorial factor, respectively. (In the super case, we shall write $`^{\mathrm{st}_1}`$ and $`^{\mathrm{st}_2}`$.)
It is easy to see that
$$(fg)^{\mathrm{T}_2}=fg^\mathrm{T},$$
for all $`f\text{Lgr}(V,V)`$ and all $`g\text{Lgr}(W,W)`$. On the other hand, under our present general assumptions, $`(fg)^{\mathrm{T}_1}`$ is not, in general, equal to $`f^\mathrm{T}g`$. Using our standard notation for $`V`$, and a similar notation for $`W`$, but with all entries overlined, we find instead that
$$(E_{ij}\overline{E}_{rs})^{\mathrm{T}_1}=\sigma (\eta _i+\eta _j,\overline{\eta }_r\overline{\eta }_s)E_{ij}^\mathrm{T}\overline{E}_{rs}.$$
In the super case, which is the case we are mainly interested in, this equation implies that
$$(fg)^{\mathrm{T}_1}=\sigma (\phi ,\gamma )f^\mathrm{T}g,$$
where $`f\text{Lgr}(V,V)`$ is homogeneous of degree $`\phi `$ and $`g\text{Lgr}(W,W)`$ is homogeneous of degree $`\gamma `$.
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# Starburst in the Ultra-luminous Galaxy Arp 220 - Constraints from Observations of Radio Recombination Lines and Continuum
## 1 Introduction
Arp 220 is the nearest ($`d`$ = 73 Mpc for $`H_0`$ = 75 km s<sup>-1</sup>Mpc<sup>-1</sup>) example of an ultra-luminous infrared galaxy (ULIRG). ULIRGs are a class of galaxies with enormous infrared luminosities of $`L_{IR}10^{12}L_{}`$, believed to have formed through mergers of two gas-rich galaxies (Sanders et al 1988). What powers the high luminosity of ULIRGs is not well understood, although there is considerable evidence that intense starburst which may be triggered by the merger process may play a dominant role (Sanders and Mirabel 1996, Gënzel et al 1998). There is also evidence in many ULIRGs for a dust-enshrouded AGN which could contribute significantly to or in some cases dominate the observed high luminosity (Gënzel et al 1998). In the case of Arp 220, there is clear evidence for a merger in the form of extended tidal tails observed in the optical band (Arp 1966) and the presence of a double nucleus in the radio and NIR bands with a separation of $`1^{\prime \prime }`$ (370 pc)(Norris 1985, Graham et al 1990). Prodigious amounts of molecular gas ($`M_{H_2}10^{10}M_{}`$; Scoville, Yun & Bryant 1997, Downes & Solomon 1998) are present in the central few hundred parsecs. Detection of CS emission (Solomon, Radford & Downes 1990) indicates that the molecular densities may be high ($`n_{H_2}10^5`$ cm<sup>-3</sup>). Recent high resolution observations reveal that the two nuclei in Arp 220 consists of a number of compact radio sources - possibly luminous radio supernovae (Smith et al 1998) - and a number of luminous star clusters in the near IR (Scoville et al 1998). High resolution observations of OH megamasers (Lonsdale et al 1998) reveal multiple maser spots with complex spatial structure. These observations have led to the suggestion that the nucleus of Arp 220 is dominated by an intense, compact starburst phenomena rather than the presence of a hidden AGN (Smith et al 1998, Downes & Solomon 1998). The presence of AGNs in the two nuclei is, however, not ruled out since large velocity gradients, $`500`$ km s<sup>-1</sup> over $`1^{\prime \prime }`$ at different position angles across each nuclei (Sakamoto et al 1999) have been observed, indicating high mass concentrations at the location of both nuclei. Obscuration due to dust in the nuclear region of Arp 220 is extremely high. Ratios of fine structure lines observed with the ISO satellite (Lutz et al 1996, Gënzel et al 1998) indicate a lower limit of $`A_V`$ 45 mag. This high extinction seriously hampers measurements even at IR wavelengths (see also Scoville et al 1998).
Radio recombination lines (RRLs) do not suffer from dust obscuration and thus may potentially provide a powerful method of studying the kinematics, spatial structure and physical properties of ionized gas in the nuclear region of galaxies. RRLs have been detected, to date, in about 15 Galaxies (Shaver, Churchwell and Rots 1977, Seaquist and Bell 1977, Puxley et al 1991, Anantharamaiah et al 1993, Zhao et al 1996, Phookun, Anantharamaiah and Goss, 1998). Arp 220 is the most distant among the galaxies that are detected in recombination lines (Zhao et al 1996). The detection of an RRL in Arp 220 (Zhao et al 1996) opens up the possibility of studying the ionized gas in the nuclear region. Zhao et al (1996) presented VLA observations of the H92$`\alpha `$ line ($`\nu _{rest}`$ =8.309 GHz) with an angular resolution of $`4^{\prime \prime }`$ and a velocity coverage of $`900`$ km s<sup>-1</sup>. This paper presents new H92$`\alpha `$ VLA observations with an angular resolution of $`1^{\prime \prime }`$ and a velocity coverage of $`1700`$ km s<sup>-1</sup>. We also present observations of millimeter wavelength RRLs H42$`\alpha `$ ($`\nu _{rest}`$ =85.688 GHz), H40$`\alpha `$($`\nu _{rest}`$ = 99.022 GHz) and H31$`\alpha `$ ($`\nu _{rest}`$ = 210.501 GHz) made using the IRAM 30 m radio telescope in Pico Veleta, Spain. Additionally, we report a sensitive upper limit to lower frequency RRLs H167$`\alpha `$ ($`\nu _{rest}`$ = 1.40 GHz) and H165$`\alpha `$($`\nu _{rest}`$ = 1.45 GHz ) obtained using the VLA. These observations together with the observed continuum spectrum are used to constrain the physical properties and the amount of ionized gas in the nuclear region of Arp 220. Since the Lyman continuum photons which ionize the gas are generated by young O and B stars with a finite life time, the derived properties of the ionized gas are used to constrain the properties of the starburst in Arp 220.
The larger velocity coverage used for the present H92$`\alpha `$ observations showed that the previous observation with half the bandwidth (Zhao et al 1996) had underestimated both the width and the strength of the H92$`\alpha `$ line. The integrated line intensity is found to be more than a factor of two larger than reported earlier and thus the deduced number of Lyman continuum photons that are required to maintain the ionization is also larger. The millimeter wavelength lines are found to be much stronger than expected based on models for the H92$`\alpha `$ line. The observed strength of the millimeter wavelength RRLs indicate the presence of a higher density ($`n_e>10^5`$ cm<sup>-3</sup>) ionized gas component and leads to information about star formation rates at recent epochs ($`t<10^5`$ yrs). The upper limits to the lower frequency lines provide significant constraints on the amount of low density ($`n_e<10^3`$ cm<sup>-3</sup>) ionized gas in the nuclear region of Arp 220. A comparison of the predictions (based on RRLs) of the strengths of NIR recombination lines (e.g. Br$`\alpha `$ and Br$`\gamma `$lines) with observations (e.g. Goldadder et al 1995, Sturm et al 1996) leads to a revised estimate of the dust extinction.
This paper is organized as follows. In Section 2, observations with the the VLA and the IRAM 30 m telescope are described and the results are presented. Section 3 summarizes the available radio continuum measurements of Arp 220 in the frequency range 150 MHz to 230 GHz including new measurements at 327 MHz, 1.4 GHz, 8.3 GHz, 99 GHz and 206 GHz. In Section 3, we discuss the various possible models to explain the observed RRLs and continuum in Arp 220 and arrive at a three-component model, consistent with the observations. Section 4 discusses the implications of the three-component model to several aspects of the starburst in Arp 220. We attempt to answer the question of whether the deduced star formation rates in Arp 220 can power the observed high luminosity without invoking the presence of an obscured active nucleus. The paper is summarized in Section 5.
## 2 Observations, Data Reduction and Results
### 2.1 VLA observations and Results
#### 2.1.1 H92$`\alpha `$ line
Observations of the H92$`\alpha `$ line ($`\nu _{rest}`$ = 8309.38 MHz) were made in Aug 1998 in the B configuration of the VLA. The parameters of the observations are given in Table 1. Phase, amplitude and frequency response of the antennas were corrected using observations of the calibrator source 1600+335 ($`S_{3.5cm}`$ 2 Jy) every 30 minutes. Standard procedures in AIPS were used for generating the continuum image and the line cube. A self-calibration in both amplitude and phase was performed on the continuum channel and the solutions were applied to the spectral line data. For subtracting the continuum, two channels, one on each end (channels 2 and 14), were used in the procedure UVLIN (Cornwel, Uson and Haddad 1992). During off-line data processing, it became necessary to apply hanning smoothing which resulted in a spectral resolution of 6.25 MHz (i.e. 230 km/s) with two spectral points per resolution element. The synthesized beam with natural weighting of the visibility function was $`1.1^{\prime \prime }\times 0.9^{\prime \prime }`$, PA = 74.
The continuum image of Arp220 is shown in Fig 1 as contours with the image of the velocity-integrated H92$`\alpha `$ line (moment 0) superposed in grey scale. The two radio nuclei of Arp 220 (Norris et al 1988) which are separated by about $`1^{\prime \prime }`$ are barely resolved. Fig 1 shows that the recombination line emission is detected from both nuclei. The peak of the line emission is slightly offset to the north with respect to the continuum peak. The position angle of the elongated line emission is tilted with respect to the line joining the two continuum peaks by $`+7^{}`$. Fig 2 shows a contour image of the velocity-integrated line emission (moment 0) with velocity dispersion (moment2) superposed in grey scale. The recombination line is stronger and wider near the western source and there appears to be a weak extension of the ionized gas towards the north-east.
Fig 3 shows the integrated H92$`\alpha `$ line profile (top frame) along with line profiles near the eastern and western continuum peaks. The line profile is wider on the western peak. The parameters of the integrated line profile are given in Table 2. The spatially-integrated H92$`\alpha `$ line has a peak flux density of 0.6 mJy/beam and a FWHM of 363 km s<sup>-1</sup>. Both these values are larger than the values reported by Zhao et al (1996). The integrated line flux is 8$`\times 10^{23}`$ W m<sup>-2</sup> which is a factor of 2.3 larger than the values reported by Zhao et al (1996). The integrated flux was underestimated by Zhao et al (1996) since a narrower bandwidth was used in their observation, thereby missing a portion of the line.
Because of the coarse velocity resolution (230 km s<sup>-1</sup>) of the H92$`\alpha `$ data, detailed information about the kinematics is limited. Fig 4a shows a position-velocity diagram along the line joining the two continuum peaks (P.A. = 98). The velocity range is wider near the western peak compared to the eastern peak. After correcting for the instrumental resolution, the full velocity range over the western peak is about 700 km s<sup>-1</sup>. On the eastern peak the velocity range is 440 km s<sup>-1</sup>. The FWHMs on the western and the eastern peaks are 500 km s<sup>-1</sup>and 340 km s<sup>-1</sup>respectively. Fig 4b shows a position-velocity diagram along the line perpendicular to the line joining the two peaks. The full velocity range of emission in Fig 4b is about 770 km s<sup>-1</sup>with a FWHM of 430 km s<sup>-1</sup>. No velocity gradients are visible in Figs 4a and 4b.
#### 2.1.2 H167$`\alpha `$ and H165$`\alpha `$ Lines
The H167$`\alpha `$ ($`\nu _{obs}`$ = 1399.37 MHz) and the H165$`\alpha `$ ($`\nu _{obs}`$ = 1450.72 MHz) lines were observed simultaneously over a 7 hour period in March 1998 in the A configuration of the VLA. Observational details are summarized in Table 1. Data were processed using standard procedures in AIPS. With natural weighting of the visibilities, the beam size is $`1.8^{\prime \prime }\times 1.6^{\prime \prime }`$. The peak continuum flux density was 215 mJy/beam and the integrated flux density is 312$`\pm `$3 mJy. After hanning smoothing, the resolution of each spectral channel was 390 kHz ($`84`$ km s<sup>-1</sup>) and the rms noise was $`130\mu `$Jy/beam. Neither line was detected. For a line width of 350 km s<sup>-1</sup>, we obtain a $`3\sigma `$ upper limit to the line strength of 0.25 mJy. The parameters are listed in Table 2. As shown in the next section, the upper limit to the RRLs near 1.4 GHz provide significant constraints on the density of ionized gas in Arp 220.
### 2.2 IRAM-30 m Observations and Results
Observations with the IRAM-30 m telescope in Pico Veleta, Spain were carried out on April 19, 1996. The pointing of the telescope was adjusted from time to time; the rms pointing error was $`3^{\prime \prime }`$. Three receivers were used simultaneously. A 1.2 mm receiver connected to a 512 MHz filter bank spectrometer was used to observe the H31$`\alpha `$ line ($`\nu _{rest}`$ = 210.5018 GHz). The other two receivers were used to observe the H40$`\alpha `$ ($`\nu _{rest}`$ = 99.0229 GHz) and the H42$`\alpha `$ ($`\nu _{rest}`$ = 85.6884 GHz) lines in the 3 mm band. The two 3 mm lines were observed using two subsets of the autocorrelation spectrometer, each with 409 channels. Instrumental parameters are summarized in Table 3. In the 1.2 mm band, the system temperature ranged between 260 K and 285 K and in the 3 mm band, the system temperatures were between 180 K and 200 K. The wobbler was used to remove the sky emission and to determine the instrumental frequency response. Spectra were recorded every 30 seconds.
The data from the IRAM-30 m telescope were reduced using a special program written by one of us (FV) to reduce single-dish spectral data in the Groningen Image Processing System (GIPSY) environment (van der Hulst et al 1992). In this new program, a baseline subtracted final spectrum as well as a continuum level are simultaneously determined from the data base which consists of a series of spectra sampled every 30 s. The temporal information is used to perform statistics on the data to measure the rms noise as a function of integration time and frequency resolution. These statistics also provide the uncertainties associated with the fitted baselines. These measured uncertainties are used in the determination of error bars in the line and continuum flux densities. In addition, the program also determines a statistical rms noise level for each spectral channel separately, which helps in judging the reality or otherwise of a spectral feature. The details of this procedure are described by Viallefond (2000, in preparation).
The final spectra obtained with the IRAM-30 m data are shown in Fig 5 and the parameters are listed in Table 4. The quoted error bars are the quadrature sums of the system noise and baseline uncertainties. To obtain the H42$`\alpha `$ spectrum shown in the bottom panel of Fig 5, linear baselines were fitted over the velocity ranges 4600-5050 km s<sup>-1</sup>and 5800-6200 km s<sup>-1</sup>. The spectrum has been smoothed to a resolution of 30 MHz and re-sampled to obtain 17 independent frequency channels out of the original 409 channels. The H42$`\alpha `$ recombination line extends from 5050 km s<sup>-1</sup>to 5800 km s<sup>-1</sup>. The emission feature at the higher velocity end of the spectrum (i.e $`v>`$ 6200 km s<sup>-1</sup>) is probably the rotational line of the molecule C<sub>3</sub>H<sub>2</sub> ($`\nu _{rest}`$ = 85.339 GHz). The continuum level is poorly determined from this data set and thus only a $`3\sigma `$ upper limit to the continuum flux density is listed in Table 4.
The H40$`\alpha `$ spectrum shown in the middle panel of Fig 5 was determined by fitting linear baselines to channels outside the velocity range 4950-5850 km s<sup>-1</sup>. The spectrum has been smoothed to a resolution of 30 MHz and re-sampled at 17 independent frequency points. The dashed lines in Fig 5 shows the rms noise across the spectrum determined as described in Viallefond (2000, in preparation). The H40$`\alpha `$ line emission is quite similar to the H42$`\alpha `$ emission in Fig 5. The consistency between the two lines confirms the reality of the detections. This data set also provided a measurement of the continuum flux density near 3 mm which is listed in Table 4
The observations of the H31$`\alpha `$ line near 202.7 GHz have a serious limitation. The 512 MHz ($``$750 km s<sup>-1</sup>) bandwidth used for the observations was too small to cover the entire extent of the line (FWZI $`>`$ 800 km s<sup>-1</sup>). To determine the continuum emission and the spectral baseline level, we used the few channels at one end of the spectrum where no line emission was expected. Although the resulting spectrum, shown as a solid line in the top panel of Fig 5, covers only a part of the line emission, it is consistent in shape with the other recombination lines in Fig 5. The statistical significance of the detection of the H31$`\alpha `$ line is $`>10\sigma `$ near the peak. Table 2 lists the ratios of the strengths of various recombination lines integrated over the partial velocity range 5220-5755 km s<sup>-1</sup>over which the H31$`\alpha `$ line has been observed. The strength of the H31$`\alpha `$ line over the full extent of the line is estimated based on these ratios.
The millimeter wavelength recombination lines (H42$`\alpha `$, H40$`\alpha `$ and H31$`\alpha `$) detected with the IRAM-30 m telescope are much stronger than the centimeter wavelength line (H92$`\alpha `$) detected with the VLA. While the peak line strength at $`\lambda `$ = 3.5 cm (H92$`\alpha `$) is 0.6 mJy, the corresponding lines at $`\lambda `$ = 3 mm (H40$`\alpha `$) and 1.2 mm (H31$`\alpha `$) are an order of magnitude more intense: 15 mJy and 80 mJy, respectively. An increase in the recombination line strength at millimeter wavelengths was expected from high density models discussed by Anantharamaiah et al (1993) and Zhao et al (1996). However, the observed increase in the line strength is much larger than expected on the basis of the models for the H92$`\alpha `$ line. The mm RRLs H40$`\alpha `$ and H42$`\alpha `$ also appear to be broader than the cm RRL H92$`\alpha `$ line. The width of the mm RRLs are comparable to that of CO lines (Downes and Solomon 1998).
## 3 Radio Continuum Spectrum of Arp 220
As a byproduct of the RRL measurements, continuum flux densities are obtained at 1.4, 8.2, 97.2 and 206.7 GHz. The measured flux densities are listed in Table 5 along with published values at other frequencies. The continuum flux density measurements around 1.5 mm are consistent with the determination by Downes and Solomon (1998), but the value near 3 mm is higher than that reported by the same authors. The increase in the flux density above 200 GHz is due to contribution by thermal dust emission.
In addition, we observed Arp 220 near 0.325 GHz in August 99 using the A-configuration of the VLA using 3C286 as the flux calibrator. The angular resolution was $`6.5^{\prime \prime }`$. Arp 220 is unresolved with this beam. The integrated flux density of Arp 220 is 380$`\pm `$15 mJy. A flux density measurement at an even lower frequency of 0.15 GHz has been reported by Sopp and Alexander (1991), which is also included in Table 5.
The continuum spectrum in the frequency range 0.1 to 100 GHz is plotted in Fig 6. The continuum spectrum is non-thermal in the frequency range 3 to 30 GHz with a spectral index $`\alpha 0.6`$ (where $`S_\nu \nu ^\alpha `$). The spectrum changes both below and above this frequency range. A break in the spectrum is observed around 2 GHz with a change in the spectral index to $`0.1`$ between 1.6 GHz and 0.325 GHz. A further change in spectrum occurs below 300 MHz. The spectral index changes to $`+0.5`$ and between 0.325 GHz and 0.15 GHz. These changes in the spectrum at lower frequencies are complex and cannot be explained by a simple free-free absorbing thermal screen or thermal gas mixed with the non-thermal component (e.g. Sopp & Alexander 1991). At the higher frequencies the change in the spectrum is gradual. The spectral index appears to become slightly flatter (i.e. $`\alpha >0.6`$) above 22.5 GHz. The solid line in Fig 6 is a fit to the continuum spectrum based on a 3-component ionized gas model (next section) developed to explain the observed recombination lines.
## 4 Models for RRL and Continuum emission
Models for RRL emission from the nuclear region of external galaxies have been discussed by Puxley et al (1991), Anantharamaiah et al (1993), Zhao et al (1996, 1998) and Phookun et al (1998). The main constraints for these models are the integrated RRL strength at one or more frequencies, the observed radio continuum spectrum and geometrical considerations. Two types of models have been considered: one based on a uniform slab of ionized gas and the other based on a collection of compact HII regions. The observed non-thermal radio continuum spectrum of the nuclear region, with spectral index $`\alpha 0.6`$ (where $`S_\nu \nu ^\alpha `$), imposes strong constraints on the nature of the ionized gas that produces the observed recombination lines from the same region. If the models are constrained by a single RRL measurement and the continuum spectrum, then the derived physical parameters ($`T_e,n_e,M_{HII}`$ etc) are not unique. In a majority of the cases that were considered (Anantharamaiah et al 1993, Zhao et al 1996, Phookun et al 1998), the model with a collection of compact, high density HII regions is favored. The uniform slab models produce an excess of thermal continuum emission at centimeter wavelengths inconsistent with the observed non-thermal spectrum. However, in the case of Arp 220, Zhao et al (1996) were able to fit both types of models; thus further observations at higher and lower frequencies were required to choose between the models. We consider below these models of RRL emission in the light of the new measurements.
### 4.1 Uniform Slab Model
In the uniform slab model, the parameters are the electron temperature ($`T_e`$), the electron density ($`n_e`$) and the thickness of the slab (l). For a given combination of $`T_e`$ and $`n_e`$, l is adjusted to account for the observed integrated H92$`\alpha `$ line strength. While calculating the line strength, stimulated emission due to the background non-thermal continuum, as well as internal stimulated emission due to the thermal continuum emission from the ionized gas are taken into account. The relevant expressions for the computation are given by Anantharamaiah et al (1993). The observed total continuum emission at two frequencies, together with the computed thermal emission from the ionized slab, is used to estimate the intrinsic non-thermal emission and its spectral index. Models in which the thermal emission from the slab at 5 GHz exceeds a substantial fraction ($``$ 30 – 50%) of the total continuum emission are rejected since the resulting spectral index will not be consistent with the observed non-thermal spectrum. Finally, the expected variation of line and continuum emissions as a function of frequency are computed and compared with the observations.
Fig 7a shows the expected variation of integrated RRL strength as a function of frequency for $`T_e`$= 7500 K and three values of electron density. Fig 7b shows the expected variation of continuum emission for the corresponding models. In these models, the slab of ionized gas is assumed to be in front of the non-thermal source. Results are qualitatively similar if the ionized gas is mixed with the non-thermal gas. The models are normalized to the H92$`\alpha `$ line and the continuum flux densities at 4.7 and 15 GHz. Fig 7 shows that while the non-thermal continuum spectrum above 2 GHz and the change in the spectrum near 1.6 GHz agree with the models, the flux densities at 0.15 and 0.325 GHz are not explained. In all the models, free-free absorption is prominent below 1.5 GHz. No uniform slab model can be found which both accounts for the H92$`\alpha `$ line and also produces a turnover in the continuum spectrum at a frequency $`<`$ 1.5 GHz. Furthermore, all the models that fit the H92$`\alpha `$ line predict a strength for RRLs near 1.4 GHz line inconsistent with the upper limit. For other values of $`T_e`$ (5000 K and 10000 K), the curves are similar to Figs 7a and 7b. No models with $`n_e>10^3`$ cm<sup>-3</sup> could be fitted to the H92$`\alpha `$ line. The parameters of the models shown in Fig 7 are given in Table 6. In these models, external stimulated emission (i.e. due to the background non-thermal radiation) accounts for more than 75% of the H92$`\alpha `$ line strength. No uniform slab model, at any density, could be fitted to the higher frequency (H42$`\alpha `$, H40$`\alpha `$ and H31a) lines. Since all the uniform slab models are inconsistent with low frequency RRL and continuum data, we consider models with a collection of HII regions.
### 4.2 Model with a Collection of HII regions
In this model, the observed RRLs are thought to arise in a number of compact, high density HII regions whose total volume filling factor in the nuclear region is small ($`<10^4`$) (Puxley et al 1991, Anantharamaiah et al 1993). The low volume filling factor ensures that the HII regions, regardless of their continuum opacities, have only a small effect on the propagation of the non-thermal continuum which originates in the nuclear region. Thus, the observed radio continuum spectrum can be non-thermal even below the frequency at which the HII regions themselves become optically thick. The low filling factor also implies that in these models there is little external stimulated emission. Since the HII regions are nearly optically thick at centimeter wavelengths, they produce only a modest amount of thermal continuum emission (typically $`<`$ 30% of S<sub>obs</sub>). The recombination line emission from these regions arises mainly through internal stimulated emission due to the continuum generated within the HII regions. If the density of the HII regions are lower ($`n_e<10^3`$ cm<sup>-3</sup>) and the area covering factor $`f_c`$ (i.e. the fraction of the area of the non-thermal emitting region that is covered by HII regions) is larger, then external stimulated emission can become dominant at centimeter wavelengths.
For simplicity, all the HII regions are considered to be characterized by the same combination of $`T_e`$, $`n_e`$ and diameter $`d_{HII}`$. For a given combination of these parameters, the expected integrated line flux density ($`S_ld_\nu `$) of a single HII region is calculated using standard expressions (e.g. Anantharamaiah et al 1993). The number of HII regions is then computed by dividing the integrated flux density of one of the observed RRLs (e.g. H92$`\alpha `$) by the expected strength from a single HII region. Since the volume filling factor of the HII regions is small, the effect of shadowing of one HII region by another is not significant. Constraints from the observed continuum flux densities at various frequencies and from several geometrical aspects, as discussed in Anantharamaiah et al (1993), are applied to restrict the acceptable combinations of $`T_e`$, $`n_e`$ and $`d_{HII}`$. Finally, the models are used to compute the expected variation of line and continuum emission with frequency and compared with observations. We show below that separate components of ionized gas are required to explain the centimeter wavelength (H92$`\alpha `$) and millimeter wavelength RRLs (H42$`\alpha `$, H40$`\alpha `$ and H31$`\alpha `$).
#### 4.2.1 Models based on the H92$`\alpha `$ line and the Continuum
Fig 8 shows three representative models that fit the observed H92$`\alpha `$ line. The parameters of the three models are given in Table 7. The nature of the curves for other successful combinations of $`T_e`$, $`n_e`$ and $`d_{HII}`$ are similar to one of the three curves shown in Fig 8, although the values of the derived parameters are different. At densities below about 500 cm<sup>-3</sup>, the area covering factor of the HII regions is unity. In other words, every line of sight through the line emitting region intersects at least one HII region (N$`{}_{}{}^{HII}{}_{los}{}^{}`$$`>1`$ and $`f_c`$ =1 in Table 7), and thus the models at these densities are similar to the uniform slab model discussed above. As seen in the dashed curves in Fig 8, these low-density models are inconsistent with the continuum spectrum below 1 GHz as well as the upper limit to the RRL emission near 1.4 GHz. Lower density models, which were considered possible by Zhao et al (1996) based only on the H92$`\alpha `$ and higher frequency continuum data, are now ruled out.
At densities above about 5000 cm<sup>-3</sup>, the area filling factor of the HII regions is very small ($`f_c<<`$ 1) and thus they do not interfere with the propagation of the non-thermal radiation. The continuum spectrum will be non-thermal even at the lowest frequencies. The dash-dot-dash curve in Fig 8 shows such a model (model C1) with the parameters given in Table 7. In this model, although there are more than $`10^5`$ HII regions, each $``$1 pc in diameter and optically thick below $``$ 3 GHz, the total continuum emission continues to have a non-thermal spectrum at lower frequencies. Thus these higher density models cannot account for the observed turnover in the continuum spectrum at $`\nu <`$ 500 MHz. These models are however consistent with the upper limit to the RRL emission near 1.4 GHz.
A continuum spectrum, which is partially consistent with observations at low frequencies is obtained by models with densities in the range $`750<n_e<1500`$ cm<sup>-3</sup>. An example is the solid curve in Fig 8 (Model A1) with the parameters as given in Table 7. For this model, the area covering factor $`f_c=`$ 0.7. In other words 30% of the non-thermal radiation propagates unhindered by the HII regions and the remaining 70% is subjected to the free-free absorbing effects of the HII regions. The net continuum spectrum thus develops a break near the frequency at which $`\tau _{HII}1`$. At much lower frequencies where 70% of the non-thermal radiation is completely free-free absorbed, the remaining 30% propagates through with its intrinsic non-thermal spectrum as shown in the solid curve in Fig 8. While this model accounts for the break in the spectrum near $``$ 2 GHz, it fails to account for the complete turnover in the continuum spectrum below 500 MHz. An additional thermal component with a covering factor close to unity and a turnover frequency near $``$ 300 MHz is required to account for the low-frequency spectrum. This component is discussed in Section 4.2.3. Model A1 in Fig 8 is also consistent with the upper limit to the RRL emission near 1.4 GHz.
In the models shown in Fig 8 and Table 7, between 25% to 70% of the line emission arises due to external stimulated emission, i.e. amplification of the background non-thermal radiation at the line frequency due to non-LTE effects. For lower densities, the fraction of external stimulated emission is increased. For $`n_e1000`$ cm<sup>-3</sup> (favored above), 70% of the H92$`\alpha `$ line is due to stimulated emission. The other parameter of interest in Table 9 which is related to the line emission mechanism is the non-LTE factor $`f_{nlte}`$, the ratio of the intrinsic line emission from an HII region if non-LTE effects are included to the line emission expected under pure LTE conditions (i.e. $`b_n=\beta _n=1`$). $`f_{nlte}`$ is in the range 1 to 3 for the models in Table 7. For lower densities, the non-LTE effect on the intrinsic line emission from an HII region are reduced. Thus, in the models with a collection of HII regions, external and internal stimulated emissions vary with density in opposite ways. If the density is increased, internal stimulated emission increases whereas external stimulated emission decreases.
Two important parameters that can be derived from these models are the mass of the ionized gas M<sub>HII</sub> and the number of Lyman continuum (Lyc) photons N<sub>Lyc</sub> . The latter can be directly related to the formation rate of massive stars if direct absorption of Lyc photons by dust is not significant. Furthermore, predictions can be made for the expected strengths of optical and IR recombination lines, which can then be compared with observed values in order to estimate the extinction. Some derived quantities are listed in Table 7. As seen in this Table, although a range of $`n_e`$, $`T_e`$, and $`d_{HII}`$ values fit the H92$`\alpha `$ data, the derived value of N<sub>Lyc</sub> varies by less than a factor of two. For the model labelled A1 in Fig 8, the total mass of ionized gas M<sub>HII</sub> = $`3\times 10^7`$ $`M_{}`$ and the number of Lyman continuum photons N<sub>Lyc</sub> = $`1.2\times 10^{55}`$ s<sup>-1</sup>. For models that satisfy all the constraints, the derived parameters M<sub>HII</sub> and N<sub>Lyc</sub> increase if either $`T_e`$ is increased or $`d_{HII}`$ is decreased. On the other hand, as the density $`n_e`$ is increased, M<sub>HII</sub> decreases, whereas N<sub>Lyc</sub> first increases and then decreases.
None of the models (at any density) that fit the H92$`\alpha `$ line can account for the observed H42$`\alpha `$, H40$`\alpha `$ and H31$`\alpha `$ lines. Although an increase in line strength towards shorter wavelengths is predicted by all the models (Fig 8), the expected integrated line flux density falls short by almost an order of magnitude, well above any uncertainty in the measurements. The model A1 in Fig 8 can explain the H92$`\alpha `$ line, the break in the continuum spectrum at 1.5 GHz and is consistent with the upper limit to the H167$`\alpha `$ line. Additional components of ionized gas are required to explain the millimeter wavelength RRLs and the turnover the continuum spectrum below 500 MHz.
#### 4.2.2 Models based on the H42$`\alpha `$ line and the Continuum
Since no model that fits the H92$`\alpha `$ line could account for the observed 3 mm and 1.2 mm lines, separate models, also based on a collection of HII regions, were fitted to the H42$`\alpha `$ line. Only models with densities above $`10^5`$ cm<sup>-3</sup> are consistent with the mm wavelength RRLs. Three successful models, labelled A2, B2, and C2 are shown in Fig 9 and the parameters of the models are listed in Table 8. As seen in Fig 9, the contribution of this high density component to RRLs at centimeter wavelengths is negligible since the HII regions are optically thick at these frequencies. This high density component is thus detectable only in RRLs at millimeter wavelengths.
In the models summarized in Table 8, the contribution to the line emission from external stimulated emission is negligible ($`<0.5\%)`$. On the other hand, enhancement of the line due to internal stimulated emission within the HII regions is pronounced and sensitive to the parameters of the HII regions ($`n_e`$, $`T_e`$, and $`d_{HII}`$). The factor $`f_{nlte}`$ range from 30 - 1300 for the models in Table 8. Because of the sensitivity of the expected line strength to parameters of the HII regions, the derived quantities (M<sub>HII</sub> and N<sub>Lyc</sub> ) can be well constrained by the observed relative strengths of the H40$`\alpha `$ and H31$`\alpha `$ lines. Fig 9 indicates that the high density models can also be distinguished by their contribution to the continuum flux density at mm wavelengths.
The Model labelled A2 in Fig 9 (and Table 8) gives a good fit to the observed relative strengths of the H31$`\alpha `$, H40$`\alpha `$, and H42$`\alpha `$ lines. The solid line in Fig 9(b) shows the expected continuum spectrum if only thermal emission from component A2 is added to the non-thermal emission. Because of the low area filling factor ($`f_c`$) of these components, there is no effect on the propagation of the non-thermal radiation, although the HII regions are optically thick below about 40 GHz. Thermal emission from component A2 at 100 GHz is only about 0.6 mJy. In this model, the mass of ionized gas M<sub>HII</sub> = $`3.6\times 10^3`$ $`M_{}`$ which is negligible compared to the mass of component A1, $`6\times 10^7`$ $`M_{}`$ (Table 7), required to explain the H92$`\alpha `$ line. The number of Lyman continuum photons N<sub>Lyc</sub> for the high density component A2 is $`3.5\times 10^{53}`$ s<sup>-1</sup>, about 3% of that for the lower density component A1 (Table 7).
The parameters of the high density component that fit the mm wavelength RRLs are not unique. The alternative model B2, shown as a dashed line in Fig 9, which is reasonably consistent with the mm RRLs has a factor of ten higher M<sub>HII</sub> and N<sub>Lyc</sub> (see Table 8). This model contributes a higher thermal continuum flux density ($`8`$ mJy) at 100 GHz (see Fig 9). A better determination of RRL and continuum flux densities are required to choose between models A2 and B2. As discussed in section 5.4, the high density component provides important information about the star formation rate during recent epochs in Arp 220.
#### 4.2.3 A Combined Model with Three Ionized Components
In this section, we combine the best-fitting models presented in the previous two sections and introduce a third component to account for all the observations. As discussed above, no single-density ionized component is consistent with all the observations. The presence of multiple components of ionized gas is to be expected since it is extremely unlikely that a complex starburst region like Arp 220 could consist of a single density ionized component. To construct a three component model, we first select two models which provide good fits to the H92$`\alpha `$ and H42$`\alpha `$ lines. We selected models which are labelled A1 and A2 in Figs 8 and 9, respectively. The parameters of these models are listed in Tables 7 and 8. The sum of the contributions from these two models to the line emission can account for both H92$`\alpha `$ and H42$`\alpha `$ and also is consistent with H40$`\alpha `$ and H31$`\alpha `$ line and the upper limit to the H167$`\alpha `$ line. Fig 10 illustrates the line and continuum emission from these two components. These two components together can also account for the continuum spectrum above 1 GHz. In this model, the intrinsic non-thermal emission has a spectral index $`\alpha 0.8`$. Because of the presence of the thermal components, the spectral index changes to $`0.6`$ in the range 2–20 GHz. A break in the spectrum is observed near 2 GHz which can be accounted for by component A1. This component, which has an area covering factor of 0.7, becomes optically thick around 2 GHz and therefore progressively absorbs about 70% of the non-thermal radiation at lower frequencies. Component A2, which becomes optically thin above $`40`$ GHz, contributes very little to the continuum emission at any wavelength. Thermal contribution to the continuum emission comes mainly from component A1 and it becomes significant in comparison to the non-thermal component above $`20`$ GHz. The thermal and non-thermal emissions contributions are about equal at $`\nu `$ 50 GHz. Since the thermal component exceeds the non-thermal component above 50 GHz, the continuum spectrum becomes flatter at millimeter wavelengths. In the continuum spectrum shown as a solid line in Fig 10b, thermal dust emission, which may have a significant contribution even at 100 GHz, has not been included. Above 200 GHz, the total continuum is dominated by dust emission (Rigopoulou et al 1996, Downes and Solomon 1998). From the sub-millimeter continuum measurements by Rigopolou et al (1996), we estimate that $``$20% of the continuum emission at 3 mm could be contributed by dust.
The observed continuum spectrum below a few hundred MHz is not accounted for by the two thermal components A1 and A2 discussed above. An additional free-free absorbing component with an emission measure of a few times $`10^5`$ pc cm<sup>-6</sup> and an area covering factor $`f_c1`$ is needed to account for the observed turnover. Furthermore, the flux densities measured at 327 MHz and 150 MHz (Table 5) indicate that the roll off in the spectrum is not steep enough to be accounted for by a foreground thermal screen. On the other hand, if the thermal gas is mixed with the non-thermal emitting gas, then a shallower roll off is expected. It was possible to obtain a good fit to the observed spectrum by adding a thermal component (mixed with the non-thermal gas) with an emission measure EM = $`1.3\times 10^5`$, $`n_e`$= 1000 cm<sup>-3</sup>, $`T_e`$= 7500 K and an area covering factor of unity. The line and continuum emission of this component (shown as model D in Fig 10) is very weak and is practically undetectable at any frequency. The only observable aspect of this component is the turnover in the continuum spectrum below a few hundred MHz. The main constraint for this component is its emission measure. Densities lower than about 500 cm<sup>-3</sup> can be ruled out as they are inconsistent with the upper limit to RRLs near 1.4 GHz. For the model shown in Fig 10, $`n_e`$=1000 cm<sup>-3</sup> and $`T_e`$=7500 K which are same as those of component A1.
All the parameters of this ionized component (model D) along with those of components A1 and A2 and the sum of the three components, where appropriate, are given in Table 9. Less than 10% of the ionized mass is in component D and $``$ 6% of the Lyman continuum photons arise from this component. Most of the mass (94%) is in component A1, accounting for $`92`$% of the total Lyman continuum photons. The mass in the high density component (A2) is negligible ($`0.01`$%) and it accounts for about 3% of the Lyman continuum photons.
Although the models given in Table 9 provide good fit to the observed line and continuum data as seen in Fig 10, there are two aspects which are not satisfactory: (1) Component A1, which contains the bulk of the ionized gas at a density around 1000 cm<sup>-3</sup>, is barely consistent with the upper limit to the H167$`\alpha `$ line. In fact, if this model is correct, then the H167$`\alpha `$ line should be detectable with a factor of 2-3 increase in sensitivity. Although an increase in density of this component will reduce the intensity of the H167$`\alpha `$ line (e.g. model C1 in Fig 8), the model will then fail to account for the observed break in the continuum spectrum around 1.5 GHz. (2) Model B2 in Fig 9, which is an alternative to A2 also provides a reasonable fit to the observations, but has an order of magnitude higher M<sub>HII</sub> and N<sub>Lyc</sub> . This model slightly overestimates the continuum flux density near 100 GHz but it provides a good fit to the mm wavelength RRLs. As explained in Section 5.4, there are significant implications to the star formation history of Arp 220 if N<sub>Lyc</sub> in the high density component is as high as in component B2. A resolution of the two difficulties mentioned here must await a firm detection or determination of a more sensitive upper limit to the H167$`\alpha `$ line and a more accurate determination of the line and continuum parameters at millimeter wavelengths. The validity or otherwise of model B2 can be determined if several mm wavelength RRLs and continuum are observed in the frequency range 100 to 300 GHz and a proper separation of continuum emission by dust and ionized gas is performed using the data.
## 5 Discussion
### 5.1 Density
As discussed in the previous section, only models with densities $`>10^3`$ cm<sup>-3</sup> can fit the recombination line data. Lower density models are inconsistent with the upper limit to the recombination line strength near 20 cm. While densities in the range $`10^32.5\times 10^4`$ cm<sup>-3</sup> fit the H92$`\alpha `$ line, even higher densities ($`n_e`$$`15\times 10^5`$ cm<sup>-3</sup>) are required to account for the millimeter wavelength recombination lines. These results confirm the point made by Zhao et al (1996) that recombination lines at different frequencies act as “density filters”. Thus multi-frequency RRL observations provides an excellent method of determining the various density components in starburst regions.
That a substantial fraction of the gas in the nuclear regions of Arp 220 is at densities $`>10^4`$ cm<sup>-3</sup> is evident from the high $`\mathrm{L}_{\mathrm{HCN}}/\mathrm{L}_{\mathrm{CO}}`$ ratio observed by Solomon, Downes and Radford (1992). The total mass of H<sub>2</sub> gas at densities higher than $`10^4`$ cm<sup>-3</sup> is $`>10^9`$ $`M_{}`$ or more than 25% of the dynamical mass of the two nuclear components (Downes & Solomon 1998). Less than $`10^4`$ of this gas needs to be in ionized form to account for the high density ionized gas deduced from millimeter wavelength recombination lines. The bulk of the ionized gas ($`3\times 10^7`$ $`M_{}`$) which is in component A1 (Table 9) is at a density $`10^3`$ cm<sup>-3</sup> or higher.
Electron densities inferred from \[S III\] and \[Ne III\] IR line ratios observed with the ISO satellite typically yield lower values ($`n_e100300`$ cm<sup>-3</sup>) (Lutz et al 1996). We suggest that this is a selection effect caused by the insensitivity of the \[S III\] line ratios to densities outside the range $`10^210^{3.5}`$ cm<sup>-3</sup> (Houck et al 1984). The upper limit to the H167$`\alpha `$ line near 1.4 GHz presented in Table 2, implies severe constraints on the amount of low-density (i.e. $`n_e`$$`<`$ 500 cm<sup>-3</sup>) ionized gas present in Arp 220. Fig 11 shows the expected strength of radio recombination lines as a function of frequency from ionized gas with $`n_e=500`$ cm<sup>-3</sup> which is ionized by a Lyman continuum luminosity equal to that deduced from Br$`\alpha `$ and Br$`\gamma `$ recombination lines observed with ISO (i.e. N<sub>Lyc</sub> = $`1.3\times 10^{55}`$ s<sup>-1</sup>, Gënzel et al 1998). The predicted strength of the H167$`\alpha `$ recombination line is five times the 3$`\sigma `$ upper limit in Table 2. Thus the only way to consistently explain the ISO and radio results is if the IR recombination lines (Br$`\alpha `$ and Br$`\gamma `$) also arise in higher density gas (i.e. $`n_e>10^3`$ cm<sup>-3</sup>). The three component model in Table 9 shows that the total number of Lyman continuum photons, N<sub>Lyc</sub> , is $`1.3\times 10^{55}`$ photons s<sup>-1</sup>. This value of N<sub>Lyc</sub> is equal to the intrinsic Lyman continuum luminosity inferred from IR recombination lines observed with the ISO satellite (Sturm et al 1996, Lutz et al 1996, Gënzel et al 1998). The simplest reason for this correspondence is that almost all of the IR recombination lines also arise in the same gas with density $`n_e`$$`>10^3`$ cm<sup>-3</sup>. That still leaves the question why \[S III\] line ratios observed with ISO indicates a density of only a few hundred cm<sup>-3</sup>, although they are sensitive to up to ten times higher densities. It is possible that if dust extinction is very high (see next section), \[SIII\] lines are not reliable indicators of density.
### 5.2 Extinction
The predicted flux of the Br$`\alpha `$ line from the three component model in Table 9 is $`1.0\times 10^{15}`$ W m<sup>-2</sup>. Sturm et al (1996) report a measured flux of $`2.1\times 10^{16}`$ W m<sup>-2</sup> from ISO data. Thus extinction of Br$`\alpha `$ is $`A_{Br\alpha }1.7`$ which correspond to a $`V`$ band extinction $`A_V45`$. Similar calculation based on the predicted Br$`\gamma `$flux of $`3.2\times 10^{16}`$ W m<sup>-2</sup> and a measured flux of $`5.9\times 10^{18}`$ W m<sup>-2</sup> (Goldader et al 1995) yield $`A_{Br\gamma }4.3`$ and $`A_V41`$. These derived values of $`A_V`$ are comparable to the values obtained using various IR line ratios by Sturm et al (1996). This similarity in the derived $`A_V`$ values provides further evidence that the IR recombination lines of hydrogen observed by ISO arise in high density gas. The radio recombination line data confirms the high extinction values found from ISO spectroscopy (Sturm et al 1996, Gënzel et al 1998) and provides support to the idea that Arp 220 is powered by a massive starburst (see below).
### 5.3 N<sub>Lyc</sub> and Star Formation Rate (SFR)
One of the derived quantities of the three component model in Table 9 is the production rate of Lyman continuum photons N<sub>Lyc</sub> . Lyman continuum (Lyc) photons are produced by massive O and B stars during their main sequence phase. Since these stars have a relatively short main sequence life time of $`<10^7`$ years, the derived value of N<sub>Lyc</sub> can be related to the formation rate of O and B stars if direct absorption of Lyc photons by dust is not significant. If an initial mass function (IMF) is adopted, then the formation rate of OB stars can be used to derive the total star formation rate (SFR) (e.g. Mezger 1985). The derived SFR is both a function of the adopted IMF and assumed upper and lower mass limits of stars that are formed ($`m_u`$ and $`m_l`$). It is thought that in the case of induced star formation, triggered for example by an external event such as a galaxy-galaxy interaction or a merger, the lower mass cut off in the IMF may be a few solar masses (Mezger 1987). In spontaneous star formation in a quiescent molecular cloud, the lower mass cut off may be $`<0.1`$ $`M_{}`$ set by theoretical considerations (Silk 1978). In the following discussion we use $`m_l`$ = 1 $`M_{}`$ and $`m_u`$ = 100 $`M_{}`$, the latter corresponding to an O3.5 star. Using these mass limits and assuming a constant rate of star formation (over the life time of O stars), we get from the formulae given by Mezger (1985)
$$N_{Lyc}=5.4\times 10^{52}\times 𝚿_{OB}\mathrm{s}^1,$$
(1)
where $`𝚿_{OB}`$ ($`M_{}`$ yr<sup>-1</sup>) is the SFR averaged over the lifetime of the OB stars. In eqn 1, the IMF proposed by Miller and Scalo (1978) has been used. Using a single power-law IMF of Salpeter (1955) results in a factor of $`3`$ increase in N<sub>Lyc</sub> . A reduction in the lower mass limit $`m_l`$ to 0.1 $`M_{}`$ decreases N<sub>Lyc</sub> by $`2`$. The value of N<sub>Lyc</sub> is less sensitive to the upper mass limit in the range 50-100 $`M_{}`$.
The sum of N<sub>Lyc</sub> for the three components in Table 9 is $`1.3\times 10^{55}`$ s<sup>-1</sup>. An examination of the various models discussed in the previous section indicate that the uncertainty in this value is likely to be less than a factor of two. This value of N<sub>Lyc</sub> is consistent with the results from ISO based on NIR recombination lines (Sturm et al 1996 and Gënzel et al 1998) and also with the lower limit of $`1\times 10^{55}`$ s<sup>-1</sup> deduced by Downes and Solomon (1998) based on their estimated thermal continuum at 113 GHz.
The derived SFR based on equation (1) and the total N<sub>Lyc</sub> in Table 9 is $``$ 240 $`M_{}`$ yr<sup>-1</sup>. The SFR in Arp 220 is thus about two orders of magnitude higher than in the Galaxy and may be the highest SFR derived in any normal, starburst or ultra-luminous galaxy (e.g. Kennicutt et al 1998). If the Salpeter-IMF is used and the upper mass limit is reduced to 60 $`M_{}`$, then the SFR is reduced to 90 $`M_{}`$ yr<sup>-1</sup>.
The high SFR in Arp 220 is likely a consequence of the large gas content in the nuclear region with high volume and surface densities. Downes and Solomon (1998) have estimated that within a radius of about 500 pc, which includes the east and west nuclei as well as a portion of a gas disk that surrounds them, the total mass of H<sub>2</sub> gas is $`4\times 10^9`$ $`M_{}`$. For the same region, Scoville et al (1997) obtain a higher mass of $`8\times 10^9`$ $`M_{}`$. For the discussion here, we take a mean of the two, $`6\times 10^9`$ $`M_{}`$. The average volume density of the H<sub>2</sub> is $`250`$ cm<sup>-3</sup> and the average gas surface density $`\mu _g`$ 6000 $`M_{}`$ pc<sup>-2</sup>. A detailed study of star formation rates in normal and ultra-luminous galaxies by Kennicutt (1998) has shown that SFR per unit area can be fitted to a Schmidt law of the form $`\mathrm{\Sigma }_{SFR}=2.5\times 10^4\mu _g^{1.4}`$ $`M_{}`$ yr<sup>-1</sup> kpc<sup>-2</sup>. This relation yields an SFR of 38 $`M_{}`$ yr<sup>-1</sup> for $`\mu _g`$ = 6000 $`M_{}`$ pc<sup>-2</sup>, which is a factor of six lower than deduced above. In his compilation of starburst properties of luminous galaxies, Kennicutt (1998) uses $`\mu _g`$ = 5.8$`\times 10^4`$ $`M_{}`$ pc<sup>-2</sup> for Arp 220 and an SFR of 955 $`M_{}`$ yr<sup>-1</sup> kpc<sup>-2</sup>. This high value of $`\mu _g`$ is the peak surface density in the inner most region obtained by Scoville et al (1997) and it is about a factor of 10 higher than the average surface density over a $``$kpc<sup>2</sup> area. Thus, when averaged over a kpc<sup>2</sup> region, the SFR predicted by Schmidt law with an exponent of 1.4 falls short of the value deduced above (240 $`M_{}`$ yr<sup>-1</sup>) by about a factor of six. Although other parameters such as the upper and lower mass limits and the shape of the IMF could be adjusted to lower the deduced SFR, it is unlikely to reduce it to the value predicted by the empirical Schmidt law obtained by Kennicutt (1998).
In a study of star forming regions in normal galaxies, Viallefond et al (1982) found a relation between $`M_{gas}`$, N<sub>Lyc</sub> and density $`n`$ of the form N<sub>Lyc</sub> = $`7.9\times 10^{44}M_{gas}n^{0.36}`$. Using $`M_{gas}=6\times 10^9`$ $`M_{}`$ and $`n=250`$ cm<sup>-3</sup> (the average molecular density), this relation predicts N<sub>Lyc</sub> = $`3.3\times 10^{55}`$ s<sup>-1</sup>. Although this value is about a factor of 2.5 higher than the derived total N<sub>Lyc</sub> in Table 9, these values are comparable given the uncertainties. This result shows that the starburst in Arp 220 behaves as expected from scaling the known properties of star forming regions in normal galaxies.
### 5.4 Star Formation at Recent Epochs
As shown in Section 4, the observed intensities of the millimeter wavelength RRLs H31$`\alpha `$, H40$`\alpha `$ and H42$`\alpha `$ in Arp 220 can only be explained by ionized gas at a high density of $`2.5\times 10^5`$ cm<sup>-3</sup>. The presence of these high density HII regions which account for about 3% of the total N<sub>Lyc</sub> of $`1.3\times 10^{55}`$ s<sup>-1</sup>, can be used to derive the star formation rates at recent epochs. Since the high density compact HII regions are relatively short-lived ($`\tau _{HII}10^5`$ yrs), the N<sub>Lyc</sub> value corresponding to these regions will indicate the SFR averaged over the life time of the HII regions rather than the main-sequence life time of O, B stars which ionize them. Based on equation (1), and approximating the Lyc photon production rate of OB stars to be constant during their life time, we can write
$$𝚿_{HII}=𝚿_{OB}\times \left[\frac{\tau _{OB}}{\tau _{HII}}\right]\times \left[\frac{N_{Lyc,HII}}{N_{Lyc,OB}}\right].$$
(2)
The compact, high density phase of an HII region is short lived because the HII region expands as it is over-pressured with respect to the surroundings. The expansion proceeds at the sound speed $`c_i`$ ($`10`$ km s<sup>-1</sup>) and approximately follows the relation
$$r(t)=r_i\left[1+\frac{7c_it}{4r_i}\right]^{4/7}$$
(3)
Spitzer (1968), where $`r_i`$ is the initial size of the HII regions and $`r(t)`$ is its size after a time $`t`$. The HII regions of component A2 in Table 9 have a density of $`2.5\times 10^5`$ cm<sup>-3</sup>. The density of these HII regions will drop to about 1000 cm<sup>-3</sup> and become indistinguishable from component A1, if they expand to $`6`$ times their initial size. Taking the initial size to be 0.1 pc, as given in Table 9, equation (3) gives $`t=10^5`$ yrs. If the initial size is smaller, then the time scale is also shorter. (For a single O5 star, the initial size of the Stromgren sphere is $``$0.05 pc if the density is $`2.5\times 10^5`$ cm<sup>-3</sup>.) Using $`\tau _{HII}=10^5`$ yrs, $`\tau _{OB}=3\times 10^6`$ yrs, $`𝚿_{OB}`$ = 240 $`M_{}`$ yr<sup>-1</sup> and using N<sub>Lyc</sub> values from Table 9, equation 2 gives the SFR averaged over the life time of HII regions $`𝚿_{HII}=`$ 194 $`M_{}`$ yr<sup>-1</sup>. This rate is similar to the SFR averaged over $`\tau _{OB}`$. This result indicates that the starburst in Arp 220 is an on-going process with a minimum age exceeding $`\tau _{OB}`$. Other evidence (see below) indicates that the age of the starburst is probably much longer.
The above calculation demonstrates that a reliable measurement of the number of Lyman continuum photons in high density HII regions can lead to a determination of the “instantaneous” SFR. Such a measurement seems possible through RRLs and continuum at millimeter wavelengths which are sensitive to the dense component. The actual value of the instantaneous SFR determined above (194 $`M_{}`$ yr<sup>-1</sup>) can be almost a factor of ten higher or about a factor of two lower since the value of $`N_{Lyc,HII}`$ in equation 2 is uncertain by those factors (see Table 8). If the ten times higher value of $`N_{Lyc,HII}4.6\times 10^{54}`$ s<sup>-1</sup> corresponding to model B2 in Table 8 is established through further measurements, then the implied SFR at recent epochs will be ten times the average value. Such a result will fundamentally alter our picture of the starburst process in Arp 220. Instead of being bursts of constant but high SFR (approximately 10-100 times the Galactic rate) over $`10^610^7`$ yrs, starburst in Arp 220 will be multiple events of short duration ($`10^510^6`$ yrs) but at extremely high SFR (100-1000 times the Galactic rate). Further continuum and RRL measurements at millimeter wavelengths are required to clarify this aspect.
### 5.5 N<sub>Lyc</sub> and IR-Excess
In dense galactic HII regions, the ratio $`\frac{L_{IR}}{L_\alpha }`$ is in the range 5-20 which is known as the IR-excess (IRE) (Mezger et al 1974, Panagia 1978, Mathis 1986). $`L_{IR}`$ is the observed infrared luminosity. $`L_\alpha `$ is the luminosity of Lyman-$`\alpha `$ photons deduced from observations of the ionized gas using recombination lines or radio continuum and assuming that all the Lyman continuum photons which ionize the gas are eventually converted to L$`\alpha `$ photons (i.e. $`N_{L_\alpha }`$=N<sub>Lyc</sub> ). These L$`\alpha `$ photons are then absorbed by dust. If this is the only source of heating the dust in an ionization bounded HII region, and since dust re-radiates most of the absorbed energy in the IR, then the expected IRE is $``$1. If IRE $`>1`$, then the dominant mechanism for heating the dust in HII regions contains a major contribution other than $`L\alpha `$ heating. Two possible sources of heating are: (1) direct absorption of Lyc photons by dust and (2) absorption of photons long ward of the Lyman limit. Lyc photons which are directly absorbed by dust are not counted in the N<sub>Lyc</sub> derived from RRL or continuum observation of the ionized gas. Therefore, if direct absorption of Lyc photons is the main contributor to the IRE, then the SFR derived using eqn 1 will be an underestimate by a factor $``$IRE. On the other hand, heating of dust by photons long ward of the Lyman limit (which are produced mainly by lower mass, non-ionizing stars) can account for the IRE without a corresponding increase in the SFR. Viallefond (1987) has shown that a population of non-ionizing stars formed continuously or in a burst with a Miller-Scalo IMF can indeed explain an IRE of $``$10. In Arp 220, $`L_{IR}=1.3\times 10^{12}L_{}`$ and $`L_{L_\alpha }=5.5\times 10^{10}L_{}`$ using $`N_{L_\alpha }`$ = N<sub>Lyc</sub> = $`1.3\times 10^{55}`$ s<sup>-1</sup>. Thus for Arp 220, IRE $``$ 24. Based on this value of IRE we make a few deductions about the star formation history in Arp 220.
#### 5.5.1 Starburst or AGN?
In giant and super-giant HII regions of the normal galaxy M33, the value of IRE is in the range 3 to 8 and the global value of IRE is $`14`$ (Rice et al 1990). In the Galaxy, IRE is typically $``$ 7 in compact HII regions and $`5`$ in extended HII regions (Mezger 1987). Gënzel et al (1998) have computed a modified form of IRE (they use $`L_{Bol}/L_{Lyc}`$) in starburst galaxies observed by ISO. Converting their ratios to the above definition of IRE (i.e.$`\frac{L_{IR}}{L_\alpha }`$), the mean value of IRE in starburst galaxies is $`25`$ with values ranging from 12 to 45. Similar computation made for AGN dominated sources show that IRE$`4565`$ (Gënzel et al 1998). Thus IRE in Arp 220 ($`24`$) is very similar to the values observed in starburst galaxies and slightly higher than in star forming regions in the Galaxy and in M33. IRE in Arp 220 is significantly lower than in AGN dominated sources. The high infrared luminosity of Arp 220 can therefore be entirely accounted for by processes that operate in starburst regions. We therefore conclude that Arp 220 is powered entirely by a starburst. There is no significant contribution from an “active” nucleus to the observed high IR luminosity in Arp 220.
#### 5.5.2 IMF, Age of Starburst and Star Formation Efficiency in Arp 220
The fact that the IRE for Arp 220 is not significantly different from the IRE in star forming regions in the Galaxy and in M33 also implies that the starburst in Arp 220 has an IMF which is not unusual. If, for example, IRE$``$1, then there is no contribution to heating of dust by non-ionizing photons which would imply that the burst is mainly producing massive stars and thus $`m_l>`$ 5 $`M_{}`$. This is not the case in Arp 220. On the other hand if IRE$`>>`$10, then it would imply one of the following: (1) the IMF is truncated at the upper end ($`m_u<1020`$ $`M_{}`$ or so) which makes the heating of dust by non-ionizing photons even more pronounced, leading to a much larger IRE or (2) the starburst ended about $`5\times 10^6`$ years ago which would reduce N<sub>Lyc</sub> in relation to non-ionizing photons, again resulting in a larger IRE or (3) an AGN is contributing predominantly to the heating and IR emission. These three possibilities are also ruled out in Arp 220.
Another implication of IRE $``$ 24 in Arp 220, which is significantly higher than in young star forming regions, is that the starburst in Arp 220 must be much longer in duration than the lifetime of OB stars (i.e. $`t_{SB}>10^7`$ yrs). As mentioned above, in younger star forming regions such as the super-giant HII region NGC 604 in M33 and in compact HII regions in the Galaxy, the IRE is significantly smaller ($`57`$, Mezger 1985, Rice et al 1990). To illustrate that a relatively high value of IRE implies a longer duration of starburst, we consider a simple model of continuous star formation discussed in Viallefond and Thuan (1983) and the tabulation of the derived parameters in Viallefond (1987). In this model, for an assumed IMF with lower ($`m_l`$) and upper ($`m_u`$) stellar mass limits, N<sub>Lyc</sub> and L<sub>bol</sub> are computed as a function of duration of star formation ($`\tau _{csf}`$) for a SFR of 1 $`M_{}`$ yr<sup>-1</sup>. These values, taken from Viallefond (1987), are tabulated in Table 10 for a Miller-Scalo IMF with $`m_u=100`$ $`M_{}`$ and $`m_l=1`$ $`M_{}`$. The results would be similar for $`m_l=0.1`$ $`M_{}`$. In this model, a steady state production rate of Lyc photons is reached after $`5\times 10^6`$ yrs. The steady state value of N<sub>Lyc</sub> is $`3.7\times 10^{52}`$ ph s<sup>-1</sup>/($`M_{}`$ yr<sup>-1</sup>) and implies a SFR of $``$ 350 $`M_{}`$ yr<sup>-1</sup> since the derived value of N<sub>Lyc</sub> = $`1.3\times 10^{55}`$ ph s<sup>-1</sup>. This SFR is consistent with the value derived in Section 5.3. Table 10 also shows that in this model, the observed IRE$``$ 24 in Arp 220 is reached at $`\tau _{csf}1.5\times 10^7`$ yrs.
Finally, a continuous SFR of 350 $`M_{}`$ yr<sup>-1</sup> for a duration of $`1.5\times 10^7`$ yrs converts $`5\times 10^9`$ $`M_{}`$ of gas into stars. Thus the mass of stars formed is approximately equal to the present mass of gas ($`6\times 10^9`$ $`M_{}`$). In a closed-box model (i.e. with no fresh gas added to the system from outside), the star formation efficiency (SFE) in Arp 220 is 50%, much higher than the average SFE of a few % in Galactic star forming regions (Elmegreen 1983, Myers et al 1986). The high SFE ($`50`$%) and the high SFR ($``$300$`M_{}`$ yr<sup>-1</sup>) in Arp 220 and its corresponding high luminosity ($`L_{IR}=1.3\times 10^{12}`$ L) must be a result of confinement of a large quantity of gas (M$`{}_{gas}{}^{}6\times 10^9`$ $`M_{}`$) in a relatively small volume ($``$kpc<sup>3</sup>) with a high average density of 250 cm<sup>-3</sup>. At the East and West peaks of Arp 220, about $`10^9`$ $`M_{}`$ of gas (i.e. about the mass of gas in the Galaxy) is confined to a region $``$100 pc in size leading to a much higher average density of $`10^4`$ cm<sup>-3</sup>. The high concentrations of large quantities of gas in Arp 220 have given rise to a starburst with an efficiency and SFR that is consistent with known properties of star forming regions in normal and starburst galaxies.
## 6 Summary
We have presented observations of radio recombination lines and continuum emission in Arp 220 at centimeter and millimeter wavelengths. We have showed that to explain both the observed variation of recombination line intensity with quantum number and also the observed continuum spectrum in the frequency range 0.15 - 113 GHz, three components of ionized gas with different densities and covering factors are required.
The bulk of the ionized gas is in a component (A1) with an electron density $`1000`$ cm<sup>-3</sup>. The total mass in this component is $`3\times 10^7`$ $`M_{}`$. This component of ionized gas consists of $`2\times 10^4`$ HII regions each $``$5 pc in diameter. This ionized component produces detectable recombination lines at centimeter wavelengths and substantially modifies the non-thermal continuum spectrum at these wavelengths. The ionized gas in this component becomes optically thick below $`1.5`$ GHz and partially absorbs the non-thermal radiation since its area covering factor is $`0.7`$. The intrinsic spectral index of the non-thermal radiation in Arp 220 is $`\alpha 0.8`$ which is modified to an observed value of $`0.6`$ in the 5-10 GHz range due to the presence of this ionized gas. At $`\nu `$ = 5 GHz, the total observed flux density of Arp 220 is $`210`$ mJy, of which 175 mJy is non-thermal and the remaining $`35`$ mJy is thermal emission from component A1. The other two thermal components (see below) produce negligible continuum emission at this frequency.
A second component (labelled A2) at a density of $`2.5\times 10^5`$ cm<sup>-3</sup> with a mass of only $`3.6\times 10^3`$ $`M_{}`$ accounts for the recombination lines observed at millimeter wavelengths and it is not detected in RRLs at centimeter wavelengths. Component A2 consists of $`10^3`$ HII regions, each 0.1 pc in diameter; these regions are optically thick below 40 GHz. The area covering factor of component A2 is small ($`10^5`$) and thus it does not affect the intrinsic spectrum of the non-thermal radiation at centimeter wavelengths. An alternative model which also fits the data has an order of magnitude higher mass at a similar density. Improved measurements of line and continuum parameters at millimeter wavelengths are required to choose the correct model.
A third component (labelled D) of ionized gas with an emission measure of 1.3$`\times 10^5`$ pc cm<sup>-6</sup>, density $`>500`$ cm<sup>-3</sup>, mixed with the non-thermal gas, is needed to account for the observed turnover in the continuum spectrum below 500 MHz. The mass in this ionized component is $`2\times 10^6`$ $`M_{}`$ and requires about $`8\times 10^{53}`$ Lyc photons s<sup>-1</sup>.
The total mass of ionized gas in the three components is $`3.2\times 10^7`$ $`M_{}`$ requiring $`3\times 10^5`$ O5 stars to maintain the ionization. The total production rate of Lyman continuum photons, N<sub>Lyc</sub> = $`1.3\times 10^{55}`$ s<sup>-1</sup> of which $``$92% is used by component A1, $``$3% by component A2 and $`5`$% by component D. N<sub>Lyc</sub> deduced from radio recombination lines is consistent with the values obtained using NIR recombination lines observed with ISO. A comparison of the predicted strengths of Br$`\alpha `$ and Br$`\gamma `$ lines summed over all the three components with observed values shows that the V band extinction due to dust, $`A_V45`$ magnitudes. This value of $`A_V`$ is consistent with the results from ISO observations.
On the assumption of continuous star formation with an IMF proposed by Miller and Scalo (1978) and assuming a stellar mass range of 1-100 $`M_{}`$, the deduced value of N<sub>Lyc</sub> implies a SFR of 240 $`M_{}`$ yr<sup>-1</sup> averaged over the mean main sequence life time of O, B stars ($`3\times 10^6`$ yrs). The dense HII regions of component A2 are short lived ($`\tau _{HII}10^5`$ yrs) since they are highly over-pressured. Thus the value of N<sub>Lyc</sub> corresponding to this component is related to very recent rate of star formation, i.e. averaged over a time scale $`\tau _{HII}10^5`$ yrs. The deduced recent SFR is similar to the average over $`10^7`$ yrs; thus the starburst in Arp 220 is an on-going process. An alternative model, which cannot be excluded on the basis of the available data, predicts an order of magnitude higher mass in the dense HII regions and correspondingly high SFR at recent epochs (i.e. $`t<10^5`$ yrs). If this alternative model is confirmed through further observations at millimeter wavelengths, then the starburst in Arp 220 consists of multiple episodes of very high SFR (several thousand $`M_{}`$ yr<sup>-1</sup>) of short durations ($`10^5`$ yrs).
Finally, based on the value of N<sub>Lyc</sub> deduced from RRL and continuum data, the IR-excess (i.e. the ratio $`L_{IR}/L_{L_\alpha }`$) in Arp 220 is $`24`$, comparable to the values found in starburst galaxies. This similarity in IR-excess implies that Arp 220 is most likely powered entirely by a starburst rather than an AGN. A comparison of the IRE in Arp 220 with the IRE in star forming regions in the Galaxy and in M33 indicates that the starburst in Arp 220 has a normal IMF and a duration much longer than $`10^7`$ yrs. If no in-fall of gas has taken place during this period, then the star formation efficiently (SFE) in Arp 220 is $`50`$%. The high SFR and SFE in Arp 220 is a consequence of concentration of a large mass in a relatively small volume and is consistent with known dependence of star forming rates on mass and density of gas deduced from observations of star forming regions in normal galaxies.
The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc.
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# Vortex shedding and drag in dilute Bose-Einstein condensates
## 1 Introduction
One enticing consequence of the discovery of Bose-Einstein condensation (BEC) in dilute alkali vapours is the potential for refining our understanding of quantum fluids. In particular, the dilute Bose gas provides a near-ideal testing ground for elucidating the role of vortices in the onset of dissipation in superfluids. Recent experiments have demonstrated the formation of quantized vortices by rotational excitation of one and two-component condensates , in analogy with the famous ‘rotating bucket’ experiments in liquid helium . In addition, by moving a far-off resonant laser beam through a condensate, Raman et al. measure a heating rate which suggests a critical velocity for dissipation characteristic of vortex shedding.
The attractive feature of experiments on dilute Bose gases is that the weakly-interacting limit permits quantitative comparison between theory and experiment. The theoretical description is based on the Gross-Pitaevskii equation, a form of non-linear Schrödinger equation (NLSE) . In the NLSE model, the critical velocity for vortex shedding by a cylindrical object was found to be, $`v_\mathrm{c}0.4c`$, where $`c`$ is the speed of sound . In a trapped (inhomogeneous) condensate the critical velocity is lower due to the reduction in density, and hence sound speed, towards the edge of the condensate .
In this paper, we study vortex shedding due to the motion of an object through homogeneous and trapped Bose-Einstein condensates using the NLSE model. The hydrodynamical properties of a NLSE fluid are reviewed in Sec. 2. In Sec. 3 and 4 we discuss the critical velocity for vortex nucleation and the link between vortex formation and drag in homogeneous condensates. In Sec. 5 we consider the properties of trapped condensates and highlight the differences from the homogeneous case.
## 2 Quantum fluid mechanics
At low temperatures and low densities, atoms interact by elastic $`s`$-wave scattering, and collisions can be parameterised by a single variable, the scattering length, $`a`$. For atoms of mass $`m`$, the wavefunction of the condensate, $`\psi (𝒓,t)`$, is given by the solution of the time-dependent Schrödinger equation:
$$i\mathrm{}_t\psi (𝒓,t)=\left[\frac{\mathrm{}^2}{2m}^2+V(𝒓,t)+g|\psi (𝒓,t)|^2\right]\psi (𝒓,t),$$
(1)
where the wavefunction is normalised to the number of atoms, $`N`$, the coefficient of the non-linear term, $`g=4\pi \mathrm{}^2a/m`$, describes the interactions within the fluid, and $`V(𝒓,t)`$ represents external potentials arising from the trap and any moving obstacle.
### 2.1 Fluid equations
The link between the nonlinear Schrödinger equation (1) and the equivalent equations of fluid mechanics is well known . However some points concerning condensate flow near penetrable objects are less familiar. In this section we gather some of the key concepts and equations that figure in the discussion of our simulations.
Classical (isentropic) fluid mechanics is based on two coupled differential equations: one describing the transport of mass, the other the transport of momentum . The relevant quantum variables can be constructed from the wavefunction: the mass density $`\rho `$ and momentum current density $`𝑱`$ are defined as,
$$\rho m\psi ^{}\psi \mathrm{and}J_k(\mathrm{}/2i)(\psi ^{}_k\psi \psi _k\psi ^{}),$$
(2)
where the index $`k`$ denotes the vector component. The fluid velocity is defined by $`v_kJ_k/\rho `$, or equivalently in terms of the phase, $`\varphi `$, of the wavefunction, $`v_k(\mathrm{}/m)_k\varphi `$. Clearly, the velocity field is a potential flow, however it is also compressible and furthermore can support circulation (vorticity) as will be seen.
The conservation of mass (probability), i.e., the continuity equation, follows from the definition of $`\rho `$ and equation (1)
$$_t\rho +_kJ_k=0.$$
(3)
The conservation of momentum equation may be found by considering the rate of change of the momentum current density,
$$_tJ_k+_jT_{jk}+\rho _k(V/m)=0,$$
(4)
where the momentum flux density tensor takes the form ,
$$T_{jk}=\frac{\mathrm{}^2}{4m}(_j\psi ^{}_k\psi \psi ^{}_j_k\psi +\mathrm{c}.\mathrm{c}.)+\frac{g}{2}\delta _{jk}|\psi |^4.$$
(5)
This can be rewritten as,
$$T_{jk}=\rho v_jv_k\sigma _{jk},$$
(6)
where the stress tensor $`\sigma _{jk}`$ is given by,
$$\sigma _{jk}=\frac{1}{2}\delta _{jk}g(\rho /m)^2+(\mathrm{}/2m)^2\rho _j_k\mathrm{ln}\rho .$$
(7)
### 2.2 Pressure, sound and drag
The form of equations (3), (4), and (6) is identical to those for classical fluid flow , the difference emerges from the nature of the stress, equation (7). A classical ideal fluid is characterised by $`\sigma _{jk}=0,`$ for all $`j,k`$. In a viscous fluid, the shear stress ($`\sigma _{jk},jk`$) is produced by velocity gradients between neighbouring streams such that, $`\sigma _{jk}=\eta (_jv_k+_kv_j)`$, where $`\eta `$ is the coefficient of viscosity. This creates a frictional force which gives rise to energy loss. In a pure dilute Bose-Einstein condensate there is no frictional viscosity, but a shear stress arises from density gradients, the second term in equation (7). This property gives rise to the possibility of vortex formation and drag without viscosity.
The pressure (normal stress, $`\sigma _{jk},j=k`$) within the quantum fluid takes the simple form,
$$p=\frac{1}{2}g(\rho /m)^2(\mathrm{}/2m)^2\rho ^2\mathrm{ln}\rho .$$
(8)
The second term, called the quantum pressure, is weak in homogeneous regions of the fluid, that is, far from obstacles or boundaries, vortex lines or shocks. The essential difference between interacting and noninteracting (ideal) fluids is the existence of interaction pressure which supports sound propagation. In the bulk of the fluid, where the quantum pressure is negligible, the speed of sound ,
$$c=\sqrt{p/\rho }=(g\rho /m^2)^{\frac{1}{2}}.$$
(9)
The force on an obstacle moving through a condensate can be calculated from the rate of momentum transfer to the fluid. By integrating Eq. (4), one finds that the $`k`$-th component of the force,
$$F_k=_t_\mathrm{\Omega }d\mathrm{\Omega }J_k=_SdSn_jT_{jk}_\mathrm{\Omega }d\mathrm{\Omega }\rho _k(V/m),$$
(10)
where $`S`$ is the surface of the object or control surface within the fluid , $`\mathrm{\Omega }`$ is the volume enclosed by $`S`$, $`n_j`$ is the $`j`$-component of the normal vector to $`S`$, and $`\mathrm{d}S`$ is a surface element. The second term on the right-hand side can be likened to the buoyancy of the fluid. In the case of homogeneous flow past an impenetrable object (Sec. 3 and 4), the wavefunction vanishes on the object surface and the potential is uniform elsewhere, therefore, only the first term contributes. Conversely, for a penetrable object in a trapped condensate (Sec. 5), $`\mathrm{\Omega }`$ may be chosen to encompass the entire fluid, and the first term is negligible compared to the second.
### 2.3 Quantisation of circulation
The quantum Euler equation follows from combining the equations describing the conservation of mass and momentum, (3) and (4), along with the identity,
$$\rho ^1_j[\rho _j_k\mathrm{ln}\rho ]=2_k[\rho ^{\frac{1}{2}}_j_j\rho ^{\frac{1}{2}}],$$
(11)
allowing the momentum equation to be written as,
$$_tv_k+v_j_jv_k+_k[g\rho /m^2(\mathrm{}^2/2m)\rho ^{\frac{1}{2}}_j_j\rho ^{\frac{1}{2}}+V/m]=0.$$
(12)
The conservation of energy (Bernoulli equation) then follows as an integral of Euler’s equation, or more directly from the real part of equation (1):
$$\mathrm{}_t\varphi +\frac{1}{2}mv^2+g\rho /m(\mathrm{}^2/2m)\rho ^{\frac{1}{2}}^2\rho ^{\frac{1}{2}}+V=0.$$
(13)
Perhaps the most significant quantum effect on the mechanics of the fluid is the quantisation of angular momentum. The circulation is given by,
$$\mathrm{\Gamma }=d𝒓𝒗=(\mathrm{}/m)2\pi ss=0,1,2,\mathrm{}$$
(14)
where the closed contour joins fluid particles. The conservation of angular momentum (Kelvin’s theorem), follows from Euler’s equation (12) and states that the circulation around a closed ‘fluid’ contour does not change in time. This means that within the fluid, vortex lines must created in pairs which emerge from a point. The exception is at boundaries, where the wavefunction is clamped to zero and no closed fluid loop can be drawn, e.g., at the surface of an impenetrable object or from the edge of a trapped condensate .
### 2.4 Units
For a homogeneous fluid flow, where the external potential is due to the obstacle only, it is convenient to rescale length and velocity in terms of the healing length $`\xi =\mathrm{}/\sqrt{mn_0g}`$ and the asymptotic speed of sound $`c=\sqrt{n_0g/m}`$, respectively. In this case, equation (1) becomes
$$\mathrm{i}_t\stackrel{~}{\psi }(𝒓,t)=\left[\frac{1}{2}^2+V(𝒓^{\mathbf{}})+|\stackrel{~}{\psi }(𝒓,t)|^2\right]\stackrel{~}{\psi }(𝒓,t),$$
(15)
where $`\stackrel{~}{\psi }=\psi /\sqrt{n_0}`$ and $`n_0`$ is the number density far from the object. The force per unit length is measured in units of $`\mathrm{}\sqrt{n_0^3g/m}`$. Unless otherwise stated we use these units throughout. For steady flow, in which $`𝒗`$ and $`n`$ are independent of time and $`\varphi =\mu t`$, where $`\mu `$ is the chemical potential, the Bernoulli equation takes the form
$$n\frac{1}{2}(\sqrt{n})^1^2\sqrt{n}+V+\frac{1}{2}v^2=\mathrm{constant}.$$
(16)
## 3 The critical velocity
The critical velocity for the breakdown of superfluidity is often expressed in terms of the Landau condition , $`v_\mathrm{c}=(ϵ/p)_{\mathrm{min}}`$, where $`ϵ`$ and $`p`$ are the energy and momentum of elementary excitations in the fluid, and $`v_c`$ is the flow velocity in the fluid bulk . In the dilute Bose gas, the long wavelength elementary excitations are sound waves and the Landau criterion predicts that $`v_\mathrm{c}=c`$. However, for flow past an object, the local velocity near the obstacle, $`v`$, can become supersonic even when flow velocity, $`U`$, is subsonic. Consequently, the critical flow velocity, $`v_\mathrm{c}`$, where laminar flow becomes unstable occurs at a fraction of the sound speed. An estimate of $`v_\mathrm{c}`$ may be found following the argument suggested by Frisch et al. . For an incompressible flow past a solid object, Bernoulli’s equation (16) (neglecting the quantum pressure) has the simple form,
$$n(v)+\frac{1}{2}v^2=1+\frac{1}{2}U^2,$$
(17)
where $`U`$ is the background flow velocity. The maximum velocity which occurs at the equator of the object is $`v=\frac{3}{2}U`$ for a sphere (or $`v=2U`$ for a cylinder), therefore, $`\frac{1}{2}(\frac{9}{4}1)U^2=1n(v)`$. The critical velocity is reached when the maximum speed, $`v=\frac{3}{2}U`$, is equal to the ‘bulk’ sound speed, $`c=\sqrt{n(v)}`$, which gives $`v_\mathrm{c}=U=\sqrt{8/23}0.59`$ (or $`\sqrt{2/11}0.44`$ for a cylinder). However, for a compressible fluid, the equatorial velocity is slightly larger due to pressure effects. The first-order correction gives $`v=\frac{3}{2}U+\frac{1}{2}U^3`$ which reduces the critical velocity to $`v_\mathrm{c}0.53`$.
The exact value of $`v_\mathrm{c}`$ may be found by solving the uniform flow equation ,
$$\mathrm{i}_t\stackrel{~}{\psi ^{}}(𝒓^{\mathbf{}},t)=\left[\frac{1}{2}^2+V(𝒓^{\mathbf{}})+|\stackrel{~}{\psi ^{}}(𝒓^{\mathbf{}},t)|^2+\mathrm{i}𝒗^{}\right]\stackrel{~}{\psi ^{}}(𝒓^{\mathbf{}},t),$$
(18)
where $`\stackrel{~}{\psi ^{}}(𝒓^{\mathbf{}},t)=\stackrel{~}{\psi }(𝒓,t)`$ is the wavefunction in the fluid rest frame written in terms of the object frame coordinates, $`𝒓^{\mathbf{}}=𝒓𝒗t`$. Stationary solutions of the form, $`\stackrel{~}{\psi ^{}}(𝒓^{\mathbf{}},t)=\varphi (𝒓^{\mathbf{}})\mathrm{e}^{\mathrm{i}\mu t}`$ are found to exist only for $`vv_\mathrm{c}`$, where $`v_\mathrm{c}`$ is the critical velocity for vortex formation.
To illustrate the behaviour of the exact solutions near the critical velocity, we solve equation (18) in 3D for an impenetrable sphere with radius $`R=50`$. The wavefunction, velocity and quantum pressure term near the object are shown in Fig. 2. Note that these parameters are related via the Bernoulli equation (16). The intersection between the velocity $`v`$ and wavefunction amplitude, $`|\psi |`$, curves defines the position where the velocity is equal to the ‘bulk’ sound speed (9). Note that close to the object the effective sound speed is increased due to the quantum pressure term, (8), therefore even though the density is low the flow is not ‘supersonic’. The critical velocity is reached when flow velocity exceeds the speed of sound in the bulk of the fluid, i.e., when the intersection between the velocity and wavefunction curves moves into the region where the quantum pressure term is zero, see Fig. 1(right).
The complication for trapped condensates is that the speed of sound and hence the critical velocity depend upon position. In addition, the object potential is typically penetrable and non-uniform. We return to these topics in Sec. 5.
## 4 Vortex shedding, drag, and dissipation
For flow faster than the critical velocity, $`v>v_\mathrm{c}`$, vortices are emitted approximately periodically. A typical vortex stream pattern for flow past an impenetrable cylinder is shown in Fig. 2. The background flow is from right to left and the vortex - anti-vortex pairs produce a flow pattern which opposes the background. Consequently, the vortex trail separates the main flow from an almost stationary wake. The momentum loss from the fluid is transferred to the object creating a drag force. The contribution of vortex shedding to the drag force can be estimated by considering the momentum transfer due to vortex emission, i.e.,
$$𝑭_\mathrm{v}=f_\mathrm{v}𝒑_\mathrm{v},$$
(19)
where $`f_\mathrm{v}`$ is the vortex shedding frequency and $`𝒑_\mathrm{v}`$ is the momentum of a vortex pair as it is created at the equatorial plane. Small fluctuations in the vortex shedding frequency occur because as the vortices move downstream they interact with each other creating fluctuations in the flow pattern around the object. The drag force is taken to be the time-average over many vortex emission cycles.
In addition to vortex shedding, drag may also arise due to sound waves. The time-average drag, evaluated using equation (10), as a function of velocity for an impenetrable cylinder is shown in Fig. 3. The drag is zero up to the critical velocity, then increases approximately quadratically with $`v`$ . Also shown is the contribution to the drag force from vortex shedding alone, i.e., equation (19). This comparison illustrates that for $`v<c`$ the drag is produced by vortex shedding, whereas for $`v>c`$, an increasingly significant contribution arises from sound waves. For $`v>c`$, the reflected matter waves create a standing wave pattern in front of the object as shown in Fig. 4.
Using similar arguments for energy loss of the flow. Consider now the condensate at rest and the obstacle moving with speed $`v`$. A drag force leads to energy transfer to the condensate. The energy transfer rate is given by
$$\frac{\mathrm{d}E}{\mathrm{d}t}=𝑭_{\mathrm{drag}}.𝒗.$$
(20)
Eqs. (19) and (20) make the important link between vortex shedding, drag and energy dissipation.
## 5 Motion in a trapped condensate
The inhomogeneous density profile and finite size of trapped condensates means that a steady, uniform flow is difficult to achieve. The MIT experiment partially overcame this problem by sweeping the object back and forth at constant velocity within the central region of the condensate, where the density is approximately uniform. In this case, the object moves through its own low-density wake, and consequently the drag law is different from the uniform flow case discussed above.
In Fig. 5 we show the time-averaged drag on a laser beam oscillating in a trapped condensate. Important differences with the uniform flow case (Fig. 3) arise at high and low velocities. The drag force tends to saturate at higher velocities as the object expels fluid from the region of oscillation and the pressure drops.
Fig. 5 also highlights important differences between two and three dimensions. Two dimensions corresponds to the limit of a ‘cylindrical’ condensate, where the density is uniform parallel to the axis of the object (which we define as the $`z`$-axis). However, in realistic three-dimensional situations the density is inhomogeneous along $`z`$, leading to a variation of the speed of sound which vanishes at the condensate edge. This results in a lower critical velocity in 3D than in 2D, as apparent in Fig. 5. In the Thomas-Fermi limit ($`ga^31`$) the density profile is parabolic, therefore the average sound speed along $`z`$ is a factor of $`\pi /4`$ smaller than that at the centre. The remaining reduction arises from the fact that even at very low velocities, vortices tend to be formed where the laser beam intersects the lower density fluid at the condensate edge. The lower densities arising in 3D also leads to enhanced sound emission, and hence enhanced drag at high velocities.
Below the critical velocity, dissipation due to sound emission occurs at the motion extrema, where the object accelerates. This is illustrated in Fig. 6 which shows a cross section through a 2D condensate cut by a moving laser beam. For constant motion, Fig. 6(a), the fluid is distributed symmetrically around the object and the drag, which is given by an overlap integral between the condensate density and the gradient of the object potential (i.e., the second term in equation (10)) is zero. When the object accelerates, Fig. 6(b), the fluid fails to respond rapidly enough to the abrupt change in velocity, and the asymmetry in the overlap between the fluid and the objects leads to a resistance force analogous to dynamic buoyancy. The system relaxes to the uniform flow case by the emission of a sound wave. This corresponds to the ‘phonon heating’ process discussed in .
Although the NLSE model describes dissipation in the sense of energy transfer between a moving object and the fluid, it does not say anything about how that energy (mostly stored within the vortex core) may be subsequently converted into heat. A complete description should include coupling of the condensate to a thermal cloud, and would describe the damping of phonon and vortex modes. Recent work on the non-equilibrium dynamics of the condensate and non-condensate predict a depletion of the condensate fraction as observed in the MIT experiment.
## 6 Summary
The motion of an object through a dilute Bose-Einstein condensate provides an ideal system to study the fundamental problem of the onset of dissipation in superfluids. In this paper we have explained the role of vortex shedding and sound emission in energy transfer between the object and the condensate. No energy transfer is observed under the condition of uniform, steady flow at speeds below a critical velocity. However, if the object accelerates there is a small dissipative effect, even below the critical velocity, due to sound emission. The critical velocity is reached when the local velocity in the bulk of the fluid (i.e., where the quantum pressure is zero) exceeds the speed of sound. Above the critical velocity vortices are emitted leading to a drag force and energy transfer to the fluid. Vortex shedding dominates the energy transfer for intermediate velocities, while sound emission becomes increasingly important for supersonic motion. We highlight some important differences between homogeneous and inhomogeneous trapped condensates. In particular, that the critical velocity is substantially reduced when the object intersects regions of lower density at the condensate edge and that the trap inhomogeneity gives rise to an additional term in the drag force analogous to a buoyancy.
This work is supported by the EPSRC.
## References
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# Dual Banach spaces which contain an isometric copy of 𝐿₁
## 1. Introduction
The duality between a Banach space containing a ‘nice’ copy of $`\mathrm{}_1`$ and its dual space containing a ‘nice’ copy of $`L_1`$ is summarized in the diagram below. Each upward implication follows straight from the definitions and the absence of a downward arrow indicates that the corresponding implication does not hold in general.
The investigation of this duality began when Pełczyński \[P\] proved that if $`X`$ contains an isomorphic copy of $`\mathrm{}_1`$ then $`X^{}`$ contains an isomorphic copy of $`L_1`$. He also proved the converse result under a technical assumption which was later removed by Hagler \[H2\]. Earlier, James \[J\] had shown that if $`X`$ contains $`\mathrm{}_1`$ isomorphically then $`X`$ contains $`\mathrm{}_1`$ almost isometrically. Recently, Dowling, N. Randrianantoanina and Turett \[DRT\] proved that a dual Banach space contains almost isometric copies of $`L_1`$ whenever it contains isomorphic copies of $`L_1`$ (see also \[H1, Corollary 2.32\] for this result). The main result of this paper, Theorem 2, shows that $`X`$ contains asymptotically isometric copies of $`\mathrm{}_1`$ if and only if $`X^{}`$ contains $`L_1`$ isometrically. In the real case, this is a hitherto unpublished result of Hagler \[H1, Theorem 2.2\].
## 2. Notation and Terminology
Henceforth, all Banach spaces are either real or complex. $`X`$, $`Y`$, and $`Z`$ will denote arbitrary (infinite-dimensional) Banach spaces. Let $`C(K)`$ be the space of continuous functions on some compact Hausdorff space $`K`$, let $`L_1`$ be the space of Lebesgue–integrable functions on , and let $`\mathrm{}_p(\mathrm{\Gamma })`$ be the space of scalar-valued functions on the set $`\mathrm{\Gamma }`$ with finite $`_p`$-norm where $`1p\mathrm{}`$, all with their usual norms. Let $`\mathrm{\Delta }`$ be the Cantor set, $`\mathrm{}_p`$ be $`\mathrm{}_p()`$, and $`C`$ be $`C([0,1])`$.
The concept of asymptotically isometric copies of $`\mathrm{}_1`$ was introduced by Hagler \[H1, pg. 14\]. It was revitalized recently by Dowling and Lennard in fixed point theory \[DL\]. A Banach space contains asymptotically isometric copies of $`\mathrm{}_1`$ provided it satisfies one (hence all) of the conditions in the lemma below.
###### Lemma 1.
For a Banach space $`X`$, the following are equivalent.
1. There exist a null sequence $`(\epsilon _n)`$ of positive numbers less than one and a sequence $`(x_n)`$ in $`X`$ such that
$$\underset{n=1}{\overset{m}{}}(1\epsilon _n)|a_n|\underset{n=1}{\overset{m}{}}a_nx_n\underset{n=1}{\overset{m}{}}|a_n|$$
for each finite sequence $`(a_n)_{n=1}^m`$ of scalars.
2. There exist a null sequence $`(\epsilon _n)`$ of positive numbers and a sequence $`(x_n)`$ in $`X`$ such that
$$\underset{n=1}{\overset{m}{}}|a_n|\underset{n=1}{\overset{m}{}}a_nx_n\underset{n=1}{\overset{m}{}}(1+\epsilon _n)|a_n|$$
for each finite sequence $`(a_n)_{n=1}^m`$ of scalars.
3. There exist a null sequence $`(\epsilon _n)`$ of positive numbers and a sequence $`(x_n)`$ in $`X`$ such that
$$\underset{n=k}{\overset{m}{}}|a_n|\underset{n=k}{\overset{m}{}}a_nx_n(1+\epsilon _k)\underset{n=k}{\overset{m}{}}|a_n|$$
for each finite sequence $`(a_n)_{n=k}^m`$ of scalars and $`k`$.
The proof of this lemma is elementary (cf. \[DLT, Theorem 1.7\] for further equivalent formulations). Note that each condition is equivalent to the variant obtained by replacing ‘There exist a’ by ‘For each’ and ‘and’ by ‘there exists.’ A sequence $`(x_n)`$ satisfying one of the conditions in the lemma is called an asymptotically isometric copy of $`\mathrm{}_1`$. See \[DLT\] for a splendid survey of this topic and its applications to fixed point theory.
The proof of James’s theorem \[J\] for $`\mathrm{}_1`$ shows that if $`X`$ contains $`\mathrm{}_1`$ almost isometrically, then for each null sequence $`(\epsilon _n)`$ of positive numbers there exists a sequence $`(x_n)`$ in $`X`$ such that
$$(1\epsilon _k)\underset{n=k}{\overset{m}{}}|a_n|\underset{n=k}{\overset{m}{}}a_nx_n\underset{n=k}{\overset{m}{}}|a_n|$$
for each finite sequence $`(a_n)_{n=k}^m`$ of scalars and $`k`$. Indeed, the line between containing $`\mathrm{}_1`$ almost isometrically and asymptotically isometrically is very fine.
A sequence $`(x_n)`$ in a Banach space $`X`$ is a ($`1+\epsilon `$)-perturbation of an isometric copy of $`\mathrm{}_1`$ (for short, a ($`1+\epsilon `$)-p.i. $`\mathrm{}_1`$-sequence) provided that there exist a Banach space $`Y`$, a linear isometric embedding $`T:XY`$, and a sequence $`(y_n)`$ in $`Y`$ such that $`(y_n)`$ is isometrically equivalent to the unit vector basis of $`\mathrm{}_1`$ and $`y_nTx_n\epsilon `$ for each $`n`$. If furthermore
$$y_nTx_n\stackrel{\text{def}}{=}\epsilon _n\stackrel{n\mathrm{}}{}0,$$
then $`(x_n)`$ is a perturbation of an isometric copy of $`\mathrm{}_1`$ (for short, a p.i. $`\mathrm{}_1`$-sequence) with respect to $`(\epsilon _n)`$. Note that if $`X`$ is separable then $`Y`$ may be taken to be separable.
If $`X`$ is a Banach space, then $`X^{}`$ is its dual space, $`B(X)`$ is its closed unit ball, and $`S(X)`$ is its unit sphere. The closed linear span of a subset $`A`$ of $`X`$ is $`[A]`$. If $`Y`$ is a subspace of $`X`$ then $`\pi :XX/Y`$ is the natural quotient mapping.
For a surjective bounded linear operator $`T:XZ`$, the corresponding bounded linear operator $`T_q`$ is defined by the following (commutative) diagram.
The operator $`T`$ is called an isometric quotient mapping provided $`T_q`$ is an isometry, which is the case if and only if $`T^{}`$ is an isometric embedding. If $`S:XZ`$ is an isomorphic embedding, then $`S^{}`$ is an isometric quotient mapping if and only if $`S`$ is an isometric embedding.
All notation and terminology, not otherwise explained, are as in \[LT\].
## 3. Main Result
Theorem 2, the main result of this paper, may be viewed as the isometric version of the theorems of Pełczyński and Hagler.
###### Theorem 2.
For a Banach space $`X`$, the following are equivalent.
1. $`X`$ contains asymptotically isometric copies of $`\mathrm{}_1`$.
2. $`X`$ contains a perturbation of an isometric copy of $`\mathrm{}_1`$.
3. $`\mathrm{}_1`$ is linearly isometric to a quotient space of a subspace $`X`$.
4. $`L_1`$ is linearly isometric to a subspace of $`X^{}`$.
5. $`C^{}`$ is linearly isometric to a subspace of $`X^{}`$.
6. $`X^{}`$ contains an infinite set $`\mathrm{\Gamma }`$ which is isometrically equivalent to the usual basis of $`\mathrm{}_1(\mathrm{\Gamma })`$ and which is dense-in-itself in the weak-star topology on $`X^{}`$.
And if $`X`$ is separable, then the following is equivalent to each of the above conditions.
* $`C(\mathrm{\Delta })`$ is isometric to a quotient space of $`X`$.
Recall that a subset $`K`$ of a topological space is dense-in-itself if $`K`$ has no isolated points in the relative topology. Our proof of Theorem 2 uses the following results.
###### Lemma 3.
If $`(x_n)`$ is a p.i. $`\mathrm{}_1`$-sequence, then $`(\lambda _nx_n)`$ is an asymptotically isometric copies of $`\mathrm{}_1`$ satisfying (A2) for some suitable choice of scalars $`(\lambda _n)`$. Conversely, an asymptotically isometric copies of $`\mathrm{}_1`$ satisfying (A2) is a p.i. $`\mathrm{}_1`$-sequence.
###### Proof.
Let $`(\stackrel{~}{x}_n)`$ be a p.i. $`\mathrm{}_1`$-sequence with respect to $`(\stackrel{~}{\epsilon }_n)`$. Then
$$\underset{n=1}{\overset{m}{}}(1\stackrel{~}{\epsilon }_n)|a_n|\underset{n=1}{\overset{m}{}}a_n\stackrel{~}{x}_n\underset{n=1}{\overset{m}{}}(1+\stackrel{~}{\epsilon }_n)|a_n|.$$
for each finite sequence $`(a_n)_{n=1}^m`$ of scalars. Define
$$\epsilon _n\stackrel{\text{def}}{=}\frac{1+\stackrel{~}{\epsilon }_n}{1\stackrel{~}{\epsilon }_n}1\text{and}\text{ }x_n\stackrel{\text{def}}{=}\frac{\stackrel{~}{x}_n}{1\stackrel{~}{\epsilon }_n}.$$
Then $`(\epsilon _n)`$ and $`(x_n)`$ satisfy (A2); thus, $`(x_n)`$ is an asymptotically isometric copy of $`\mathrm{}_1`$.
Conversely, let $`(\epsilon _n)`$ and $`(x_n)`$ satisfy (A2). Then $`(x_n)`$ is a p.i. sequence. To see this, let $`X_0=[x_n]`$ and
$$W=\{\left(w_n\right)_{n=1}^{\mathrm{}}:w_n\text{ and }|w_n|=1\text{ for each }n\}.$$
For each $`\omega =(w_n)W`$, define $`f_\omega B(X_0^{})`$ by $`f_\omega (x_n)=w_n`$; for indeed,
$$\left|f_\omega \left(\underset{n=1}{\overset{m}{}}a_nx_n\right)\right|=\left|\underset{n=1}{\overset{m}{}}a_nw_n\right|\underset{n=1}{\overset{m}{}}\left|a_n\right|\underset{n=1}{\overset{m}{}}a_nx_n_{X_0}$$
for each finite sequence $`(a_n)_{n=1}^m`$ of scalars. For each $`\omega W`$, let $`\stackrel{~}{f}_\omega B(X^{})`$ be a norm-preserving Hahn-Banach extension of $`f_\omega `$.
Let
$$Y\stackrel{\text{def}}{=}C(B\left(X^{}\right),\sigma (X^{},X)),$$
endowed with the usual sup norm, and consider the isometric embedding
$$T:XY$$
given by
$$(Tx)(x^{})\stackrel{\text{def}}{=}x^{}(x).$$
Let $`y_nB(Y)`$ be the ‘truncation’ of $`Tx_n`$; specifically,
$$y_n(x^{})=\{\begin{array}{cc}(Tx_n)(x^{})\hfill & \text{if}\text{ }|(Tx_n)(x^{})|1\hfill \\ \frac{(Tx_n)(x^{})}{|(Tx_n)(x^{})|}\hfill & \text{if}\text{ }|(Tx_n)(x^{})|>1.\hfill \end{array}$$
(1)
For each $`n`$, condition (A2) gives that $`x_n1+\epsilon _n,`$ and so by (1)
$$y_nTx_n_Y\epsilon _n.$$
Since for each $`n`$ and $`\omega =(w_j)W`$
$$(Tx_n)(\stackrel{~}{f}_\omega )=\stackrel{~}{f}_\omega (x_n)=f_\omega (x_n)=w_n=y_n(\stackrel{~}{f}_\omega ),$$
it follows that
$$\underset{n=1}{\overset{m}{}}a_ny_n_Y\underset{\omega W}{sup}\left|\underset{n=1}{\overset{m}{}}a_ny_n(\stackrel{~}{f}_\omega )\right|=\underset{(w_n)W}{sup}\left|\underset{n=1}{\overset{m}{}}a_nw_n\right|=\underset{n=1}{\overset{m}{}}\left|a_n\right|.$$
for each finite sequence $`(a_n)_{n=1}^m`$ of scalars. Also, $`y_n1`$ for each $`n`$. Thus $`(y_n)`$ is isometrically equivalent to the unit vector basis of $`\mathrm{}_1`$. ∎
###### Remark 4.
Minor modifications to the above proof give an isomorphic version of Lemma 3. Indeed, if $`(\stackrel{~}{x}_n)`$ be a ($`1+\stackrel{~}{\epsilon }`$)-p.i. $`\mathrm{}_1`$-sequence with $`\stackrel{~}{\epsilon }<1`$ and
$$\epsilon \stackrel{\text{def}}{=}\frac{1+\stackrel{~}{\epsilon }}{1\stackrel{~}{\epsilon }}1\text{and}\text{ }x_n\stackrel{\text{def}}{=}\frac{\stackrel{~}{x}_n}{1\stackrel{~}{\epsilon }},$$
then
$$\underset{n=1}{\overset{m}{}}|a_n|\underset{n=1}{\overset{m}{}}a_nx_n(1+\epsilon )\underset{n=1}{\overset{m}{}}|a_n|$$
(2)
for each finite sequence $`(a_n)_{n=1}^m`$ of scalars. Conversely, if $`(x_n)`$ satisfies (2) for each finite sequence $`(a_n)_{n=1}^m`$ of scalars, then $`(x_n)`$ is a ($`1+\epsilon `$)-p.i. $`\mathrm{}_1`$-sequence.
###### Lemma 5.
If $`X`$ satisfies (f) of Theorem 2, then there exists a separable subspace $`X_0`$ of $`X`$ and a countable subset $`\mathrm{\Gamma }^{}`$ of $`X_0^{}`$ which satisfies (f) of Theorem 2.
###### Proof.
Let $`X`$ be a Banach space satisfying (f) of Theorem 2. We shall inductively construct a sequence $`(\mathrm{\Lambda }_i)`$ of countably infinite subsets of $`\mathrm{\Gamma }`$ and a sequence $`(Z_i)`$ of separable subspaces of $`X`$ which satisfy, for each $`n`$,
1. $`Z_nZ_{n+1}`$,
2. $`Z_n`$ norms $`[_{i=1}^n\mathrm{\Lambda }_i]`$, i.e., if $`x^{}\left[_{i=1}^n\mathrm{\Lambda }_i\right]`$ then
$$x^{}=\underset{zB(Z_n)}{sup}|x^{}(z)|,$$
3. $`_{i=1}^n\mathrm{\Lambda }_i`$ is contained in the $`Z_n`$-cluster points of $`\mathrm{\Lambda }_{n+1}`$,
i.e., if $`x^{}_{i=1}^n\mathrm{\Lambda }_i`$ and $`(w_i)_{i=1}^k`$ are from $`Z_n`$ and $`\epsilon >0`$ then
$$\{y^{}\mathrm{\Lambda }_{n+1}:\left|\left(y^{}x^{}\right)\left(w_i\right)\right|<\epsilon \text{for}\text{ }1ik\}\left\{x^{}\right\}\mathrm{}.$$
For the first step of the induction choose a countably infinite subset $`\mathrm{\Lambda }_1`$ of $`\mathrm{\Gamma }`$ and find a separable subspace $`Z_1`$ of $`X`$ which satisfies (2). Suppose that we have chosen $`(\mathrm{\Lambda }_i)_{i=1}^n`$ and $`(Z_i)_{i=1}^n`$ satisfying the three conditions. Since $`Z_n`$ is separable and elements of $`\mathrm{\Gamma }`$ are of norm one, there is a countable subset $`\mathrm{\Lambda }_{n+1}`$ of $`\mathrm{\Gamma }`$ satisfying (3). Next we find a separable subspace $`Z_{n+1}`$ of $`X`$ which satisfies (1) and (2). This completes the inductive step.
Now let $`X_0=[_{n=1}^{\mathrm{}}Z_n]`$ and
$$\mathrm{\Gamma }^{}\stackrel{\text{def}}{=}\{x^{}|_{X_0}:x^{}_{n=1}^{\mathrm{}}\mathrm{\Lambda }_n\}.$$
Condition (2) gives that $`\mathrm{\Gamma }^{}`$ is isometrically equivalent to the usual basis of $`\mathrm{}_1(\mathrm{\Gamma }^{})`$. Conditions (1) and (3), along with the fact that $`\mathrm{\Gamma }S(X^{})`$, give that $`\mathrm{\Gamma }^{}`$ is dense-in-itself in the weak-star topology on $`X_0^{}`$. ∎
###### Fact 6.
(cf. \[HS, Lemma 4\]) Let $`N`$ and $`M`$ be compact Hausdorff spaces with $`M`$ perfect and suppose that $`\varphi :NM`$ is continuous and onto. Then there exists a subset $`Q`$ of $`N`$ such that $`Q`$ is dense-in-itself and $`\varphi |_Q:QM`$ is a bijection.
###### Fact 7.
(Haskell P. Rosenthal \[R, Proposition 3 and its Remark 2\]) Let $`X_0`$ be a separable Banach space satisfying (f) of Theorem 2. Then there exists
$$KB(X_0^{}),$$
which is homeomorphic to $`\mathrm{\Delta }`$, such that the restriction operator
$$R:X_0C(K)$$
given by $`(Rx_0)(x_0^{})=x_0^{}(x_0)`$ is an isometric quotient mapping.
###### Proof of Theorem 2.
We shall assume that $`X`$ is a complex Banach space as the proof in the real case is easier.
The equivalence of (a) and (b) follows directly from Lemma 3.
To see that (a) implies (c), let $`(\epsilon _n)`$ and $`(x_n)`$ be sequences satisfying condition (A3). Partition $``$ into infinite sets $`\{J_n\}_n`$ and let $`T:[x_n]\mathrm{}_1`$ be the bounded linear operator that maps $`x_j`$ to the $`n^{\text{th}}`$ unit vector of $`\mathrm{}_1`$ when $`jJ_n`$. Then $`T`$ is an isometric quotient mapping.
To see that (c) implies (a), let
$$T:X_0/X_1\mathrm{}_1$$
be an isometry from a quotient space of a subspace $`X_0`$ of $`X`$ onto $`\mathrm{}_1`$. Fix a null sequence $`(\epsilon _n)`$ of positive numbers. Find a sequence $`(x_n)`$ in $`X_0`$ such that $`T(x_n+X_1)`$ is the $`n^{\text{th}}`$ unit vector of $`\mathrm{}_1`$ and
$$1x_n_X1+\epsilon _n.$$
Then $`(\epsilon _n)`$ and $`(x_n)`$ satisfy (A2).
To see that (a) implies (e), let $`(\epsilon _n)`$ and $`(x_n)`$ satisfy (A2). We shall define the bounded linear operators in the (commutative) diagram below
(3)
as follows. Let $`(z_n)`$ be dense in the unit sphere of $`C`$. Define $`T`$ by $`Tx_n=z_n`$. Condition (A2) gives that $`T`$ is a surjective norm-one bounded linear operator. Furthermore, $`T`$ is an isometric quotient mapping for if $`fS(C)`$ then there is a subsequence $`(z_{k_n})`$ converging in norm to $`f`$ and so
$$1T_q^1f_{Y/\text{ker }T}\underset{¯}{\mathrm{lim}}_n\mathrm{}x_{k_n}_X=1.$$
Let $`j`$ be the natural embedding and let $`\widehat{ı}`$ be the canonical isometric embedding given by point evaluation. Since $`C^{}`$ has the Hahn-Banach Extension Property, $`T`$ admits a norm-preserving extension $`\stackrel{~}{T}`$. Dualizing gives the commutative diagram
where $`h`$ is the canonical isometric embedding given by point evaluation. To see that $`R\stackrel{\text{def}}{=}\stackrel{~}{T}^{}h`$ is the desired isometric embedding, let $`\mu C^{}`$. Then, since $`T^{}`$ is an isometric embedding and $`\widehat{ı}^{}h`$ is the identity mapping,
$$\mu _C^{}=T^{}\mu _Y^{}=j^{}R\mu _Y^{}R\mu _X^{}\mu _C^{}.$$
Clearly, (e) implies (d).
To see that (d) implies (a), let $`T:L_1X^{}`$ be an isometric embedding and let $`(\epsilon _n)`$ be a null sequence of positive numbers. Then $`T^{}:X^{}L_{\mathrm{}}`$ is a weak-star continuous isometric quotient mapping. By Goldstine’s Theorem,
$$W\stackrel{\text{def}}{=}T^{}(B(X))$$
is weak-star dense in $`B(L_{\mathrm{}})`$. For each $`n`$, let
$$F_n\stackrel{\text{def}}{=}\{z_j^n:1jM(n)\}$$
be an $`(\epsilon _n/2)`$–net for $`\{z:|z|=1\}`$.
Let $`𝒯`$ be the tree
$$𝒯=\underset{n}{}𝒯_n$$
where $`𝒯_n`$, the $`n^{\text{th}}`$-level of $`𝒯`$, is
$$\begin{array}{c}𝒯_n\stackrel{\text{def}}{=}\{(m_0,m_1,m_2,\mathrm{}m_{n1})^n:\hfill \\ \hfill m_0=1\text{and}\text{ }1m_jM(j)\text{for each}\text{ }j\}.\end{array}$$
(4)
If $`\alpha =(m_0,m_1,m_2,\mathrm{}m_{n1})𝒯`$ then
$$(\alpha ,j)\stackrel{\text{def}}{=}(m_0,m_1,m_2,\mathrm{}m_{n1},j);$$
thus, for each $`n`$
$$𝒯_{n+1}=\{(a,j):\alpha 𝒯_n\text{and}\text{ }1jM(n)\}.$$
We will define inductively, for each $`n`$, a collection $`\{A_\alpha \}_{\alpha 𝒯_n}`$ of disjoint sets of positive (Lebesgue) measure and a function $`f_nW`$ such that, for each $`n`$ and $`\alpha 𝒯_n`$,
$$\underset{j=1}{\overset{M(n)}{}}A_{(\alpha ,j)}A_\alpha [0,1]$$
(5)
and, for each $`1jM(n)`$,
$$|f_nz_j^n|<\frac{\epsilon _n}{2}\text{on}\text{ }A_{(\alpha ,j)}.$$
(6)
To start the induction, let
$$A_{(m_0)}=[0,1].$$
For the inductive step, let $`n`$ and suppose that we have constructed disjoint sets
$$\{A_\alpha :\alpha 𝒯_n\}$$
of positive measure. For each $`\alpha 𝒯_n`$, partition $`A_\alpha `$ into sets $`\{D_{(\alpha ,j)}\}_{j=1}^{M(n)}`$ of positive measure. Consider the function $`g_nB(L_{\mathrm{}})`$ defined by
$$g_n(t)=\{\begin{array}{cc}z_j^n\hfill & \text{if}\text{ }tD_{(\alpha ,j)}\text{and}\text{ }\alpha 𝒯_n\hfill \\ 0\hfill & \text{otherwise}\text{ }.\hfill \end{array}$$
Since $`W`$ is weak-star dense in $`B(L_{\mathrm{}})`$ there exists $`f_nW`$ approximating $`g_n`$ closely enough to ensure that the sets
$$A_{(\alpha ,j)}\stackrel{\text{def}}{=}\{|f_nz_j^n|<\epsilon _n/2\}D_{(\alpha ,j)}$$
all have positive measure. This completes the proof of the inductive step.
For each $`n`$, select $`x_nB(X)`$ such that $`T^{}(x_n)=f_n`$. To see that $`(x_n)`$ is an asymptotically isometric copy of $`\mathrm{}_1`$, let $`(a_n)_{n=1}^m`$ be a finite complex sequence. Define $`(\stackrel{~}{a}_n)_{n=1}^m`$ from $`\{z:|z|=1\}`$ by
$$\stackrel{~}{a}_n=\{\begin{array}{cc}\frac{\overline{a}_n}{\left|a_n\right|}\hfill & \text{if}\text{ }a_n0\hfill \\ 1\hfill & \text{if}\text{ }a_n=0;\hfill \end{array}$$
thus, $`a_n\stackrel{~}{a}_n=\left|a_n\right|`$. For each $`1nm`$, find $`1j_nM(n)`$ so that
$$|\stackrel{~}{a}_nz_{j_n}^n|<\frac{\epsilon _n}{2}.$$
Then $`\alpha \stackrel{\text{def}}{=}(1,j_1,\mathrm{},j_m)𝒯_{m+1}`$ and so by (5) and (6)
$$|f_nz_{j_n}^n|<\frac{\epsilon _n}{2}\text{on}\text{ }A_\alpha $$
for each $`1nm`$. Thus
$`{\displaystyle \underset{n=1}{\overset{m}{}}}a_nx_n`$ $`T^{}\left({\displaystyle \underset{n=1}{\overset{m}{}}}a_nx_n\right)_L_{\mathrm{}}={\displaystyle \underset{n=1}{\overset{m}{}}}a_nf_n_L_{\mathrm{}}`$
$`{\displaystyle \underset{n=1}{\overset{m}{}}}a_n\stackrel{~}{a}_n1_{A_\alpha }_L_{\mathrm{}}{\displaystyle \underset{n=1}{\overset{m}{}}}a_n\left(\stackrel{~}{a}_nf_n\right)1_{A_\alpha }_L_{\mathrm{}}`$
$`{\displaystyle \underset{n=1}{\overset{m}{}}}|a_n|{\displaystyle \underset{n=1}{\overset{m}{}}}\epsilon _n|a_n|={\displaystyle \underset{n=1}{\overset{m}{}}}(1\epsilon _n)|a_n|.`$
So $`(\epsilon _n)`$ and $`(x_n)`$ do indeed satisfy (A1).
To show that (a) implies (f), let (a) hold. Then we have the situation depicted in (3). For $`t[0,1]`$, let $`\delta _tC^{}`$ denote the point mass measure at $`t`$. Since $`T^{}`$ is a weak-star continuous isometric embedding
$$M\stackrel{\text{def}}{=}\{T^{}(\delta _t):t[0,1]\}B(Y^{}),$$
equipped with the weak-star topology of $`Y^{}`$, is homeomorphic to $`[0,1]`$ and is isometrically equivalent to the usual basis of $`\mathrm{}_1([0,1])`$. By the Hahn-Banach Theorem $`j^{}(B(X^{}))=B(Y^{})`$ and so
$$N\stackrel{\text{def}}{=}j^1(M)B(X^{})$$
is weak-star compact and satisfies $`j^{}(N)=M`$. By Fact 6 there exists
$$\mathrm{\Gamma }=\{n_t:t[0,1]\}N$$
such that $`j^{}(n_t)=T^{}(\delta _t)`$ and such that $`\mathrm{\Gamma }`$ is dense-in-itself in the weak-star topology on $`X^{}`$. Since for any finite set $`\{a_t\}_{tF}`$ of scalars
$`{\displaystyle \underset{tF}{}}\left|a_t\right|`$ $`={\displaystyle \underset{tF}{}}a_tT^{}\delta _t_Y^{}`$ $`={\displaystyle \underset{tF}{}}a_tj^{}n_t_Y^{}`$
$`{\displaystyle \underset{tF}{}}a_tn_t_X^{}`$ $`{\displaystyle \underset{tF}{}}\left|a_t\right|,`$
the set $`\mathrm{\Gamma }`$ is isometrically equivalent to the usual basis of $`\mathrm{}_1([0,1])`$.
To see that (f) implies (e), let $`X`$ satisfy (f). Then by Lemma 5, there is a separable subspace $`X_0`$ of $`X`$ which satisfies (f). From Fact 7 and the fact that $`C^{}(\mathrm{\Delta })`$ is linearly isometric to $`C^{}`$, it follows that $`X_0`$ satisfies (e). The equivalence of (a) and (e) gives that $`X`$ also satisfies (e).
Thus (a) through (f) are equivalent. The fact that $`C^{}(\mathrm{\Delta })`$ is linearly isometric to $`C^{}`$ gives that (g) implies (e). That (f) implies (g) when $`X`$ is separable is due to Rosenthal: Fact 7. ∎
###### Remark 8.
Without the added assumption of separability, (g) is not equivalent to the other conditions. Clearly, $`\mathrm{}_{\mathrm{}}`$ satisfies conditions (a) through (f). But by a result of Grothendieck \[G\], a separable quotient of $`\mathrm{}_{\mathrm{}}`$ is reflexive and so $`\mathrm{}_{\mathrm{}}`$ does not satisfy (g).
###### Remark 9.
A complemented isomorphic version of Theorem 2 is due to Hagler and Stegall \[HS, Theorem 1\]. A $`K`$-complemented isometric version of Theorem 2 is due to Hagler (\[H1, Theorem 2.13\] or \[H3\]).
###### Remark 10.
Many Banach spaces (and their subspaces) which arise naturally in analysis contain an abundance of asymptotically isometric copies of $`\mathrm{}_1`$: for example, Carothers, Dilworth and Lennard \[CDL\] proved that every nonreflexive subspace of the Lorentz space $`L_{w,1}(0,\mathrm{})`$ contains asymptotically isometric copies of $`\mathrm{}_1`$ whenever the weight $`w`$ satisfies very mild regularity conditions. On the other hand, $`L_{w,1}`$ does not contain an isometric copy of the $`2`$-dimensional space $`\mathrm{}_1^2`$ whenever $`w`$ is strictly decreasing \[CDT\].
###### Remark 11.
Theorem 2 improves a recent result of Shutao Chen and Bor-Luh Lin \[CL\] who proved that $`X`$ contains an asymptotically isometric copy of $`\mathrm{}_1`$ whenever $`X^{}`$ contains an isometric copy of $`\mathrm{}_{\mathrm{}}`$.
Dowling, Johnson, Lennard and Turett \[DJLT\] gave some concrete examples of equivalent norms on $`\mathrm{}_1`$ such that the corresponding Banach spaces do not contain asymptotically isometric copies of $`\mathrm{}_1`$. Theorem 2 easily yields other equivalent norms on $`\mathrm{}_1`$ with this property.
###### Corollary 12.
Let $`(\gamma _n)`$ be a sequence of non-zero scalars with $`(\gamma _n)_2<\epsilon `$. Then the Banach space $`(\mathrm{}_1,_1^{})`$, where
$$(a_n)_1^{}\stackrel{\text{def}}{=}inf\{[(a_n+b_n)_1^2+\left(\gamma _n^1b_n\right)_2^2]^{\frac{1}{2}}:\left(\gamma _n^1b_n\right)\mathrm{}_2\},$$
does not contain asymptotically isometric copies of $`\mathrm{}_1`$ and
$$\left(1+\epsilon ^2\right)^{\frac{1}{2}}\left(a_n\right)_1\left(a_n\right)_1^{}\left(a_n\right)_1$$
(7)
for each $`\left(a_n\right)\mathrm{}_1`$.
###### Proof.
We exhibit $`(\mathrm{}_1,_1^{})`$ as a quotient space of $`X=\mathrm{}_1_2\mathrm{}_2`$ with its usual norm
$$((a_n),(b_n))_X=\left[\left(a_n\right)_1^2+\left(b_n\right)_2^2\right]^{1/2}.$$
Let
$`Y`$ $`\stackrel{\text{def}}{=}`$ $`\{((a_n),(b_n))X:b_n=\gamma _n^1a_n\}`$
$`\stackrel{\text{note}}{=}`$ $`\{((a_n),(\gamma _n^1a_n)):(\gamma _n^1a_n)\mathrm{}_2\}.`$
Since each element of $`X/Y`$ has a representative of the form $`(\left(a_n\right),0)`$,
$$\begin{array}{c}((a_n),0)+Y_{X/Y}=\hfill \\ \hfill inf\{(\left(a_n\right),0)+(\left(b_n\right),(\gamma _n^1b_n))_X:\left(\gamma _n^1b_n\right)\mathrm{}_2\}\\ \hfill =(a_n)_1^{^{}}.\end{array}$$
Thus $`X/Y`$ is isometrically isomorphic to $`(\mathrm{}_1,_1^{})`$.
Observe that $`X^{}=\mathrm{}_{\mathrm{}}_2\mathrm{}_2`$ with its usual norm
$$((c_n),(d_n))_X^{}=\left[\left(c_n\right)_{\mathrm{}}^2+\left(d_n\right)_2^2\right]^{1/2}$$
and
$$Y^{}=\{((c_n),(d_n))X^{}:d_n=\gamma _nc_n\}=\{((c_n),(\gamma _nc_n)):(c_n)\mathrm{}_{\mathrm{}}\}.$$
It follows that $`Y^{}`$ is isometrically isomorphic to $`(\mathrm{}_{\mathrm{}},_{\mathrm{}}^{})`$ where
$$(c_n)_{\mathrm{}}^{}=\left[(c_n)_{\mathrm{}}^2+\left(\gamma _nc_n\right)_2^2\right]^{1/2}$$
(8)
and the mapping
$$i:(\mathrm{}_{\mathrm{}},_{\mathrm{}}^{})(\mathrm{}_1,_1^{})^{}$$
given by $`\left(i\left(c_n\right)\right)\left(a_n\right)=_na_nc_n`$ is an isometry. Since the norm $`_{\mathrm{}}^{}`$ is strictly convex, the space $`(\mathrm{}_{\mathrm{}},_{\mathrm{}}^{})`$ does not contain an isometric copy of $`L_1`$. Thus by Theorem 2 the space $`(\mathrm{}_1,_1^{})`$ does not contain asymptotically isometric copies of $`\mathrm{}_1`$. From (8) it follows that
$$\left(c_n\right)_{\mathrm{}}\left(c_n\right)_{\mathrm{}}^{}\sqrt{1+\epsilon ^2}\left(c_n\right)_{\mathrm{}}$$
for each $`\left(c_n\right)\mathrm{}_{\mathrm{}}`$ and so (7) holds by duality. ∎
Finally, Alspach’s \[A\] example of an isometry $`T:KK`$ without a fixed point, where $`K`$ is a certain weakly compact convex subset of $`L_1`$, yields an obvious corollary.
###### Corollary 13.
If $`X`$ contains asymptotically isometric copies of $`\mathrm{}_1`$, then there exists a (nonempty) weakly compact convex subset $`K`$ of $`X^{}`$ and an isometry $`T:KK`$ without a fixed point.
###### Acknowledgement.
The authors thank Patrick Dowling, William B. Johnson, and Haskell P. Rosenthal for their helpful comments. H. P. Rosenthal suggested the term p.i. $`\mathrm{}_1`$-sequence and his helpful conversations \[HPR\] led us to the proof of Lemma 3.
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# 1 Introduction
## 1 Introduction
The problem of the Bose-Fermi duality in two-dimensional space-time has a long history. Since the pioneering works in the thirties and till now, various of its aspects, conceptual, as well as technical, remain attractive due to the ambitious idea to extend this remarkable feature also beyond one space dimension. In recent papers , starting with bare fermions which form the $`C^{}`$–algebra $`𝒜=\mathrm{CAR}(𝐑)`$, we have constructed a chain of algebraic inclusions to substanciate our understanding of this phenomenon:
$$\mathrm{CAR}(\mathrm{𝑏𝑎𝑟𝑒})\pi _\beta (𝒜)^{\prime \prime }𝒜_c\overline{𝒜}_c\overline{\pi }_\beta (\overline{𝒜}_c)^{\prime \prime }\mathrm{CAR}(\mathrm{𝑑𝑟𝑒𝑠𝑠𝑒𝑑})$$
(1.1)
To get the first extension, we note that the shift $`\tau _t`$ is an automorphism of $`𝒜`$ which has KMS-states $`\omega _\beta `$ and associated representations $`\pi _\beta `$. One then identifies in $`\pi _\beta (𝒜)^{\prime \prime }`$ bosonic modes — the currents, which form the current algebra $`𝒜_c`$ with a $`\beta `$-independent structure for $`0<\beta <\mathrm{}`$. The crucial ingredient needed at this step is the appropriately chosen state. We choose the KMS-state which is unique for the shift over the CAR algebra. Another possibility would be to chose the Dirac vacuum. This is what has originally been done in , in order to achieve stability for a fermionic system, and recovered later by Mattis and Lieb in the context of the Luttinger model.
The essential result is the appearance of an anomalous (Schwinger) term in the quantum current commutator
$$[j(x),j(x^{})]=\frac{i}{2\pi }\delta ^{}(xx^{})$$
(1.2)
based on which Mandelstam proposed $`e^{i2\pi _{\mathrm{}}^xj(y)𝑑y}`$ as a fermion field . In this note we construct the anyonic fields
$$\mathrm{\Psi }_\alpha (x)e^{i2\pi \sqrt{\alpha }\underset{\mathrm{}}{\overset{x}{}}j(y)𝑑y}$$
(1.3)
through their action on the bosonic current and study them as operators in a Hilbert space by exhibiting their n-point function in a $`\tau `$-KMS state $`\omega `$, $`\tau `$ being the shift automorphism. This corresponds to taking the crossed product of $`𝒜_c`$ with an outer automorphism or, equivalently, augmenting $`𝒜_c`$ by an unitary operator $`U_\alpha =e^{i2\pi \sqrt{\alpha }j_{\phi _\eta }}(\alpha 𝐑^+,`$ and $`j_{\phi _\eta }`$ being the smeared current with an appropriately chosen test function $`\phi _\eta `$, which converges to a constant for $`\eta 0`$, $`x\mathrm{}`$) to $`\overline{𝒜}_c`$. Then one discovers in $`\overline{\pi }_\beta (\overline{𝒜}_c)^{\prime \prime }`$ anyonic modes which satisfy Heisenberg’s Urgleichung in a distributional sense
$$\frac{1}{i}\frac{}{x}\mathrm{\Psi }_\alpha (x)=\pi \sqrt{\alpha }[j(x),\mathrm{\Psi }_\alpha (x)]_\alpha .$$
Thus we do not introduce the anyons as fields with deformed commutation relations as e.g. in but we obtain these relations.
The important point here is that the Hilbert space $`\overline{}_\beta `$ assumes a sectorial structure, being for fixed $`\alpha `$ a countable orthogonal sum of sectors with $`n`$ particles created by $`U_\alpha `$
$$\overline{}_\beta =\overline{}_\beta ^n,\overline{}_\beta ^n=𝒜_c\underset{i=1}{\overset{n}{}}\mathrm{\Psi }_\eta (x_i)|\mathrm{\Omega }$$
(1.4)
Here the following is to be observed:
1. $`\mathrm{\Psi }_\alpha `$ as in (1.3) has an infrared and an ultraviolet problem. The infrared divergence actually shows that admitting the (smeared) step function as a test function, one creates new elements in the field algebra which lead to orthogonal sectors in a larger Hilbert space, Eq.(1.4). The ultraviolet divergence does not lead out of $`\overline{\pi }_\beta `$ if we smear $`j(y)`$ over a region of size $`\eta `$ to get $`\mathrm{\Psi }_{\alpha ,\eta }`$ and consider the renormalized field
$$\underset{\eta 0^+}{lim}c_\alpha (\eta )𝑑xf(x)\mathrm{\Psi }_{\alpha ,\eta }(x)=\mathrm{\Psi }_\alpha (f)$$
with a suitable $`c_\alpha (\eta )`$. This limit exists in a strong sense and $`\mathrm{\Psi }(f)`$ has finite n-point functions.
2. When the statistical parameter $`\alpha `$ is an integer, two special families of such renormalized operators are distinguished: for odd $`\alpha `$’s we get fermions and for even $`\alpha `$’s — bosons. However, only the field $`\mathrm{\Psi }_1`$ turns out to be a canonical Fermi field,
$$[\mathrm{\Psi }_1^{}(x),\mathrm{\Psi }_1(x^{})]_+=\delta (xx^{}),$$
with an n-point function of the familiar determinant form. $`\mathrm{\Psi }_2`$ is a non-canonical Bose field, whose commutator is not a $`c`$-number
$$[\mathrm{\Psi }_2^{}(x),\mathrm{\Psi }_2(x^{})]=\delta ^{}(xx^{})+ij(x)\delta (xx^{}).$$
Similarly, the operator $`\mathrm{\Psi }_3`$ describes a non-canonical (unbounded) Fermi field. For $`\alpha 𝐙`$ the anyonic commutator vanishes.
The algebraic chain (1.1) means that the dressed fermions obtained for special values of $`\alpha `$ can be constructed either from bare fermions or directly from the current algebra, so in this case it cannot be decided whether fermions or bosons are more fundamental. Moreover, we shall argue that the fermions at both “ends” are actually equivalent, so the corresponding algebras do coincide. To make this statement precise, the correlation functions arrising in both cases have to be compared and this will be done subsequently.
## 2 The bosonic algebra and its states
Consider the smeared currents $`j(f)=j(x)f(x)𝑑x`$ with real-valued functions $`f(x)𝒞_0^{\mathrm{}}`$, so that neither infrared nor ultraviolet problems occur. $`f`$’s form a real pre-Hilbert space to be defined subsequently. For the bosonic algebra built by the Weyl operators $`e^{ij(f)}`$, (1.2) is replaced by the multiplication law for the unitaries
$$e^{ij(f)}e^{ij(g)}=e^{\frac{i}{2}\sigma (g,f)}e^{ij(f+g)},$$
(2.1)
with a symplectic form
$$\sigma (f,g)=\frac{1}{4\pi }_{\mathrm{}}^{\mathrm{}}𝑑x(f^{}(x)g(x)f(x)g^{}(x)).$$
(2.2)
As an integral kernel the symplectic form reads $`\sigma (xy)=\delta ^{}(xy)/2\pi `$. Therefore the bosonic algebra satisfies the requirement of a local field theory.
We shall be interested in states that are invariant under a time evolution $`\tau _t`$, in particular — KMS states. An equilibrium state of a quantum system at a finite temperature $`T=\beta ^1`$ is characterized by the KMS-condition
$$\omega _\beta \left(\tau _t(A)B\right)=\omega _\beta (B\tau _{t+i\beta }A)$$
with the time evolution $`\tau _t`$ as an automorphism of the algebra of observables analytically continued for imaginary times. In the bosonic (current) algebra $`𝒜_c`$ this role is played by the shift
$$\tau _tj(f(x))=j(f(xt)).$$
The time-invariant states are in general defined by the two-point function
$$\omega (j(f)j(g))=𝑑x𝑑yw(xy)f(x)g(y),$$
so for the $`\tau `$-KMS state at a temperature $`\beta =\pi `$ the kernel $`w(xy)`$ reads
$$w(xy)=\underset{\epsilon 0^+}{lim}\frac{1}{(2\pi )^2\mathrm{sh}^2(xyi\epsilon )}.$$
(2.3)
If we want to avoid distributions being involved into the calculations, we can consider the smeared functions
$$f_\eta (x)=f(y)h_\eta (xy)𝑑y,j(f_\eta )=j_\eta (f),$$
(2.4)
with some $`h_\eta 𝒞_0^{\mathrm{}}`$. This leads to a new symplectic form
$$\sigma (f_\eta ,g_\eta )\sigma _\eta (f,g),$$
(2.5)
$$\begin{array}{ccc}\hfill \sigma _\eta (xy)& =& \frac{1}{2\pi }\delta ^{}(zz^{})h_\eta (zx)h_\eta (z^{}y)𝑑z𝑑z^{}\hfill \\ & =& \frac{1}{2\pi }h_\eta (zx)h_\eta ^{}(zy)𝑑z.\hfill \end{array}$$
In this case $`\epsilon `$ of (2.3) need not be finite but just indicates which distribution is understood. We shall henceforth omit keep writing $`lim_{\epsilon 0}`$.
$`𝒞_0^{\mathrm{}}`$ functions and functions with steps are mapped by the smearing into $`𝒞_0^{\mathrm{}}`$ functions. Therefore the bosonic algebra remains unchanged. The essential point is that the fields $`j_\eta (x)`$ are not local any more in the sense that they do not commute if they are spatially separated. On the other hand, we gain that they become (unbounded) operators and not only operator-valued distributions. Their time evolution is given by
$$\tau _tj_\eta (x)=j_\eta (x+t).$$
Again we can express the state by
$$\omega (j(f_\eta )j(g_\eta ))=\omega (j_\eta (f)j_\eta (g)):=f|g_\eta $$
where the corresponding integral kernel
$$w_\eta (xy)=w(zz^{})h_\eta (xz)h_\eta (yz^{})𝑑z𝑑z^{}$$
satisfies
$$w_\eta (xy)w_\eta (yx)=\sigma _\eta (xy).$$
The expectation of the Weyl operators is given by
$$\omega (e^{ij_\eta (f)})=e^{\frac{1}{2}f|f_\eta }.$$
(2.6)
Eqs.(2.1),(2.6) imply
$$\omega (\underset{k}{}e^{ij_\eta (f_k)})=\mathrm{exp}\left\{\frac{1}{2}\left[\underset{k}{}f_k|f_k_\eta +2\underset{k<m}{}f_k|f_m_\eta \right]\right\}.$$
(2.7)
If $`lim_{\eta 0}h_\eta (x)=\delta (x)`$, then $`lim_{\eta 0}f|g_\eta =f|g`$ and the operator $`j_\eta (x)`$ approaches the distribution $`j(x)`$.
With the scalar product $`f|g`$ the one-particle Hilbert space $`h`$ is defined as the closure of the pre-Hilbert space $`𝒞_0^{\mathrm{}}`$, that determines the bosonic von Neumann algebra $`\overline{\pi }_\beta (𝒜)^{\prime \prime }=:𝒜_B`$. Of course, $`f|g_\eta `$ defines the same Hilbert space, but infrared and ultraviolet divergencies are kept apart.
## 3 The extension to anyons
The field operators we wish to construct in the spirit of Mandelstam are of the form $`\mathrm{\Psi }_\alpha (x)e^{i2\pi \sqrt{\alpha }_{\mathrm{}}^xj(y)𝑑y}`$, so they are defined by the function $`F_x(y)=2\pi \mathrm{\Theta }(xy)`$, whereas the functions from $`h`$ have to vanish for $`x\pm \mathrm{}`$. The symplectic form (2.2) is defined for functions that tend to a constant, however they cannot be reached as limits of functions from $`h`$. For instance,
$$F_{x,\delta }(y)=\mathrm{\Theta }(xy)\mathrm{\Theta }(xy\delta ),F_{x,\delta ,\eta }𝒞_0^{\mathrm{}}.$$
(3.1)
does not work, since $`\sigma (F_{x,\delta },F_{x^{},\delta ^{}})`$ depends on the order in which the limits $`\delta ,\delta ^{}\mathrm{}`$ are taken and only for $`\delta =\delta ^{}\mathrm{}`$ we get the desired result $`i\mathrm{sgn}(xx^{})`$. Since this appears in the $`c`$-number part, in no representation can $`j(F_{x,\delta })`$ converge strongly. Nevertheless, for functions with the same (nontrivial) asymptotics at, say, $`x\mathrm{}`$ and whose difference $`h`$ one can succeed in getting the expectation values as limits.
The desired extension of the algebra $`𝒜_B`$ can be achieved in two equivalent ways. The one-particle Hilbert space can be enlarged allowing also for (appropriately smeared) step functions, e.g. $`F_x^\alpha =\sqrt{\alpha }F_x`$, such that only $`F^{}𝒞_0^{\mathrm{}}`$. For the enlarged algebra the symplectic form (2.2) is kept. If the localization of the current is given by support $`F^{}`$, then fields that are localized in different regions satisfy the following exchange relation ($`\alpha `$-commutator)
$$e^{ij(F^\alpha )}e^{ij(G^\alpha )}=e^{i\pi \alpha }e^{ij(G^\alpha )}e^{ij(F^\alpha )},\text{supp }F^{}<\text{supp }G^{}$$
(3.2)
$$\underset{x\mathrm{}}{lim}F(x)=\underset{x\mathrm{}}{lim}G(x)=2\pi ,$$
$$\underset{x\mathrm{}}{lim}F(x)=\underset{x\mathrm{}}{lim}G(x)=0.$$
The same algebra can be obtained by the automorphism $`\gamma _F^{}`$ of the initial bosonic algebra $`𝒜_c`$
$$\gamma _F^{}e^{ij(f)}=e^{i{\scriptscriptstyle F^{}(x)f(x)𝑑x}}e^{ij(f)}.$$
This automorphism is not inner but allows the construction of a crossed product in which $`\gamma _F^{}`$ is implemented by $`e^{ij(F)}`$ (compare ). A given quasifree automorphism $`\rho `$ on the initial algebra can be uniquely extended to an automorphism on the enlarged algebra provided $`\gamma _F^{}\rho \gamma _F^{}^1\rho ^1`$ is an inner automorphism of $`𝒜_c`$. For instance, $`\tau _t`$ is extendable and acts again as a shift on $`j(F)`$. Also a state on the bosonic algebra can be extended according to
$$\omega (e^{ij(F)})=0\text{if }\underset{x\mathrm{}}{lim}F(x)0.$$
For higher products we combine
$$\omega (e^{ij(f)}e^{ij(F)}e^{ij(G)}e^{ij(g)})=\omega (e^{ij(f)}e^{ij(FG)}e^{ij(g)})e^{i\sigma (F,G)/2}.$$
The operator $`e^{ij(FG)}𝒜_B`$ and the expectation value is well defined.
As already mentioned, whereas $`e^{ij(F)}`$ itself cannot be obtained as a limit of operators in $`𝒜_B`$, the situation with the expectation value is different. We concentrate on the currents $`j_\eta (F_x)`$, where $`\eta `$ indicates that we have smeared as in (2.4),(2.5) and choose the test function (3.1). With this, we get for $`\delta >|x\overline{x}|+|\text{supp }h_\eta |`$
$`e^{ij_\eta (F_{x,\delta }^\alpha )}e^{ij_\eta (F_{\overline{x},\delta }^\alpha )}=e^{ij_\eta (F_{x,\delta }^\alpha F_{\overline{x},\delta }^\alpha )+i\alpha \sigma _\eta (F_{x,\delta },F_{\overline{x},\delta })/2}`$
$`=`$ $`e^{ij_\eta (F_x^\alpha F_{\overline{x}}^\alpha )}e^{ij_\eta (F_{x\delta }^\alpha F_{\overline{x}\delta }^\alpha )}e^{i\alpha \sigma _\eta (F_x,F_{\overline{x}})},`$
noting that for sufficiently large $`\delta `$,
$$\sigma _\eta (F_xF_{\overline{x}},F_{x\delta }F_{\overline{x}\delta })=0$$
$$\sigma _\eta (F_x,F_{\overline{x}})=\sigma _\eta (F_{x,\delta },F_{\overline{x},\delta })/2.$$
If we take into account that $`\omega `$ is translation invariant and clustering
$`\underset{\delta \mathrm{}}{lim}\omega (e^{ij_\eta (F_x^\alpha F_{\overline{x}}^\alpha )}e^{ij_\eta (F_{x\delta }^\alpha F_{\overline{x}\delta }^\alpha )})`$
$`=`$ $`\omega (e^{ij_\eta (F_x^\alpha F_{\overline{x}}^\alpha )})\omega (e^{ij_\eta (F_{x\delta }^\alpha F_{\overline{x}\delta }^\alpha )})=e^{\alpha F_xF_{\overline{x}}|F_xF_{\overline{x}}_\eta },`$
we can conclude
$$\omega (e^{ij_\eta (f)}e^{ij_\eta (F_x^\alpha )}e^{ij_\eta (F_{\overline{x}}^\alpha )}e^{ij_\eta (g)})=\underset{\delta \mathrm{}}{lim}\omega (e^{ij_\eta (f)}e^{ij_\eta (F_{x,\delta }^\alpha )/\sqrt{2}}e^{ij_\eta (F_{\overline{x},\delta }^\alpha )/\sqrt{2}}e^{ij_\eta (g)}).$$
Now we can apply (2.7) and obtain
$`\omega ({\displaystyle \underset{k}{}}e^{ij_\eta (F_{x_k}^\alpha )})=\underset{\delta \mathrm{}}{lim}\omega ({\displaystyle \underset{k}{}}e^{ij_\eta (F_{x_k,\delta }^\alpha )/\sqrt{2}})`$
$`=`$ $`\underset{\delta \mathrm{}}{lim}\mathrm{exp}\left\{{\displaystyle \frac{\alpha }{2}}\left[{\displaystyle \frac{1}{2}}{\displaystyle \underset{k}{}}F_{x_k,\delta }|F_{x_k,\delta }_\eta +{\displaystyle \underset{k<m}{}}F_{x_k,\delta }|F_{x_m,\delta }_\eta \right]\right\}.`$
We can evaluate the scalar products involved in (3.3):
$`F_{x_k,\delta }|F_{x_m,\delta }_\eta ={\displaystyle 𝑑z\underset{\delta +x_k}{\overset{x_k}{}}𝑑y\underset{\delta +x_m}{\overset{x_m}{}}𝑑y^{}\frac{\widehat{h}_\eta (z)}{(2\pi )^2\mathrm{sh}^2(yy^{}zi\epsilon )}}`$
$`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle \widehat{h}_\eta (z)𝑑z\mathrm{ln}\frac{\mathrm{sh}^2(x_kx_mzi\epsilon )}{\mathrm{sh}(x_kx_m+\delta zi\epsilon )\mathrm{sh}(x_kx_m\delta zi\epsilon )}}.`$
To get something finite for $`\delta \mathrm{}`$ we have to take operators of the form
$$e^{\pm ij(F_{x_k}^\alpha )}:=e^{ij(\overline{F}_{x_k}^\alpha )},$$
where $`\overline{F}_{x_k}^\alpha =s_kF_{x_k}^\alpha `$, $`s_k=\pm 1`$. The individual expressions diverge with $`\delta \mathrm{}`$. If therefore the anyon contributions do not neutralize, i.e. if $`_ks_k0`$, the expectation value (3.3) vanishes. On the other hand, if $`_ks_k=0`$, we have as many positive contributions as negative ones, this means for the products $`2r+2r(r1)`$ positive contributions and $`2r^2`$ negative ones. Those that contain $`\delta `$ can be combined in pairs to
$$\underset{\delta \mathrm{}}{lim}\left[\mathrm{ln}\frac{\mathrm{sh}(x+\delta i\epsilon )}{\mathrm{sh}(y+\delta i\epsilon )}\mathrm{ln}\frac{\mathrm{sh}(x\delta i\epsilon )}{\mathrm{sh}(y\delta i\epsilon )}\right]=\underset{\delta \mathrm{}}{lim}\left(\mathrm{ln}\frac{e^{x+\delta }}{e^{y+\delta }}\mathrm{ln}\frac{e^{\delta x}}{e^{\delta y}}\right)=0.$$
Therefore we remain in the limit $`\delta \mathrm{}`$ with
$$\omega (\underset{k}{}e^{ij_\eta (\overline{F}_{x_k}^\alpha )})=[c^2(\eta )]^{n\alpha }\mathrm{exp}\left\{𝑑z\widehat{h}_\eta (z)\alpha \mathrm{ln}\left(\underset{k<m}{}s_ks_m\mathrm{sh}(x_kx_mzi\epsilon )\right)\right\},$$
$$c^1(\eta )=\mathrm{exp}\left[\frac{1}{2}\widehat{h}_\eta (z)\mathrm{ln}\mathrm{sh}(zi\epsilon )𝑑z\right].$$
(3.5)
the latter taking care for the ultraviolet divergence in $`\mathrm{\Psi }_\alpha `$.
In terms of $`\mathrm{\Psi }`$’s this means that the expectation value of a product of $`\mathrm{\Psi }`$’s and $`\mathrm{\Psi }^{}`$’s is different from zero only if there are as many $`\mathrm{\Psi }`$’s as $`\mathrm{\Psi }^{}`$’s, equivalently — if the “total” statistic parameter of creation operators equals the one of annihilation operators, for instance $`\mathrm{\Psi }_1^{}\mathrm{\Psi }_1^{}\mathrm{\Psi }_40`$, or otherwise that they lead to orthogonal sectors of the enlarged Hilbert space.
Performing the limit $`\eta 0`$, where $`\widehat{h}_\eta (x)\delta (x)`$ we achieve our aim — construction of the local anyonic field operators. Note that these are now strong limits of the anyonic Weyl operators as discussed in .
The divergence for $`\eta 0`$ remains and determines the necessary renormalization of the operators $`\mathrm{\Psi }_\alpha `$
$$\text{s-}\underset{\eta 0}{lim}f(x)c(\eta )e^{ij_\eta (F_x^\alpha )}𝑑x=f(x)\mathrm{\Psi }_\alpha (x)𝑑x.$$
(3.6)
From (3.5), (3.6) we obtain for the $`n`$-point function
$`\omega \left(\mathrm{\Psi }_\alpha ^{}(x_1)\mathrm{}\mathrm{\Psi }_\alpha ^{}(x_n)\mathrm{\Psi }_\alpha (y_n)\mathrm{}\mathrm{\Psi }_\alpha (y_1)\right)`$ (3.7)
$`=`$ $`\underset{\eta 0}{lim}[c^2(\eta )]^{n\alpha }\omega \left(e^{ij_\eta (F_{x_1}^\alpha )}\mathrm{}e^{ij_\eta (F_{x_n}^\alpha )}e^{ij_\eta (F_{y_n}^\alpha )}\mathrm{}e^{ij_\eta (F_{y_1}^\alpha )}\right)`$
$`=`$ $`{\displaystyle \frac{{\displaystyle \underset{k>l}{}}(\mathrm{sh}(x_kx_li\epsilon ))^\alpha {\displaystyle \underset{k>l}{}}(\mathrm{sh}(y_ky_li\epsilon ))^\alpha }{{\displaystyle \underset{k,l}{}}\left(2\pi i\mathrm{sh}(x_ky_li\epsilon )\right)^\alpha }}.`$
The exact exchange relations are hidden in the factor $`i^\alpha `$ in (3.7) and we shall return to their detailed analysis later on.
For all $`\alpha `$’s the two-point function (for $`x>x^{}`$ and $`\beta =\pi `$)
$$\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (x^{})_\beta =\mathrm{\Psi }_\alpha (x)\mathrm{\Psi }_\alpha ^{}(x^{})_\beta =\left(\frac{i}{2\pi \mathrm{sh}(xx^{})}\right)^\alpha =:S_\alpha (xx^{})$$
has the desired properties
Hermiticity:
$$S_\alpha ^{}(x)=S_\alpha (x)\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (x^{})_\beta ^{}=\mathrm{\Psi }_\alpha ^{}(x^{})\mathrm{\Psi }_\alpha (x)_\beta ;$$
$`\alpha `$-commutativity:
$$S_\alpha (x)=e^{i\pi \alpha }S_\alpha (x)\mathrm{\Psi }_\alpha (x^{})\mathrm{\Psi }_\alpha ^{}(x)_\beta =e^{i\pi \alpha }\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (x^{})_\beta ;$$
KMS-property:
$$S_\alpha (x)=S_\alpha (x+i\pi )\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (x^{})_\beta =\mathrm{\Psi }_\alpha (x^{})\mathrm{\Psi }_\alpha ^{}(x+i\pi )_\beta .$$
(3.8)
For $`\alpha =2`$ and an arbitrary temperature $`\beta ^1`$ we get like for the $`j`$’s
$$\mathrm{\Psi }_2^{}(x)\mathrm{\Psi }_2(x^{})_\beta =\frac{1}{\left(2\beta \mathrm{sh}\frac{\pi (xx^{}i\epsilon )}{\beta }\right)^2},$$
similarly, for $`\alpha =3`$ we get a different kind of fermions
$$\mathrm{\Psi }_3^{}(x)\mathrm{\Psi }_3(x^{})_\beta =\frac{i}{\left(2\beta \mathrm{sh}\frac{\pi (xx^{}i\epsilon )}{\beta }\right)^3}.$$
These fields, though locally (anti)commuting, are not canonical and this becomes transparent by analyzing temperature dependence and operator structure of their exchange relations.
## 4 The $`\alpha `$-commutator
As a direct consequence of the Weyl relations (3.2) for the anyonic Weyl operators (1.3) and with (3.6) in mind, it follows that
$$\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (y)e^{i\pi \alpha \mathrm{sgn}(xy)}\mathrm{\Psi }_\alpha (y)\mathrm{\Psi }_\alpha ^{}(x)=0\text{for }xy.$$
It remains to calculate the distribution that emerges by bringing the field arguments together.
Evidently, the case $`\alpha 𝐙`$ plays a special role since we then deal with Fermi- or Bose- commutation relations. We first concentrate on $`\alpha =1`$. Then we observe
$$\omega (\mathrm{\Psi }^{}(x_1)\mathrm{}\mathrm{\Psi }(y_n))=\mathrm{Det}\omega (\mathrm{\Psi }^{}(x_i)\mathrm{\Psi }(y_k))=\mathrm{Det}w(x_iy_k).$$
Proof: From Cauchy’s determinant formula
$$\mathrm{Det}(x_iy_k)^1=\frac{{\displaystyle \underset{i>k}{}}(x_ix_k){\displaystyle \underset{i>k}{}}(y_iy_k)}{{\displaystyle \underset{i,k}{}}(x_iy_k)}$$
one gets
$$\mathrm{Det}\frac{\sqrt{x_iy_j}}{(x_iy_j)}=\sqrt{x_i}\sqrt{y_j}\frac{{\displaystyle \underset{i>j}{}}(x_ix_j){\displaystyle \underset{i>j}{}}(y_iy_j)}{{\displaystyle \underset{i,j}{}}(x_iy_j)}=\mathrm{Det}\frac{1}{\sqrt{x_i/y_j}\sqrt{y_j/x_i}}.$$
This corresponds by replacing $`\sqrt{x}`$ by $`e^x`$, $`\sqrt{y}`$ by $`e^y`$, to
$$\mathrm{Det}\frac{1}{\mathrm{sh}(x_iy_ki\epsilon )}=\frac{{\displaystyle \underset{i>k}{}}\mathrm{sh}(x_ix_ki\epsilon ){\displaystyle \underset{i>k}{}}\mathrm{sh}(y_iy_ki\epsilon )}{{\displaystyle \underset{i,k}{}}\mathrm{sh}(x_iy_ki\epsilon )}.$$
(4.1)
Theorem
For $`\alpha =1`$ the renormalized fields $`\mathrm{\Psi }_1^{}(x),\mathrm{\Psi }_1(x)`$ canonically anticommute
$$\mathrm{\Psi }_1^{}(x)\mathrm{\Psi }_1(y)+\mathrm{\Psi }_1(y)\mathrm{\Psi }_1^{}(x)=\delta (xy).$$
The state over the algebra of fermions is the quasifree state given by the two-point function
$$\omega (\mathrm{\Psi }_1^{}(x)\mathrm{\Psi }_1(y))=\frac{i}{2\pi \mathrm{sh}(xyi\epsilon )}.$$
It satisfies the KMS condition with respect to the shift for temperature $`\beta =\pi `$ (comp. (3.8)). For arbitrary temperature by scaling arguments it follows
$$\omega _\beta (\mathrm{\Psi }_1^{}(x)\mathrm{\Psi }_1(y))=\frac{i}{2\beta \mathrm{sh}\frac{\pi (xx^{}i\epsilon )}{\beta }},$$
with the same commutation relations.
For arbitrary $`\alpha `$ we have to analyse in detail the expectation
$`\omega _\beta (\mathrm{\Psi }_\alpha ^{}(x_1)\mathrm{}\mathrm{\Psi }_\alpha (y_n))`$ $`=`$ $`{\displaystyle \frac{{\displaystyle \underset{k>l}{}}[\mathrm{sh}(x_kx_li\epsilon )]^\alpha {\displaystyle \underset{k>l}{}}[\mathrm{sh}(y_ky_li\epsilon )]^\alpha }{{\displaystyle \underset{i,j}{}}\mathrm{sh}(x_iy_ji\epsilon )^\alpha }}`$ (4.2)
$`=`$ $`\left(\mathrm{Det}\mathrm{sh}^1(x_iy_ji\epsilon )\right)^\alpha .`$
Evidently, for $`\alpha 1`$ the state is determined again by the two-point function but not in a way that corresponds to a truncation. In order to deduce the commutation relations it is preferable to evaluate
$`{\displaystyle \omega \left(e^{ij(f_1)}\left[\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (y)e^{i\pi \alpha \mathrm{sgn}(xy)}\mathrm{\Psi }_\alpha (y)\mathrm{\Psi }_\alpha ^{}(x)\right]e^{ij(f_2)}\right)g_1(x)g_2(y)𝑑x𝑑y}`$
$`=`$ $`{\displaystyle \omega (e^{ij(f_1)}e^{ij(f_2)})\left(\frac{1}{\mathrm{sh}^\alpha (xyi\epsilon )}\frac{(1)^\alpha }{\mathrm{sh}^\alpha (yxi\epsilon )}\right)(i/2\pi )^\alpha }`$
$`\times `$ $`e^{2\pi \sqrt{\alpha }[f_1|\mathrm{\Theta }(x)f_1|\mathrm{\Theta }(y)\mathrm{\Theta }(x)|f_2+\mathrm{\Theta }(y)|f_2]}g_1(x)g_2(y)dxdy.`$
Therefore the distribution to be considered is
$$\underset{\epsilon 0^+}{lim}\left(\frac{1}{\mathrm{sh}^\alpha (xi\epsilon )}\frac{1}{\mathrm{sh}^\alpha (x+i\epsilon )}\right),\alpha 𝐙^+.$$
This expression tends to 0 for $`x0`$. As for the singularity at $`x=0`$, after a partial integration one realizes that it suffices to assume $`0<\alpha <1`$ and to evaluate the remaining integral around the origin, so with $`\mathrm{sh}(xi\epsilon )^\alpha (xi\epsilon )^\alpha `$
$`\underset{\epsilon 0}{lim}{\displaystyle \underset{\delta }{\overset{\delta }{}}}\left({\displaystyle \frac{1}{(xi\epsilon )^\alpha }}{\displaystyle \frac{1}{(x+i\epsilon )^\alpha }}\right)f(x)𝑑x`$
$`=`$ $`\underset{\epsilon 0}{lim}{\displaystyle \underset{\delta /\epsilon }{\overset{\delta /\epsilon }{}}}\left({\displaystyle \frac{1}{(yi)^\alpha }}{\displaystyle \frac{1}{(y+i)^\alpha }}\right)f(\epsilon y)\epsilon ^{1\alpha }𝑑y=0`$
Therefore for anyonic fields with noninteger statistic parameter the $`\alpha `$-commutator vanishes
$$\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (y)e^{i\pi \alpha \mathrm{sgn}(xy)}\mathrm{\Psi }_\alpha (y)\mathrm{\Psi }_\alpha ^{}(x)=0\alpha 𝐙^+.$$
For $`\alpha `$ integer we can perform again partial integration. This gives us for $`\alpha =2`$
$$\underset{\epsilon 0}{lim}\left(\frac{1}{\mathrm{sh}^2(xi\epsilon )}\frac{1}{\mathrm{sh}^2(x+i\epsilon )}\right)=2\pi i\delta ^{}(x).$$
Therefore
$`{\displaystyle \omega \left(e^{ij(f_1)}[\mathrm{\Psi }_2^{}(x)\mathrm{\Psi }_2(y)\mathrm{\Psi }_2(y)\mathrm{\Psi }_2^{}(x)]e^{ij(f_2)}g_1(x)g_2(y)\right)𝑑x𝑑y}`$
$`=`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle \delta ^{}(xy)\omega (e^{ij(f_1)}e^{ij(f_2)})e^{2\pi \sqrt{2}[f_1|\mathrm{\Theta }(x)f_1|\mathrm{\Theta }(y)\mathrm{\Theta }(x)|f_2+\mathrm{\Theta }(y)|f_2]}}`$
$`\times g_1(x)g_2(y)dxdy+i\sqrt{2}{\displaystyle g_1(x)g_2(x)\omega (e^{ij(f_1)}e^{ij(f_2)})𝑑x}`$
or
$$[\mathrm{\Psi }_2^{}(x),\mathrm{\Psi }_2(y)]=\frac{i}{2\pi }\delta ^{}(xy)+i\sqrt{2}\delta (xy)j(x).$$
This follows from comparing the contributions in the expectation value
$$\delta ^{}(xy)e^{[f_1|\mathrm{\Theta }(x)f_1|\mathrm{\Theta }(y)\mathrm{\Theta }(x)|f_2+\mathrm{\Theta }(y)|f_2]}\delta (xy)[f_1|\delta (x)\delta (x)|f_2]$$
with
$$\omega (e^{ij(f_1)}j(x)e^{ij(f_2)})=\frac{d}{d\gamma }\omega (e^{ij(f_1)}e^{ij(f_2)})e^{\gamma [f_1|\delta (x)\delta (x)|f_2]}.$$
Thus the bosonic fields of the first level do not satisfy canonical commutation relations. Though the two-point functions look similar, for the initial bosonic algebra this was a two-point function of currents whereas now it is a two-point function of fields, that are not invariant under gauge automorphisms but are adjoint of each other.
For $`\alpha =3`$ again fermionic fields are obtained. Partial integration yields (we do not fix the temperature)
$`\omega \left(e^{ij(f_1)}[\mathrm{\Psi }_3^{}(x)\mathrm{\Psi }_3(y)+\mathrm{\Psi }_3(y)\mathrm{\Psi }_3^{}(x)]e^{ij(f_2)}\right)`$
$`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left(\delta ^{\prime \prime }(xy){\displaystyle \frac{\pi ^2}{\beta ^2}}\delta (xy)\right)e^{2\pi \sqrt{3}[f_1|\mathrm{\Theta }(x)f_1|\mathrm{\Theta }(y)\mathrm{\Theta }(x)|f_2+\mathrm{\Theta }(y)|f_2]}.`$
Here the $`\delta ^{\prime \prime }(xy)`$ splits again in a distribution acting only on the smearing function of the fermion fields and additional terms, so that we can write the commutator as an operator-valued distribution
$$[\mathrm{\Psi }_3^{}(x),\mathrm{\Psi }_3(y)]_+=\frac{1}{8\pi ^2}[\delta ^{\prime \prime }(xy)2\pi \sqrt{3}\delta ^{}(xy)j(x)+12\pi ^2\delta (xy)j^2(x)]+\frac{1}{8\beta ^2}\delta (xy)$$
or more precisely,
$`[\mathrm{\Psi }_3^{}(f),\mathrm{\Psi }_3(g)]_+`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}{\displaystyle }dx\{{\displaystyle \frac{d^2}{dx^2}}f(x)g(x)2\pi \sqrt{3}j(x)\left({\displaystyle \frac{d}{dx}}f(x)g(x)\right)`$
$`+`$ $`12\pi ^2f(x)g(x)j^2(x)(\pi /\beta )^2f(x)g(x)\}.`$
Thus we are faced with a new feature: the commutation relations become temperature dependent, though they remain local. The operators $`\mathrm{\Psi }_3(f)`$ are obviously not bounded. The fermions of the third level form a subalgebra of the whole fermionic algebra (the fermions of the first level). This demonstrates that representations corresponding to different temperatures are inequivalent not only globally, but already also locally.
Back to the current algebra
The current algebra can be considered as the subalgebra of any anyonic algebra $`\{\mathrm{\Psi }_\alpha (x)\}^{\prime \prime }`$ that is invariant under the gauge automorphism $`\gamma _s\mathrm{\Psi }_\alpha (x)=e^{i\alpha s}\mathrm{\Psi }_\alpha (x)`$. Furthermore, the local fields $`j(x)`$ can be directly constructed out of the anyonic fields $`\mathrm{\Psi }_\alpha (x)`$. Starting with the fermions $`\mathrm{\Psi }_1(x)`$ this procedure is well known . In fact, $`\mathrm{\Psi }^{}(x)\mathrm{\Psi }(x)`$ is at best a quadratic form and by point splitting it becomes the operator-valued distribution $`\mathrm{\Psi }^{}(x+\epsilon )\mathrm{\Psi }(x\epsilon )`$. We can relate this to the bounded operator $`e^{ij(f_{x,\epsilon })}`$ with $`f_{x,\epsilon }(y)=1`$ for $`x\epsilon <y<x+\epsilon `$, otherwise $`0`$. But then we produce by the sharp edge of $`f`$ an ultraviolet problem with diverging expectation values. To start with a well defined expression with a zero expectation value we consider
$$j(x)=\underset{\genfrac{}{}{0pt}{}{f_{x,\epsilon }(x)=1}{{\scriptscriptstyle f_{x,\epsilon }(y)𝑑y{\scriptscriptstyle \genfrac{}{}{0pt}{}{}{\epsilon 0}}0}}}{lim}\frac{e^{i\alpha j(f_{x,\epsilon })}e^{i\alpha j(f_{x,\epsilon })}}{2i\alpha f_{x,\epsilon }(y)𝑑y}.$$
(4.3)
Now we choose
$$f_{x,\epsilon }(y):=\{\begin{array}{cc}0& \text{ for }y<x\epsilon ,y>x+\epsilon ,\hfill \\ 1\left|\frac{xy}{\epsilon }\right|& \text{ for }x\epsilon <y<x+\epsilon .\hfill \end{array}$$
This can be interpreted as a smearing in the sense of (2.4) so that
$$e^{i\alpha j(f_{x,\epsilon })}=\mathrm{\Psi }_{\alpha ,\epsilon }^{}(x+\epsilon )\mathrm{\Psi }_{\alpha ,\epsilon }(x\epsilon ),$$
whereas
$$e^{i\alpha j(f_{x,\epsilon })}=\mathrm{\Psi }_{\alpha ,\epsilon }^{}(x\epsilon )\mathrm{\Psi }_{\alpha ,\epsilon }(x+\epsilon )$$
and $`f_{x,\epsilon }(y)𝑑y`$ corresponds to the renormalization factor needed to pass from the anyonic Weyl operators to the anyonic fields. Of course, $`j(x)`$ is not an operator but has to be smeared by a suitable $`f`$. Then (4.3) becomes a strong resolvent limit. In a sense we thus provide a replacement of the familiar expression
$$j(x)=:\mathrm{\Psi }^{}(x)\mathrm{\Psi }(x):$$
that has been used to construct a bosonic current algebra with Schwinger commutation relations out of the fermions. In fact, since we obtain the anyonic fields as limit of anyonic Weyl operators that are not quadratic in the Fermi field, the current emerges out of a kind of a Dirac sea filling. This procedure is state dependent and the states under consideration are not locally normal, therefore the structure of the Dirac sea, the higher-level fields exhibit a nontrivial temperature dependence. Nevertheless, the structure of the currents defined by (4.3) is $`\alpha `$ and $`\beta `$ independent.
## 5 The Luttinger model
The Luttinger model has been designed to describe a one-dimensional interacting electron system. The model spectrum in the ground state consists of plasmons with a well-defined energy . The essential input in these considerations is the fact that the free part of the Hamiltonian makes a Bogoliubov transformation necessary, so that the Bose operators, which replace (appropriately) the products of fermionic creation and annihilation operators $`a^{}(x)a(x)`$, satisfy the Schwinger commutation relations (1.2).
Several interaction potentials can be considered
1. $`H=H_0+\lambda j_1(x)v(xy)j_2(y)`$;
2. $`H=H_0+\lambda (j_1(x)+j_2(x))v(xy)(j_1(y)+y_2(y))`$;
3. $`H=H_0+\lambda j_1(x)v(xy)j_1(y)`$.
With $`\rho _i(p)=j_i(x)e^{ipx}𝑑p`$, (1.2) leads to
1. $`(\ddot{\rho }_1(p)+\ddot{\rho }_2(p))=p^2(1\lambda ^2\stackrel{~}{v}^2(p))(\rho _1(p)+\rho _2(p))`$;
2. $`(\ddot{\rho }_1(p)\pm \ddot{\rho }_2(p))=(p^2+2\lambda p\stackrel{~}{v}(p))(\rho _1(p)\pm \rho _2(p))`$;
3. $`\dot{\rho }(p)=p(1+\lambda \stackrel{~}{v}(p))\rho (p)`$.
In all cases we obtain a quasifree evolution for the currents. There exist the corresponding KMS-states, provided they are positiv definite (which might bring in some restrictions on the coupling constant). With $`\tau _tj(f)=j(e^{i\epsilon (p)t}f)`$,
$`\omega (j(f)j(g))`$ $`=`$ $`f|g_{\beta ,v}={\displaystyle f(x)g(y)w(xy)_{\beta ,v}}`$
$`=`$ $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{p}{1e^{\beta \epsilon (p)}}}\stackrel{~}{f}(p)\stackrel{~}{g}(p)𝑑p`$
$`=`$ $`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{p}{1e^{\beta \epsilon (p)}}}(\stackrel{~}{f}(p)\stackrel{~}{g}(p)+\stackrel{~}{f}(p)\stackrel{~}{g}(p)e^{\beta \epsilon (p)}).`$
As for the free evolution, the state can be extended to the anyonic Weyl operators as in (3.3), resulting in a KMS-state for the extended algebra. For sufficiently smooth short-range potentials
$$\underset{p\mathrm{}}{lim}\frac{\epsilon (p)}{p}=1,\underset{x\pm \mathrm{}}{lim}\frac{w(x)_{\beta ,v}}{w(x)_\beta }=1.$$
Therefore the infrared behaviour in $`F_{x,\delta }|F_{\overline{x},\delta }`$ for $`\delta \mathrm{}`$ can be controled as for the free case. Since $`\stackrel{~}{w}_\beta (p)\stackrel{~}{w}_{\beta ,v}(p)𝒮`$ also the ultraviolet behaviour is unchanged. It is essential again that the bosonic one-particle Hamiltonian is positiv definite and does not produce an additional singularity at some $`p0`$. Also here the $`n`$-point functions become products of two-point functions (4.2). However, already for the $`\alpha =1`$ fermions $`\epsilon (p)p`$ and they do not allow the decomposition (4.1), so that the state over the fermionic algebra is not quasifree any more. Nevertheless, the resulting fermions satisfy canonical anticommutation relations because the leading singularity in $`w_{\beta ,v}(x)`$ coincides with that of $`w_\beta (x)`$ and determines the commutation relations
$$\omega _{\beta ,v}\left(e^{ij(f_1)}[\mathrm{\Psi }^{}(x)\mathrm{\Psi }(y)+\mathrm{\Psi }(y)\mathrm{\Psi }^{}(x)]e^{ij(f_2)}\right)=\omega _{\beta ,v}(e^{ij(f_1)}e^{ij(f_2)})\delta (xy).$$
This shows that $`[.,.]_+`$ is a $`c`$-number and that $`\omega _{\beta ,v}`$ is in fact a KMS-state for fermions with a point pair interaction, provided we accept that it results from a regularization procedure involving both delocalization and filling of the Dirac sea. Due to the fact that only the leading singularity remains unchanged, one should be aware of the indispensable dependence of the exchange relations of the anyonic fields (especially of the higher-level Bose- and Fermi- fields) both on the temperature and interaction.
## 6 Conclusions
We have represented the two-dimensional chiral anyonic field $`\mathrm{\Psi }_\alpha `$ as an operator in a Hilbert space and have studied its exchange relations and its thermal correlation functions. The latter are still determined by the two-point function
$$\omega _\beta (\mathrm{\Psi }_\alpha ^{}(x)\mathrm{\Psi }_\alpha (y))=\left(\frac{i}{2\beta \mathrm{sh}\frac{\pi (xx^{}i\epsilon )}{\beta }}\right)^\alpha ,$$
however not via simple truncation, but as
$$\omega _\beta (\mathrm{\Psi }_\alpha ^{}(x_1)\mathrm{}\mathrm{\Psi }_\alpha (y_n))=\left(\mathrm{Det}\frac{i}{2\beta \mathrm{sh}\frac{\pi (x_iy_ki\epsilon )}{\beta }}\right)^\alpha .$$
For $`\alpha `$ not integer the $`\alpha `$-commutator vanishes, for $`\alpha `$ odd the fields obtained are fermions and for $`\alpha `$ even they are bosons. However, for $`\alpha >1`$ the fields still being local, are not canonical, as it follows from the operator structure of their (anti)commutators. Moreover, already the first noncanonical fermions are not bounded and their commutation relations exhibit essential temperature dependence, so representations corresponding to different temperatures are not locally normal. Along this line of considerations, in the Luttinger model we have found an example of a not quasi-free KMS-state.
## 7 Acknowledgments
We are grateful to A. Alekseev and I. Todorov for stimulating our interest into the problem and to H.J. Borchers, E. Lieb and B. Schroer for suggestive discussions .
N.I. thanks the International Erwin Schrödinger Institute for Mathematical Physics where the research has been performed, for hospitality and financial support. This work has been supported in part also by “Fonds zur Förderung der wissenschaftlichen Forschung in Österreich” under grant P11287–PHY.
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# Untitled Document
hep-ph/0004128
A quasi unitarity bound on the radion mass in the Randall-Sundrum model
Uma Mahanta
Mehta Research Institiute
Chhatnag Road, Jhusi
Allahabad-211019, India
Abstract
In this paper we derive a quasi unitarity bound on the radion mass ($`m_\varphi `$) from the process $`hh\varphi \varphi `$. We find that at sufficiently high energies (i.e. when $`sm_h^2,m_\varphi ^2`$) the J=0 partial wave amplitude for the above process violates unitarity if $`m_\varphi ^2>m_h^2+16\pi \varphi ^2`$. Combining this result with the unitarity bound $`m_h^2<\frac{16\pi v\varphi }{3}`$ obtained from the process $`hhh\varphi `$ we get an upper bound of $`4\sqrt{\pi }\varphi [1+\frac{v}{3\varphi }]^{\frac{1}{2}}`$ on the radion mass. This bound is however valid only in the low energy effective theory where we can ignore the effects of the gravitational KK modes and the string/M theoretic excitations.
Introduction
Recently several radical proposals based on extra dimensions have been put forward to explain the large hierarchy between the Planck scale and the weak scale. Among them the Randall-Sundrum \[RS\] model is most interesting because it propses a five dimensional world with a non-factorizable metric
$$ds^2=e^{2kr_c|\theta |}\eta _{\mu \nu }dx^\mu dx^\nu r_c^2d\theta ^2.$$
$`(1)`$
Here $`r_c`$ measures the size of the extra dimension which is an $`\frac{S^1}{Z_2}`$ orbifold. k is a mass parameter of the order of the fundamental five dimensional Planck mass M. $`x^\mu `$ are the coordinates of the usual four dimensional space time. $`\pi \theta \pi `$ is the coordinate of the extra dimension with $`\theta `$ and $`\theta `$ identified. Two three-branes extending in the space time directions are placed at the orbifold fixed points $`\theta =0`$ and $`\theta =\pi `$. The 3 brane located at $`\theta =0`$ is called the hidden (Planck) brane and the 3 brane located at $`\theta =\pi `$ is called the visible brane. Randall and Sundrum showed that any field with a fundamental mass parameter $`m_0`$ gets an effective four dimensional mass given by $`m=m_0e^{k\pi r_c}`$. Thus for $`kr_c14`$ the electro-weak scale is generated from the fundamental Planck scale by the exponential warp factor of the metric.
In the original proposal of Randall and Sundrum the compactification radius $`r_c`$ was assumed to be determined by the vacuum expectation value (vev) of scalar field $`T(x)`$. However the modulus field was massless and therefore its vev was not stabilized by some dynamics. Goldberger and Wise showed that by introducing a scalar field $`\chi (x,\theta )`$ in the bulk with interactions localized on the two branes it is possible to generate a potential for $`T(x)`$. They also showed that the parameters of the potential can be adjusted to yield a minimum at $`kr_c`$ 14 without any extreme fine tuning.
The Randall-Sundrum model assumes that the Standard Model \[SM\] fields are localized on the the visible brane at $`\theta =\pi `$. However the SM action is modified due to the exponential warp factor of the metric. Fluctuations of the modulus field $`T(x)`$ about its vev $`r_c`$ give rise to non-trivial couplings of the modulus field with the SM fields on the visible brane. In fact it was shown in Ref. that linear fluctuations in the radion field $`\varphi `$ couples to the SM fields on the visible brane through the Lagrangian
$$L_I=\frac{T_\mu ^\mu }{\varphi }\widehat{\varphi }.$$
$`(2)`$
Here $`T_\mu ^\mu `$ is the trace of the energy-momentum tensor of the SM fields localized on the visible brane. $`\widehat{\varphi }`$ is a small fluctuation of the radion field from its vev and is given by $`\varphi =fe^{k\pi T(x)}=\varphi +\widehat{\varphi }`$. $`\varphi `$ is the vev of $`\varphi `$ and $`f\sqrt{\frac{24M^3}{k}}`$ is a mass parameter of the order of M.
In this paper we shall derive the couplings of a light stabilized radion with the SM higgs scalar up to quadratic order in $`\widehat{\varphi }`$. We shall then use these couplings to calculate the tree level transition amplitude for the processes $`hh\varphi \varphi `$ and $`hhh\varphi `$ at very high energies i.e. when $`sm_h^2,m_\varphi ^2`$. We shall however assume that $`\sqrt{s}`$ remains much below the lowest gravitational KK mode and the string/M theoretic excitation. This will enable us to neglect the contributions of the tower of KK gravitons and stringy excitations to the transition amplitude for the process $`hh\varphi \varphi `$. An exact computation of such contributions which involves summing over a tower of states is also quite difficult. In the Goldberger-Wise stabilzation mechanism the radion mass ( $`m_\varphi `$ 10 Gev) indeed turn out to be much smaller than that of the gravitational KK modes and string excitations which usually lie in the several Tev range. Hence an effective field theory involving the SM particles and the light radion is quite realistic. By requiring the J=0 partial wave amplitude for the process $`hh\varphi \varphi `$ in the context of this low energy effective theory satisfies the unitarity constraint we obtain an upper bound on $`m_\varphi ^2m_h^2`$. Combining this result with the unitarity bound on $`m_h`$ obtained from the process $`hhh\varphi `$ we finally obtain an upper bound on the radion mass.
Radion couplings to the higgs scalar
The couplings of the radion field to the SM higgs field localized on the brane at $`\theta =\pi `$ is completely determined by general covariance. The action for the SM higgs field in the Randall-Sundrum model is given by
$$S=d^4x\sqrt{g_v}[\frac{1}{2}g_v^{\mu \nu }_\mu h_\nu hV(h)].$$
$`(3)`$
Here $`g_v^{\mu \nu }`$ is the induced metric on the visible brane. In the abscence of graviton fluctuations about the background metric it is given by $`g_v^{\mu \nu }=e^{2k\pi T(x)}\eta ^{\mu \nu }=(\frac{\varphi }{f})^2\eta ^{\mu \nu }`$ where $`\eta ^{\mu \nu }`$ is the Minkowski metric. $`\sqrt{g_v}=\sqrt{det(g_v)}=e^{4k\pi T(x)}=(\frac{\varphi }{f})^4`$. The higgs potential expanded about the classical minimum is given by
$$V(h)=\frac{\lambda }{4}(h^4+4h^3v+4h^2v^2).$$
$`(4)`$
In the above expression we have subtracted the vacuum energy from $`V(h)`$. The mass of the higgs scalar can be determined from the above potential and is given by $`m_h^2=2\lambda v^2`$. The couplings of the RS radion to the higgs field is therefore given by
$$S=d^4x[\frac{1}{2}\eta ^{\mu \nu }(\frac{\varphi }{f})^2_\mu h_\nu h\frac{\lambda }{4}(\frac{\varphi }{f})^4(h^4+4h^3v+4h^2v^2)].$$
$`(5)`$
Rescaling h and v according to $`h\frac{f}{\varphi }h`$ and $`v\frac{f}{\varphi }v`$ the action becomes
$$S=d^4x[\frac{1}{2}\eta ^{\mu \nu }(\frac{\varphi }{\varphi })^2_\mu h_\nu h\frac{\lambda }{4}(\frac{\varphi }{\varphi })^4(h^4+4h^3v+4h^2v^2)].$$
$`(6)`$
Expanding $`\varphi `$ about its vev and keeping terms only up to quadratic order in $`\widehat{\varphi }`$ we get
$$\begin{array}{ccc}\hfill S& =d^4x[\frac{1}{2}\eta ^{\mu \nu }_\mu h_\nu h\frac{\lambda }{4}(h^4+4h^3v+4h^2v^2)]\hfill & \\ & +d^4x[\eta ^{\mu \nu }_\mu h_\nu h\lambda (h^4+4h^3v+4h^2v^2)]\frac{\widehat{\varphi }}{\varphi }\hfill & \\ & +d^4x\frac{1}{2}[\eta ^{\mu \nu }_\mu h_\nu h3\lambda (h^4+4h^3v+4h^2v^2)]\frac{\widehat{\varphi }^2}{\varphi ^2}+..\hfill & (7)\hfill \end{array}$$
The dots stand for the couplings of cubic and higher order fluctuations in the radion field to the higgs scalar. These couplings will not be needed for the processes under consideration in this paper.
Unitarity bound from the process $`hh\varphi \varphi `$
We shall now derive the unitarity bound that follows from the process $`hh\varphi \varphi `$ in the context of the low energy effective field theory where we can legitimately neglect the contributions from gravitational KK modes and string excitations. At tree level the process $`hh\varphi \varphi `$ gets contributions from three distinct Feynman diagrams. Let $`M_1`$ and $`M_2`$ be transition amplitudes for t and u channel higgs exchanges. $`M_3`$ be the transition amplitude for the contact interaction diagram. It can then be shown that
$$M_1=\frac{1}{2\varphi ^2}[(tm_h^2)2(2m_h^2+m_\varphi ^2)+\frac{(2m_h^2+m_\varphi ^2)^2}{(tm_h^2)}].$$
$`(8)`$
$$M_2=\frac{1}{2\varphi ^2}[(um_h^2)2(2m_h^2+m_\varphi ^2)+\frac{(2m_h^2+m_\varphi ^2)^2}{(um_h^2)}].$$
$`(8)`$
and
$$M_3=\frac{1}{2\varphi ^2}[s+10m_h^2)].$$
$`(9)`$
At sufficiently high energies i.e. when $`sm_h^2,m_\varphi ^2`$ but $`s<m_{0g}^2,m_{0s}^2`$ ($`m_{0g}`$ and $`m_{0s}`$ being the masses of the lightest gravitational KK mode and string excitation respectively) the total transition amplitude exhibits the following asymptotic behaviour
$$\begin{array}{ccc}\hfill M& =M_1+M_2+M_3=\frac{1}{2\varphi ^2}[2(m_h^2m_\varphi ^2)\hfill & \\ & +(2m_h^2+m_\varphi ^2)^2(\frac{1}{(tm_h^2)}+\frac{1}{(um_h^2)})]\hfill & \\ & \frac{(m_h^2m_\varphi ^2)}{\varphi ^2}+O(\frac{1}{s})\hfill & (10)\hfill \end{array}$$
Following Lee, Quigg and Thacker if we require that the J=0 partial wave amplitude must satisfy the constraint $`|a_0|<1`$ we get the unitarity bound $`|m_\varphi ^2m_h^2|<16\pi \varphi ^2`$.
Unitarity bound on $`m_h`$ from the process $`hhh\varphi `$
The transition amplitude for the process $`hhh\varphi `$ gets contributions only from radion couplings in the linear approximation. Since the gravitational KK modes do not contribute to the process $`hhh\varphi `$ the bound on $`m_h`$ that we shall present below will be valid more generally in the complete model of Randall and Sundrum. Let $`M_1`$, $`M_2`$ and $`M_3`$ denote the transition amplitudes due to s, t and u channel higgs exchanges for the process $`hhh\varphi `$. $`M_4`$ be the transition amplitude for the conatct interaction diagram. We then find that at very high energies to leading order in $`\frac{m_h^2}{s}`$ and $`\frac{m_\varphi ^2}{s}`$
$$M_1=\frac{6\lambda v}{\varphi }[1\frac{(2m_h^2+m_\varphi ^2)}{(sm_h^2)}]\frac{6\lambda v}{\varphi }(1\frac{2m_h^2+m_\varphi ^2}{s}).$$
$`(11)`$
$$M_2=\frac{6\lambda v}{\varphi }[1\frac{(2m_h^2+m_\varphi ^2)}{(tm_h^2)}]\frac{6\lambda v}{\varphi }(1+\frac{2}{s}\frac{2m_h^2+m_\varphi ^2}{(\alpha \beta x)}).$$
$`(11)`$
$$M_3=\frac{6\lambda v}{\varphi }[1\frac{(2m_h^2+m_\varphi ^2)}{(um_h^2)}]\frac{6\lambda v}{\varphi }(1+\frac{2}{s}\frac{2m_h^2+m_\varphi ^2}{(\alpha +\beta x)}).$$
$`(12)`$
and
$$M_4=24\lambda v\varphi .$$
$`(13)`$
where $`\alpha =1\frac{m_h^2+m_\varphi ^2}{s}`$ and $`\beta =1\frac{3m_h^2+m_\varphi ^2}{s}`$. $`x=\mathrm{cos}\theta `$ where $`\theta `$ is the angle between the outgoing h and one of the incoming h. At sufficiently high energies i.e. for $`sm_h^2,m_\varphi ^2`$ the J=0 partial wave amplitude exhibits the following asymptotic behaviour
$$a_0\frac{3\lambda v}{8\pi \varphi }[1+\frac{2m_h^2+m_\varphi ^2}{s}\frac{2}{\beta }\frac{2m_h^2+m_\varphi ^2}{s}\mathrm{ln}\frac{\alpha +\beta }{\alpha \beta }].$$
$`(14)`$
The unitarity constraint then gives rise to the bound $`m_h^2<\frac{16\pi }{3}\varphi v`$. We would like to note first that in the purely SM a similar unitarity bound on $`m_h`$ can be derived from the process $`hhhh`$. Second the radion mass $`m_\varphi `$ decouples from the asymptotic form of the transition amplitude for the process $`hhh\varphi `$. This causes the unitarity bound on $`m_h^2`$ derived in this paper to depend linearly on v in contrast to its quadratic dependence on v in the SM.
Results and Conclusion
Combining the unitarity bound on $`m_h`$ derived from the process $`hhh\varphi `$ with the result $`|m_\varphi ^2m_h^2|<16\pi \varphi ^2`$ we obtain $`m_\varphi <4\sqrt{\pi }\varphi [1+\frac{v}{3\varphi }]^{\frac{1}{2}}`$. Since the process $`hh\varphi \varphi `$ can receive contributions from gravitational KK modes and string excitations and we have ignored such contributions the above bound on $`m_\varphi `$ is valid only in the low energy effective theory involving the SM particles and the light radion. It is in this sense that we have called the bound a quasi unitarity bound. The bound on $`m_h`$ is however valid more generally in the context of the complete Randall-Sundrum model. Note that the results obtained in this paper should be valid irrespective of the value of $`\varphi `$. Thus for GUT size compactification where $`\varphi 10^{15}`$ Gev we would expect unitarity violation in the process $`hhh\varphi `$ to occur if $`m_h>10^9`$ Gev. However for the same type of compactification unitarity violation would occur in the process $`hh\varphi \varphi `$ if $`m_\varphi >5\times 10^{15}`$ Gev. Hence for $`\varphi v`$ the unitarity bound on $`m_\varphi `$ is determined mainly by $`\varphi `$. On the other hand if $`\varphi v`$ the uniatity bounds on both $`m_h`$ and $`m_\varphi `$ will be much small compared to the EW symmetry breaking scale.
References
1. L. Randall and R. Sundrum, Phys. Rev. Lett 83, 3370 (1999).
2. W. D. Goldberger and M. B. Wise, Phys. Rev. Lett. 83, 4922 (1999).
3. C. Csaki, M. Graesser, L. Randall and J. Terning, hep-ph/9911406; W. D. Goldberger and M. B. Wise, hep-ph/ 9911457.
4. B. W. Lee, C. Quigg and H. B. Thacker, Phys. Rev. D, 16, 1519 (1977); D. Dicus and V. Mathur, Phys. Rev. D 7, 3111 (1973).
5. S. Dawson and and S. Willenbrock, Phys. Rev. Lett. 62, 1232 (1989); W. Marciano, G. Valencia and S. Willenbrock, Phys. Rev. D 40, 1725 (1989); L. Durand, J. Johnson and J. Lopez, Phys. Rev. Lett. 64, 1215 (1990).
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# ON 7Li ENRICHMENT BY LOW MASS METAL POOR RED GIANT BRANCH STARS
## 1 INTRODUCTION
The main purpose of this letter is to discuss the possibility of <sup>7</sup>Li enrichment in some components of the Galactic halo in view of recent discoveries of very strong and strong Li giant stars in globular clusters M3 (Kraft et al., 1999), NGC 362 (Smith, Shetrone, & Keane, 1999) and M5 (Carney, Fry, & Conzalez, 1998) In general, Galactic chemical evolution models contain the following sources of <sup>6</sup>Li and <sup>7</sup>Li enrichment in the halo: Galactic cosmic rays $`\alpha \alpha `$ nucleosynthesis, SN II $`\nu `$-process and the <sup>7</sup>Be process in intermediate mass ($`3M_{}<M<7M_{}`$) AGB stars (Deliyannis, Boesgaard, & King, 1995; Boesgaard et al., 1998; D’Antona & Matteucci, 1991; Romano et al., 1999). The <sup>6</sup>Li detection in two stars HD 84937 (Smith, Lambert, & Nissen, 1993, 1998) and BD$`+26^{}3578`$ (Hobbs & Thorburn, 1994, 1997) indicated for the first time that a Galactic source other than a cosmologic one was at play, at least in the metallic range of the interstellar matter preceding the formation of stars having \[Fe/H\]<sup>1</sup><sup>1</sup>1here we use the following notation \[X/H\]$`=\mathrm{log}(\mathrm{X}/\mathrm{H})_{}\mathrm{log}(\mathrm{X}/\mathrm{H})_{}`$ and $`\mathrm{log}ϵ_\mathrm{X}=\mathrm{log}(\mathrm{X}/\mathrm{H})+12`$, where X and H are the number of atoms of the element X and hydrogen respectively $`2.3`$. Cosmic rays have produced <sup>6</sup>Li in a pre-stellar interstellar medium by $`\alpha \alpha `$ process or directly on stars. This last possibility has been shown, however, to be inefficient (Lambert, 1995). Also, the creation of Li on stars’ surfaces appears to be uncertain (Montes & Ramsey, 1998). Difficulties related to <sup>7</sup>Li formation by SN II $`\nu `$ and AGB giants are discussed by (Boesgaard et al., 1998). They concern the absence of spectral signatures in the Li enriched dwarf and sub-giant stars proving the action of these processes such as the absence of excesses of Mg in the case of SN II and of the $`s`$-elements in the case of AGB stars. These Li-rich objects as the field star BD$`+23^{}3912`$ (King, Deliyannis, & Boesgaard, 1996) present an overabundance in respect to the mean <sup>7</sup>Li abundance of the Spite plateau (Spite, & Spite, 1982) Observational <sup>7</sup>Li enrichment in AGB giants which have been discovered in the LMC (Smith & Lambert, 1989, 1990) are, however, in a much sounder theoretical basis (Sackmann & Boothroyd, 1992; Mazzitelli et al., 1999). It appears interesting to explore the possibilities of a new source of <sup>7</sup>Li enrichment in population II stars systems. This source is the metal poor low mass $`(M<2.5M_{})`$ first-ascent red giant branch stars (RGB).
## 2 THE METAL DEFICIENT GIANTS <sup>7</sup>Li ENRICHMENT PROCESS
Low metal RGB stars are expected to lose a reasonable quantity of mass $`(0.10.2M_{})`$ during their first ascent on the giant branch. Considering a time scale to reach the tip of this branch as $`2\times 10^6`$ y (Rood, 1972; Kraft et al., 1993) we obtain mean stellar mass losses of the order of $`5\times 10^8`$ up to $`10^7M_{}`$/y. It is interesting to note that the same mass loss rates are present, nevertheless in a discontinued way in the mass loss – <sup>7</sup>Li enrichment scenario proposed by de la Reza, Drake, & da Silva (1996) and de la Reza et al. (1997). In this scenario, all low mass giants $`(M<2.5M_{}`$) suffer a prompt <sup>7</sup>Li enrichment in the upper part of the RGB, after the first dredge-up and before the RGB tip. The internal mechanism producing new <sup>7</sup>Li forms a circumstellar shell (CS) enriched with <sup>7</sup>Li, which detaches from the star when the internal process ceases. In this way, the interstellar medium is enriched with <sup>7</sup>Li. The synchronized expansion of the dusty CS enriched with <sup>7</sup>Li and the subsequent <sup>7</sup>Li depletion in star photospheres can be followed by means of closed loops in an IRAS color-color diagram and compared with observed positions of the stars (de la Reza, Drake, & da Silva, 1996; de la Reza et al., 1997). These times of expansion measure the “lithium cycles” which are the periods when strong Li lines are observed. Li cycles of the order of $`10^3`$ up to $`10^5`$ y have been estimated for an CS expanding velocity of 2 km/s (de la Reza et al., 1997). Maybe the best mechanism for producing the <sup>7</sup>Li photospheric enrichment from internal origin for these low mass giants is the “cool bottom process” (CBP) (Sackmann & Boothroyd, 1999) based on the <sup>7</sup>Be production by means of the excess of <sup>3</sup>He in the H-burning shell which is characteristic of these low mass stars. The fresh <sup>7</sup>Be is transported by a conveyor circulating mechanism up to the base of the convective layer to be then taken to the stellar surface where it is transformed into <sup>7</sup>Li. Because the mentioned conveyor mechanism attains deeper and hotter regions in the case of metal deficient giants, very large <sup>7</sup>Li surface abundances ($`\mathrm{log}ϵ_{\mathrm{Li}}4.2`$) have been obtained for \[Fe/H\]$`=2.3`$ in a short episode (see Fig. 10 in Sackmann & Boothroyd (1999). In the above mentioned scenario which links a discontinuous <sup>7</sup>Li enrichment to mass loss, the process can be repeated depending on the available quantity of <sup>3</sup>He. The nature of the internal process which provokes the formation of the CS is not specified however; some mechanisms can be considered, such as a conversion of a rapidly rotating core (Fekel & Balachandran, 1993; Fekel et al., 1996) or an outward extension of large convection cells as those which appear to be producing the extended atmosphere of Betelgeuse (Lim et al., 1998).
Another way to examine the Li enrichment is by relating it to the luminosity bump in the RGB stage. This bump corresponds to the evolutionary stage when the hydrogen burning shell erases the chemical discontinuity left behind by the first dredge-up at the moment the convective envelope was at its maximum extent or deepest penetration. Because the CBP is assumed to start at this stage until the tip of the RGB is reached, we expect, if the CBP is responsible for the <sup>7</sup>Li enrichment, to observe Li-rich giants at luminosities equal or higher than the RGB bump. This is the case for the giant IV-101 in M3 where the star is observed at $`V=13.2^\mathrm{m}`$ and the very precise RGB bump of M3 is $`V=15.45^\mathrm{m}\pm 0.05^\mathrm{m}`$ (Ferraro et al. 1999). Concerning population I Li K giants, and especially those stars having Hipparcos parallaxes, de la Reza (2000) has shown that these giants have higher luminosities than those corresponding to the bumps proposed by Charbonnel (1994) at least for stars with masses between 1.0 to $`1.5M_{}`$.
The relation of the <sup>7</sup>Li and <sup>13</sup>C enrichment is less clear, principally because of the scarce number of Li rich stars among the metal poor giants. In general, among population I Li K giants da Silva, de la Reza, & Barbuy (1995) and Drake (1998) have not found a clear relation between the Li abundances and the <sup>12</sup>C/<sup>13</sup>C ratios indicating that probably both enrichments are not time correlated. As far as the RGB bump is concerned, Charbonnel et al. (1998) have shown that the <sup>13</sup>C enrichment is found only at the bump or higher luminosities when considering the <sup>12</sup>C/<sup>13</sup>C ratios measured by Shetrone et al. (1993) in stars with $`1.0[\mathrm{Fe}/\mathrm{H}]0.5`$. Depending on the metallicity, population II giants can be sources of <sup>7</sup>Li only in the RGB stage and not in the following AGB stage (Sackmann & Boothroyd, 1999). Due to their higher hydrogen burning temperatures, very metal poor giants at \[Fe/H\] $`2.3`$ for example, not only produce higher quantities of fresh <sup>7</sup>Li, but also destroy all the <sup>3</sup>He during the RGB stage so that no more of it remains to produce <sup>7</sup>Li in the AGB stage. As far as the very Li-rich giant star IV-101 in M3 is concerned, it does not show any particular enhancement of the $`s`$-process elements, presenting then a typical RGB scenario. Due to incomplete <sup>3</sup>He burning, low mass mild deficient stars will maintain <sup>7</sup>Li creation by CBP during the AGB stage.
Concerning population I giants, the distinction between RGB and AGB contributions of <sup>7</sup>Li will depend more on the stellar mass. Low mass stars, below $`2.5M_{}`$, produce <sup>7</sup>Li in the RGB and the AGB by CBP (Sackmann & Boothroyd, 1999), whereas giants with masses between 4 and $`7M_{}`$ contribute in the AGB by means of the Hot Bottom Process (Sackmann & Boothroyd, 1992; Mazzitelli et al., 1999). The CBP needs to be explored to see if it can explain the intermediate mass zone between 2.5 and $`4M_{}`$.
To evaluate the interstellar Li enrichment in the prompt <sup>7</sup>Li scenario presented before, it is fundamental to know the maximum Li abundances for stars presenting fresh <sup>7</sup>Li before depletion begins to act. We consider that this is the case for IV-101 and for estimating the <sup>7</sup>Li enrichment in M3 as will be presented in a next section, we have calculated the Li abundance of IV-101 in Non-LTE (NLTE) by means of the $`\lambda 6708`$ and $`\lambda 6104`$ Li I lines using a new self-consisting methodology for chromospheric treating. Details of this method can be found in Terra (1997) and will be submitted elsewhere (Terra, de la Reza, & Batalha, 2000). The obtained <sup>7</sup>Li abundance for IV-101 is log $`ϵ_{\mathrm{Li}}`$ = 4.0. This NLTE value is an order of magnitude larger than the LTE Li abundance proposed by Kraft et al. (1999) based on the $`\lambda 6708`$ resonant Li line alone.
## 3 IS THERE A <sup>7</sup>Li - Na/Al ENRICHMENT CONNECTION?
An extensive literature exists on the CNO, Na and Al variations among globular cluster stars (a review can be seen, for instance, in Kraft (1994)). The most remarkable are the Na and Al versus O anticorrelation and the Na and Al versus N correlation. These variations indicate that a relatively rapid mixing is taking place. Considering this, we can ask if there is a <sup>7</sup>Li – Na and a <sup>7</sup>Li – Al simultaneous surface enrichment. Kraft et al. (1999) have suggested that a Li – Al correlation could be present in some giants of metallicity $`[\mathrm{Fe}/\mathrm{H}]=1.5`$ in M3. We did not find a similar correlation for mild deficient giants having extremely high <sup>7</sup>Li abundances as is the case of the high velocity star PDS 68 (\[Fe/H\]$`=0.4`$) where no substantial Na enrichment was found (Drake, 1998). Substantial Li – Al/Na correlations must exist in very low metal RGB stars because at \[Fe/H\] $`2.3`$ relatively similar internal star regions produce large increases of <sup>23</sup>Na and <sup>27</sup>Al (from the seed elements <sup>20</sup>Ne and <sup>24</sup>Mg respectively) (Cavallo, Sweigart, & Bell 1996) and that of <sup>7</sup>Li (from the seed element <sup>3</sup>He) (Sackmann & Boothroyd, 1999). The cool bottom processing proposed by Sackmann & Boothroyd for low mass giants producing the large surface Li enrichments could also play an important role concerning the Na and Al enrichment variations. If this is the case, a <sup>7</sup>Li – Al/Na connection will be independent from any stellar interactions due to high stellar density in globular clusters and will be also valid for field stars. Fujimoto, Aikawa, & Kato (1999) have suggested another scenario to explain Na/Al variations (not <sup>7</sup>Li) in globular cluster stars by means of shell flashes induced by a deep H mixing provoked by star – star interactions in a dense cluster.
## 4 ON THE <sup>7</sup>Li PRODUCTION IN M3
We calculate here the enrichment of new <sup>7</sup>Li in the interstellar matter of M3 produced by low mass giant stars $`(M<2.5M_{}`$). These stars are considered being a second generation stars formed by matter already enriched in a large part of heavy elements by a first generation of high mass, short life, stars. The prompt <sup>7</sup>Li enrichment - mass loss scenario used here will get into action only when the first giants of mass $`2.5M_{}`$ appear in M3. That is after $`10^8`$ y, which is a small fraction of the lifetime of the globular cluster. The <sup>7</sup>Li production by the $`2.0M_{}`$ stars will continue during an important part of the life of M3. Later, the <sup>7</sup>Li production will increase in time due to the rise of population of low mass giants ($`1M_{}`$). Due to time evolution constraints, the <sup>7</sup>Li prompt enrichment – mass loss mechanism will very probably never be able to operate during any initial self-enrichment of the gas cloud from which the present M3 cluster was formed. Some recent results (Drake, 1998) indicate that <sup>7</sup>Li photospheric depletion, following a strong enrichment, depends on the value of the stellar mass for masses less than $`2.5M_{}`$. RGB stars with masses $`2M_{}`$ have larger depletion times ($`10^4`$ y) than those of stellar masses of $`1M_{}`$ ($`3\times 10^3`$ y). These results were obtained for mass losses between $`10^810^7M_{}/`$y and CS expansion velocities equal to 2 km/s and for \[Fe/H\] between $`0.5`$ and $`0.2`$. In the case of population II giants $`(M<2.5M_{}`$), where stellar masses around $`1M_{}`$ are of interest, we obtain short Li cycles resulting in lower probabilities of detection. Considering the time to reach the RGB tip as $`2\times 10^6`$ y, the probability to detect a Li-rich RGB star will be $`3\times 10^3\mathrm{y}/2\times 10^6\mathrm{y}`$, that is 0.15%. However, if the <sup>7</sup>Li enrichment process is repeated by a recurrence factor of 10 for example, we obtain 1.5%. A similar result is found by Kraft et al. (1999) considering that already two Li-rich RGB giants have been found among near 100 observed globular clusters giants. We must note that the physical basis for the actual value of the recurrence factor depends on the quantity of <sup>3</sup>He that remains to be burned. Let us estimate the <sup>7</sup>Li enrichment in the interstellar medium in the globular cluster M3 by means of mass loss of super Li-rich RGB stars, following the scenario of de la Reza, Drake, & da Silva (1996); de la Reza et al. (1997). Considering that all RGB stars are potential sources of <sup>7</sup>Li during the short time in which they form and eject a Li-rich CS, the total production of <sup>7</sup>Li will be $`P_{\mathrm{Li}}=N_{\mathrm{RGB}}N_{\mathrm{Li}}ft_{\mathrm{CS}}\dot{M}/t_{\mathrm{RGB}}`$. Here $`N_{\mathrm{RGB}}`$ is the total number of RGB stars in the globular cluster M3, $`N_{\mathrm{Li}}=(n_{\mathrm{Li}}/n_\mathrm{H})(m_{\mathrm{Li}}/m_\mathrm{H})`$ where $`(n_{\mathrm{Li}}/n_\mathrm{H})`$ is the ratio of the number of Li and H atoms equal to $`10^8`$ for $`\mathrm{log}ϵ_{\mathrm{Li}}=4.0`$ and $`m_{\mathrm{Li}}=7m_\mathrm{H}`$, $`f`$ is the recurrence factor, $`t_{\mathrm{CS}}=200`$ y is the time of CS formation (de la Reza et al. 1996), $`\dot{M}`$ is the mass loss equal to $`10^7M_{}/`$y and $`t_{\mathrm{RGB}}`$ is the time necessary to attain the RGB tip ($`2\times 10^6`$ y). The number of RGB stars ($`N_{\mathrm{RGB}}`$) can be estimated in the following way: we assume that the present observed $`N_{\mathrm{RGB}}`$ will represent a reasonable mean value of the number of RGB stars during the entire life of the globular cluster. (We consider that RGB stars have typical masses around the turn-off mass and similar to the one assumed by Kraft et al. (1999) for IV-101 ($`0.85M_{}`$)). Probably a better evaluation of the production of <sup>7</sup>Li during the life of M3 can be made, if we are able to distinguish the Li yields produced by giants with masses around $`2.0M_{}`$ from those with $`1.0M_{}`$, which is not unfortunately the case. Even in a crude way, $`N_{\mathrm{RGB}}`$ can be estimated by extrapolating the already known number of RGB stars (424) counted among almost 19 000 stars in M3 (Ferraro et al., 1997). Maintaining the same proportion for the total estimated number of stars in M3 equal to $`3.44\times 10^5`$ stars (Lang 1992) we obtain approximately 7700 RGB stars in M3. Taking first $`f=1`$ we obtain $`P_{\mathrm{Li}}=5.4\times 10^{15}M_{}`$/y of new <sup>7</sup>Li in the gas of the globular cluster interstellar medium. If we multiply this value by the mean age of the globular clusters (13 Gyrs, Mould (1998)) and if we consider this age to be that of M3, we obtain a total <sup>7</sup>Li mass production of $`7.0\times 10^5`$$`M_{}`$. To have a better idea of what represents this quantity of new <sup>7</sup>Li let us compare the latter with the initial content of <sup>7</sup>Li in M3. To estimate this initial quantity of <sup>7</sup>Li we can consider a Jeans mass of $`10^6M_{}`$ at the earliest time of formation for M3 (Peebles, 1993). Supposing for the sake of simplicity, that almost all matter consists of H atoms and considering an initial <sup>7</sup>Li abundance of $`n_{\mathrm{Li}}/n_\mathrm{H}=1.6\times 10^{10}`$ corresponding to the Spite plateau (Molaro, Primas, & Bonifacio, 1995; Bonifacio & Molaro, 1997), we obtain an initial mass of <sup>7</sup>Li of $`10^3M_{}`$ in M3. The RGB production will represent only the $``$ 7% of this quantity. But if we consider $`f=10`$ this value will increase up to $`70`$%. If we consider alternatively possible higher <sup>7</sup>Li abundances as $`\mathrm{log}ϵ_{\mathrm{Li}}=4.5`$, as obtained for some giants such as PDS 68, and if we use a larger, but yet realistic, mass loss value of $`5\times 10^7M_{}/`$y, we will obtain increased enrichment factors. It is interesting to note, that this new <sup>7</sup>Li in the interstellar gas of M3 will, very probably, not be maintained for long time in the cluster and will be transferred, by ram pressure, to the Galactic disk when M3 crosses this one. Scholz, Odenkirchen, & Irwin (1993) have calculated that M3 crossed the plane nearly 34 times during the last 10 Gyrs. Globular clusters appear then as new potential sources of <sup>7</sup>Li enrichment in the Galactic disk!
We thanks the referee for very useful suggestions which clarified some aspects of this letter. N.A.D. and M.A.T. thank FAPERJ for the financial support under the contracts E-26/151.172/98 and E-26/150.571/98 respectively. RdlR and LdS thank CNPq by the financial support grants 301375/86-0 and 200580/97-0 respectively.
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# Irreducible Hamiltonian BRST symmetry for reducible first-class systems
## 1 Introduction
It is well known that the Hamiltonian BRST formalism stands for one of the strongest and most popular quantization methods for theories with first-class constraints. This method can be applied to irreducible, as well as reducible first-class systems. In the irreducible case the ghosts can be regarded as one-forms dual to the vector fields associated with the first-class constraints. This interpretation fails in the reducible framework, being necessary to introduce ghosts for ghosts together with their canonical conjugated momenta (antighosts). The ghosts for ghosts ensure the incorporation of the reducibility relations within the cohomology of the exterior derivative along the gauge orbits, while their canonical momenta are required in order to kill the higher-resolution-degree nontrivial co-cycles from the homology of the Koszul-Tate differential. Similar considerations apply to the antifield BRST treatment .
Another way of approaching reducible systems is based on the idea of replacing such a system with an irreducible one , . This idea has been enforced in the framework of the irreducible quantization of reducible theories from the point of view of the antifield-BRST formulation , as well as of the antifield BRST-anti-BRST method . Some applications of these treatments can be found in . At the level of the Hamiltonian symmetry, this type of procedure has been developed so far only in the case of some peculiar models , but a general theory covering on-shell reducible first-class systems has not yet been given.
In view of this, here we propose an irreducible Hamiltonian BRST procedure for quantizing on-shell reducible first-class theories. In this light, we enforce the following steps: (i) we transform the original reducible first-class constraints into some irreducible ones on a larger phase-space in a manner that allows the substitution of the BRST quantization of the reducible system by that of the irreducible theory; (ii) we quantize the irreducible system accordingly the Hamiltonian BRST formalism. As a consequence, the ghosts for ghosts and their antighosts do not appear in our formalism.
The paper is organized in five sections. Section 2 deals with the derivation of an irreducible set of first-class constraints associated with the original reducible one by means of constructing an irreducible Koszul-Tate complex. The irreducible Koszul-Tate complex is obtained by requiring that all the antighost number one co-cycles from the reducible case become trivial under an appropriate redefinition of the antighost number one antighosts. This procedure leads to the introduction of new canonical variables and antighosts. In section 3 we infer the irreducible BRST symmetry corresponding to the irreducible constraint set deduced in section 2 and prove that we can replace the Hamiltonian BRST quantization of the reducible system with that of the irreducible theory. Section 4 illustrates our method in the case of a two-stage reducible model involving three-form gauge fields. In section 5 we present the main conclusions of the paper.
## 2 Derivation of an irreducible first-class constraint set. Irreducible Koszul-Tate differential
### 2.1 Setting the problem
Our starting point is a Hamiltonian system with the phase-space locally described by $`N`$ canonical pairs $`z^A=(q^i,p_i)`$, subject to the first-class constraints
$$G_{a_0}(q^i,p_i)0,a_0=1,\mathrm{},M_0,$$
(1)
which are assumed to be on-shell $`L`$-stage reducible. We suppose that there are no second-class constraints in the theory (if any, they can be eliminated with the help of the Dirac bracket). The first-class property of the constraints (1) is expressed by
$$[G_{a_0},G_{b_0}]=C_{a_0b_0}^{c_0}G_{c_0},$$
(2)
while the reducibility relations are written as
$$Z_{a_1}^{a_0}G_{a_0}=0,a_1=1,\mathrm{},M_1,$$
(3)
$$Z_{a_2}^{a_1}Z_{a_1}^{a_0}=C_{a_2}^{a_0b_0}G_{b_0},a_2=1,\mathrm{},M_2,$$
(4)
$$\mathrm{}$$
$$Z_{a_L}^{a_{L1}}Z_{a_{L1}}^{a_{L2}}=C_{a_L}^{a_{L2}b_0}G_{b_0},a_L=1,\mathrm{},M_L,$$
(5)
with the symbol $`[,]`$ denoting the Poisson bracket. For definiteness we approach here the bosonic case, but our analysis can be easily extended to fermions modulo introducing some appropriate sign factors. The first-order structure functions $`C_{a_0b_0}^{c_0}`$ may involve the phase-space coordinates (open gauge algebra) and are antisymmetric in the lower indices. The reducibility functions $`\left(Z_{a_{k+1}}^{a_k}\right)_{k=0,\mathrm{},L1}`$ and the coefficients $`\left(C_{a_{k+2}}^{a_kb_0}\right)_{k=0,\mathrm{},L2}`$ appearing in the right hand-side of the relations (45) may also depend on the canonical variables, and, in addition, $`C_{a_2}^{a_0b_0}`$ should be antisymmetric in the upper indices.
The standard Hamiltonian BRST symmetry for the above on-shell reducible first-class Hamiltonian system, $`s_R=\delta _R+D_R+\mathrm{}`$, contains two crucial operators. The Koszul-Tate differential $`\delta _R`$ realizes an homological resolution of smooth functions defined on the first-class constraint surface (1), while the model of longitudinal exterior derivative along the gauge orbits $`D_R`$ is a differential modulo $`\delta _R`$ and accounts for the gauge invariances. The degree of $`\delta _R`$ is called antighost number ($`antigh`$), the degree of $`D_R`$ is named pure ghost number ($`pgh`$), while the overall degree of the BRST differential is called ghost number ($`gh`$) and is defined like the difference between the pure ghost number and the antighost number ($`antigh\left(\delta _R\right)=1`$, $`pgh\left(D_R\right)=1`$, $`gh\left(s_R\right)=1`$). In the case of a first-stage reducible Hamiltonian system ($`L=1`$) the proper construction of $`\delta _R`$ relies on the introduction of the generators (antighosts) $`𝒫_{a_0}`$ and $`P_{a_1}`$, with the Grassmann parity ($`ϵ`$) and the antighost number given by
$$ϵ\left(𝒫_{a_0}\right)=1,ϵ\left(P_{a_1}\right)=0,antigh\left(𝒫_{a_0}\right)=1,antigh\left(P_{a_1}\right)=2,$$
(6)
on which $`\delta _R`$ is set to act like
$$\delta _R𝒫_{a_0}=G_{a_0},$$
(7)
$$\delta _RP_{a_1}=Z_{a_1}^{a_0}𝒫_{a_0},$$
(8)
being understood that
$$\delta _Rz^A=0.$$
(9)
The antighosts $`P_{a_1}`$ are required in order to “kill” the antighost number one nontrivial co-cycles
$$\mu _{a_1}=Z_{a_1}^{a_0}𝒫_{a_0},$$
(10)
(which are due to the definitions (7) and the reducibility relations (3)) in the homology of $`\delta _R`$, and thus restore the acyclicity of the Koszul-Tate differential at nonvanishing antighost numbers. For a two-stage reducible Hamiltonian system ($`L=2`$), one should supplement the antighost spectrum from the first-stage case with the antighosts $`P_{a_2}`$, displaying the characteristics $`ϵ\left(P_{a_2}\right)=1`$, $`antigh\left(P_{a_2}\right)=3`$, and define the action of $`\delta _R`$ on them through $`\delta _RP_{a_2}=Z_{a_2}^{a_1}P_{a_1}\frac{1}{2}C_{a_2}^{a_0b_0}𝒫_{b_0}𝒫_{a_0}`$. The acyclicity of $`\delta _R`$ is thus achieved by making exact the antighost number two nontrivial co-cycles $`\mu _{a_2}=Z_{a_2}^{a_1}P_{a_1}+\frac{1}{2}C_{a_2}^{a_0b_0}𝒫_{b_0}𝒫_{a_0}`$, which are present on behalf of the definitions (8) and the reducibility relations (4). The process goes along the same lines at higher antighost numbers. In the general situation of an $`L`$-stage reducible Hamiltonian theory, the antighost spectrum contains the variables $`𝒫_{a_0}`$ and $`\left(P_{a_k}\right)_{k=1,\mathrm{},L}`$, with $`ϵ\left(P_{a_k}\right)=k+1mod\mathrm{\hspace{0.33em}2}`$, $`antigh\left(P_{a_k}\right)=k+1`$, the actions of $`\delta _R`$ on higher antighost number antighosts being taken in such a way to ensure the acyclicity of the Koszul-Tate differential at nonvanishing antighost numbers.
The problem to be investigated in the sequel is the derivation of an irreducible set of first-class constraints associated with the $`L`$-stage reducible one (1). Our basic idea is to redefine the antighosts $`𝒫_{a_0}`$ in such a way that the new co-cycles of the type (10) are trivial. Then, the antighosts $`P_{a_1}`$ ensuring the triviality of the co-cycles (10) are no longer necessary, so they will be discarded from the Koszul-Tate complex. Moreover, the absence of these antighosts implies the absence of nontrivial co-cycles at antighost number greater that one, hence the antighosts $`\left(P_{a_k}\right)_{k=2,\mathrm{},L}`$ will also be discarded from the antighost spectrum. The enforcement of this idea leads to an irreducible set of first-class constraints underlying some physical observables that coincide with those deriving from the original reducible system. In order to simplify the presentation we initially approach the cases $`L=1,2`$, and further generalize the results to an arbitrary $`L`$.
### 2.2 The case $`L=1`$
We begin with the definitions (78) and the reducibility relations (3). As we have previously mentioned, we redefine the antighosts $`𝒫_{a_0}`$ like
$$𝒫_{a_0}𝒫_{a_0}^{}=𝒫_{a_0}Z_{b_1}^{b_0}\overline{D}_{c_1}^{b_1}A_{a_0}^{c_1}𝒫_{b_0},$$
(11)
where $`\overline{D}_{c_1}^{b_1}`$ stands for the inverse of $`D_{c_1}^{b_1}=Z_{c_1}^{c_0}A_{c_0}^{b_1}`$, while $`A_{a_0}^{c_1}`$ are some functions that may involve at most the variables $`z^A`$ and are chosen to satisfy $`rank\left(D_{c_1}^{b_1}\right)=M_1`$. Next, we replace (7) with
$$\delta 𝒫_{a_0}^{}=G_{a_0}.$$
(12)
The definitions (12) imply some co-cycles of the type (10)
$$\mu _{a_1}^{}=Z_{a_1}^{a_0}𝒫_{a_0}^{},$$
(13)
which are found trivial on behalf of (11), namely, $`\mu _{a_1}^{}0`$. Thus, the definitions (12) do not lead to nontrivial co-cycles at antighost number one, therefore the antighosts $`P_{a_1}`$ will be discarded. Moreover, formulas (12) are helpful at deriving some irreducible first-class constraints corresponding to (1). For this reason in (12) we used the notation $`\delta `$ instead of $`\delta _R`$. The derivation of the irreducible first-class constraints relies on enlarging the phase-space with the bosonic canonical pairs $`(y^{a_1},\pi _{a_1})`$, where the momenta $`\pi _{a_1}`$ are demanded to be nonvanishing solutions to the equations
$$D_{a_1}^{b_1}\pi _{b_1}=\delta \left(Z_{a_1}^{a_0}𝒫_{a_0}\right).$$
(14)
As $`D_{a_1}^{b_1}`$ is invertible, the nonvanishing solutions to (14) implement the irreducibility. This is because the equations (14) possess nonvanishing solutions if and only if
$$\delta \left(Z_{a_1}^{a_0}𝒫_{a_0}\right)0,$$
(15)
hence if and only if (10) are not co-cycles. Inserting (14) in (12) and using (11), we arrive at
$$\delta 𝒫_{a_0}=G_{a_0}A_{a_0}^{a_1}\pi _{a_1},$$
(16)
which emphasize the irreducible constraints
$$\gamma _{a_0}G_{a_0}+A_{a_0}^{a_1}\pi _{a_1}0.$$
(17)
From (17) it is easy to see that
$$\pi _{a_1}=\overline{D}_{a_1}^{b_1}Z_{b_1}^{b_0}\gamma _{b_0},G_{a_0}=\left(\delta _{a_0}^{b_0}Z_{b_1}^{b_0}\overline{D}_{a_1}^{b_1}A_{a_0}^{a_1}\right)\gamma _{b_0},$$
(18)
so
$$[\gamma _{a_0},\gamma _{b_0}]=\overline{C}_{a_0b_0}^{c_0}\gamma _{c_0},$$
(19)
for some $`\overline{C}_{a_0b_0}^{c_0}`$. Thus, the irreducible constraints (17) are first-class. In the meantime, if we take the standard action of $`\delta `$ on the phase-space coordinates
$$\delta z^A=0,\delta z^{A_1}=0,$$
(20)
with $`z^{A_1}=(y^{a_1},\pi _{a_1})`$, then formulas (16) and (20) completely define an irreducible Koszul-Tate complex corresponding to an irreducible system based on the first-class constraints (17).
At this point we underline two important observations. First, the number of physical degrees of freedom of the irreducible system coincides with the original one. Indeed, in the reducible case there are $`N`$ canonical pairs and $`M_0M_1`$ independent first-class constraints, hence $`NM_0+M_1`$ physical degrees of freedom. In the irreducible situation there are $`N+M_1`$ canonical pairs and $`M_0`$ independent first-class constraints, therefore as many physical degrees of freedom as in the reducible case. Second, from (14) it results (due to the invertibility of $`D_{a_1}^{b_1}`$) that the momenta $`\pi _{a_1}`$ are $`\delta `$-exact. These results represent two desirable features, which will be requested also in connection with higher-order reducible Hamiltonian systems. Anticipating a bit, we notice that for higher-order reducible theories it is necessary to further add some supplementary canonical variables and antighost number one antighosts. While the former request indicates the number of new canonical pairs and first-class constraints to be added within the irreducible framework, the latter ensures, as it will be further seen, that there exists a proper redefinition of the antighost number one antighosts that makes trivial the co-cycles from the reducible approach.
### 2.3 The case $`L=2`$
Now, we start from the reducibility relations (34) and intend to preserve the definitions (16) and (20) with respect to the irreducible Koszul-Tate differential $`\delta `$. Nevertheless, there appear two difficulties. First, the matrix $`D_{c_1}^{b_1}`$ is not invertible now due to the second-stage reducibility relations (4). Rather, it possesses $`Z_{a_2}^{c_1}`$ as on-shell null vectors
$$D_{c_1}^{b_1}Z_{a_2}^{c_1}=A_{c_0}^{b_1}Z_{c_1}^{c_0}Z_{a_2}^{c_1}=C_{a_2}^{c_0d_0}A_{c_0}^{b_1}G_{d_0}0,$$
(21)
hence the transformations (11) fail to be correct. We will still maintain the definitions (16) and will subsequently show that there exists a redefinition of the antighosts $`𝒫_{a_0}`$ that brings the constraint functions $`\gamma _{a_0}`$ under the form (17). Second, the irreducible first-class constraints (17) are not enough in order to maintain the initial number of physical degrees of freedom in connection with the irreducible theory, so they should be supplemented with $`M_2`$ new constraints in such a way to ensure a first-class irreducible behaviour of the overall constraint set. Thus, we must introduce some fermionic antighost number one antighosts $`𝒫_{a_2}`$ and set
$$\delta 𝒫_{a_2}=\gamma _{a_2},$$
(22)
where $`\gamma _{a_2}0`$ denote the new first-class constraints. On the one hand, we should infer the concrete form of $`\gamma _{a_2}`$ such that the new Koszul-Tate complex based on the definitions (16), (20) and (22) is irreducible, and, on the other hand, we should prove that it is possible to perform a redefinition of the antighosts $`𝒫_{a_0}`$ (eventually involving also the new antighosts) such that (16) is gained. This goes as follows. We begin with the relations (16) on which we apply $`Z_{b_1}^{a_0}`$, and obtain
$$\delta \left(Z_{b_1}^{a_0}𝒫_{a_0}\right)=D_{b_1}^{a_1}\pi _{a_1}.$$
(23)
Formulas (21) make permissible a representation of $`D_{b_1}^{a_1}`$ under the form
$$D_{b_1}^{a_1}=\delta _{b_1}^{a_1}Z_{a_2}^{a_1}\overline{D}_{b_2}^{a_2}A_{b_1}^{b_2}+A_{a_0}^{a_1}C_{c_2}^{a_0b_0}\overline{D}_{b_2}^{c_2}A_{b_1}^{b_2}G_{b_0},$$
(24)
where $`\overline{D}_{b_2}^{a_2}`$ is the inverse of $`D_{b_2}^{a_2}=Z_{b_2}^{a_1}A_{a_1}^{a_2}`$ and $`A_{a_1}^{a_2}`$ are some functions that may depend at most on $`z^A`$ and are taken to fulfill $`rank\left(D_{b_2}^{a_2}\right)=M_2`$. Inserting (24) in (23), we infer that
$$\delta \left(Z_{b_1}^{a_0}𝒫_{a_0}\right)=\pi _{b_1}+Z_{a_2}^{a_1}\overline{D}_{b_2}^{a_2}A_{b_1}^{b_2}\pi _{a_1}C_{c_2}^{a_0b_0}\overline{D}_{b_2}^{c_2}A_{b_1}^{b_2}A_{a_0}^{a_1}\pi _{a_1}G_{b_0},$$
(25)
which turns into
$$\delta \left(Z_{b_1}^{a_0}𝒫_{a_0}C_{c_2}^{a_0b_0}\overline{D}_{b_2}^{c_2}A_{b_1}^{b_2}G_{b_0}𝒫_{a_0}\right)=\pi _{b_1}+\overline{D}_{b_2}^{a_2}A_{b_1}^{b_2}Z_{a_2}^{a_1}\pi _{a_1},$$
(26)
on account on the one hand of the antisymmetry of $`C_{c_2}^{a_0b_0}`$, which implies the relations $`C_{c_2}^{a_0b_0}A_{a_0}^{a_1}\pi _{a_1}G_{b_0}=C_{c_2}^{a_0b_0}\gamma _{a_0}G_{b_0}`$, and, on the other hand, of the definitions (16) and (20). The second term in the right hand-side of (26) implies that $`\pi _{a_1}`$ are not $`\delta `$-exact. An elegant manner of complying with the requirement on the $`\delta `$-exactness of these momenta is to take the functions $`\gamma _{a_2}`$ like
$$\gamma _{a_2}Z_{a_2}^{a_1}\pi _{a_1},$$
(27)
so $`\delta 𝒫_{a_2}=Z_{a_2}^{a_1}\pi _{a_1}`$, which then leads, via (26), to
$$\pi _{b_1}=\delta \left(Z_{b_1}^{a_0}𝒫_{a_0}+C_{c_2}^{a_0b_0}\overline{D}_{b_2}^{c_2}A_{b_1}^{b_2}G_{b_0}𝒫_{a_0}\overline{D}_{b_2}^{a_2}A_{b_1}^{b_2}𝒫_{a_2}\right).$$
(28)
On behalf of (27) and of the fact that the functions $`Z_{a_2}^{a_1}`$ have no null vectors, we find that (22) imply no nontrivial co-cycles. In this way the first task is achieved. Introducing (28) in (16), we arrive at
$$\delta 𝒫_{a_0}^{\prime \prime }=G_{a_0},$$
(29)
where
$$𝒫_{a_0}^{\prime \prime }=𝒫_{a_0}Z_{b_1}^{b_0}A_{a_0}^{b_1}𝒫_{b_0}+C_{c_2}^{c_0b_0}\overline{D}_{b_2}^{c_2}A_{b_1}^{b_2}A_{a_0}^{b_1}G_{b_0}𝒫_{c_0}\overline{D}_{b_2}^{a_2}A_{b_1}^{b_2}A_{a_0}^{b_1}𝒫_{a_2}.$$
(30)
Formula (30) expresses a redefinition of the antighost number one antighosts that is in agreement with (16). It is precisely the requirement on the $`\delta `$-exactness of the momenta $`\pi _{a_1}`$ that allows us to deduce (30). Thus, the second task is also attained.
It remains to be proved that (29) also gives no nontrivial co-cycles. From (29) we obtain the co-cycles
$$\mu _{a_1}^{\prime \prime }=Z_{a_1}^{a_0}𝒫_{a_0}^{\prime \prime },$$
(31)
at antighost number one. After some computation, we find that they are trivial as
$$\mu _{a_1}^{\prime \prime }=\delta \left(\frac{1}{2}C_{c_2}^{a_0b_0}\overline{D}_{b_2}^{c_2}A_{a_1}^{b_2}𝒫_{a_0}^{\prime \prime }𝒫_{b_0}^{\prime \prime }\right).$$
(32)
In conclusion, we associated an irreducible Koszul-Tate complex based on the definitions (16), (20) and (22) with the starting second-stage reducible one. This complex underlies the irreducible constraint functions (17) and (27), which can be shown to be first-class. Indeed, with the help of (17) and (27) we get that
$$\pi _{b_1}=\left(Z_{b_1}^{a_0}C_{c_2}^{a_0b_0}\overline{D}_{b_2}^{c_2}A_{b_1}^{b_2}G_{b_0}\right)\gamma _{a_0}+\overline{D}_{b_2}^{a_2}A_{b_1}^{b_2}\gamma _{a_2},$$
(33)
$$G_{a_0}=\left(\delta _{a_0}^{c_0}Z_{b_1}^{c_0}A_{a_0}^{b_1}+C_{c_2}^{c_0b_0}\overline{D}_{b_2}^{c_2}A_{b_1}^{b_2}A_{a_0}^{b_1}G_{b_0}\right)\gamma _{c_0}\overline{D}_{b_2}^{a_2}A_{b_1}^{b_2}A_{a_0}^{b_1}\gamma _{a_2}.$$
(34)
Evaluating the Poisson brackets among the constraint functions $`(\gamma _{a_0},\gamma _{a_2})`$ with the help of (3334), we find that they weakly vanish on the surface $`\gamma _{a_0}0`$, $`\gamma _{a_2}0`$, hence (17) and (27) are first-class. This completes the case $`L=2`$.
### 2.4 Generalization to an arbitrary $`L`$
Now, we are in the position to generalize the above discussion to an arbitrary $`L`$. Acting like in the previous cases, we enlarge the phase-space with the bosonic canonical pairs $`z^{A_{2k+1}}=(y^{a_{2k+1}},\pi _{a_{2k+1}})_{k=0,\mathrm{},a}`$ and construct an irreducible Koszul-Tate complex based on the definitions
$$\delta z^A=0,\delta z^{A_{2k+1}}=0,k=0,\mathrm{},a,$$
(35)
$$\delta 𝒫_{a_0}=\gamma _{a_0},$$
(36)
$$\delta 𝒫_{a_{2k}}=\gamma _{a_{2k}},k=1,\mathrm{},b,$$
(37)
which emphasize the irreducible constraints
$$\gamma _{a_0}G_{a_0}+A_{a_0}^{a_1}\pi _{a_1}0,$$
(38)
$$\gamma _{a_{2k}}Z_{a_{2k}}^{a_{2k1}}\pi _{a_{2k1}}+A_{a_{2k}}^{a_{2k+1}}\pi _{a_{2k+1}}0,k=1,\mathrm{},b.$$
(39)
The antighosts $`\left(𝒫_{a_{2k}}\right)_{k=0,\mathrm{},b}`$ are all fermionic with antighost number one, while the notations $`a`$ and $`b`$ signify
$$a=\{\begin{array}{c}\frac{L1}{2},\mathrm{for}L\mathrm{odd},\\ \frac{L}{2}1,\mathrm{for}L\mathrm{even},\end{array}b=\{\begin{array}{c}\frac{L1}{2},\mathrm{for}L\mathrm{odd},\\ \frac{L}{2},\mathrm{for}L\mathrm{even}.\end{array}$$
(40)
In order to avoid confusion, we use the conventions $`f^{a_k}=0`$ if $`k<0`$ or $`k>L`$. The matrices of the type $`A_{a_k}^{a_{k+1}}`$ implied in (3839) may involve at most the original variables $`z^A`$ and are taken to fulfill the relations
$$rank\left(D_{b_k}^{a_k}\right)\underset{i=k}{\overset{L}{}}()^{k+i}M_i,k=1,\mathrm{},L1,$$
(41)
$$rank\left(D_{b_L}^{a_L}\right)=M_L,$$
(42)
where $`D_{b_k}^{a_k}=Z_{b_k}^{a_{k1}}A_{a_{k1}}^{a_k}`$. We remark that the choice of the functions $`A_{a_{k1}}^{a_k}`$ is not unique. Moreover, for a definite choice of these functions, the equations (4142) are unaffected if we modify $`A_{a_{k1}}^{a_k}`$ as
$$A_{a_{k1}}^{a_k}A_{a_{k1}}^{a_k}+\mu _{b_{k2}}^{a_k}Z_{a_{k1}}^{b_{k2}},$$
(43)
hence these functions carry some ambiguities.
Acting like in the cases $`L=1,2`$, after some computation we find the relations
$$\pi _{a_{2k+1}}=m_{a_{2k+1}}^{a_{2j}}\gamma _{a_{2j}},k=0,\mathrm{},a,$$
(44)
$$G_{a_0}=m_{a_0}^{a_{2j}}\gamma _{a_{2j}},$$
(45)
for some functions $`m_{a_{2k+1}}^{a_{2j}}`$ and $`m_{a_0}^{a_{2j}}`$. Computing the Poisson brackets among the constraint functions in (3839), we find that they weakly vanish on the surface (3839), hence they form a first-class set. This ends the general construction of an irreducible Koszul-Tate complex associated with the original on-shell $`L`$-stage reducible Hamiltonian system.
## 3 Irreducible BRST symmetry for on-shell reducible Hamiltonian systems
### 3.1 Derivation of the irreducible BRST symmetry
The first step in deriving an irreducible BRST symmetry for the investigated on-shell $`L`$-stage reducible Hamiltonian system has been implemented by constructing an irreducible Koszul-Tate complex based on the irreducible first-class constraints (3839). The BRST symmetry corresponding to this irreducible first-class constraint set can be decomposed like
$$s=\delta +D+\mathrm{},$$
(46)
where $`\delta `$ is the irreducible Koszul-Tate differential generated in the previous section, $`D`$ represents the longitudinal exterior derivative along the gauge orbits, and the other pieces, generically denoted by “$`\mathrm{}`$”, ensure the nilpotency of $`s`$. The ghost spectrum of the longitudinal exterior complex includes only the fermionic ghosts $`\eta ^\mathrm{\Delta }\left(\eta ^{a_{2k}}\right)_{k=0,\mathrm{},b}`$ with pure ghost number one ($`pgh\left(\eta ^{a_{2k}}\right)=1`$), respectively associated with the irreducible first-class constraints (3839), to be redenoted by $`\gamma _\mathrm{\Delta }\left(\gamma _{a_{2k}}\right)_{k=0,\mathrm{},b}`$. The standard definitions of $`D`$ acting on the generators from the longitudinal exterior complex read as
$$DF=[F,\gamma _\mathrm{\Delta }]\eta ^\mathrm{\Delta },$$
(47)
$$D\eta ^\mathrm{\Delta }=\frac{1}{2}C_{\mathrm{\Delta }^{}\mathrm{\Delta }^{\prime \prime }}^\mathrm{\Delta }\eta ^\mathrm{\Delta }^{}\eta ^{\mathrm{\Delta }^{\prime \prime }},$$
(48)
where $`F`$ can be any function of the variables $`z^A`$, $`\left(z^{A_{2k+1}}\right)_{k=0,\mathrm{},a}`$, and $`C_{\mathrm{\Delta }^{}\mathrm{\Delta }^{\prime \prime }}^\mathrm{\Delta }`$ stand for the first-order structure functions corresponding to the first-class constraint functions
$$[\gamma _\mathrm{\Delta },\gamma _\mathrm{\Delta }^{}]=C_{\mathrm{\Delta }\mathrm{\Delta }^{}}^{\mathrm{\Delta }^{\prime \prime }}\gamma _{\mathrm{\Delta }^{\prime \prime }}.$$
(49)
The action of $`\delta `$ can be extended to the ghosts through
$$\delta \eta ^\mathrm{\Delta }=0,$$
(50)
with $`antigh\left(\eta ^\mathrm{\Delta }\right)=0`$, such that both the nilpotency and acyclicity of the irreducible Koszul-Tate differential are maintained, while $`D`$ can be appropriately extended to the antighosts in such a way to become a differential modulo $`\delta `$ on account of the first-class behaviour of the irreducible constraints. On these grounds, the homological perturbation theory guarantees the existence of a nilpotent BRST symmetry $`s`$ of the form (46) associated with the irreducible first-class constraints (3839).
### 3.2 Link with the standard reducible BRST symmetry
In the sequel we establish the correlation between the standard Hamiltonian BRST symmetry corresponding to the starting reducible first-class system and that associated with the irreducible one, investigated in the above subsection. In this light we show that the physical observables of the two theories coincide. Let $`F`$ be an observable of the irreducible system. Consequently, it satisfies the equations
$$[F,\gamma _{a_{2k}}]0,k=0,\mathrm{},b,$$
(51)
where the weak equality refers to the surface (3839). Using the relations (4445), we then find that $`F`$ also fulfills the equations
$$[F,G_{a_0}]=[F,m_{a_0}^{a_{2j}}]\gamma _{a_{2j}}+[F,\gamma _{a_{2j}}]m_{a_0}^{a_{2j}}0,$$
(52)
$$[F,\pi _{a_{2k+1}}]=[F,m_{a_{2k+1}}^{a_{2j}}]\gamma _{a_{2j}}+[F,\gamma _{a_{2j}}]m_{a_{2k+1}}^{a_{2j}}0,k=0,\mathrm{},a,$$
(53)
on this surface. So, every observable of the irreducible theory should verify the equations (5253) on the surface (3839). Now, we observe that the first-class surface described by the relations (3839) is equivalent with that described by the equations
$$G_{a_0}0,\pi _{a_{2k+1}}0,k=0,\mathrm{},a.$$
(54)
Indeed, it is clear that when (54) hold, (3839) hold, too. The converse results from (4445), which show that if (3839) hold, then (54) also hold. This proves the equivalence between the first-class surfaces (3839) and (54). As a consequence, we have that every observable of the irreducible theory, which we found that verifies the equations (5253) on the surface (3839), will check these equations also on the surface (54). This means that every observable of the irreducible system is an observable of the theory based on the first-class constraints (54). Conversely, if $`F`$ represents a physical observable of the system underlying the constraints (54), then it should check the equations
$$[F,G_{a_0}]0,[F,\pi _{a_{2k+1}}]0,k=0,\mathrm{},a,$$
(55)
on the surface (54). Then, it follows that $`F`$ satisfies the relations
$$[F,\gamma _{a_0}]=[F,G_{a_0}]+[F,A_{a_0}^{a_1}]\pi _{a_1}+[F,\pi _{a_1}]A_{a_0}^{a_1}0,$$
(56)
$`[F,\gamma _{a_{2k}}]=[F,Z_{a_{2k}}^{a_{2k1}}]\pi _{a_{2k1}}+[F,\pi _{a_{2k1}}]Z_{a_{2k}}^{a_{2k1}}+`$
$`[F,A_{a_{2k}}^{a_{2k+1}}]\pi _{a_{2k+1}}+[F,\pi _{a_{2k+1}}]A_{a_{2k}}^{a_{2k+1}}0,k=1,\mathrm{},b,`$ (57)
on the same surface. Recalling once again the equivalence between this surface and the one expressed by (3839), we find that $`F`$ will verify the equations (563.2) also on the surface (3839), being therefore an observable of the irreducible system. From the above discussion we conclude that the physical observables of the irreducible theory coincide with those associated with the system subject to the first-class constraints (54). Next, we show that the physical observables of the system possessing the constraints (54) and the ones corresponding to the original reducible theory coincide. In this light, we remark that the surface (54) can be inferred in a trivial manner from (1) by adding the canonical pairs $`(y^{a_{2k+1}},\pi _{a_{2k+1}})_{k=0,\mathrm{},a}`$ and requiring that their momenta vanish. Thus, the observables of the original redundant theory are unaffected by the introduction of the new canonical pairs. In fact, the difference between an observable $`F`$ of the system subject to the constraints (54) and one of the original theory, $`\overline{F}`$, is of the type $`F\overline{F}=_{k=0}^af^{a_{2k+1}}\pi _{a_{2k+1}}`$. As any two observables that differ through a combination of first-class constraint functions can be identified, we find that the physical observables of the initial theory coincide with those of the system described by the constraints (54). So far, we proved that the observables of the system with the constraints (54) coincide on the one hand with those of the irreducible theory, and, on the other hand, with those of the original reducible one. In conclusion, the physical observables associated with the irreducible system also coincide with those of the starting on-shell reducible first-class theory. In turn, this result will have a strong impact at the level of the BRST analysis.
In the above we have shown that starting with an arbitrary on-shell reducible first-class Hamiltonian system displaying the standard Hamiltonian BRST symmetry $`s_R`$ we can construct a corresponding irreducible first-class theory, whose BRST symmetry $`s`$ complies with the basic requirements of the Hamiltonian BRST formalism. The previous result on the physical observables induces that the zeroth order cohomological groups of the corresponding BRST symmetries are isomorphic
$$H^0\left(s_R\right)H^0\left(s\right).$$
(58)
In addition, both symmetries are nilpotent
$$s_R^2=0=s^2.$$
(59)
Then, from the point of view of the fundamental equations of the BRST formalism, namely, the nilpotency of the BRST operator and the isomorphism between the zeroth order cohomological group of the BRST differential and the algebra of physical observables, it follows that it is permissible to replace the Hamiltonian BRST symmetry of the original on-shell $`L`$-stage reducible system with that of the irreducible theory. Thus, we can substitute the path integral of the reducible system in the Hamiltonian BRST approach by that of the irreducible theory.
However, it would be convenient to infer a covariant path integral with respect to the irreducible system. The present phase-space coordinates may not ensure the covariance. For instance, if we analyze the gauge transformations of the extended action of the irreducible system, we remark that those corresponding to the Lagrange multipliers of the constraint functions $`\gamma _{a_0}`$ will not involve the term $`Z_{a_1}^{a_0}ϵ^{a_1}`$, which is present within the reducible context with respect to the constraint functions $`G_{a_0}`$. For all known models, the presence of this term is essential in arriving at some covariant gauge transformations at the Lagrangian level. For this reason it is necessary to gain such a term also within the irreducible setting. Moreover, it is possible that some of the newly added canonical variables lack covariant Lagrangian gauge transformations. This signalizes that we need to add more phase-space variables to be constrained in an appropriate manner. In view of this, we introduce the additional bosonic canonical pairs
$$(y^{(1)a_{2k+1}},\pi _{a_{2k+1}}^{(1)}),(y^{(2)a_{2k+1}},\pi _{a_{2k+1}}^{(2)}),k=0,\mathrm{},a,$$
(60)
$$(y^{a_{2k}},\pi _{a_{2k}}),k=1,\mathrm{},b,$$
(61)
subject to some constraints of the type
$$\gamma _{a_{2k+1}}^{(1)}\pi _{a_{2k+1}}\pi _{a_{2k+1}}^{(1)}0,\gamma _{a_{2k+1}}^{(2)}\pi _{a_{2k+1}}^{(2)}0,k=0,\mathrm{},a,$$
(62)
$$\gamma _{a_{2k}}\pi _{a_{2k}}0,k=1,\mathrm{},b.$$
(63)
In this manner we do not affect in any way the properties of the irreducible theory as (6263) still form together with (3839) an irreducible first-class set. The newly added constraints implies the introduction of some supplementary ghosts and antighosts and the extension of the action of the BRST operator on them in the usual manner. Then, there exists a consistent Hamiltonian BRST symmetry with respect to the new irreducible theory, described by the constraints (3839) and (6263). Now, if we choose the first-class Hamiltonian with respect to the above first-class constraints in an adequate manner, we can in principle generate a gauge algebra that leads to some covariant Lagrangian gauge transformations. From (44) it results that the former set of constraints in (62) reduces to $`\pi _{a_{2k+1}}^{(1)}0`$. Thus, we observe that the surface (3839), (6263) results in a trivial way from (3839) by adding the canonical variables (6061) and demanding that their momenta vanish. Then, the difference between an observable $`F`$ of the new irreducible theory and one of the previous irreducible system, $`\overline{F}`$, is of the type $`F\overline{F}=_{k=0}^af^{a_{2k+1}}\pi _{a_{2k+1}}^{(1)}+_{k=0}^ag^{a_{2k+1}}\pi _{a_{2k+1}}^{(2)}+_{k=1}^bh^{a_{2k}}\pi _{a_{2k}}`$, hence $`F`$ and $`\overline{F}`$ can be identified. Therefore, the physical observables corresponding to the two irreducible systems coincide, such that the supplementary constraints (6263) do not afflict the previously established equivalence with the physical observables of the original redundant theory. In consequence, we can replace the Hamiltonian BRST symmetry of the original reducible system with that of the latter irreducible theory, and similarly with regard to the associated path integrals.
With all these elements at hand, the Hamiltonian BRST quantization of the irreducible theory goes along the standard lines. Defining a canonical action for $`s`$ in the usual way as $`s=[,\mathrm{\Omega }]`$, with $`\mathrm{\Omega }`$ the canonical generator (the BRST charge), the nilpotency of $`s`$ implies that $`\mathrm{\Omega }`$ should satisfy the equation
$$[\mathrm{\Omega },\mathrm{\Omega }]=0.$$
(64)
The existence of the solution to the equation (64) is then guaranteed by the acyclicity of the irreducible Koszul-Tate differential at positive antighost numbers. Once constructed the BRST charge, we pass to the construction of the BRST-invariant Hamiltonian, $`H_B=H^{}+\mathrm{`}\mathrm{`}\mathrm{more}\mathrm{"}`$, that satisfies $`[H_B,\mathrm{\Omega }]=0`$, where $`H^{}`$ stands for the first-class Hamiltonian with respect to the constraints (3839) and (6263). In order to fix the gauge, we have to choose a gauge-fixing fermion $`K`$ that implements some irreducible gauge conditions. The possibility to construct some irreducible gauge conditions is facilitated by the introduction of the pairs $`(y,\pi )`$, which, at this level, play the same role like the auxiliary variables in the non-minimal approach. Actually, these variables are more relevant than the corresponding non-minimal ones appearing in the gauge-fixing process from the reducible case because they prevent the appearance of the reducibility (via the irreducible first-class constraints), while the non-minimal coordinates in the reducible situation are mainly an effect of the redundancy. The gauge-fixed Hamiltonian $`H_K=H_B+[K,\mathrm{\Omega }]`$ will then produce a correct gauge-fixed action $`S_K`$ with respect to the irreducible theory. In this way, we showed how a reducible first-class Hamiltonian system can be approach along an irreducible BRST procedure, hence without using either ghosts for ghosts or their antighosts.
Finally, a word of caution. The appearance of the inverse of the matrices $`D_{b_k}^{a_k}`$ in various formulas implies that our approach may have problems with locality. On the other hand, locality is required for all concrete applications in field theory. However, taking into account the ambiguities present in the choice of the reducibility functions and of the functions $`A_{a_k}^{a_{k+1}}`$ (see (43)), it might be possible to obtain a local formulation. This completes our treatment.
## 4 Example
In this section we illustrate the general formalism exposed in the above in the case of a second-stage reducible model. We start with the Lagrangian action
$$S_0^L\left[A^{\mu \nu \rho }\right]=d^7x\left(\frac{1}{48}F_{\mu \nu \rho \lambda }F^{\mu \nu \rho \lambda }+\xi \epsilon _{\mu \nu \rho \lambda \alpha \beta \gamma }F^{\mu \nu \rho \lambda }A^{\alpha \beta \gamma }\right),$$
(65)
where
$$F_{\mu \nu \rho \lambda }=_\mu A_{\nu \rho \lambda }_\nu A_{\mu \rho \lambda }+_\rho A_{\lambda \mu \nu }_\lambda A_{\rho \mu \nu }_{[\mu }A_{\nu \rho \lambda ]},$$
(66)
$`\epsilon _{\mu \nu \rho \lambda \alpha \beta \gamma }`$ represents the completely antisymmetric seven-dimensional symbol and $`\xi `$ is a constant. From the canonical analysis of action (65) we deduce the first-class constraints
$$\gamma _{ij}^{(1)}\pi _{0ij}0,$$
(67)
$$G_{ij}^{(2)}3\left(^k\pi _{kij}+\xi \epsilon _{0ijklmn}F^{klmn}\right)0,$$
(68)
and the first-class Hamiltonian
$$H=d^6x\left(\frac{1}{48}F_{ijkl}^23\left(\pi _{ijk}4\xi \epsilon _{0ijklmn}A^{lmn}\right)^2+A^{0ij}G_{ij}^{(2)}\right),$$
(69)
where $`\pi _{\mu \nu \rho }`$ stand for the canonical momenta conjugated with $`A^{\mu \nu \rho }`$. The notation $`F_{ijkl}^2`$ signifies $`F_{ijkl}F^{ijkl}`$, and similarly for the other square. The constraint functions in (68) are second-stage reducible, the first-, respectively, second-stage reducibility relations being expressed by
$$Z_k^{ij}G_{ij}^{(2)}=0,Z^kZ_k^{ij}=0,$$
(70)
where the reducibility functions have the form
$$Z_{a_1}^{a_0}Z_k^{ij}=^{[i}\delta _k^{j]},Z_{a_2}^{a_1}Z^k=^k.$$
(71)
Acting like in section 2, we enlarge the original phase-space with the bosonic canonical pairs $`(y^{a_1},\pi _{a_1})(H^i,\pi _i)`$ and construct an irreducible first-class constraint set corresponding to (68) of the type
$$\gamma _{a_0}\gamma _{ij}^{(2)}3\left(^k\pi _{kij}+\xi \epsilon _{0ijklmn}F^{klmn}\right)_{[i}\pi _{j]}0,$$
(72)
$$\gamma _{a_2}\gamma ^{(2)}^i\pi _i0.$$
(73)
In inferring the above irreducible constraints we took the functions of the type $`A_{a_k}^{a_{k+1}}`$ like
$$A_{a_0}^{a_1}A_{ij}^k=_{[i}\delta _{j]}^k,A_{a_1}^{a_2}A_k=_k,$$
(74)
such that they comply with the requirements from the general theory. As we mentioned in section 3, in order to generate a covariant gauge-fixed action as a result of the Hamiltonian BRST quantization of the irreducible model it is still necessary to enlarge the phase-space with some canonical pairs subject to additional constraints such that the overall constraint set is first-class and irreducible. In this light, we introduce the canonical pairs $`(H^0,\pi _0)`$, $`(H^{(1)i},\pi _i^{(1)})`$, $`(H^{(2)i},\pi _i^{(2)})`$ constrained like
$$\gamma ^{(1)}\pi _00,\gamma _i^{(1)}\pi _i\pi _i^{(1)}0,\gamma _i^{(2)}2\pi _i^{(2)}0.$$
(75)
Obviously, the constraint set (67), (7273), (75) is irreducible and first-class. We take the first-class Hamiltonian with respect to this set under the form
$`H^{}={\displaystyle }d^6x({\displaystyle \frac{1}{48}}F_{ijkl}^23(\pi _{ijk}4\xi \epsilon _{0ijklmn}A^{lmn})^2+A^{0ij}\gamma _{ij}^{(2)}+`$
$`H^0\gamma ^{(2)}+H^i\pi _i^{(2)}+H^{(2)i}(^j\gamma _{ji}^{(2)}+_i\gamma ^{(2)})),`$ (76)
such that the irreducible gauge algebra will lead to a covariant path integral.
Next, we approach the Hamiltonian BRST quantization of the irreducible model. In view of this, we add the minimal fermionic canonical pairs ghost-antighost $`(\eta ^{(\mathrm{\Delta })ij},𝒫_{ij}^{(\mathrm{\Delta })})`$, $`(\eta ^{(\mathrm{\Delta })i},𝒫_i^{(\mathrm{\Delta })})`$, $`(\eta ^{(\mathrm{\Delta })},𝒫^{(\mathrm{\Delta })})`$, with $`\mathrm{\Delta }=1,2`$, associated with the corresponding constraints in (67), (7273) and (75). All the ghosts possess ghost number one, while their antighosts display ghost number minus one. Moreover, we take a non-minimal sector organized as $`(P_b^{ij},b_{ij})`$, $`(P_b,b)`$, $`(P_{b^1}^{ij},b_{ij}^1)`$, $`(P_{b^1},b^1)`$, $`(P_{\overline{\eta }}^{ij},\overline{\eta }_{ij})`$, $`(P_{\overline{\eta }},\overline{\eta })`$, $`(P_{\overline{\eta }^1}^{ij},\overline{\eta }_{ij}^1)`$, $`(P_{\overline{\eta }^1},\overline{\eta }^1)`$. The first four sets of non-minimal variables are bosonic with ghost number zero, while the others are fermionic, with the $`P_{\overline{\eta }}`$’s and $`\overline{\eta }`$’s of ghost number one, respectively, minus one. Consequently, the non-minimal BRST charge reads as
$`\mathrm{\Omega }={\displaystyle }d^6x({\displaystyle \underset{\mathrm{\Delta }=1}{\overset{2}{}}}(\gamma _{ij}^{(\mathrm{\Delta })}\eta ^{(\mathrm{\Delta })ij}+\gamma _i^{(\mathrm{\Delta })}\eta ^{(\mathrm{\Delta })i}+\gamma ^{(\mathrm{\Delta })}\eta ^{(\mathrm{\Delta })})+`$
$`P_{\overline{\eta }}^{ij}b_{ij}+P_{\overline{\eta }}b+P_{\overline{\eta }^1}^{ij}b_{ij}^1+P_{\overline{\eta }^1}b^1).`$ (77)
The BRST-invariant Hamiltonian corresponding to (4) is given by
$`H_B={\displaystyle }d^6x(\eta ^{(1)ij}𝒫_{ij}^{(2)}+\eta ^{(1)}𝒫^{(2)}{\displaystyle \frac{1}{2}}\eta ^{(1)i}𝒫_i^{(2)}+{\displaystyle \frac{1}{2}}\eta ^{(2)ij}_{[i}𝒫_{j]}^{(2)}+`$
$`{\displaystyle \frac{1}{2}}\eta ^{(2)}^i𝒫_i^{(2)}2\eta ^{(2)i}(^j𝒫_{ji}^{(2)}+_i𝒫^{(2)}))+H^{}.`$ (78)
In order to determine the gauge-fixed action, we work with the gauge-fixing fermion
$`K={\displaystyle }d^6x(𝒫_{ij}^{(1)}(_kA^{kij}+{\displaystyle \frac{1}{2}}^{[i}H^{(1)j]})+𝒫^{(1)}_iH^{(1)i}+`$
$`𝒫_i^{(1)}(2_jA^{ji0}+^iH^0)+P_{b^1}^{ij}(𝒫_{ij}^{(1)}\overline{\eta }_{ij}+\underset{ij}{\overset{1}{\stackrel{.}{\overline{\eta }}}})+`$
$`P_{b^1}(𝒫^{(1)}\overline{\eta }+\stackrel{1}{\stackrel{.}{\overline{\eta }}})+P_b^{ij}(\overline{\eta }_{ij}^1+\underset{ij}{\overset{.}{\overline{\eta }}})+P_b(\overline{\eta }^1+\stackrel{.}{\overline{\eta }})).`$ (79)
It is clear that the above gauge-fixing fermion implements some irreducible canonical gauge conditions. Introducing the expression of the gauge-fixed Hamiltonian $`H_K=H_B+[K,\mathrm{\Omega }]`$ in the gauge-fixed action and eliminating some auxiliary variables on their equations of motion, we finally find the path integral
$$Z_K=𝒟A^{\mu \nu \rho }𝒟H^{(1)\mu }𝒟b_{\mu \nu }𝒟b𝒟\eta ^{(2)\mu \nu }𝒟\eta ^{(2)}𝒟\overline{\eta }_{\mu \nu }𝒟\overline{\eta }\mathrm{exp}\left(iS_K\right),$$
(80)
where
$`S_K=S_0^L\left[A^{\mu \nu \rho }\right]+{\displaystyle }d^7x(b_{\mu \nu }(_\rho A^{\rho \mu \nu }+{\displaystyle \frac{1}{2}}^{[\mu }H^{(1)\nu ]})+`$
$`b_\mu H^{(1)\mu }\overline{\eta }_{\mu \nu }\mathrm{}\eta ^{(2)\mu \nu }\overline{\eta }\mathrm{}\eta ^{(2)}).`$ (81)
In deriving (4) we performed the identifications
$$H^{(1)\mu }=(H^0,H^{(1)i}),b_{\mu \nu }=(\pi _i^{(1)},b_{ij}),$$
(82)
$$\overline{\eta }_{\mu \nu }=(𝒫_i^{(1)},\overline{\eta }_{ij}),\eta ^{(2)\mu \nu }=(\eta ^{(2)i},\eta ^{(2)ij}),$$
(83)
and used the symbol $`\mathrm{}=_\rho ^\rho `$. It is clear that the gauge-fixed action (4) is covariant and displays no residual gauge invariances although we have not used any ghosts for ghosts. It is precisely the introduction of the supplementary canonical pairs constrained accordingly to (75) and the choice (4) of the first-class Hamiltonian that generates a Hamiltonian gauge algebra implementing the covariance. Actually, the gauge-fixed action (4) can be alternatively inferred in the framework of the antifield BRST treatment if we start with the action
$$S_0^L[A^{\mu \nu \rho },H^{(1)\mu }]=S_0^L\left[A^{\mu \nu \rho }\right],$$
(84)
subject to the irreducible gauge transformations
$$\delta _ϵA^{\mu \nu \rho }=^{[\mu }ϵ^{\nu \rho ]},\delta _ϵH^{(1)\mu }=2_\nu ϵ^{\nu \mu }+^\mu ϵ,$$
(85)
and employ an adequate non-minimal sector and gauge-fixing fermion. In fact, the irreducible and covariant gauge transformations (85) are the Lagrangian result of our irreducible Hamiltonian procedure inferred via the extended and total actions and their gauge transformations.
In the end, let us briefly compare the above results with those derived in the standard reducible BRST approach. The gauge-fixed action in the usual reducible BRST framework can be brought to the form
$`S_\psi =S_0^L\left[A^{\mu \nu \rho }\right]+{\displaystyle }d^7x(_{\mu \nu }(_\rho A^{\rho \mu \nu }+{\displaystyle \frac{1}{2}}^{[\mu }\phi ^{\nu ]})+`$
$`_\mu \phi ^\mu \overline{C}_{\mu \nu }\mathrm{}C^{\mu \nu }\overline{C}^2\mathrm{}\overline{C}^1\overline{C}_\mu \mathrm{}C^\mu \overline{C}\mathrm{}C),`$ (86)
where $`C^{\mu \nu }`$ stand for the ghost number one ghosts, $`C^\mu `$ represent the ghost number two ghosts for ghosts, and $`C`$ is the ghost number three ghost for ghost for ghost. The remaining variables belong to the non-minimal sector and have the properties
$$ϵ\left(_{\mu \nu }\right)=0,gh\left(_{\mu \nu }\right)=0,ϵ\left(\phi ^\mu \right)=0,gh\left(\phi ^\mu \right)=0,$$
(87)
$$ϵ\left(\right)=0,gh\left(\right)=0,ϵ\left(\overline{C}_{\mu \nu }\right)=1,gh\left(\overline{C}_{\mu \nu }\right)=1,$$
(88)
$$ϵ\left(\overline{C}_\mu \right)=0,gh\left(\overline{C}_\mu \right)=2,ϵ\left(\overline{C}\right)=1,gh\left(\overline{C}\right)=3,$$
(89)
$$ϵ\left(\overline{C}^2\right)=1,gh\left(\overline{C}^2\right)=1,ϵ\left(\overline{C}^1\right)=1,gh\left(\overline{C}^1\right)=1.$$
(90)
By realizing the identifications
$$C^{\mu \nu }\eta ^{(2)\mu \nu },_{\mu \nu }b_{\mu \nu },\phi ^\mu H^{(1)\mu },$$
(91)
$$b,\overline{C}_{\mu \nu }\overline{\eta }_{\mu \nu },\overline{C}^2\overline{\eta },\overline{C}^1\eta ^{(2)},$$
(92)
among the variables involved with the gauge-fixed actions inferred within the irreducible and reducible approaches, (4), respectively, (4), the difference between the two gauge-fixed actions is given by
$$S_KS_\psi =d^7x\left(\overline{C}_\mu \mathrm{}C^\mu +\overline{C}\mathrm{}C\right).$$
(93)
We remark that $`S_KS_\psi `$ is proportional to the ghosts for ghosts $`C^\mu `$ and the ghost for ghost for ghost $`C`$, which are some essential compounds of the reducible BRST quantization. Although identified at the level of the gauge-fixed actions, the fields $`\phi ^\mu `$ and $`H^{(1)\mu }`$ play different roles within the two formalisms. More precisely, the presence of $`H^{(1)\mu }`$ within the irreducible model prevents the appearance of the reducibility, while $`\phi ^\mu `$ is nothing but one of its effects. In fact, the effect of introducing the fields $`H^{(1)\mu }`$ is that of forbidding the appearance of zero modes which exist within the original reducible theory by means of the first term present in the gauge transformations of $`H^{(1)\mu }`$ from (85). Indeed, if we take $`ϵ^{\mu \nu }=^{[\mu }ϵ^{\nu ]}`$ in (85), then $`_\mu ϵ^{\mu \nu }`$ is non-vanishing, such that the entire set of gauge transformations is irreducible. In consequence, all the quantities linked with zero modes, like the ghosts with ghost number greater than one or the pyramid-like structured non-minimal sector, are discarded when passing to the irreducible setting. This completes the analysis of the investigated model.
## 5 Conclusion
In conclusion, in this paper we have shown how systems with reducible first-class constraints can be quantized by using an irreducible Hamiltonian BRST formalism. The key point of our approach is given by the construction of a Hamiltonian Koszul-Tate complex that emphasizes an irreducible set of first-class constraints. As the physical observables associated with the irreducible and reducible versions coincide, the main equations underlining the Hamiltonian BRST formalism make legitimate the substitution of the BRST quantization of the reducible theory by that of the irreducible one. The canonical generator of the irreducible BRST symmetry exists due to the acyclicity of the irreducible Koszul-Tate differential, while the gauge-fixing procedure is facilitated by the enlargement of the phase-space with the canonical pairs of the type $`(y,\pi )`$. The general formalism is exemplified on a second-stage reducible model involving three-form gauge fields, which is then compared with the results from the standard reducible BRST approach.
## Acknowledgment
This work has been supported by a Romanian National Council for Academic Scientific Research (CNCSIS) grant.
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# The response of a dwarf nova disc to real mass transfer variations
## 1 Introduction
Cataclysmic variable systems (CV’s), in which a white dwarf accretes material from a Roche lobe filling secondary star (see Warner 1995 for an enyclopaedic review) can be broadly divided into three subclasses: in non-magnetic systems, the white dwarf accretes via an accretion disc; in the magnetic polars, the infalling matter couples directly onto the strong magnetic field of the white dwarf before it can build a disc and is funnelled to accretion region(s) near the magnetic pole(s) of the white dwarf; finally, in CV’s containing a weekly magnetic white dwarf (intermediate polars), a partial disc may exist, with the mass flowing from the inner edge of the disrupted disc through magnetically funnelled accretion curtains onto the white dwarf.
Dwarf novae are non-magnetic cataclysmic variable stars which show characteristic eruptions with photometric amplitudes in the range of 2–8 magnitudes <sup>1</sup><sup>1</sup>1Also intermediate polars are known to show outbursts in their partial discs; e.g. EX Hya, GK Per. These eruptions typically last a few days to several weeks and recur quasi-periodically on timescales of weeks up to many years. They are generally thought to result from thermal instabilities associated with hydrogen ionization in the accretion disc (see Cannizzo 1993a for a review and Ludwig et al. 1994 for a detailed parameter study).
An important boundary condition of the disc-limit-cycle model is the mass transfer rate from the secondary to the accretion disc. In most studies of dwarf nova outbursts this mass transfer rate is kept constant (Cannizzo 1993b; Ludwig & Meyer 1997; Hameury et al. 1998, 1999). Duschl and Livio (1989) were the first to examine combined mass transfer and disc outbursts, though within the framework of the mass transfer instability model where individual and short-lived mass transfer events are capable of producing single outbursts. Smak (1991, 1999) discussed mass transfer variations in the context of the superoutburst phenomenon in dwarf novae of the SU UMa type and dwarf nova outbursts with enhanced mass transfer during outburst.
King & Cannizzo (1998) and Leach et al. (1999) tested how the accretion disc in a dwarf nova system behaves if the mass transfer from the secondary varies abruptly between different levels. They found that these mass transfer variations produce only subtle effects on normal dwarf novae, including variations in the outburst shape, and that the dwarf novae keep on having outbursts even if the transfer rate is reduced close to zero.
In spite of these efforts we are left with the question of what the real variations of the mass tranfer rates from the secondaries in dwarf nova systems are and how they can influence the outburst behaviours of the accretion discs. Fortunately, nature provides an answer to the first question in the form of discless cataclysmic variables, the polars or AM Her systems. In these systems, the mass transfer rate can be estimated directly from the observations because there is no accretion disc acting as a mass buffer. Thus, we can use the long-term light curve of an AM Her system as a measure for the mass transfer variations in a fictitious but realistic dwarf nova system with the same system parameters. In the present paper, we present the results of such a numerical experiment.
In the following section, we briefly review the equations for the viscous and thermal evolution of the accretion disc, discuss the numerical methods used in our code, and compare the results of two standard models with other fine mesh computations. Thereafter we derive the mass transfer rate in AM Her as a function of time, $`\dot{M}_{\mathrm{tr}}(t)`$. Finally, we apply this mass transfer rate to a fictitious dwarf nova and discuss the effects that can be observed in the outburst behaviour.
## 2 Finite Element Methods in the context of disc evolution
The classical equation describing the viscous evolution of the surface density $`\mathrm{\Sigma }`$ in a geometrically thin, axisymmetric accretion disc is obtained by combining the vertically averaged Navier-Stokes and mass-conservation equations:
$`{\displaystyle \frac{\mathrm{\Sigma }}{t}}`$ $`=`$ $`{\displaystyle \frac{3}{R}}{\displaystyle \frac{}{R}}\left(R{}_{}{}^{\frac{1}{2}}{\displaystyle \frac{}{R}}(R^{\frac{1}{2}}\nu \mathrm{\Sigma })\right)+{\displaystyle \frac{\dot{M}_{\mathrm{tr}}(t,R)}{2\pi R}}`$ (1)
where $`\nu =(2/3)\alpha R_\mathrm{g}T/\mu \mathrm{\Omega }`$ is the kinematic viscosity, $`\alpha `$ denotes the viscosity parameter, $`R_\mathrm{g}`$ the gas constant and $`\mu `$ the mean molecular weight. The first term on the right hand side is the classical disc diffusion term and the second term describes the external mass-deposition. For the purposes of this study, we follow Cannizzo (1993b) and simply add the mass lost from the secondary $`\dot{M}_{\mathrm{tr}}`$(t) in a Gaussian distribution at a fixed radius near the outer edge.
Similarly, one can derive an energy equation (in the central temperature $`T`$:
$$\frac{T}{t}=\frac{2(HC+J)}{c_\mathrm{p}\mathrm{\Sigma }}\frac{R_\mathrm{g}T}{\mu c_\mathrm{p}}\frac{1}{R}\frac{(Rv_\mathrm{R})}{R}v_\mathrm{R}\frac{T}{R},$$
(2)
where $`c_\mathrm{p}`$ denotes the specific heat,
$$v_\mathrm{R}=\frac{3\nu }{R}\frac{\mathrm{log}(\nu \mathrm{\Sigma }r^{\frac{1}{2}})}{\mathrm{log}R}$$
is the local radial flow velocity, $`H=\frac{9}{8}\nu \mathrm{\Omega }^2\mathrm{\Sigma }`$ represents viscous heating, $`C=\sigma T_{\mathrm{eff}}^4`$ is the radiative cooling and
$$J=\frac{3}{2}c_\mathrm{p}\nu \frac{\mathrm{\Sigma }}{R}\frac{}{R}(R\frac{T}{R})$$
the radial energy flux carried by viscous processes (see Smak 1984; Mineshige & Osaki 1983; Mineshige 1986; Ichikawa & Osaki 1992 and Cannizzo 1993b, for discussions).
We solve Eqs. (1) and (2) using a combined Finite-Element / Finite-Difference algorithm (FE for the spatial part and FD for the time-evolution). Apart from our own work (Schreiber & Hessman 1998) the method of Finite Elements has not been used in this context. As this method proved to be extremely robust, it warrants a somewhat more detailed description.
The idea of FE is to divide the region of interest (the disc radii between $`R_{\mathrm{in}}`$ and $`R_{\mathrm{out}}`$) into $`n1`$ elements and to expand the function $`u(x)`$ which is supposed to solve the differential equation with suitable functions $`u(x)=_{i=1}^na_i\phi _i(x)`$ for every element. In order to get a continuous solution over all elements, the functions ($`\phi `$) of every element have to be transformed to the so-called local basis functions $`N_i`$ and the coefficients to the so-called nodes $`(c_i)`$ before collected together (Gruber & Rappaz 1985, Schwarz 1991). To solve the differential equation, the function
$$u(x)=\underset{k=1}{\overset{n}{}}c_iN_i(x)$$
(3)
has to meet the requirement formulated by Garlerkin: the integral of the residuum (which one gets by inserting Eq. (3) into the differential equation) weighted with the functions $`N_j(x)(j=1,\mathrm{},n)`$ has to vanish. This requirement, the interchange of integration and summation and partial integration lead to matrix-equations of the form
$$B\dot{c}+Ac=D,$$
(4)
with $`A=(a_{ij}),B=b_{ij}`$ and $`D=d_i`$ ($`i=1,..,n;j=1,..,n`$ for $`n`$ nodes).
To solve the differential Eqs. (1) and (2), we have to fill $`A`$, $`B`$, $`D`$ in the sense mentioned above and calculate $`c`$ from Eq. (4). After transforming to the variables $`X=2R^{\frac{1}{2}}`$ and $`S=X\mathrm{\Sigma }`$ we derive for the surface density from Eq. (1):
$`a_{ij}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}\nu _k({\displaystyle }{\displaystyle \frac{12}{X^2}}{\displaystyle \frac{N_i}{X}}N_k{\displaystyle \frac{N_j}{X}}dX`$
$`+{\displaystyle }{\displaystyle \frac{12}{X^2}}N_i{\displaystyle \frac{N_k}{X}}{\displaystyle \frac{N_j}{X}}dX)`$
$`b_{ij}`$ $`=`$ $`{\displaystyle N_iN_j𝑑X}`$
$`d_i`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{2\dot{M}_{\mathrm{tr},k}N_k}{\pi X^2}N_j𝑑X}.`$
Similarly it is easy to obtain for the central temperature Eq. (2):
$`a_{ij}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}({\displaystyle }p_k^{(1)}{\displaystyle \frac{N_i}{R}}N_k{\displaystyle \frac{N_j}{R}}dR`$
$`+{\displaystyle }p_k^{(2)}N_k{\displaystyle \frac{N_i}{R}}N_jdR+{\displaystyle }p_k^{(3)}N_iN_kN_jdR)`$
$`b_{ij}`$ $`=`$ $`{\displaystyle N_iN_j𝑑R}`$
$`d_j`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle p_k^{(4)}N_kN_j𝑑R}.`$
The coefficients $`p_k^{(i)}`$ are given by
$`p_k^{(1)}`$ $`=`$ $`3c_\mathrm{p}\nu \mathrm{\Sigma }`$
$`p_k^{(2)}`$ $`=`$ $`{\displaystyle \frac{3c_\mathrm{p}\nu \mathrm{\Sigma }}{R}}v_\mathrm{R}`$
$`p_k^{(3)}`$ $`=`$ $`{\displaystyle \frac{R_\mathrm{g}}{\mu c_\mathrm{p}R}}{\displaystyle \frac{Rv_\mathrm{R}}{R}}`$
$`p_k^{(4)}`$ $`=`$ $`{\displaystyle \frac{9/4\nu \mathrm{\Sigma }\mathrm{\Omega }^2\sigma T_{\mathrm{eff}}}{c_\mathrm{p}\mathrm{\Sigma }}}.`$
The index $`k`$ refers to the value of the quantities at node number $`k`$, which is equivalent to the radius $`R_k`$. Notice that we only used linear basis functions in this paper.
As mentioned above, Eq. (4) has to be solved to get $`c(t+\mathrm{\Delta }t)`$ from $`c(t)`$. Using simple finite differences leads to
$`c(t+\mathrm{\Delta }t)`$ $`=`$ $`c(t){\displaystyle \frac{1}{2}}\mathrm{\Delta }t(B^1Ac(t)`$ (5)
$`+B^1Ac(t+\mathrm{\Delta }t)(+B^1D)).`$
In order to test our code, we carried out two sets of calculations using the binary parameters and cooling functions from Cannizzo (1993b)
$`M_1=1M_{},\alpha _\mathrm{h}=0.1,\alpha _\mathrm{c}=0.02,`$
$`R_{\mathrm{in}}=5.0\times 10^8\mathrm{cm},R_{\mathrm{out}}=4.0\times 10^{10}\mathrm{cm},`$
$`\dot{M}_{\mathrm{tr}}=1.5\times 10^9M_{}/\mathrm{yr}`$
and Ludwig & Meyer (1997)
$`M_1=0.63M_{},\alpha _h=0.2,\alpha _c=0.04,`$
$`R_{\mathrm{in}}=8.4\times 10^8\mathrm{cm},R_{\mathrm{out}}=1.7\times 10^{10}\mathrm{cm},`$
$`\dot{M}_{\mathrm{tr}}=5\times 10^{15}\mathrm{g}/\mathrm{s}.`$
The resulting light curves are shown in Figs. 1 and 2. Our code reproduces the sequence of only relatively long outbursts found by Cannizzo (1993b) for the parameter of SS Cygni as well as the sequence of one long outburst followed by two short outbursts found by Ludwig & Meyer (1997) to describe VW Hydri. The short outbursts arise when there is not enough mass stored in the disc and therefore the heating wave gets reflected before it has reached the outer edge of the disc.
We find that at least $`100`$ nodes are necessary for long-term convergence. The outburst and quiescence duration decreases with an increasing number of nodes because - with finer zoning - the length of time spent on the viscous plateau becomes shorter (Cannizzo 1993b). This effect is smaller in Fig. 2 because the disc is smaller in this system.
We conclude that our FE-code produces results which are in excellent agreement with those of other fine-mesh computations.
## 3 The mass loss rate of the secondary star in AM Herculis
The strong magnetic field of the white dwarf primary in polars prevents the formation of an accretion disc. Without an accretion disc acting as a buffer for the transferred mass, the mass loss rate from the secondary equals the mass accretion rate on the white dwarf at every moment, $`\dot{M}_{\mathrm{acc}}=\dot{M}_{\mathrm{tr}}`$ (the free-fall time is $`1`$ h). As an observational consequence, any variation in rate at which the secondary star loses mass through the $`L_1`$ point will result in a quasi-immediate change of the observed accretion luminosity. The brightest polar, AM Her, has been intensely monitored at optical wavelengths by observers of the AAVSO for more than 20 years and shows an irregular long-term variability, switching back and forth between high- and low states of accretion on timescales of days to months.
The problem of deriving the mass loss history of the secondary star is then equivalent to that of determining the accretion luminosity $`L_{\mathrm{acc}}`$ as a function of time. As the bulk of the accretion luminosity is emitted in the X-ray regime, a bolometric correction relating the densly monitored optical magnitude to the total luminosity has to be derived. This approach has been followed in detail by Hessman et al. (A&A, submitted) using X-ray observations obtained at multiple epochs. We summarize only briefly the results here. The accretion luminosity is computed from the observed accretion-induced flux $`F_{\mathrm{acc}}`$ as
$$L_{\mathrm{acc}}(t)=4\pi d^2F_{\mathrm{acc}}(t)$$
(6)
with $`d=90pc`$ (Gänsicke et al. 1995). We further use $`F_{\mathrm{acc}}F_{\mathrm{SX}}+3\times F_{\mathrm{HX}}`$ with $`F_{\mathrm{SX}}`$ and $`F_{\mathrm{HX}}`$ the observed soft and hard X-ray fluxes respectively. The factor three accounts for the additional cyclotron radiation emitted from the accretion column and for the thermal reprocession of bremsstrahlung and cyclotron radiation intercepted by the white dwarf and emitted in the ultraviolet (Gänsicke et al. 1995). The mass loss ( = transfer) rate is then
$$\dot{M}_{\mathrm{tr}}(t)=\frac{L_{\mathrm{acc}}(t)R_{\mathrm{wd}}}{GM_{\mathrm{wd}}}$$
(7)
where $`G`$ is the gravitational constant and $`R_{\mathrm{wd}}`$ and $`M_{\mathrm{wd}}`$ are the white dwarf radius and mass, respectively. As the actual properties of the white dwarf in AM Her are still the subject of controversial discussions (Gänsicke et al. 1998, Cropper et al. 1999), we use the parameters of an average white dwarf, $`M_{\mathrm{wd}}=0.6M_{}`$ and $`R_{\mathrm{wd}}=8.4\times 10^8\mathrm{cm}`$. $`\dot{M}_{\mathrm{tr}}(t)`$ is shown in Fig. 3. The average value of the mass transfer rate in AM Her is $`\dot{M}_{\mathrm{av}}=7.88\times 10^{15}\mathrm{g}\mathrm{s}^1=1.24\times 10^{10}M_{}\mathrm{yr}^1`$. The derived mass transfer rates of AM Her are in general agreement with results published in the literature. For example, Beuermann & Burwitz (1995) found transfer rates between $`0.8`$ and $`2.0\times \mathrm{\hspace{0.17em}10}^{10}\dot{M_{}}\mathrm{yr}^1`$ and Greeley et al. (1999) estimated a mass transfer rate of $`2\times \mathrm{\hspace{0.17em}10}^{16}\mathrm{g}\mathrm{s}^1`$ for the high state of AM Her from far ultraviolet spectra.
## 4 Results
### 4.1 The fictitious dwarf nova
We devised a fictitious dwarf nova with a non-magnetic primary of mass $`M_{\mathrm{wd}}=0.6M_{}`$ and an orbital period of $`P=3.08\mathrm{hr}`$, i.e. a non-magnetic twin of AM Her. For these binary parameters, we obtain $`R_{\mathrm{in}}R_{\mathrm{wd}}=8.4\times 10^8`$ cm for the inner disc radius and $`R_{\mathrm{out}}=2.2\times 10^{10}`$ cm for the outer edge of the disc. We then used our FE code to follow the structure of the accretion disc in our fictitious dwarf nova for 7000 d, applying the variable mass transfer rate $`\dot{M}_{\mathrm{tr}}(t)`$ derived above and standard viscosity parameters $`\alpha _\mathrm{h}=0.2`$ and $`\alpha _\mathrm{c}=0.04`$.
We show $`500`$ day-long samples of our calculations in Figs. 4–8. In each figure, the top panel show the mass transfer rate as a function of time (solid line) and the average mass transfer rate $`\mathrm{log}(\dot{M}_{\mathrm{av}}[\mathrm{g}\mathrm{s}^1])=15.90`$ (dotted line). The panel below displays the disc mass $`M_{\mathrm{disc}}`$ normalized with the averaged disc mass $`\overline{M}_{\mathrm{disc}}=1.64\times \mathrm{\hspace{0.17em}10}^{23}\mathrm{g}`$. The two lower panels display the light curves calculated with the varying mass transfer rate and the light curves calculated with the constant average mass transfer rate respectively. For the constant average mass transfer rate the disc goes through a $`\mathrm{\hspace{0.17em}60}`$ day-long cycle including one long outburst followed by two short outbursts. The long outbursts are those in which the entire disc is transformed into the hot state while the short outbursts arise when the outward moving heating wave is reflected as a cooling wave before it has reached the outer edge of the disc.
Our numerical experiment clearly demonstrates that the outburst light curve of the fictitious system is strongly affected by the variations of the mass transfer rate. Even in the case of relatively small fluctuations effects on the outburst behaviour ($`16.0<\mathrm{log}(\dot{M}_{\mathrm{tr}}[\mathrm{g}\mathrm{s}^1])<\mathrm{\hspace{0.17em}16.4}`$) are clearly present in the light curve. In Fig. 4, the mass transfer rate is always high during the 500 days but the disc switches between an accretion state with only long outbursts and states where one or two short outbursts follow a long outburst. When the tranfer rate decreases somewhat ($``$ day 200), the disc does not save enough mass to create consecutive long outbursts until the mass transfer rate increases again ($``$ day 350).
In addition to this effect, our experiment shows that a sharp decrease in the mass transfer rate instantaneously changes the outburst behaviour of the accretion disc. In Fig. 5 the disc first behaves as in Fig. 4 but when the transfer rate drops sharply (day 4340), the long outbursts immediately vanish and the duration of the quiescent phase increases somewhat.
Another remarkable point is that even during relatively long periods of very low transfer rates (Fig. 6, days 4700– 5000) the disc does not stop its outburst activity but produces only short outbursts with slowly decreasing amplitudes and increasing quiescence intervals. This confirms the findings of King & Cannizzo (1998).
Fig. 7 shows 500 days of our simulation in which the adopted mass transfer from the secondary varies strongly on a timescale of roughly 20 days. This is the most frequent case during our calculation but there are rare periods of nearly constant mass transfer. Fig. 8 shows the light curve of the fictitious system during 500 days in which the mass transfer rate nearly exactly equals its averaged value (top). As a result the light curves computed with the real mass transfer rate and the averaged value look equal.
In summary, one can say that the variations of the mass transfer rate leads the disc to switch between three states in which only long outbursts occur ($`\mathrm{log}(\dot{M}_{\mathrm{tr}})16.3`$), one long outburst is followed by one or two short outbursts ($`16.3>\mathrm{log}(\dot{M}_{\mathrm{tr}})>15.7`$), and only short outbursts occur ($`\mathrm{log}(\dot{M}_{\mathrm{tr}})15.7`$).
In order to understand the described behaviour of the disc we take into account the viscous timescale $`t_\mathrm{v}R^2/\nu `$ which gives an estimate of the timescale for a disc annulus to move a radial distance $`R`$. For the quiescent state $`t_{\mathrm{v},\mathrm{c}}`$ and the outburst state $`t_{\mathrm{v},\mathrm{h}}`$ and with $`R=R_{\mathrm{out}}`$ (where the mass transferred from the secondary is added to the disc) we obtain for the viscous timescale
$`t_{\mathrm{v},\mathrm{c}}{\displaystyle \frac{R^2}{\nu }}\mathrm{\hspace{0.17em}2000}\mathrm{d},`$ (8)
$`t_{\mathrm{v},\mathrm{h}}{\displaystyle \frac{R^2}{\nu }}\mathrm{\hspace{0.17em}15}\mathrm{d}.`$ (9)
The mass added to the disc during quiecence is stored in the disc because it moves inward on the long timescale $`t_{\mathrm{v},\mathrm{c}}`$ whereas during an outburst even mass from outer regions can reach the white dwarf within a viscous timescale. This makes it possible that the mass accreted onto the white dwarf during a long outburst can be up to roughly one third of the disc mass ($`\mathrm{day}\mathrm{\hspace{0.17em}4220}`$). Therefore the disc can relax to equillibrium with the mass transfer rate in only one outburst in the case of high transfer rates (Fig. 4).
The prompt response of the disc to the sharp decline in the mass transfer rate (Fig. 5) can be understood in the same way: due to the short viscous time ($`t_{\mathrm{v},\mathrm{h}}`$) the disc accretes a substantial fraction of the disc mass ($`\mathrm{\hspace{0.17em}1}/4M_{\mathrm{disc}}`$) during the last long outburst which immediately prevents long outbursts when the mass transfer rate becomes low ($`\mathrm{day}\mathrm{\hspace{0.17em}4340}`$).
Finally, the long period of low transfer rates (Fig. 6) do not prevent outbursts because the mass accreted during the short outbursts is only a few percent of the disc mass. Therefore the disc relaxes to the mass transfer rate on a longer timescale.
### 4.2 The mass transfer and mass accretion rates
An important point in understanding the physics of accreting binaries is to know how far the outburst behaviour and hence the resulting light curves depend on real variations of the mass tranfer rate.
To answer this question, we compare the averaged mass accretion rate onto the white dwarf with the mass transfer rate. Figure 9 shows that the time in which the disc relaxes to an equilibrium with the mass transfer rate depends on the occurence of long outbursts: when the accretion rate is averaged over 20 days (dotted line in Fig. 9) the mass transfer and the accretion rate correspond roughly only during periods where only long outbursts occur (days 40 – 150, see also Fig. 4) but for the periods where the disc goes through a cycle of short and long outbursts the accretion rate has to be averaged over 60 days to match the mass transfer rate (days 150 – 500). If the mass transfer rate drops steeply and stays in a low–state (days 4700 – 5000 in Fig. 10, see also Fig. 6), the accretion rate needs more than 60 days to follow this behaviour, because the disc produces only short outbursts in which only a small percentage of the disc mass is involved.
In order to make this plausible we give a relaxation timescale $`t_\mathrm{r}`$ as the ratio of the viscous timescale in outburst with the relative mass fraction accreted during an outburst:
$$t_\mathrm{r}t_{\mathrm{v},\mathrm{h}}\frac{M_{\mathrm{disc}}}{\mathrm{\Delta }M_{\mathrm{disc}}}$$
(10)
For high mass transfer rates this timescale is around $`\mathrm{\hspace{0.17em}70}\mathrm{days}`$ and so the correspondence of the averaged accretion rate and the mass transfer rate in Fig. 9. is not surprising. In the case of low transfer rates (Fig. 10, day 4700–5000) this timescale is longer ($`\mathrm{\hspace{0.17em}300}\mathrm{days}`$) because only roughly 5 percent of the disc mass are accreted during an outburst.
### 4.3 Dependence on the primary mass
In the numerical experiment above, we have assumed an average white dwarf mass for the primary in AM Her. The literature holds a large spectrum of white dwarf mass estimates for AM Her, $`M_{\mathrm{wd}}=0.39M_{\mathrm{}}`$ (Young & Schneider 1981), $`M_{\mathrm{wd}}=0.69M_{\mathrm{}}`$ (Wu et al.1995), $`M_{\mathrm{wd}}=0.75M_{\mathrm{}}`$ (Mukai & Charles 1987), $`M_{\mathrm{wd}}=0.91M_{\mathrm{}}`$ (Mouchet 1993) and $`M_{\mathrm{wd}}=1.22M_{\mathrm{}}`$ (Cropper et al. 1998). Based on the observed ultraviolet spectrum of AM Her and on its well-established distance, Gänsicke et al. (1998) estimated $`R_{\mathrm{wd}}1.1\times 10^9`$ cm, and, using the Hamada-Salpeter (1961) mass-radius relation for carbon cores, $`0.35M_{\mathrm{}}<M_{\mathrm{wd}}<\mathrm{\hspace{0.17em}0.53}M_{\mathrm{}}`$. As the Hamada-Salpeter mass-radius relation is valid for cold white dwarfs, the finite temperature $`\mathrm{20\hspace{0.17em}000}`$ K of the white dwarf in AM Her would allow also somewhat higher masses, $`M_{\mathrm{wd}}0.65M_{\mathrm{}}`$, which is very close to the average mass of field white dwarfs, $`0.6M_{\mathrm{}}`$, that we used.
Even though we exclude a massive white dwarf based on the observational evidences, we repeated our simulation with $`M_{\mathrm{wd}}=1.0M_{\mathrm{}}`$ in order to test the sensitivity of our results.
In a first step, we recompute $`\dot{M}_{\mathrm{tr}}(t)`$ from Eq. (7) with $`M_{\mathrm{wd}}=1.0M_{\mathrm{}}`$ and a corresponding $`R_{\mathrm{wd}}=5.4\times 10^8`$ cm. The resulting mean accretion rate is $`3.0\times 10^{15}\mathrm{g}\mathrm{s}^1`$, a factor 2.6 lower than before. Then, we simulated once more 7000 days of disc evolution with the new $`\dot{M}_{\mathrm{tr}}(t),\alpha _\mathrm{h}=0.1,\alpha _\mathrm{c}=0.02`$ and $`R_{\mathrm{out}}=2.8\times 10^{10}\mathrm{cm}`$.
In Fig. 11 we show 500 days of our calculation with $`M_{\mathrm{wd}}=1.0M_{\mathrm{}}`$. The disc produces only short outbursts and the outburst cycle of four outbursts with decreasing amplitude is hardly changed even by drastic variations of the mass transfer rate (day 4700).
The different responses that our fictitious dwarf novae with $`0.6M_{\mathrm{}}`$ and $`1.0M_{\mathrm{}}`$ white dwarfs show to the variable mass transfer rate are easy to understand: both the increased primary mass and the decreased radius of the white dwarf reduce - as mentioned above - the derived average mass transfer rate.
In addition, the outer disc radius of the fictitious system with $`M_{\mathrm{wd}}=1.0M_{\mathrm{}}`$ increases. The disc becomes more massive, i.e. $`\overline{M}_{\mathrm{disc}}=4.59\times 10^{23}\mathrm{g}`$, because of its increased size.
Due to the reduced mass transfer rate and the increased disc size, the heating waves are able to reach the outer edge of the disc only during the first outburst of the outburst cycle and only in cases where the mass transfer rate is high ($`\mathrm{log}\dot{M}_{\mathrm{tr}}\mathrm{\hspace{0.17em}16.1}`$). Even in this rare situation the disc stays only a few days in the hot state and accretes only a small fraction of the disc mass ($`1/8M_{\mathrm{disc}}`$). For lower transfer rates (the extremly more frequent case shown in Fig. 11) the heating front gets reflected before it has reached the outer edge of the disc. Hence, only a small percentage of the disc mass is involved in any outburst. Therefore, and due to the longer viscous timescale, the relaxation timescale given in Eq. (10) is always larger than a few years.
Summing up, the adopted primary mass and the average accretion rate play an important role on the influence that the variable mass transfer rate has on the outburst behaviour.
## 5 A note on the disc limit-cycle model
A number of new features have been added to the well known limit-cycle model (Meyer & Meyer-Hofmeister 1981; Smak 1982; Cannizzo, Gosh & Wheeler 1982). Hameury et al. (1998, 1999, 2000) have shown how the light curves change if the disc size is allowed to vary and which effects irradiation has on the outburst behaviour. Additionally these authors discussed other sources of uncertainties such as the tidal torque and the evaporation of the inner parts of the disc (Meyer & Meyer-Hofmeister 1994). They concluded that all these effects must be included in order to obtain meaningful physical information on e.g. the viscosity from the comparison of predicted and observed light curves. Another point of importance is the interaction of the accretion stream leaving the secondary star and the accretion disc. Schreiber & Hessman (1998) tested the influence of stream overflow on the disc evolution in dwarf novae. They found that significant stream overflow can lead to reversion of the inward-travelling cooling front and create an outward-travelling heating front. This behaviour would produce small dips in the light curve during the declining phase.
The model we used here contains neither irradiation nor allows it the disc radi to vary. The stream mass is simply deposited in a small Gaussian distribution near the outer edge of the disc. Therefore we do not attempt at present a comparison with observed light curves. Nevertheless, all calculated light curves with constant mass transfer from the secondary relax in a quasi-stationary outburst cycle (of one or more outbursts) which repeat periodically. Our model shows in a qualitative way that real mass transfer variations may have a dominant influence on the outburst behaviour at least in systems with relatively small discs and strong mass transfer.
## 6 Discussion and conclusions
There is no reason to assume that the mass transfer variations of the secondary observed in AM Her are not present in non-magnetic systems. Our numerical experiment including realistic variations of the mass transfer rate in a dwarf nova system is, therefore, a significant step towards a better understanding of dwarf nova light curves, and, thereby, of the underlying disc limit cycle.
The light curve produced by our fictitious system switches between three states depending on the actual mass transfer rate. High transfer rates lead to only long outbursts where the entire disc is transformed into the hot state. If the transfer rate is near the average value the disc goes through a cycle of three outbursts, one long outburst followed by two short ones. Even long periods of low transfer rates do not force the disc to stop its outburst activity: long outbursts are suppressed and the duration of quiescence increases but the disc always produces short outbursts. From this follows that the low-states of VY Sculptoris stars (a subgroup of novalike variables) could not be caused by low transfer rates alone (see also Leach et al. 1999).
We find that in our fictitious system the mass accreted during an outburst cycle is dominated by the course of the mass transfer rate if the mass transfer rate varies significantly. The disc always relaxes to equilibrium with the mass input from the secondary. Thus, our experiment strongly supports King & Cannizzo’s (1998) claim that dwarf nova accretion discs are probably never in a stationary state but are constantly adjusting to the prevailing value of $`\dot{M}_{\mathrm{tr}}`$. Only during periods where the mass transfer rate is nearly exactly constant the disc periodically repeat the quasy stationary outburst cycle. Such periods are rare but occur in AM Her.
The strong influence of the mass transfer rate on the outburst behaviour of the fictitious system clearly indicates that probably most (if not all) the deviations from periodic outburst cycles seen in the light curves of dwarf novae are caused by variations of the mass transfer rate.
###### Acknowledgements.
We thank Daisaku Nagami and the referee for their helpful comments and suggestions. MRS would like to thank the Deutsche Forschungsgemeinschaft for financial support (Ma 1545 2-1). BTG thanks for support from the DLR under grant 50 OR 99 03 6.
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# Primordial magnetic fields from inflation?
## I Formulation of the Problem
The idea that our galaxy could possess a magnetic field dates back to the (independent but simultaneous) works of Fermi (motivated by the origin of high energy cosmic rays) and Schwinger (motivated by the origin of galactic synchrotron emission). Since then, a lot of work has been done both from the experimental and theoretical side. Listening to observations large scale cosmological magnetic fields can be estimated from Faraday rotation effects and Zeeman splitting of (hyperfine) spectral lines. Listening to theory the main puzzle is connected with the dynamical origin of large scale magnetic fields.
It is usually assumed that the observed magnetic fields were (exponentially) amplified by galactic dynamo mechanism from pre-existing seed magnetic fields, coherent over scales of the order of 100 pc at the time of galaxy formation. The amplitude of the necessary seed fields is quite uncertain and depends on many details of the dynamo mechanism and on cosmological parameters . Typical vaules of the seeds range from $`10^{15}`$ to $`10^{25}`$ Gauss at the decoupling epoch.
In general, two types of ideas are considered for the explanation of the seed magnetic fields. The first one is related to astrophysics. For instance, a Biermann battery mechanism can be postulated at the level of protogalaxy and this will possibly lead to a mean current able to provide a source term for the evolution equations of the magnetic fields in the galactic plasma. The observation of large scale magnetic fields associated with clusters of galaxies also suggested the appealing possibility of connecting (inter) galactic magnetic fields with active galactic nuclei (in this picture the dynamo mechanism is not essential since the rotation of a cluster is much smaller than the rotation of a galaxy). The second type of ideas relies on the the interplay between particle physics and cosmology at different moments in the life of the Universe. In this framework various models were discussed such as cosmological defects , phase transitions electroweak anomaly , inflation , string cosmology , temporary electric charge non-conservation or breaking of the gauge invariance .
The main problem of most particle physics mechanisms of the origin of seed fields is how to produce them coherently on cosmological scales. All causal proposals (related, for example, to the bubble collisions at the phase transitions) may produce sufficiently large magnetic fields only on sub-horizon scales. Inflation, quite effective in producing density fluctuations at super-horizon scales, fails to amplify directly the vacuum fluctuations of the electromagnetic field because of its conformally invariant coupling to gravity.
An attractive and very economical idea on the possible primordial origin of the galactic magnetic fields was suggested in . In short, it is based on the following observation. While the coupling of electromagnetic field to the metric and to the charged fields is conformally invariant (this is not necessarily true in the models with dynamical dilaton ), the coupling of the charged scalar field to gravity is not. Thus, vacuum fluctuations of the charged scalar field can be amplified during inflation at super-horizon scales, leading potentially to non-trivial correlations of the electric currents and charges over cosmological distances. The fluctuations of electric currents, in their turn, may induce magnetic fields through Maxwell equations at the corresponding scales. The role of the charged scalar field may be played by the Higgs boson which couples to the hypercharge field above the electroweak phase transition; the generated hypercharged field is converted into ordinary magnetic field at the temperatures of the order of electroweak scale.
No detailed computations were carried out in in order to substantiate this idea. The suggestion of was further developed quite recently in for the standard electroweak theory with an optimistic conclusion that large scale magnetic fields can be indeed generated. In a supersymmetric model was studied.
The aim of the present paper is to re-analyze this proposal and compute the amplitude and the spectrum of seed magnetic fields that may be generated because of amplification of zero-point fluctuations of the scalar field during inflation. Our set-up resembles the one of Ref. . We suppose that an inflationary phase was followed by a radiation dominated phase and we compute the charged particle production associated with the change of the metric. This allows to define the spectrum and magnitude of the current fluctuations at the beginning of the radiation era. Taking the current distribution as an initial condition, we study the plasma-physics problem of the relaxation of such an initial condition. We compute finally the magnetic fields, which survived possibly until the time of galaxy formation.
Depending upon the mass of the charged scalar and upon various cosmological (critical fraction of energy density in matter, Hubble constant) and particle physics parameters we found that the magnetic fields generated in this way are much larger than vacuum fluctuations, in agreement with qualitative conclusions of refs. . However, in contrast with , their amplitude on cosmological distances is found to be too small, by many orders of magnitude, in order to seed the galactic magnetic fields. We trace back this difference in the conclusions to the discrepancy in the results obtained for the Bogoliubov coefficients (appearing in the problem of scalar particle production) and to the treatment of the relaxation of electric currents in conducting media during the radiation dominated epoch.
The plan of our paper is the following. In Section II we will discuss the amplification of the charged scalar field in an expanding cosmological background. In Section III we will study the connection between the amplification of the charged scalar and the production of charge and current density fluctuations. In Section IV we will develop a curved space description of the evolution of charge and current fluctuations within a kinetic (Landau-Vlasov) approach. We will apply our analysis to the physical case of an ultra-relativistic plasma prior to decoupling. Section V contains some phenomenological applications of our formalism. We will mainly be concerned with the generation of large scale magnetic fields in various models of cosmological evolution and with the possible occurrence of charge density and current fluctuations at large scales. Section VI contains our concluding remarks.
## II Amplification of a complex scalar field during inflation
In this paper we will only consider scalar electrodynamics. Possible generalizations of our results to theories containing more scalar fields (for example, for electroweak theory or its extensions) are straightforward. Since fermions are conformally coupled to gravity, their gravitational production is too small to generate any substantial seed magnetic fields .
Consider the action of a (massive) charged scalar field minimally coupled to the background geometry and to the electromagnetic field (hypercharge field, if the standard model is assumed):
$$S=d^4x\sqrt{g}\left[(𝒟_\mu \varphi )^{}𝒟^\mu \varphi m^2\varphi ^{}\varphi \frac{1}{4}_{\alpha \beta }^{\alpha \beta }\right],$$
(2.1)
where $`𝒟_\mu =_\mu ie𝒜_\mu `$, $`_{\mu \nu }=_{[\mu }𝒜_{\nu ]}`$, $`g_{\mu \nu }`$ is the four-dimensional metric and $`g`$ its determinant.
We will suppose that the background geometry is described, in conformal coordinates, by a Friedmann-Robertson-Walker (FRW) line element
$$g_{\mu \nu }=a^2\eta _{\mu \nu },ds^2=a^2(\eta )[d\eta ^2d\stackrel{}{x}^2],$$
(2.2)
where $`\eta _{\mu \nu }`$ is the usual Minkowski metric.
It is convenient to introduce rescaled fields $`\mathrm{\Phi }=a\varphi `$ and $`A_\mu =a𝒜_\mu `$. Correspondingly, we denote by $`\stackrel{}{E}`$ and $`\stackrel{}{B}`$ the electric and magnetic fluctuations in curved space. They are related to the usual flat space fields $`\stackrel{}{}`$ and $`\stackrel{}{}`$ by a time-dependent rescaling involving the scale factor:
$$\stackrel{}{E}=a^2\stackrel{}{},\stackrel{}{B}=a^2\stackrel{}{}.$$
(2.3)
In terms of the rescaled fields the action is
$$S=d^3x𝑑\eta [\eta ^{\mu \nu }D_\mu \mathrm{\Phi }^{}D_\nu \mathrm{\Phi }+\left(\frac{a^{\prime \prime }}{a}m^2a^2\right)\mathrm{\Phi }^{}\mathrm{\Phi }\frac{1}{4}F_{\alpha \beta }F^{\alpha \beta }],$$
(2.4)
where the prime denotes differentiation with respect to the conformal time coordinate $`\eta `$ and $`D_\mu =_\mu ieA_\mu `$, $`F_{\mu \nu }=_{[\mu }A_{\nu ]}`$. This is simply the action of electrodynamics in flat space-time with a time dependent mass term for the scalar field. From this form of the action it is obvious that there is no direct amplification of electromagnetic fields during inflation. Moreover, for conformal coupling of the scalar field to gravity the term proportional to $`a^{\prime \prime }`$ is absent and the scalar particle production is supressed by the charged boson mass.
To compute the magnetic field fluctuations we will use a perturbative approach. Namely, we will firstly compute the scalar particle production omitting completely the coupling of the scalar field to the gauge field (i.e. neglecting the back reaction, as it can be checked to be, a posteriori, self-consistent).
The classical evolution equations for $`\mathrm{\Phi }`$ in the case $`A_\mu =0`$ are simply given by:
$$\mathrm{\Phi }^{\prime \prime }^2\mathrm{\Phi }\frac{a^{\prime \prime }}{a}\mathrm{\Phi }+m^2a^2\mathrm{\Phi }=0.$$
(2.5)
We will often use also decomposition of the complex scalar field in terms of two real fields, $`\mathrm{\Phi }=(\mathrm{\Phi }_1+i\mathrm{\Phi }_2)/\sqrt{2}`$.
Once the background geometry is specified, the amplified inhomogeneities in the field $`\mathrm{\Phi }`$ can be computed. Suppose, as in , that the history of the Universe consists of two different epochs. A primordial phase, whose background evolution is not exactly known, and a radiation dominated phase where the scale factor $`a(\eta )`$ evolves linearly in conformal time. A continuous (and differentiable) choice of scale factors is then represented by
$`a_i(\eta )=\left({\displaystyle \frac{\eta }{\eta _1}}\right)^\alpha ,\eta <\eta _1,`$ (2.6)
$`a_r(\eta )={\displaystyle \frac{\alpha \eta +(\alpha +1)\eta _1}{\eta _1}},\eta \eta _1,`$ (2.7)
where $`\alpha `$ is some effective exponent parametrizing the dynamics of the primordial phase of the Universe. Notice that if $`\alpha =1`$ we have that the primordial phase coincides with a de Sitter inflationary epoch. In practice we will consider $`\alpha =1`$ or slightly deviating from it.
As a result of the change in the behaviour of the scale factor occurring at $`\eta _1`$ the modes of $`\mathrm{\Phi }`$ will be parametrically amplified. Defining $`x_1=k\eta _1`$ and $`\mu =m\eta _1`$ we can write the Bogoliubov coefficients for $`\alpha =1`$ (standard inflation) and for cosmologically interesting case $`x_11`$ (long ranged fluctuations) and $`\mu 1`$ (small scalar mass, potentially giving rise to large scalar fluctuations):
$`\alpha _k=e^{i\frac{\pi }{8}}\sqrt{\pi }\left\{{\displaystyle \frac{x_1^{\frac{3}{2}}}{2\mathrm{\Gamma }(\frac{3}{4})\mu ^{\frac{1}{4}}}}+{\displaystyle \frac{ix_1^{\frac{1}{2}}}{2\mathrm{\Gamma }(\frac{3}{4})\mu ^{\frac{1}{4}}}}+\left[{\displaystyle \frac{1}{2\mathrm{\Gamma }(\frac{3}{4})\mu ^{\frac{1}{4}}}}+{\displaystyle \frac{(i1)\mu ^{1/4}}{\sqrt{2}\mathrm{\Gamma }(\frac{1}{4})}}\right]\sqrt{x_1}\right\}+𝒪(\mu ^{\frac{5}{4}}),`$ (2.8)
$`\beta _k=e^{i\frac{\pi }{8}}\sqrt{\pi }\left\{{\displaystyle \frac{x_1^{\frac{3}{2}}}{2\mathrm{\Gamma }(\frac{3}{4})\mu ^{\frac{1}{4}}}}{\displaystyle \frac{ix_1^{\frac{1}{2}}}{2\mathrm{\Gamma }(\frac{3}{4})\mu ^{\frac{1}{4}}}}+\left[{\displaystyle \frac{1}{2\mathrm{\Gamma }(\frac{3}{4})\mu ^{\frac{1}{4}}}}+{\displaystyle \frac{(i+1)\mu ^{1/4}}{\sqrt{2}\mathrm{\Gamma }(\frac{1}{4})}}\right]\sqrt{x_1}\right\}+𝒪(\mu ^{\frac{5}{4}}).`$ (2.9)
More general expressions can be found in Appendix A where we also give details about the calculation. The terms kept in the expansion for small $`x_1`$ and $`\mu `$ is such that the unitarity condition for the Bogoluibov coefficients is satisfied: $`|\alpha _k|^2|\beta _k|^2=1`$. This property of the truncation is quite important: when discussing charge density fluctuations we will see that interesting cancellations between the leading terms arise.
In the opposite limit, for $`k>1/\eta _1`$ the mixing coefficients are exponentially suppressed as $`\mathrm{exp}[\lambda k\eta _1]`$ . The coefficient $`\lambda `$ depends upon the details of the transition between the inflationary and radiation dominated phases. The existence of an exponential suppression in the mixing coefficients is quite important from a general point of view: it ensures gentle ultra-violet properties for the physical quantities we ought to compute.
We should mention that the leading contribution to the amplification coefficients has been computed in different contexts . However, we are not only interested in the leading behaviour of the Bogoliubov coefficients but also in the corrections whose contribution is relevant when the leading contribution cancels, as in the case of charge density fluctuations (see the following Section).
## III Charge and current density fluctuations
The Bogoliubov coefficients obtained in the previous Section specify the probability of charged particle creation. The fluctuations in the scalar field induce also fluctuations in the electric current associated with the $`U(1)`$ symmetry of our action. The fluctuations in the charge and current density act as a source term for the evolution of the fluctuations in the gauge fields.
We will ignore, for the moment, the effect of the plasma conductivity which is essential for the calculation of induced magnetic fields (it will be discussed in the following section) and we will simply consider the structure of the current-current correlators. In the curved space, the conservation equation for the current is given by:
$$\frac{1}{\sqrt{g}}_\mu \left(\sqrt{g}j^\mu \right)=0,$$
(3.1)
where
$$j^\mu =ie(\varphi ^{}^\mu \varphi \varphi ^\mu \varphi ^{}).$$
(3.2)
It is convenient to introduce the rescaled current
$$J^\mu =\sqrt{g}g^{\mu \nu }j_\nu ,$$
(3.3)
which can be expressed as
$$J^\mu =e\left[\mathrm{\Phi }_2^\mu \mathrm{\Phi }_1\mathrm{\Phi }_1^\mu \mathrm{\Phi }_2\right]$$
(3.4)
in terms of conformal fields. Notice that in this last expressions the index appearing in the derivatives is raised (or lowered) using the Minkowski metric and not the curved metric.
If we average $`J^\mu `$ on the vacuum we have clearly that $`J^\mu =0`$. However, the fluctuations of the same quantity for two space-time separated points are not zero. By defining the two-point functions of the field operators $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$
$$𝒢(x,y)=\mathrm{\Phi }_1(x)\mathrm{\Phi }_1(y)=\mathrm{\Phi }_2(x)\mathrm{\Phi }_2(y),$$
(3.5)
the two-point function of the charge and current density fluctuations can be written as
$$J_\mu (x)J_\nu (y)=2e^2\left\{𝒢(x,y)\frac{^2}{x^\mu y^\nu }𝒢(x,y)\frac{}{y^\nu }𝒢(x,y)\frac{}{x^\mu }𝒢(x,y)\right\},$$
(3.6)
where $`x(\stackrel{}{x},\eta )`$ and $`y(\stackrel{}{y},\tau )`$ ($`\eta `$ and $`\tau `$ are two different conformal times). In the case of the vacuum (no amplification took place) the two-point functions are simply
$$𝒢(x,y)=\frac{1}{2(2\pi )^3}\frac{d^3k}{k_0}e^{ik(xy)}.$$
(3.7)
In the case of flat space-time the current density fluctuations can be expressed as
$$J_\mu (x)J_\nu (y)=\frac{e^2}{4(2\pi )^6}\frac{d^3p}{p_0}\frac{d^3k}{k_0}(p_\mu k_\nu )^2e^{i(k+p)(xy)}.$$
(3.8)
When the background passes from an inflationary phase to a radiation dominated phase we have instead that the two-point functions of the field operators can be written as
$`𝒢(x,y)`$ $`=`$ $`{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}[|\alpha _k|^2g_k(\eta )g_k^{}(\tau )+|\beta _k|^2g_k^{}(\eta )g_k(\tau )`$ (3.9)
$`+`$ $`\alpha _k\beta _k^{}g_k(\eta )g_k(\tau )+\alpha _k^{}\beta _kg_k^{}(\eta )g_k^{}(\tau )]e^{i\stackrel{}{k}\stackrel{}{r}},`$ (3.10)
where $`g_k(\eta )`$ are given by Eqs. (A.15) and where $`\stackrel{}{r}=\stackrel{}{x}\stackrel{}{y}`$ .
In order to define properly the charge and current density fluctuations at different scales it is helpful to introduce an averaging procedure both over space and over time. In the case of the charge density fluctuations we will define the charge fluctuations inside a patch of volume $`L^3`$ and within a typical time $`𝒯`$ as
$$Q_{L,𝒯}^2=\frac{1}{𝒯^2}d^3xd^3y𝑑\eta 𝑑\tau J_0(\stackrel{}{x},\eta )J_0(\stackrel{}{y},\tau )W_{L,𝒯}(\stackrel{}{x},\eta \xi _0)W_{L,𝒯}(\stackrel{}{y},\tau \xi _0),$$
(3.11)
where $`W_{L,𝒯}`$ are the smearing functions selecting the contribution of the correlator inside a given space-time region. We can choose the smearing functions to be, for instance, the so-called top hat function which has a sharp edge. It is defined as $`W_{L,𝒯}(\stackrel{}{x},\eta )=1`$ for $`|\stackrel{}{x}|L`$, $`|\eta |<𝒯`$ and it is equal to zero otherwise. This choice is, however, not so useful for our case. In fact, one can show that the Fourier transform of the top hat profile goes to zero, for large $`k`$, as $`(kr_0)^3`$. Unfortunately this behaviour (ultimately related to the shrap edge of the profile) could create spurious effects in the charge and current density fluctuations. A better choice is, for our purposes, the gaussian smearing function
$$W_{L,𝒯}(\stackrel{}{x},\eta )=e^{\frac{|\stackrel{}{x}|^2}{2L^2}\frac{\eta ^2}{2𝒯^2}}.$$
(3.12)
This expression smears the high frequency modes more efficiently than the top hat function. By inserting Eqs. (3.10) into Eq. (3.11) we obtain, after regulatization over time, in the limit $`Lm^1`$ and for $`\xi _0\eta _1`$
$$Q_{L,𝒯}^2\frac{e^2}{(2\pi )^5}d^3xd^3yd^3kd^3p|\alpha _p\beta _k\alpha _k\beta _p|^2e^{\frac{|\stackrel{}{x}|^2}{2L^2}\frac{|\stackrel{}{y}|^2}{2L^2}}e^{i(\stackrel{}{k}+\stackrel{}{p})(\stackrel{}{x}\stackrel{}{y})}.$$
(3.13)
A similar procedure can be carried on in the case of the current density fluctuation.
As we will see in the next Section, the relevant quantity to be computed in order to estimate the size of the gauge field fluctuations is given by the trace of
$$(\stackrel{}{}\times \stackrel{}{J})_k(\stackrel{}{}\times \stackrel{}{J})_l=ϵ_{ika}ϵ_{jlb}\frac{^2}{x^ay^b}J_i(\stackrel{}{x},\eta )J_j(\stackrel{}{y},\tau ).$$
(3.14)
Therefore, the regularized quantity we are interested in is
$$(\stackrel{}{}\times \stackrel{}{J})_{L,𝒯}^2=\frac{e^2}{𝒯^2}d^3xd^3y𝑑\eta 𝑑\tau (\stackrel{}{}\times \stackrel{}{J})_k(\stackrel{}{}\times \stackrel{}{J})_kW_{L,𝒯}(\stackrel{}{x},\eta \xi _0)W_{L,𝒯}(\stackrel{}{y},\tau \xi _0).$$
(3.15)
After regularization over time we find, using Eq. (3.10) into Eq. (3.15) in the limit $`m\xi _0^2/\eta _11`$, that for $`Lm^1`$ and for $`\xi _0\eta _1`$ our correlator becomes
$$(\stackrel{}{}\times \stackrel{}{J})_{L,𝒯}^2\frac{e^2}{(2\pi )^5}\frac{1}{m^2}\left(\frac{\eta _1}{\xi _0}\right)^2d^3xd^3yd^3kd^3p𝒟(k,p)e^{\frac{|\stackrel{}{x}|^2}{2L^2}\frac{|\stackrel{}{y}|^2}{2L^2}}e^{i(\stackrel{}{k}+\stackrel{}{p})(\stackrel{}{x}\stackrel{}{y})},$$
(3.16)
where
$$𝒟(k,p)=[k^2p^2(\stackrel{}{k}\stackrel{}{p})^2]|\alpha _k+\beta _k|^2|\alpha _p+\beta _p|^2.$$
(3.17)
## IV Magnetic fields from charge and current density fluactuations
Once we know the charge and current fluctuations at the beginning of the radiation dominated epoch we can compute the induced fluctuations in the electromagnetic fields from the Maxwell equations. The complete solution of this problem is hardly possible in a general case because of the complicated (and model dependent) dynamics of the inflaton decay and of the reheating of the Universe after inflation. We shall use an extremely simplified picture of this process. Namely, we will suppose that the plasma right after inflation is in thermal equilibrium with some temperature $`T`$, up to small (at super-horizon scales) perturbations of the distribution functions leading to charge and current fluctuations. Though this picture may be, in general, far from reality we expect it to reproduce the physics of the magnetic field generation quite accurately, because the rate of equilibration of the electromagnetic reaction is large compared with the rate of the Universe expansion after inflation and may be bigger than the rate of the inflaton decay. In fact, our result depends essentially on the density of the charged particles and their typical momentum only. Thus they may be qualitatively applicable even to situations with large deviations from thermal equilibrium.
Within this set-up the problem can be related to the relaxation of an initial perturbation in the distribution function. Consider an equilibrium homogeneous and isotropic conducting plasma, characterized by a distribution function $`f_0(p)`$ common for both positively and negatively charged ultrarelativistic particles (for example, electrons and positrons) . Suppose now that this plasma is slightly perturbed, so that the distribution functions are
$$f_+(\stackrel{}{x},\stackrel{}{p},\eta )=f_0(p)+\delta f_+(\stackrel{}{x},\stackrel{}{p},\eta ),f_{}(\stackrel{}{x},\stackrel{}{p},\eta )=f_0(p)+\delta f_{}(\stackrel{}{x},\stackrel{}{p},\eta ),$$
(4.1)
where $`+`$ refers to positrons and $``$ to electrons, and $`\stackrel{}{p}`$ is the conformal momentum. The Vlasov equation defining the curved-space evolution of the perturbed distributions can be written as
$`{\displaystyle \frac{f_+}{\eta }}+\stackrel{}{v}{\displaystyle \frac{f_+}{\stackrel{}{x}}}+e(\stackrel{}{E}+\stackrel{}{v}\times \stackrel{}{B}){\displaystyle \frac{f_+}{\stackrel{}{p}}}=\left({\displaystyle \frac{f_+}{\eta }}\right)_{\mathrm{coll}},`$ (4.2)
$`{\displaystyle \frac{f_{}}{\eta }}+\stackrel{}{v}{\displaystyle \frac{f_{}}{\stackrel{}{x}}}e(\stackrel{}{E}+\stackrel{}{v}\times \stackrel{}{B}){\displaystyle \frac{f_{}}{\stackrel{}{p}}}=\left({\displaystyle \frac{f_{}}{\eta }}\right)_{\mathrm{coll}},`$ (4.3)
where the two terms appearing at the right hand side of each equation are the collision terms. This system of equation represents the curved space extension of the Vlasov-Landau approach to plasma fluctuations . All particle number densities here are related to the comoving volume.
Notice that, in general $`\stackrel{}{v}=\stackrel{}{p}/\sqrt{\stackrel{}{p}^2+m_e^2a^2}`$. In the ultra-relativistic limit $`\stackrel{}{v}=\stackrel{}{p}/|\stackrel{}{p}|`$ and the Vlasov equations are conformally invariant. This implies that, provided we use conformal coordinates and rescaled gauge fields, the system of equations which we would have in flat space looks exactly the same as the one we are discussing in a curved FRW (spatially flat) background .
The evolution equations of the gauge fields coupled to the two Vlasov equations can be written as
$`\stackrel{}{}\stackrel{}{E}=e{\displaystyle d^3p[f_+(\stackrel{}{x},\stackrel{}{p},\eta )f_{}(\stackrel{}{x},\stackrel{}{p},\eta )]},`$ (4.4)
$`\stackrel{}{}\times \stackrel{}{E}+\stackrel{}{B}^{}=0,`$ (4.5)
$`\stackrel{}{}\stackrel{}{B}=0,`$ (4.6)
$`\stackrel{}{}\times \stackrel{}{B}\stackrel{}{E}^{}={\displaystyle d^3p\stackrel{}{v}[f_+(\stackrel{}{x},\stackrel{}{p},\eta )f_{}(\stackrel{}{x},\stackrel{}{p},\eta )]}.`$ (4.7)
Now, if $`\delta f_\pm (\stackrel{}{x},\stackrel{}{p},\eta )`$ at the beginning of the radiation dominated epoch $`\eta _0`$ are known (their magnitude follows from our computation of the Bogoluibov coefficients), and $`E(\stackrel{}{x},\eta _0)=B(\stackrel{}{x},\eta _0)=0`$ initially, the magnetic field at later times can be found from Eqs. (4.2)–(4.7).
This problem is solved in the Appendix B in a relaxation time approximation for the collision integral and under the assumption of small fluctuations in the distribution functions. Electric fields, as well as the charge density perturbations quickly relax to zero during the time given by the inverse conductivity of the plasma, $`\sigma n_q/(p\nu )`$, where $`n_qT^3`$ is the density of the charged particles and $`pT`$ is their average momentum, $`\nu `$ is the frequency of collisions All quantities here are given for conformal time. For instance the curved space conductivity $`\sigma `$ is related to the flat space one, $`\sigma _c`$ , through a rescaling which involves the scale factor : $`\sigma =a\sigma _c`$. Once we use our rescaled system $`n,T`$ and $`\sigma `$ do not change for adiabatic Universe expansion, provided that the effective number of massless degrees of freedom is constant. This is different from Ref. where ordinary (flat space) conductivity scales as temperature squared.. As for the magnetic fields, for the most interesting case of large scales ($`k^2\sigma /\eta `$ the result does not depend on the frequency of collisions $`\nu `$ and reads
$$\stackrel{}{B}\frac{p}{\alpha n_q}\stackrel{}{}\times \stackrel{}{J},$$
(4.8)
where $`\alpha `$ is the fine structure constant.
## V Estimates of magnetic field
We are now ready to estimate the produced current and charge density fluctuations. We will firstly discuss the case of magnetic fields assuming that the primordial background is of de Sitter or quasi-de Sitter type. We will then move to estimate the charge density fluctuations.
### A De Sitter case
From Eq. (3.16) we see that the correlation function of current density fluctuations is determined by the sum of the Bogoliubov coefficients. From Eq. (2.9) we have
$$|\alpha _k+\beta _k|^2|\alpha _p+\beta _p|^2\left(\frac{\pi ^2}{\mathrm{\Gamma }(3/4)^4}\mathrm{sin}^4\frac{\pi }{8}\right)\mu ^1|k\eta _1|^3|p\eta _1|^3,$$
(5.1)
which implies
$$(\stackrel{}{}\times \stackrel{}{J})_{L,𝒯}^2=\left(\frac{\eta _1}{\xi _0}\right)^2\left(\frac{e^2}{(2\pi )^5}\frac{\pi ^2}{\mathrm{\Gamma }(3/4)^2}\mathrm{sin}^4\frac{\pi }{8}\right)\frac{k_1^4}{\mu ^3}(L),$$
(5.2)
where
$$(L)=d^3xd^3yd^3kd^3p\frac{[k^2p^2(\stackrel{}{k}\stackrel{}{p})^2]}{k^3p^3}e^{\frac{|\stackrel{}{x}|^2}{2L^2}}e^{\frac{|\stackrel{}{y}|^2}{2L^2}}e^{i(\stackrel{}{k}+\stackrel{}{p})(\stackrel{}{x}\stackrel{}{y})}0.9\times (2\pi )^5L^2.$$
(5.3)
The final result for our correlator is
$$\frac{(\stackrel{}{}\times \stackrel{}{J})_{L,𝒯}^2}{V^2}\left(0.9\frac{e^2\pi ^2}{\mathrm{\Gamma }(3/4)^4}\mathrm{sin}^4\frac{\pi }{8}\right)\frac{k_1^4}{L^4}\mu ^3,$$
(5.4)
where $`VL^3`$ is the typical volume of a region of size $`L`$; we take $`\xi _0\eta _1`$ (larger values of $`\xi _0`$ give even smaller values of the magnetic fields).
In order to get an estimate of the magnetic field, one should specify the density of the charged particles and their average momentum, see (4.8). A most realistic estimate would be to take $`nT^3`$ and $`pT`$, where $`T`$ is the reheating temperature. Then we obtain for the gauge field fluctuations
$$\frac{B^2}{T^4}=6.26\left(\frac{m}{M_P}\right)^3\left(\frac{H_1}{M_P}\right)^5\frac{1}{(LT)^4}.$$
(5.5)
In Eq. (5.5) $`L`$ is the coherence scale of the magnetic field. The ratio of the magnetic field to the $`T^2`$ is roughly constant during the Universe expansion (if dynamo effects and the galaxy collapse are discarded as well as annihilation of heavy particles) and may be taken at any time, provided the coherence length $`L`$ is taken at the same epoch. We will take $`L1`$ Mpc at the present microwave background temperature <sup>§</sup><sup>§</sup>§Equally, in agreement with the analysis in Ref. , one can take $`L100`$ pc at the galaxy formation time, which gives essentially the same value for $`LT`$., which gives $`LT3\times 10^{25}`$. For a galactic mass perturbation (of the order of $`10^{12}`$ solar masses, including dark matter) the typical length scale is of the order of $`1.9\times (\mathrm{\Omega }_0h^2)^{1/2}`$ Mpc . The gravitational collapse of the protogalaxies enhances the magnetic flux frozen in the primeval galactic patch by roughly a factor of the order of $`10^3`$.
The obtained value for the magnetic field should be compared with the values of the magnetic field necessary to seed the galactic dynamo mechanism. The differential rotation of the galaxy introduces a parity violating term in the MHD equations (the dynamo term). The effect of this term exponentially amplify the seed magnetic field by a factor $`e^{\mathrm{\Gamma }t}`$ . In this amplification factor enter two numbers: the galactic age (of the order of $`10`$ Gyrs and the dynamo amplification rate ($`\mathrm{\Gamma }`$) whose estimate is rather uncertain: values of the order $`0.3`$$`0.5`$ Gyr for $`\mathrm{\Gamma }^1`$ are present in the literature . Following a recent analysis we have that the required value for the seed field can be expressed as
$`{\displaystyle \frac{B_{\mathrm{dec}}}{T_{\mathrm{dec}}^2}}5\times 10^{17},\mathrm{for}\mathrm{\Gamma }^1=0.5\mathrm{Gyr},`$ (5.6)
$`{\displaystyle \frac{B_{\mathrm{dec}}}{T_{\mathrm{dec}}^2}}2\times 10^{25},\mathrm{for}\mathrm{\Gamma }^1=0.3\mathrm{Gyr}.`$ (5.7)
where $`T_{\mathrm{dec}}`$ is the decoupling temperature. These values have been obtained in the case of $`\mathrm{\Omega }_00.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$ and $`h=0.65`$ as a fiducial set of cosmological parameters . Notice that in the case of a flat Universe with $`\mathrm{\Omega }_0=1`$ we would get values of $`B_{\mathrm{dec}}/T_{\mathrm{dec}}^2`$ close to $`10^{15}`$.
We want now to compare these values with the parameter space described by our estimate. If we take $`m=100`$ GeV (the lower mass limit for the Higgs boson) and if we assume $`H_1/M_P10^6`$ (i.e. the maximal $`H_1`$ compatible with microwave background anisotropies) we have that
$$\frac{B_{\mathrm{dec}}}{T_{\mathrm{dec}}^2}5.77\times 10^{40},$$
(5.8)
twenty orders of magnitude smaller than the required value in order to seed the magnetic field of our galaxy. One can argue that by lowering the mass this estimate will improve. This is not the case. By taking $`m1`$ MeV (which is not realistic at all) we get $`B_{\mathrm{dec}}/T_{\mathrm{dec}}^210^{32.5}`$, still too small to be relevant. Notice, finally, that to tune $`H_1`$ does not help: since it can only get smaller that $`10^6`$ it can only make the value of the seed even smaller than the values we just mentioned.
The most conservative estimate of the magnetic field can be obtained with the asumption that the number density of the charged particles in the plasma is the number of gravitationally created scalars,
$$n_{gr}=\frac{d^3k}{(2\pi )^3}|\beta _k|^2\frac{1}{8\pi \mathrm{\Gamma }^2(3/4)}\left(\frac{H_1}{m}\right)^{\frac{1}{2}}H_1^3\mathrm{log}\left(\frac{H_1}{m}\right)$$
(5.9)
and their average momentum is simply $`pH_1100\mathrm{G}\mathrm{e}\mathrm{V}`$, specific for the gravitational production. If this is indeed the case, the value of the magnetic field $`B`$ is larger by a factor
$$8\pi \mathrm{\Gamma }^2(3/4)\left(\frac{m}{M_P}\right)^{\frac{1}{2}}\left(\frac{M_P}{H_1}\right)^{\frac{3}{2}},$$
(5.10)
if compared with the estimate (5.8). Following these considerations the obtained magnetic field becomes
$$\frac{B}{T^2}10^2\left(\frac{H_1}{m}\right)\frac{1}{(LT)^2}.$$
(5.11)
Even with the most optimistic numbers (i.e. $`m100`$ GeV and $`H_110^6M_P`$ ) we get that $`B/T^2`$ can be, at most, $`10^{38}`$ over scales relevant for the dynamo action, with a minute gain, with respect to the estimate of Eq. (5.8). In priciple, one can think that even if the produced fields are too weak to turn-on the galactic dynamo they could be of some relevance for other processes occurring at different epochs in the life of the Universe. For instance, the electroweak horizon at the time of the electroweak phase transition gives $`L_{\mathrm{ew}}T_{\mathrm{ew}}10^{16}`$. This would imply from Eq. (5.11) that $`B_{\mathrm{ew}}/T_{\mathrm{ew}}^210^{19}`$. At the electroweak epoch the smallest coherence lenght of magnetic fields is set by the diffusivity. Around the diffusivity scale $`L_{\mathrm{diff}}T_{\mathrm{ew}}10^8`$ . Thus the obtained magnetic field can be, at most, $`B_{\mathrm{ew}}10^3T_{\mathrm{ew}}^2`$. In order to have sizable effects on the phase diagram of the electroweak transition we should have, at least, that $`B_{\mathrm{ew}}/T_{\mathrm{ew}}^20.2`$ so the produced fields are also too small in this context.
### B Quasi de Sitter case
Up to now we assumed that the primordial phase in the evolution of the Universe was of pure de Sitter type. In this subsection we are going to study what happens is this requirement is relaxed. In principle we know that during the inflationary phase the scalar field slightly decreases its value and, consequently, the inflationary curvature scale is not exactly constant but mildly decreasing as a function of time. In order to account for the decrease in the curvature and in the scalar field we can define the two slow-rolling parameters
$$ϵ=\frac{\dot{H}}{H^2},\lambda =\frac{\ddot{\varphi }}{H\dot{\varphi }},$$
(5.12)
where $`\varphi `$ denotes, in this Section, the inflaton and $`H=\dot{a}/a`$ is the Hubble parameter (the over-dot denotes differentiation with respect to the cosmic time). The important point for our purposes is that the slight decrease in the curvature corrects the evolution equation for the charged scalar. More specifically, we know that the dependence on the curvature appears in the mode function as $`a^{\prime \prime }/a`$. It is a simple exercise to show that
$$\frac{a^{\prime \prime }}{a}=2a^2H^2(1\frac{ϵ}{2}),$$
(5.13)
which implies that the slow-rolling corrections will make the index appearing in the Hankel functions slightly larger
$$\rho =\frac{3}{2}+ϵ.$$
(5.14)
the case $`ϵ=0`$ corresponds to the pure de Sitter case. Clearly, from our general expressions of the Bogoliubov coefficients (see Eq.(A.24)) an increase in $`\rho `$ implies an increase of the Bogoliubov coefficients in the infra-red part of the spectrum.
Therefore we want to estimate the magnetic fields in the case where the index $`\nu `$ is kept generic. For this purpose, from Eq. (A.24) we have that
$$|\alpha _k+\beta _k|^2|\alpha _p+\beta _p|^2\left(2^{4\rho 6}(\rho \frac{1}{2})^4\frac{\mathrm{\Gamma }(\rho )^4}{\mathrm{\Gamma }(3/4)^4}\mathrm{sin}^4\frac{\pi }{8}\right)\mu ^1|k\eta _1|^{2\nu }|p\eta _1|^{2\rho }.$$
(5.15)
Following the same steps as in the pure de Sitter case we can estimate that
$$\frac{B^2}{T^4}\left(\frac{2^{4\rho 6}(\rho \frac{1}{2})^4\mathrm{\Gamma }^4(\rho )}{2\mathrm{\Gamma }^4(3/4)e^2}\mathrm{sin}^4\frac{\pi }{8}\right)\left(\frac{m}{M_P}\right)^3\left(\frac{H_1}{M_P}\right)^{2+2\rho }(LT)^{4\rho 10}.$$
(5.16)
If $`\rho `$ could be larger than $`1.5`$ the magnetic fields would also be larger at the relevant scales.
There are two relevant bounds on $`\rho `$. The first and obvious one comes from the energy density stored in the charged scalar modes. The energy density stored in the scalar field modes goes as $`mk^3|\beta _k|^2`$. This means, in the case of generic $`\rho `$ that the energy density (in critical units) is
$$\left(\frac{m}{H_1}\right)^{\frac{1}{2}}\left(\frac{H_1}{M_P}\right)^{(\rho +\frac{1}{2})}\left(\frac{k}{T}\right)^{32\rho }.$$
(5.17)
If we take $`m100`$ GeV and $`\rho 2`$ we see that this expression still gives a value $`10^6`$ at the decoupling scale. Larger values of $`\rho `$ could induce further anisotropies in the microwave background. So we will assume $`1.5<\rho <2`$ which would be already enough to increase magnetic fields according to Eq. (5.16).
The second class of bounds stems from the fact that $`\rho `$ is connected with the slow rolling parameters which are constrained. In fact, the contribution to the scalar spectral index deduced from the COBE data could be written, in terms of the slow-rolling parameters, as
$$n=14ϵ+2\lambda .$$
(5.18)
The values of $`ϵ`$ and $`\lambda `$ are different depending upon the different models of background evolution, namely upon the different analytical forms of the inflationary potential driving inflation. This can be appreciated by writing the equations of motion of the inflaton in the slow-rolling approximation
$`3H\dot{\varphi }+{\displaystyle \frac{V}{\varphi }}0,`$ (5.19)
$`M_P^2H^2V.`$ (5.20)
By using these two equations we can re-express $`ϵ`$ and $`\lambda `$ as
$`ϵ={\displaystyle \frac{M_P^2}{6}}\left({\displaystyle \frac{\mathrm{ln}V}{\varphi }}\right)^2,`$ (5.21)
$`\lambda ={\displaystyle \frac{M_P^2}{6}}\left({\displaystyle \frac{\mathrm{ln}V}{\varphi }}\right)^2+{\displaystyle \frac{M_P^2}{3V^2}}\left({\displaystyle \frac{^2V}{\varphi ^2}}\right).`$ (5.22)
Clearly the values of $`ϵ`$ and $`\lambda `$ depend upon the value of $`\varphi `$. For instance, we could estimate the value of $`ϵ`$ coinciding with the value of the field $`50`$ e-folds before the end of inflation (corresponding to the moment where the large scales went out of the horizon) . In this case, for a power-law potential $`V\varphi ^q`$ we have that
$$ϵ=\frac{q}{q+200},\lambda =\frac{q2}{q+200}.$$
(5.23)
In the case of an exponential potential of the form $`V=\mathrm{exp}[6\varphi ^2/(pM_P^2)]`$ we have that $`ϵ=\lambda =1/p`$. Consequently, from Eq. (5.18), we have that $`n=1(2+q)/100`$ for power-law potentials, $`n=12/p`$ for exponential potentials. The scale-invariant spectrum as it has been observed by the DMR experiment aboard the COBE satellite the spectral index can lie in the range $`0.9n1.5`$. In order to have more magnetic fields we should increase $`\rho `$, namely we should go for large $`ϵ`$. The variation of the spectral constrains the maximal value of $`ϵ`$. So if we take $`0.9`$ as the minimal value for $`n`$ we would have from Eq. (5.23) that the $`q`$ (for a power-law potential) is $`q=8`$. But this would imply that the maximal $`ϵ`$ is $`0.03`$ and our effective $`\rho `$ will be $`1.53`$. Too small to give relevant consequences in Eq. (5.16) for the magnetic fields generation. Similar conclusions could be reached in the case of power-law potentials. By playing with the value of $`\rho `$ it is not possible to enhance the value obtained for large scale magnetic fields.
### C Charge density fluctuations
On the basis of the kinetic discussion the modes which survive in the plasma are the transverse ones. The charge density fluctuations, being associated with the longitudinal modes will be dissipated quite quickly in a typical time of the order of the inverse temperature. Still it is interesting to check if the charge density fluctuations are small. From our expression of the Bogoliubov coefficients given in Eq. (2.9) we have that
$$|\alpha _p\beta _k\alpha _k\beta _p|^2=\frac{\pi ^2}{2\mathrm{\Gamma }(1/4)^2\mathrm{\Gamma }(3/4)^2}\frac{|p\eta _1|}{|k\eta _1|^3}.$$
(5.24)
We can insert this last expression into Eq. (3.13) and perform the integrations. The integrations over $`x`$ and $`y`$ are trivial. The integration over the moduli over the momenta leads to a non trivial integral which can be exactly computed:
$$_0^{\mathrm{}}\frac{dz}{z}e^{z^2}_0^{\mathrm{}}w^2𝑑we^{w^2}\mathrm{sinh}2wz=1/2.$$
(5.25)
Defining then the electric charge density as
$$n_e=\frac{Q_{L,T}}{L^6}$$
(5.26)
we obtain that
$$\frac{n_e}{n_\gamma }10^2\left(\frac{H_1}{M_P}\right)^{1/2}(LT)^2,$$
(5.27)
where $`n_\gamma T^3`$ is the photon density. This value, for a length scale corresponding to the horizon at decoupling, would give $`n_e/n_\gamma 10^{58}`$. We want to stress that the bounds on the electric charge fluctuations were derived by assuming that charge fluctuations would induce electric fields coherent over the whole horizon. These fields would cause some observable effect in the microwave background so that a constraint on the charge density could be derived. Again, on the basis of kinetic treatment of plasma fluctuations, we can say that electric fields dissipate as soon as the the plasma becomes conducting. Therefore the effects on the microwave background are not present.
## VI Conclusions
In this paper we discussed the amplification and the fate of the fluctuations of a charged scalar field in the inflationary Universe scenario. These fluctuations may lead eventually to the generation of some magnetic fields in the Universe. We found that the produced magnetic field are always much smaller than the most optimistic lower bounds required in order to seed the galactic dynamo mechanism. Thus, the inflationary production of charged scalars is unlikely to be responsible for the observed galactic magnetic fields.
## Acknowledgments
We thank P. Tinyakov and C. Wagner for helpful discussions.
## A Bogoliubov Coefficients
Defining $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ as the real and imaginary part of $`\mathrm{\Phi }`$ as
$$\mathrm{\Phi }=\frac{1}{\sqrt{2}}(\mathrm{\Phi }_1+i\mathrm{\Phi }_2),$$
(A.1)
we assume that the background geometry evolves from $`\eta \mathrm{}`$ to $`\eta +\mathrm{}`$ for instance, according to Eq. (2.7). In both limits we can define a Fourier expansion of $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ in terms of two distinct orthonormal sets of modes. By then promoting the classical fields to quantum mechanical operators in the Heisenberg representation we can write, for $`\eta \mathrm{}`$
$`\mathrm{\Phi }_1^{\mathrm{in}}(\stackrel{}{x},\eta )={\displaystyle \frac{d^3k}{(2\pi )^{3/2}}\left[a_kf_k(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}+a_k^{}f_k^{}(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}\right]},`$ (A.2)
$`\mathrm{\Phi }_2^{\mathrm{in}}(\stackrel{}{x},\eta )={\displaystyle \frac{d^3p}{(2\pi )^{3/2}}\left[b_pf_p(\eta )e^{i\stackrel{}{p}\stackrel{}{x}}+b_p^{}f_p^{}(\eta )e^{i\stackrel{}{p}\stackrel{}{x}}\right]}.`$ (A.3)
where the two sets of creation and annihilation operators (i.e. $`[a_k,a_k^{}]`$ and $`[b_p,b_p^{}]`$) are mutually commuting. As $`\eta +\mathrm{}`$ $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ can be expanded in a second orthonormal set of modes
$`\mathrm{\Phi }_1^{\mathrm{out}}(\stackrel{}{x},\eta )={\displaystyle \frac{d^3k}{(2\pi )^{3/2}}\left[\stackrel{~}{a}_kg_k(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}+\stackrel{~}{a}_k^{}g_k^{}(\eta )e^{i\stackrel{}{k}\stackrel{}{x}}\right]},`$ (A.4)
$`\mathrm{\Phi }_2^{\mathrm{out}}(\stackrel{}{x},\eta )={\displaystyle \frac{d^3p}{(2\pi )^{3/2}}\left[\stackrel{~}{b}_pg_p(\eta )e^{i\stackrel{}{p}\stackrel{}{x}}+\stackrel{~}{b}_p^{}g_p^{}(\eta )e^{i\stackrel{}{p}\stackrel{}{x}}\right]}.`$ (A.5)
Since both sets of modes are complete the old modes can be expressed in terms of the new ones
$$f_k(\eta )=\alpha _kg_k(\eta )+\beta _kg_k^{}(\eta ).$$
(A.6)
This transformation, once inserted back into Eq. (A.3), implies that
$$\stackrel{~}{a}_k=\alpha _ka_k+\beta _k^{}a_k^{}.$$
(A.7)
Notice that in order to preserve the scalar products in the old and new sets of orthonormal modes we have that the two complex numbers $`\alpha _k`$ and $`\beta _k`$ are subjected to the constraints $`|\alpha _k|^2|\beta _k|^2=1`$. Exactly the same discussion applies for the field operator $`\mathrm{\Phi }_2`$. Eq. (A.7) is nothing but the well known Bogoliubov-Valatin transformation; $`\alpha _k`$ and $`\beta _k`$ are the Bogoliubov coefficients parametrizing the mixing between positive and negative frequency modes.
In order to ensure the continuity of the operators $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ we have to match continuously the old mode functions with the new ones. During the primordial phase of the Universe the evolution equations satisfied by the mode functions is
$$\frac{d^2f_k}{d\eta ^2}+\left[k^2\frac{\alpha (\alpha +1)}{\eta ^2}+\frac{\mu }{\eta _1^2}\left(\frac{\eta _1}{\eta }\right)^{2\alpha }\right]f_k=0,$$
(A.8)
where $`\alpha =1`$ for a pure de Sitter background $`\mu =m/(H_1a_1)=m\eta _1`$. Notice that with $`H`$ we will denote the Hubble factor in cosmic time ( as usual the relation between the Hubble factor in conformal time, e.g. $``$ and the Hubble factor in cosmic time is $`H=/a`$; during the de Sitter epoch $`\eta ^1`$). Notice that if the mass of the charged scalar is not of Planckian magnitude, $`\mu 1`$. Moreover, in the limit where $`\mu 1`$ one can argue that the amplified fluctuations will be exponentially suppressed. Since we want to explore the situation where the mass of the scalar field is of electroweak order we will always be (quite safely) in the limit $`\mu 1`$.
The exact solution of Eq. (A.8) which reduces, in the limit $`\eta \mathrm{}`$, to the usual positive frequency Minkowski space solution is given by
$$f_k(\eta )=\frac{1}{\sqrt{2k}}p\sqrt{x}H_\rho ^{(1)}(x),$$
(A.9)
where $`x=k\eta `$ and $`H_\rho ^{(1)}`$ is the first order Hankel function. Again, in the pure de Sitter case, $`\rho =3/2\sqrt{1(4/9)\mu ^2}`$. Since $`\mu `$ is typically small, $`\rho 3/2`$ in the pure de Sitter case. We denoted with $`p`$ a phase factor which we choose such that
$$p=\sqrt{\frac{\pi }{2}}e^{i\frac{\pi }{4}(1+2\rho )}.$$
(A.10)
With this choice of $`p`$ we have that $`f_k(\eta )e^{ik\eta }/\sqrt{2k}`$ for $`\eta \mathrm{}`$ .
After the radiation dominated phase sets in (for $`\eta >\eta _1`$ ) the evolution equation obeyed by the mode functions $`g_k(\eta )`$ is given by
$$\frac{d^2g_k}{d\eta ^2}+\left[k^2+\frac{\mu ^2(\eta +2\eta _1)^2}{\eta _1^4}\right]g_k=0.$$
(A.11)
This last equation can be easily recast in the form of a parabolic cylinder equation . Defining $`\gamma =\mu /\eta _1^2`$ we can introduce two new quantities, namely
$$z=\sqrt{2\gamma }(\eta +2\eta _1),q=\frac{k^2}{2\gamma }.$$
(A.12)
Consequently, Eq. (A.12) becomes
$$\frac{d^2g_k}{dz^2}+\left[q+\frac{z^2}{4}\right]g_k=0,$$
(A.13)
which is one of the canonical forms of the parabolic cylinder equation . The exact solutions which reduce to positive and negative frequency modes for $`\eta +\mathrm{}`$ are
$`g_k(\eta )={\displaystyle \frac{1}{(2\gamma )^{1/4}}}e^{i\frac{\pi }{8}}D_{iq\frac{1}{2}}(ie^{i\frac{\pi }{4}}z),`$ (A.14)
$`g_k^{}(\eta )={\displaystyle \frac{1}{(2\gamma )^{1/4}}}e^{i\frac{\pi }{8}}D_{iq\frac{1}{2}}(e^{i\frac{\pi }{4}}z),`$ (A.15)
where $`D_\sigma `$ are the parabolic cylinder functions in the Whittaker’s notation . Notice that with our choice of normalizations we have, in the limit $`z|q|`$ and for $`k^2\eta _1m`$, that
$$g_k(\eta )\sqrt{\frac{\eta _1}{2m\eta }}e^{\frac{i}{2}\frac{m\eta ^2}{\eta _1}}.$$
(A.16)
In view of the actual calculation it is worth recalling the exact expressions of the parabolic cylinder functions (in the Whittaker form) in terms of confluent hypergeometric (Kummer) functions $`{}_{1}{}^{}F_{1}^{}(a,b,x)`$:
$`D_{iq\frac{1}{2}}(e^{i\frac{\pi }{4}}z)`$ $`=`$ $`\sqrt{\pi }2^{\frac{iq}{2}\frac{1}{4}}e^{\frac{iz^2}{4}}[{\displaystyle \frac{1}{\mathrm{\Gamma }(\frac{3}{4}\frac{iq}{2})}}{}_{1}{}^{}F_{1}^{}({\displaystyle \frac{1}{4}}{\displaystyle \frac{iq}{2}},{\displaystyle \frac{1}{2}},{\displaystyle \frac{iz^2}{2}})`$ (A.18)
$`{\displaystyle \frac{z(1i)}{\mathrm{\Gamma }(\frac{1}{4}\frac{iq}{2})}}{}_{1}{}^{}F_{1}^{}({\displaystyle \frac{3}{4}}{\displaystyle \frac{iq}{2}},{\displaystyle \frac{3}{2}},{\displaystyle \frac{iz^2}{2}})],`$
$`D_{iq\frac{1}{2}}(ize^{i\frac{\pi }{4}})`$ $`=`$ $`\sqrt{\pi }2^{\frac{iq}{2}\frac{1}{4}}e^{i\frac{z^2}{4}}[{\displaystyle \frac{1}{\mathrm{\Gamma }(\frac{3}{4}+\frac{iq}{2})}}{}_{1}{}^{}F_{1}^{}({\displaystyle \frac{1}{4}}+{\displaystyle \frac{iq}{2}},{\displaystyle \frac{1}{2}},{\displaystyle \frac{iz^2}{2}})`$ (A.20)
$`{\displaystyle \frac{z(1+i)}{\mathrm{\Gamma }(\frac{1}{4}+\frac{iq}{2})}}{}_{1}{}^{}F_{1}^{}({\displaystyle \frac{3}{4}}+{\displaystyle \frac{iq}{2}},{\displaystyle \frac{3}{2}},{\displaystyle \frac{iz^2}{2}})].`$
The Bogoliubov coefficients are obtained from
$`f_k(\eta _1)=\alpha _kg_k(\eta _1)+\beta _kg_k^{}(\eta _1),`$ (A.21)
$`f_k^{}(\eta _1)=\alpha _kg_k^{}(\eta _1)+\beta _kg_{k}^{}{}_{}{}^{}(\eta _1),`$ (A.22)
which is a system of two equations in the two unknowns $`\alpha _k`$ and $`\beta _k`$.
By solving this system we obtain an exact expression for the Bogoliubov coefficients which is, in general a function of two variables : $`\mu =m\eta _1`$ and $`x_1=k\eta _1`$. Since $`\mu 1`$ we can expand the exact result in this limit and we obtain, in the case of a generic Bessel index $`\rho `$,
$`\alpha _k=\pi e^{i\frac{\pi }{8}}\left\{{\displaystyle \frac{i}{\sqrt{2}\mathrm{\Gamma }(\frac{3}{4})}}S_2(x_1,\rho )\mu ^{\frac{1}{4}}+{\displaystyle \frac{(1+i)}{2\mathrm{\Gamma }(\frac{1}{4})}}[S_1(x_1,\rho )+S_2(x_1,\rho )]\mu ^{\frac{1}{4}}\right\}+𝒪(\mu ^{\frac{5}{4}}),`$ (A.23)
$`\beta _k=\pi e^{i\frac{\pi }{8}}\left\{{\displaystyle \frac{i}{\sqrt{2}\mathrm{\Gamma }(\frac{3}{4})}}S_2(x_1,\rho )\mu ^{\frac{1}{4}}+{\displaystyle \frac{(i1)}{2\mathrm{\Gamma }(\frac{1}{4})}}[S_1(x_1,\rho )+S_2(x_1,\rho )]\mu ^{\frac{1}{4}}\right\}+𝒪(\mu ^{\frac{5}{4}}),`$ (A.24)
where $`S_1(x_1,\rho )`$ and $`S_2(x_1,\rho )`$ contain the explicit dependence upon the Hankel’s functions:
$`S_1(x_1,\rho )=e^{i\frac{\pi }{4}(1+2\rho )}H_\rho ^{(1)}(x_1),`$ (A.25)
$`S_2(x_1,\rho )=\sqrt{x_1}e^{i\frac{\pi }{4}(1+2\rho )}\left[\left(\rho +{\displaystyle \frac{1}{2}}\right){\displaystyle \frac{H_\rho ^{(1)}(x_1)}{\sqrt{x_1}}}\sqrt{x_1}H_{\rho +1}^{(1)}(x_1)\right].`$ (A.26)
If we are now specifically interested in the pure de Sitter phase we can insert the value $`\rho =3/2`$ into Eq. (A.26). Then, we can insert the obtained expressions into Eq. (A.24) and we obtain the wanted Bogoliubov coefficients reported in Eqs. (2.9).
Notice that it would not be correct to use the asymptotic solutions (like the one reported in Eq. (A.16)) in order to compute the Bogoliubov coefficientsWe disagree with the calculation of Ref. where the matching has been performed using approximate mode functions. In our case, for small $`\mu `$ the leading behaviour of the Bogoliubov coefficient goes as $`\mu ^{1/4}`$. In the case of it goes as $`\mu ^{1/2}`$, an artifact of the WKB approximation.. In fact Eq. (A.16) can be viewed as the a WKB-type solution of Eq. (A.13) which can be also written as
$$\frac{d^2g_k}{d\eta ^2}+\omega _k^2(\eta )g_k=0,\omega _k^2(\eta )=k^2+m^2a^2(\eta ).$$
(A.27)
By now postulating a WKB-type solution we have that
$$g_k(\eta )=\frac{1}{2W(\eta )}e^{i^\eta W(\eta ^{})𝑑\eta ^{}}.$$
(A.28)
By now inserting the trial solution back to Eq. (A.27) we get $`W(\eta )`$ is specified by the following non-linear relation
$$W^2(\eta )=\omega _k^2(\eta )\frac{1}{2}\left[\frac{W^{\prime \prime }}{W}\frac{3}{2}\left(\frac{W^{}}{W}\right)^2\right].$$
(A.29)
This equation can be solved by iteration. If we keep the lowest order we get that
$$W_0(\eta )\omega _k(\eta )$$
(A.30)
and by using the explicit expression of the scale factor during the radiation dominated epoch we exactly get Eq. (A.16). This solution is valid provided the corrections to the exact expression of $`W(\eta )`$ are small, namely, provided, from eq, (A.29)
$$\omega _k^2(\eta )\frac{1}{2}\left[\frac{W_0^{\prime \prime }}{W_0}\frac{3}{2}\left(\frac{W_0^{}}{W_0}\right)^2\right].$$
(A.31)
This last inequality, using the explicit expression of $`\omega _k(\eta )`$, implies that
$$k^2\eta ^2+m^2\eta ^2\left(\frac{\eta +2\eta _1}{\eta _1}\right)1.$$
(A.32)
Now we can see that this inequality is clearly satisfied for $`\eta +\mathrm{}`$. However, for $`\eta \eta _1`$, this inequality would imply $`m\eta _1>1`$ (since $`k\eta _1<1`$). Therefore, in order to be consistent with the requirement that $`m\eta _1<1`$ we have to use the WKB-type solution only for large (positive) $`\eta `$.
## B Vlasov-Landau approach to electromagnetic field fluctuations.
The purpose of this appendix is to give details concerning the derivation of the relation between the electromagnetic field fluctuations and the initial fluctuations in the current (or charge) density profile.
By subtracting Eqs. (4.2) and (4.3) we obtain the equations relating the fluctuations of the distributions functions of the charged particles present in the plasma to the induced gauge field fluctuations:
$`{\displaystyle \frac{}{\eta }}f(\stackrel{}{x},\stackrel{}{p},t)+\stackrel{}{v}{\displaystyle \frac{}{\stackrel{}{x}}}f(\stackrel{}{x},\stackrel{}{p},t)+2e\stackrel{}{E}{\displaystyle \frac{f_0}{\stackrel{}{p}}}=\nu (p)f,`$ (B.1)
$`\stackrel{}{}\stackrel{}{E}=e{\displaystyle d^3pf(\stackrel{}{x},\stackrel{}{p},\eta )},`$ (B.2)
$`\stackrel{}{}\times \stackrel{}{E}+\stackrel{}{B}^{}=0,`$ (B.3)
$`\stackrel{}{}\stackrel{}{B}=0,`$ (B.4)
$`\stackrel{}{}\times \stackrel{}{B}\stackrel{}{E}^{}={\displaystyle d^3p\stackrel{}{v}f(\stackrel{}{x},\stackrel{}{p},\eta )},`$ (B.5)
where $`f(\stackrel{}{x},\stackrel{}{p},\eta )=\delta f_+(\stackrel{}{x},\stackrel{}{p},\eta )\delta f_{}(\stackrel{}{x},\stackrel{}{p},\eta )`$ and $`\nu (p)`$ is a typical frequency of collisions.
We can solve this system by taking the Fourier transform of the space-dependent quantities and the Laplace transform of the time-dependent quantities:
$`\stackrel{}{E}_{\stackrel{}{k}\omega }={\displaystyle _0^{\mathrm{}}}𝑑\eta e^{i\omega \eta }{\displaystyle d^3xe^{i\stackrel{}{k}\stackrel{}{x}}\stackrel{}{E}(\stackrel{}{x},\eta )},`$ (B.6)
$`\stackrel{}{B}_{\stackrel{}{k}\omega }={\displaystyle _0^{\mathrm{}}}𝑑\eta e^{i\omega \eta }{\displaystyle d^3xe^{i\stackrel{}{k}\stackrel{}{x}}\stackrel{}{B}(\stackrel{}{x},\eta )},`$ (B.7)
$`f_{\stackrel{}{k}\omega }={\displaystyle _0^{\mathrm{}}}𝑑\eta e^{i\omega \eta }{\displaystyle d^3xe^{i\stackrel{}{k}\stackrel{}{x}}f(\stackrel{}{x},\stackrel{}{p},\eta )}.`$ (B.8)
We have now to specify, at the initial time, the form of the perturbed distribution function which can be derived from the amplification studied in the previous Section. We will call $`g_\stackrel{}{k}(\stackrel{}{p})`$ the initial profile of the distribution function. Eq. (B.5) can then be re-written as
$`g_\stackrel{}{k}(\stackrel{}{p})+i(\stackrel{}{k}\stackrel{}{v}\omega )f_{\stackrel{}{k}\omega }(\stackrel{}{p})+2e\stackrel{}{E}_{\stackrel{}{k}\omega }{\displaystyle \frac{f_0}{\stackrel{}{p}}}=\nu f,`$ (B.9)
$`i\stackrel{}{k}\stackrel{}{E}_{\stackrel{}{k}\omega }=e{\displaystyle f_{\stackrel{}{k}\omega }(\stackrel{}{p})d^3p},`$ (B.10)
$`i\stackrel{}{k}\stackrel{}{B}_{\stackrel{}{k}\omega }=0,`$ (B.11)
$`\stackrel{}{B}_{\stackrel{}{k}\omega }={\displaystyle \frac{1}{\omega }}\stackrel{}{k}\times \stackrel{}{E}_{\stackrel{}{k}\omega },`$ (B.12)
$`i\omega \left(1{\displaystyle \frac{k^2}{\omega ^2}}\right)\stackrel{}{E}_{\stackrel{}{k}\omega }+{\displaystyle \frac{i}{\omega }}\stackrel{}{k}(\stackrel{}{k}\stackrel{}{E}_{\stackrel{}{k}\omega })={\displaystyle d^3p\stackrel{}{v}f_{\stackrel{}{k}\omega }(\stackrel{}{p})},`$ (B.13)
where eq. (B.13) has been obtained by using Eq. (B.12) in the (Fourier and Laplace) transformed of the last of Eqs. (B.5). The Gauss constraint at $`\eta =0`$ implies that
$$i\stackrel{}{k}\stackrel{}{E}_0(\stackrel{}{k})=ed^3pg_\stackrel{}{k}(\stackrel{}{p}).$$
(B.14)
If we start, at the initial time, with a given profile of fluctuations fluctuation the Gauss constraint determines the initial value of the electric field. The magnetic field fluctuations are consistently equal to zero.
We can now separate the electric field in its polarizations parallel and transverse to the direction of propagation of the fluctuation. The transverse current provides a source for the evolution of transverse electric field fluctuations
$$i\omega \left(1\frac{k^2}{\omega ^2}\right)\stackrel{}{E}_{\stackrel{}{k}\omega }^{}=ed^3pf_{\stackrel{}{k}\omega }(\stackrel{}{p})\stackrel{}{v}_{},$$
(B.15)
whereas the charge fluctuations provide a source for the evolution of longitudinal electric field fluctuations
$$i\stackrel{}{k}\stackrel{}{E}_{\stackrel{}{k}\omega }^{}=ed^3pf_{\stackrel{}{k}\omega }(\stackrel{}{p}).$$
(B.16)
In Eqs. (B.15) and (B.16) we defined the longitudinal part of the electric field fluctuations and the transverse electric field as
$$\stackrel{}{E}_{\stackrel{}{k}\omega }^{}=\stackrel{}{E}_{\stackrel{}{k}\omega }^{}\frac{\stackrel{}{k}}{|\stackrel{}{k}|^2}(\stackrel{}{k}\stackrel{}{E}_{\stackrel{}{k}\omega }),\stackrel{}{E}_{\stackrel{}{k}\omega }^{}=\frac{\stackrel{}{k}}{|\stackrel{}{k}|^2}(\stackrel{}{E}_{\stackrel{}{k}\omega }\stackrel{}{k}).$$
(B.17)
The solution of Eq. (B.9) is given by
$$f_{\stackrel{}{k}\omega }(\stackrel{}{p})=\frac{1}{i(\stackrel{}{k}\stackrel{}{v}\omega i\nu )}\left[g_\stackrel{}{k}(\stackrel{}{p})2e\stackrel{}{v}\stackrel{}{E}_{\stackrel{}{k}\omega }\frac{f_0}{p}\right],$$
(B.18)
where we used that $`f_0/\stackrel{}{p}\stackrel{}{v}f_0/p`$. The longitudinal and transverse components of the electric fluctuations can be obtained by inserting Eq. (B.18) into Eqs. (B.15)-(B.16)
$`|\stackrel{}{E}_{\stackrel{}{k}\omega }^{}|={\displaystyle \frac{e}{ikϵ_{}}}{\displaystyle d^3p\frac{g_\stackrel{}{k}(\stackrel{}{p})}{i(\stackrel{}{k}\stackrel{}{v}\omega i\nu )}},`$ (B.19)
$`\stackrel{}{E}_{\stackrel{}{k}\omega }^{}={\displaystyle \frac{e\omega }{\omega ^2ϵ_{}k^2}}{\displaystyle d^3p\stackrel{}{v}^{}\frac{g_\stackrel{}{k}(\stackrel{}{p})}{(\stackrel{}{k}\stackrel{}{v}\omega i\nu )}},`$ (B.20)
where $`ϵ_{}`$ and $`ϵ_{}`$ are, respectively, the longitudinal and transverse part of the polarization tensor
$`ϵ_{}(k,\omega )=1{\displaystyle \frac{2e^2}{k^2}}{\displaystyle d^3p\frac{\stackrel{}{k}\stackrel{}{v}}{(\stackrel{}{\stackrel{}{k}}\stackrel{}{v}\omega i\nu )}\frac{f_0}{p}},`$ (B.21)
$`ϵ_{}(k,\omega )=1{\displaystyle \frac{e^2}{\omega }}{\displaystyle d^3p\frac{\stackrel{}{v}_{}^2}{(\stackrel{}{k}\stackrel{}{v}\omega i\nu )}\frac{f_0}{p}}.`$ (B.22)
Now, the general expression for the generated magnetic field is
$$\stackrel{}{B}_{\stackrel{}{k}\omega }=\frac{e}{\omega ^2ϵ_{}(k,\omega )k^2}d^3p[\stackrel{}{v}\times \stackrel{}{k}]\frac{g_\stackrel{}{k}(\stackrel{}{p})}{(\stackrel{}{k}\stackrel{}{v}\omega i\nu )}.$$
(B.23)
The space-time evolution of the magnetic fluctuations can be determined by performing the inverse Laplace and Fourier transforms:
$$\stackrel{}{B}(\stackrel{}{x},\eta )=e^{i\omega \eta }\frac{ed\omega }{\omega ^2ϵ_{}(k,\omega )k^2}d^3ke^{i\stackrel{}{k}\stackrel{}{x}}[\stackrel{}{v}\times \stackrel{}{k}]d^3p\frac{g_\stackrel{}{k}(\stackrel{}{p})}{(\stackrel{}{v}\stackrel{}{k}\omega i\nu )}.$$
(B.24)
In order to perform this integral, the explicit relations for the polarization tensors should be given. They depend on the equilibrium distribution function $`f_0(p)`$, which we take to be Notice that most of our considerations can be easily extended to the case of a Bose-Einstein or Fermi-Dirac distribution. What is important, in our context, is the analytical structure of the polarization tensors and this is the same for different distributions .
$$f_0(p)=\frac{n_q}{8\pi T^3}e^{\frac{p}{T}},$$
(B.25)
where $`T`$ is the equilibrium temperature, $`p`$ is the modulus of the momentum and $`n`$ is the equilibrium (thermal) density of charged particles in the plasma. The normalization is chosen in such a way that $`d^3pf_0(p)=n_q`$.
Then we have for transverse polarization
$$ϵ_{}(k,\omega )=1+\frac{e^2n_q}{2\omega kT}\left\{\left[1\frac{(\omega +i\nu )^2}{k^2}\right]\mathrm{ln}\frac{k\omega i\nu }{k+\omega +i\nu }2\frac{\omega +i\nu }{k}\right\},$$
(B.26)
and for the longitudinal polarization:
$$ϵ_{}(k,\omega )=1+\frac{e^2n_q}{k^2T}\left\{2+\frac{\omega +i\nu }{k}\mathrm{ln}\frac{k\omega i\nu }{k+\omega +i\nu }\right\}.$$
(B.27)
Consider now the case of very small momenta $`k\omega `$ and $`\omega \nu `$, relevant for long-ranged magnetic fields. Then the computation of the integral (B.24) in the large time limit and with the use of explicit form of the transverse polarization tensor in (B.26) gives <sup>\**</sup><sup>\**</sup>\**For small $`k`$, the equation $`\omega ^2ϵ_{}(k,\omega )k^2=0`$ defining the poles of the inverse Laplace transform implies $`\omega ik^2/\sigma `$.:
$$B(\stackrel{}{x},\eta )\frac{T}{4\pi \alpha n_q}\mathrm{exp}(k^2\eta /\sigma )\stackrel{}{}\times \stackrel{}{J},$$
(B.28)
where $`\sigma `$ is the plasma conductivity in the relaxation time approximation,
$$\sigma =\frac{2e^2n_q}{\nu T}$$
(B.29)
and initial electric current is given by
$$\stackrel{}{J}(\stackrel{}{x})=d^3p\stackrel{}{v}g_\stackrel{}{k}(\stackrel{}{p}).$$
(B.30)
In closing our discussion of the Vlasov equation we want to briefly comment about the validity of our approach. The obtained results assumed that the linearization of the Vlasov equation is consistent with the physical assumptions of our problem. This is indeed the case. In order to safely linearize the Vlasov equation we have to make sure that the perturbed distribution function of the charge fluctuations is always smaller than the first order of the perturbative expansion (given by the distribution of Eq. B.25). In other words we have to make sure that
$$|\delta f_+(\stackrel{}{x},\stackrel{}{p},\eta )|f_0(\stackrel{}{p}),|\delta f_{}(\stackrel{}{x},\stackrel{}{p},\eta )|f_0(\stackrel{}{p}).$$
(B.31)
These conditions imply that
$$\frac{e\stackrel{}{E}_{\stackrel{}{k}\omega }}{|\stackrel{}{k}\stackrel{}{v}\omega i\nu |}\frac{f_0(\stackrel{}{p})}{\stackrel{}{p}}f_0(\stackrel{}{p}).$$
(B.32)
If we now define the relativistic plasma frequency as
$$\omega _p^2=\frac{2e^2n_q}{3T},$$
(B.33)
we can see that the condition expressed by Eq. (A.29) can be restated, for modes $`k\omega _p`$, as $`|\stackrel{}{E}_{\stackrel{}{k},\omega }|^2<n_qT`$ (where we essentially took the square modulus of Eq. (B.32)). This last inequality expresses the fact that the energy density associated with the gauge field fluctuations should always be smaller than the critical energy density stored in radiation. The linear treatment of the Vlasov equation is certainly accurate provided the typical modes of the the field are smaller than the plasma frequency and provided the energy density in electric and magnetic fields is smaller than $`T^4`$, i.e. the energy density stored in the radiation background.
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# 1 Introduction
## 1 Introduction
A chiral $`p`$-form $`A`$ is defined by the equation
$`F={}_{}{}^{}F.`$ (1.1)
where $`FdA`$ is the corresponding fieldstrength. From this, it is clear that the dimension, $`d`$, of space-time is given by $`d=2(p+1)`$. Furthermore, $`p`$ should be even in the Minkowski case and odd in the Euclidean case since only in those cases is the square of the Hodge $``$-operator equal to the identity. Throughout this paper we maintain a Minkowski signature. Chiral $`p`$-forms naturally appear in string or M-theory. Chiral bosons are essential in the worldsheet formulation of the heterotic string and correspond to $`p=0`$. Chiral two-forms, which, as we will explain further, constitute the main motivation for the present study, are central in the description of the M5-brane. Finally, chiral four-forms appear in type IIB string theory where they signal the presence of D3-branes.
The strongest motivation for studying deformations of chiral forms arises from the study of coinciding M5-branes. The solitonic objects in M theory (viewed here as eleven dimensional supergravity) are M2- and M5-branes. These soliton solutions break half of the supersymmetries, reducing them from 32 to 16. Their effective worldbrane actions contain therefore 16 Goldstinos, which correspond to 8 propagating fermionic degrees of freedom. This should be matched by 8 bosonic degrees of freedom. Obvious candidates for the bosonic degrees of freedom of a $`p`$-brane living in $`d`$ dimensions are the $`dp1`$ transversal positions of the brane. For the M2 brane ($`p=2`$ and $`d=11`$), this saturates the number of bosonic degrees of freedom. For the M5-brane, however, one needs three additional bosonic degrees of freedom. The little group of the worldvolume theory is $`Spin(6)=SU(2)\times SU(2)`$, which means we need a (3,1) representation of this. This is precisely a chiral two-form in six dimensions.
In the low energy limit where bulk gravity decouples, a single M5-brane is described by a six dimensional $`N=(2,0)`$ superconformal field theory , . Its field content consists of five scalar fields and a single chiral two-form <sup>1</sup><sup>1</sup>1Throughout this paper we ignore the fermionic degrees of freedom which does not change any of our conclusions.. A Lorentz non-covariant action was constructed in , and . A covariant action was obtained in and . The covariant action contains appropriate extra auxiliary fields and gauge symmetries. Partial gauge fixing of the covariant action yields the non-covariant action.
Once $`n`$ M5-branes coincide, the situation changes. This can be seen by compactifying one direction on a circle. For small radius, the resulting theory is weakly coupled type IIA string theory. When the M5-branes are transversal to the circle, they appear in the type IIA theory as $`n`$ coinciding NS5-branes. Not much is explicitely known about this system. However, when the M5-branes are longitudinal to the circle, they emerge as $`n`$ coinciding D4-branes. The effective action for such a system is a $`U(n)`$ non-Abelian Born-Infeld action . Its leading and next to leading terms are well understood but discussion about the subleading terms remains , . Ignoring higher derivative terms and focussing on the leading term, one gets that the dynamics of the D4 system is governed by 5 scalar fields in the adjoint representation of $`U(n)`$ coupled to a 5-dimensional $`U(n)`$ gauge theory. Going back to the supergravity description, this observation suggests the existence of a non-abelian extension of chiral 2-forms.
Genuine non-abelian extensions of non-chiral $`p`$-forms, for $`p2`$ have not yet been constructed. Viewing a 2-form as a connection over loopspace, one can show that no straightforward non-abelian extension exists (see also ). Dropping geometric prejudices, all local deformations continously connected to the free action were constructed in . Though both known and novel deformations were discovered, none of them had the required property that the $`p`$-form gauge algebra becomes truly non-abelian.
Turning back to chiral 2-forms, one finds that M-theoretical considerations indicate that $`n`$ coinciding M5-branes constitute a highly unusual physical system. Indeed, the supergravity description of $`n`$ M5-branes predicts that both the entropy and the two-point function for the stress-energy tensor scale as $`n^3`$ in the large $`n`$ limit. Anomaly considerations lead to a similar behaviour , . So this suggests that a non-abelian extension of chiral two-forms falls outside the scope of finite dimensional semi-simple Lie groups as none of those have a dimension growing as fast as $`n^3`$ (where $`n`$ would be the dimension of the Cartan sub-algebra). It has been argued that “gerbes” could provide the appropriate mathematical framework .
In , we announced the result that no local field theory is able to describe a system of coinciding M5-branes. This result was obtained by showing that local deformations of the action cannot modify the abelian nature of the algebra of the $`2`$-form gauge symmetries. It holds under the assumption that the deformed action is continuous in the coupling constant (i.e., possible non-perturbative “miracles” are not investigated) and reduces, in the limit of vanishing coupling constant, to the action describing free chiral $`2`$-forms. In particular, no assumption was made on the polynomial order (cubic, quartic …) of the interaction terms.
In the present paper we present detailed proofs of that assertion. The techniques used in this paper can be applied in a straightforward fashion to prove the results in as well. There, deformations of chiral four-forms in ten dimensions were analyzed with as conclusion that the only consistent deformation was the type IIB coupling of the chiral four-form to the NS-NS and the R-R two-forms familiar from IIB supergravity .
The outline of this paper is as follows. In the next section, we review how the problem of consistent couplings can be reformulated as a cohomological problem . We then recall the non-covariant formalism for chiral $`2`$-forms, and their BRST formulation (sections 3, 4 and 5). In particular, we point out that the BRST differential $`s`$ naturally splits as the sum $`s=\delta +\gamma `$ of simpler building blocks. After a brief section in which we recall the so-called “algebraic Poincaré lemma”, which provides an important tool for our investigations, we turn to the calculation of the BRST cohomology. First we compute the cohomology of $`\gamma `$ (section 7). Next, we compute the cohomology of $`\gamma `$ modulo $`d`$, where $`d`$ is the spacetime exterior derivative (sections 8 and 9). In section 10, we compute the same cohomologies for the other piece involved in $`s`$, namely $`\delta `$. In section 11, we put together the calculations of the previous sections to derive the announced result that the gauge symmetries for a set of free chiral $`2`$-forms are rigid and cannot be deformed continuously in the local field theoretical context. Our paper ends with a short, concluding section.
## 2 Constructing consistent couplings as a deformation problem
The theoretical problem of determining consistent interactions for a given gauge invariant system has a long history. It has been formulated in general terms in (see also ).
The equations for the consistent interactions are rather intricate because they are non linear and involve simultaneously not only the deformed action, but also the deformed structure functions of the deformed gauge algebra, as well as the deformed reducibility coefficients if the gauge transformations are reducible. The problem is further complicated by the fact that one has to factor out the “trivial” interactions that are simply induced by a change of variables.
As we now review, one can reformulate the problem as a cohomological problem . This approach systematizes the recursive construction of the consistent interactions and, furthermore, enables one to use the powerful tools of homological algebra.
Starting with a “free” action $`\stackrel{(0)}{S_0}[\phi ^i]`$ with “free” gauge symmetries
$$\delta _\epsilon \phi ^i=\underset{\alpha }{\overset{i}{\stackrel{(0)}{R}}}\epsilon ^\alpha ,$$
(2.1)
leading to the Noether identities
$$\frac{\delta \stackrel{(0)}{S}}{\delta \phi ^i}\underset{\alpha }{\overset{i}{\stackrel{(0)}{R}}}=0,$$
(2.2)
we introduce a coupling constant $`g`$ and modify $`\stackrel{(0)}{S_0}`$,
$$\stackrel{(0)}{S_0}S_0=\stackrel{(0)}{S_0}+g\stackrel{(1)}{S_0}+g^2\stackrel{(2)}{S_0}+\mathrm{}$$
(2.3)
We consider only consistent deformations, meaning that the deformed action should be gauge invariant as well. In the generic case this requires a deformation of the gauge transformation rules,
$$\underset{\alpha }{\overset{i}{\stackrel{(0)}{R}}}R_\alpha ^i=\underset{\alpha }{\overset{i}{\stackrel{(0)}{R}}}+g\underset{\alpha }{\overset{i}{\stackrel{(1)}{R}}}+g^2\underset{\alpha }{\overset{i}{\stackrel{(2)}{R}}}+\mathrm{}.$$
(2.4)
Consistency is then translated into the requirement that the Noether identities should hold to all orders
$$\frac{\delta S}{\delta \phi ^i}R_\alpha ^i=0,$$
(2.5)
where,
$$\delta _\epsilon \phi ^i=R_\alpha ^i\epsilon ^\alpha .$$
(2.6)
Expanding Eq. (2.5) order by order in the coupling constant gives consistency condition of increasing complexity.
For reducible theories, which is the case relevant to chiral $`2`$-forms, there is an additional constraint. The gauge transformations of the free theory are not independent,
$$\underset{\alpha }{\overset{i}{\stackrel{(0)}{R}}}\underset{A}{\overset{\alpha }{\stackrel{(0)}{Z}}}=0$$
(2.7)
(possibly on-shell). One must then also impose that the gauge transformations remain reducible, possibly in a deformed way. This yields additional conditions on the coefficients $`R`$’s in Eq. (2.4).
The deformations of an action fall into three classes. In the first one, gauge invariant terms are added to the original lagrangian and therefore no modification of the gauge transformations is required. Examples of this are functionals of the field strength and its derivatives, as well as Chern-Simons-like terms . In the second class, both the action and the transformation rules are modified. However, the terms added to the transformation rules are invariant under the original gauge transformations. As a consequence, the gauge algebra is not modified to first order in the coupling constant. An example of this is the Freedman-Townsend model for two-forms in four dimensions. Finally, in the last class, the additional terms in the deformed transformation rules are not gauge invariant. Therefore the gauge algebra itself gets modified as well. The best known example of this is the deformation of an abelian Yang-Mills theory to a non-abelian theory.
The key to translating the problem of consistent interactions into a cohomological problem is the antifield formalism (for reviews, see ). Let us assume that we solved the master equation for the undeformed theory. Its solution is denoted by $`\stackrel{(0)}{S}`$, which satisfies $`(\stackrel{(0)}{S},\stackrel{(0)}{S})=0`$. The existence of a consistent deformation of the original gauge invariant action implies the existence of a deformation of $`\stackrel{(0)}{S}`$, which we denote by $`S`$,
$$\stackrel{(0)}{S}S=\stackrel{(0)}{S}+g\stackrel{(1)}{S}+g^2\stackrel{(2)}{S}+\mathrm{}$$
(2.8)
Expanding the master equation for $`S`$, $`(S,S)=0`$, order by order in the coupling constant yields various consistency relations,
$`(\stackrel{(0)}{S},\stackrel{(0)}{S})`$ $`=0`$ (2.9)
$`(\stackrel{(0)}{S},\stackrel{(1)}{S})`$ $`=0`$ (2.10)
$`2(\stackrel{(0)}{S},\stackrel{(2)}{S})+(\stackrel{(1)}{S},\stackrel{(1)}{S})`$ $`=0`$
$`\mathrm{}.`$
The first equation is satisfied by assumption. As $`(\stackrel{(0)}{S},(\stackrel{(0)}{S},))=0`$, the second equation implies that $`\stackrel{(1)}{S}`$ is a cocycle for the free differential $`\stackrel{(0)}{s}(\stackrel{(0)}{S},)`$. If $`\stackrel{(1)}{S}`$ is a coboundary, $`\stackrel{(1)}{S}=(\stackrel{(1)}{T},\stackrel{(0)}{S})`$, one can show that this corresponds to a trivial deformation (i.e. a deformation which amounts to a simple redefinition of the fields).
In practice, we consider deformations which are local in spacetime, i.e., we impose that $`\stackrel{(1)}{S},\stackrel{(2)}{S},\mathrm{}`$ be local functionals. Reformulating the equations in terms of the Lagrange densities takes care of this problem. E.g., rewriting equation (2.10) as
$$\stackrel{(0)}{s}\stackrel{(1)}{S}=\stackrel{(0)}{s}(\stackrel{(1)}{𝒮})=0(\stackrel{(0)}{s}\stackrel{(1)}{𝒮})=0,$$
(2.12)
we obtain the following condition on the Lagrange density $`\stackrel{(1)}{𝒮}`$,
$$\stackrel{(0)}{s}\stackrel{(1)}{𝒮}+d=0,$$
(2.13)
where $``$ is a local form of degree $`n1`$, where $`n`$ is the dimensionality of space-time and $`d`$ is the spacetime exterior derivative<sup>2</sup><sup>2</sup>2Throughout this paper we ignore boundary contributions. Again one can show that BRST-exact terms modulo $`d`$ are trivial solutions of (2.13) and corresponds to trivial deformations. In the local context, the proper cohomology to evaluate is thus $`H^{0,n}(\stackrel{(0)}{s}d)`$ where the first and second superscripts denote the ghost number and form degree, respectively.
Note that when all the representatives of $`H^{0,n}(\stackrel{(0)}{s}d)`$ can be taken not to depend on the antifields, one may take the first-order deformations $`\stackrel{(1)}{𝒮}`$ to be antifield-independent. In this case Eq. (2.10) reduces to $`(\stackrel{(0)}{S},\stackrel{(2)}{S})=0`$ and implies that the deformation at order $`g^2`$ defines also an element of $`H^{0,n}(\stackrel{(0)}{s}d)`$. One can thus take $`\stackrel{(2)}{S}`$ not to depend on the antifields either. Proceeding in this manner order by order in the coupling constant, we conclude that the additional terms in $`S`$ are all independent of the antifields. Since the antifield-dependent terms in the deformation of the master equation are related to the deformations of the gauge transformations, this means that there is no deformation of the gauge transformations. Summarizing, if there is no non-trivial dependence on the antifields in $`H^{0,n}(\stackrel{(0)}{s}d)=0`$, the only possible consistent interactions are of the first class and do not modify the gauge symmetry. This is the situation met for a system of chiral $`2`$-forms, as we now pass to discuss.
## 3 System of free chiral 2-forms in 6 dimensions
The non-covariant action for a system of N free chiral 2-forms is ,
$$S_0[A_{ij}^A]=\underset{A}{}𝑑td^5xB^{Aij}(\dot{A}_{ij}^AB_{ij}^A),(A=1,\mathrm{},N),$$
(3.1)
where
$$B^{Aij}=\frac{1}{6}ϵ^{ijklm}F_{klm}^A=\frac{1}{2}ϵ^{ijklm}_kA_{lm}^A.$$
(3.2)
The integer $`N`$ can be any function of the number $`n`$ of coincident M5-branes (e.g., $`Nn^3`$). The action (3.1) differs from the one in - where a space-like dimension was singled out. Here we take time as the distinguished direction; from the point of view of the PST formulation , the two approaches simply differ in the gauge fixation. We work in Minkowski spacetime. This implies, in particular, that the topology of the spatial sections $`R^5`$ is trivial. Most of our considerations would go unchanged in a curved background of the product form $`R\times \mathrm{\Sigma }`$ provided the De Rham cohomology groups $`H_{DeRham}^2(\mathrm{\Sigma })`$ and $`H_{DeRham}^1(\mathrm{\Sigma })`$ of the spatial sections $`\mathrm{\Sigma }`$ vanish. \[If $`H_{DeRham}^2(\mathrm{\Sigma })`$ is non-trivial, there are additional gauge symmetries besides (3.3) below, given by time-dependent spatially closed $`2`$-forms; similarly, if $`H_{DeRham}^1(\mathrm{\Sigma })`$ is non-trivial, there are additional reducibility identities besides (3.4) below. One would thus need additional ghosts and ghosts of ghosts. These, however, would not change the discussion of local Lagrangians because they would be global in space (and local in $`t`$).\]
The action $`S_0`$ is invariant under the following gauge transformations
$$\delta _\mathrm{\Lambda }A_{ij}^A=_i\mathrm{\Lambda }_j^A_j\mathrm{\Lambda }_i^A,$$
(3.3)
because $`B^{Aij}`$ is gauge-invariant and identically transverse ($`_iB^{Aij}0`$) <sup>3</sup><sup>3</sup>3Since $`A_{0i}^A`$ does not occur in the action – even if one replaces $`_0A_{ij}^A`$ by $`_0A_{ij}^A_iA_{0j}^A_jA_{i0}^A`$ (it drops out because $`B^{Aij}`$ is transverse) –, the action is of course invariant under arbitrary shifts of $`A_{0i}^A`$.. As $`\delta A_{ij}^A=0`$ for
$$\mathrm{\Lambda }_i^A=_i\epsilon ^A,$$
(3.4)
this set of gauge transformations is reducible. This exhausts completely the redundancy in $`\mathrm{\Lambda }_i^A`$ since $`H_{DeRham}^1(R^5)=0`$.
The equations of motion obtained from $`S_0[A_{ij}^A]`$ by varying $`A_{ij}^A`$ are
$$ϵ^{ijklm}_k\dot{A}_{lm}^A2_kF^{Aijk}=0ϵ^{ijklm}_k(\dot{A}_{lm}^AB_{lm}^A)=0.$$
(3.5)
Using $`H_{DeRham}^2(R^6)=0`$, one finds that the general solution of (3.5) is
$$\dot{A_{ij}^A}B_{ij}^A=_i\mathrm{\Lambda }_j^A_j\mathrm{\Lambda }_i^A.$$
(3.6)
The ambiguity in the solutions of the equations of motion is thus completely accounted for by the gauge freedom (3.3). Hence the set of gauge transformations is complete.
We can view $`\mathrm{\Lambda }_i^A`$ as $`A_{0i}^A`$, so the equation (3.6) can be read as the self-duality equation
$$F_{0ij}^AF_{0ij}^A=0,$$
(3.7)
where $`F_{0ij}^A=\dot{A_{ij}^A}+_iA_{j0}^A+_jA_{0i}^A`$. Alternatively, one may use the gauge freedom to set $`\mathrm{\Lambda }_i^A=0`$, which yields the self-duality condition in the temporal gauge.
## 4 Fields - Antifields - Solution of the master equation
The solution of the master equation is easy to construct in this case because the gauge transformations are abelian. We refer to for the general construction.
The fields in presence here are
$$\{\mathrm{\Phi }^M\}=\{A_{ij}^A,C_i^A,\eta ^A\}.$$
(4.1)
The ghosts $`C_i^A`$ corresponds to the gauge parameters $`\mathrm{\Lambda }_i^A`$, and the ghosts of ghosts $`\eta ^A`$ corresponds to $`ϵ^A`$.
Now, to each field $`\mathrm{\Phi }^M`$ we associate an antifield $`\mathrm{\Phi }_M^{}`$. The set of antifields is then
$$\{\mathrm{\Phi }_M^{}\}=\{A^{Aij},C^{Ai},\eta ^A\}.$$
(4.2)
The fields and antifields have the respective parities
$`ϵ(A_{ij}^A)=ϵ(\eta ^A)=ϵ(C^{Ai})=0`$ (4.3)
$`ϵ(C_i^A)=ϵ(A^{Aij})=ϵ(\eta ^A)=1.`$ (4.4)
The antibracket is defined as
$$(X,Y)=d^nx\left(\frac{\delta ^RX}{\delta \mathrm{\Phi }^M(x)}\frac{\delta ^LY}{\delta \mathrm{\Phi }_M^{}(x)}\frac{\delta ^RX}{\delta \mathrm{\Phi }_M^{}(x)}\frac{\delta ^LY}{\delta \mathrm{\Phi }^M(x)}\right)$$
(4.5)
where $`\delta ^R/\delta Z(x)`$ and $`\delta ^L/\delta Z(x)`$ denote functional right- and left-derivatives.
Because the set of gauge transformations is complete and defines a closed algebra, the (minimal, proper) solution of the master equation $`(S,S)=0`$ takes the general form
$$S=S_0+\underset{M}{}()^{ϵ(M)}\mathrm{\Phi }_M^{}s\mathrm{\Phi }^M,$$
(4.6)
where $`ϵ(M)`$ is the Grassmann parity of $`\mathrm{\Phi }^M`$. More explicitly, we have
$$S=S_0+\underset{A}{}𝑑td^5x(A^{Aij}_iC_j^AC^{Ai}_i\eta ^A)$$
(4.7)
The solution $`S`$ of the master equation captures all the information about the gauge structure of the theory : the Noether identities, the closure of the gauge transformations and the higher order gauge identities are contained in the master equation. The existence of $`S`$ reflects the consistency of the gauge transformations.
## 5 BRST operator
The BRST operator $`s`$ is obtained by taking the antibracket with the proper solution $`S`$ of the classical master equation,
$$sX=(S,X).$$
(5.1)
The BRST operator can be decomposed as
$$s=\delta +\gamma $$
(5.2)
where $`\delta `$ is the Koszul–Tate differential . What distinguishes $`\delta `$ and $`\gamma `$ is the antighost number ($`antigh`$) defined through
$`antigh(A_{ij}^A)=antigh(C_i^A)=antigh(\eta ^A)=0,`$ (5.3)
$`antigh(A^{Aij})=1,antigh(C^{Ai})=2,antigh(\eta ^A)=3.`$ (5.4)
The ghost number ($`gh`$) is related to the antighost number by
$$gh=pureghantigh$$
(5.5)
where $`puregh`$ is defined through
$`puregh(A_{ij}^A)=0,puregh(C_i^A)=1,puregh(\eta ^A)=2,`$ (5.6)
$`puregh(A^{Aij})=puregh(C^{Ai})=puregh(\eta ^A)=0.`$ (5.7)
The differential $`\delta `$ is characterized by $`antigh(\delta )=1`$, i.e. it lowers the antighost number by one unit and acts on the fields and antifields according to
$`\delta A_{ij}^A`$ $`=`$ $`\delta C_i^A=\delta \eta ^A=0,`$ (5.8)
$`\delta A^{Aij}`$ $`=`$ $`2_kF^{Akij}ϵ^{ijklm}_k\dot{A}_{lm}^A,`$ (5.9)
$`\delta C^{Ai}`$ $`=`$ $`_jA^{Aij},`$ (5.10)
$`\delta \eta ^A`$ $`=`$ $`_iC^{Ai}.`$ (5.11)
The differential $`\gamma `$ is characterized by $`antigh(\gamma )=0`$ and acts as
$`\gamma A_{ij}^A`$ $`=`$ $`_iC_j^A_jC_i^A,`$ (5.12)
$`\gamma C_i^A`$ $`=`$ $`_i\eta ^A,`$ (5.13)
$`\gamma \eta ^A`$ $`=`$ $`0,`$ (5.14)
$`\gamma A^{Aij}`$ $`=`$ $`\gamma C^{Ai}=\gamma \eta ^A=0.`$ (5.15)
Furthermore we have,
$$sx^\mu =0,s(dx^\mu )=0.$$
(5.16)
## 6 Local forms - Algebraic Poincaré lemma
A $`\mathrm{l}\mathrm{o}\mathrm{c}\mathrm{a}\mathrm{l}`$ $`\mathrm{f}\mathrm{u}\mathrm{n}\mathrm{c}\mathrm{t}\mathrm{i}\mathrm{o}\mathrm{n}`$ is a function of the fields, the ghosts, the antifields, and their derivatives up to some finite order $`k`$ (which depends on the function),
$$f=f(\mathrm{\Phi },_\mu \mathrm{\Phi },\mathrm{},_{\mu _1}\mathrm{}_{\mu _k}\mathrm{\Phi }).$$
(6.1)
A *local function* is thus a function over a finite dimensional vector space $`J^k`$ called “jet space”. A *local form* is an exterior polynomial in the $`dx^\mu `$’s with local functions as coefficients. The algebra of local forms will be denoted by $`𝒜`$. In practice, the local forms are polynomial in the ghosts and the antifields, as well as in the differentiated fields, so we shall from now on assume that the local forms under consideration are of this type. One can actually show that polynomiality in the ghosts, the antifields and their derivatives follows from polynomiality in the derivatives of the $`A_{ij}`$ by an argument similar to the one used in for $`1`$-forms; and polynomiality in the derivatives is automatic in our perturbative approach where we work order by order in the coupling constant(s).
Note also that we exclude an explicit $`x`$-dependence of the local forms. One could allow for one without change in the conclusions. In fact, as we shall indicate below, allowing for an explicit $`x`$-dependence simplifies some of the proofs. We choose not to do so here since the interaction terms in the Lagrangian should not depend explicitly on the coordinates in the Poincaré-invariant context.
The following theorem describes the cohomology of $`d`$ in the algebra of local forms, in degree $`q<n`$.
###### Theorem 6.1
The cohomology of $`d`$ in the algebra of local forms of degree $`q<n`$ is given by
$`H^0(d)`$ $``$ $`R,`$
$`H^q(d)`$ $`=`$ $`\{\text{Constant Forms}\},\mathrm{\hspace{0.33em}0}<q<n.`$
Constant forms are by definition polynomials in the $`dx^\mu `$’s with constant coefficients. This theorem is called the algebraic Poincaré lemma (for $`q<n`$). There exist many proofs of this lemma in the literature. One of the earliest can be found in .
Constant $`q`$-forms are trivial in degree $`0<q<n`$ in the algebra of local forms with an explicit $`x`$-dependence; e.g., $`dx^0=df`$, where $`f`$ is the $`x^0`$-dependent function $`f=x^0`$. Thus, in this enlarged algebra, the cohomology of $`d`$ is simpler and vanishes in degrees $`0<q<n`$. This is the reason that the calculations are somewhat simpler when one allows for an expicit $`x`$-dependence.
We work in a formalism where the time direction is privileged. For this reason, it is useful to introduce the following notation : the $`l`$-th time derivative of a field $`\mathrm{\Phi }`$ (including the ghosts and antifields) is denoted by $`\mathrm{\Phi }^{(l)}`$ ($`=_0^l\mathrm{\Phi }`$), and the spatial differential is denoted by $`\stackrel{~}{d}=dx^i_i`$.
A *local spatial form* is an exterior polynomial in the spatial $`dx^k`$’s with coefficients that are local functions. If we write the set of the generators of the jet space $`J^k`$ as
$$\{\mathrm{\Phi }^{(l_0)},_{i_1}\mathrm{\Phi }^{(l_1)},\mathrm{},_{i_1}\mathrm{}_{i_k}\mathrm{\Phi }^{(0)};l_j=0,\mathrm{},kj\},$$
(6.2)
it is clear that
###### Theorem 6.2
The cohomology of $`\stackrel{~}{d}`$ in the algebra of local spatial forms of degree $`q<n1`$ is given by
$`H^0(\stackrel{~}{d})`$ $``$ $`R,`$
$`H^q(\stackrel{~}{d})`$ $`=`$ $`\{\text{Constant spatial forms}\},\mathrm{\hspace{0.33em}0}<q<n1.`$
A similar decomposition of space and time derivatives occurs of course in the Hamiltonian formalism. A discussion of the problem of consistent deformations of a gauge invariant action has been carried out in the Hamiltonian context in .
## 7 Cohomology of $`\gamma `$
The following theorem completely gives $`H(\gamma )`$.
###### Theorem 7.1
The cohomology of $`\gamma `$ is given by,
$$H(\gamma )=V.$$
(7.1)
Here, the algebra $``$ is the algebra of the local forms with coefficients that depend only on the variables $`F_{ijk}^A`$, the antifields $`\varphi _M^{}`$, and all their partial derivatives up to a finite order (“gauge-invariant” local forms). These variables are collectively denoted by $`\chi `$. The algebra $`V`$ is the polynomial algebra in the ghosts $`\eta ^A`$ of ghost number two and their time derivatives.
* The generators of $`𝒜`$ can be grouped in three sets:
$`T=\{t^i\}=\{_{\mu _1\mathrm{}\mu _k}F_{ijk}^A,_{\mu _1\mathrm{}\mu _k}\varphi _M^{},\eta ^{A(l)},dx^\mu \}`$ (7.2)
$`U=\{u^\alpha \}=\{_{(i_1\mathrm{}i_k}A_{[i)_2j]_1}^{A(l)},_{(i_1\mathrm{}i_{k1}}C_{i_k)}^{A(l)}\}`$ (7.3)
$`V=\{v^\alpha \}=\{_{i_1\mathrm{}i_k}_{[i}C_{j]}^{A(l)},_{i_1\mathrm{}i_k}\eta ^{A(l)}\}`$ (7.4)
($`k,l=0,\mathrm{}`$) where $`[`$ $`]`$ and $`(`$ $`)`$ mean respectively antisymmetrization and symmetrization; the subscript indicates the order in which the operations are made.
The differential $`\gamma `$ acts on these three sets in the following way
$$\gamma T=0,\gamma U=V,\gamma V=0.$$
(7.5)
The elements of $`U`$ and $`V`$ are in a one-to-one correspondence and are linearly independent with respect to each other, so they constitute a manifestly contractible part of the algebra and can thus be removed from the cohomology.
No element in the algebra of generated by $`T`$ is trivial in the cohomology of $`\gamma `$, except $`0`$. Indeed, let us assume the existence of a local form $`F(t^i)0`$ which is $`\gamma `$-exact, then
$`F(t^i)=\gamma G(t^i,u^\alpha ,v^\alpha )=v^\alpha {\displaystyle \frac{^LG}{u^\alpha }}(t^i,u^\alpha ,v^\alpha ).`$ (7.6)
But this implies that
$$F(t^i)=F(t^i)_{v^\alpha =0}=0,$$
(7.7)
as announced. $`\mathrm{}`$
Note that contrary to what happens in the non-chiral case, the temporal derivatives of the ghosts $`\eta ^A`$ are non-trivial in cohomology. There is thus an infinite number of generators in ghost number two for $`H(\gamma )`$, namely, all the $`\eta ^{A(l)}`$’s. In contrast, in the non-chiral case, one has $`_0\eta ^A=\gamma C_0^A`$ and so $`_0\eta ^A`$ (and all the subsequent derivatives) are $`\gamma `$-exact. In the chiral case, there is no $`C_0^A`$.
Let$`\{\omega ^I\}`$ be a basis of the vector space $`V`$ of polynomials in the variables $`\eta ^A`$ and all their time derivatives. Theorem 7.1 tells us that
$$\gamma \alpha =0,\alpha 𝒜\alpha =\underset{}{}𝒫_{}(\chi )\omega ^{}+\gamma \beta .$$
(7.8)
Furthermore, because $`\omega ^I`$ is a basis of $`V`$
$$\underset{I}{}P_I(\chi )\omega ^I=\gamma \beta P_I(\chi )=0.$$
(7.9)
It will be useful in the sequel to choose a special basis $`\{\omega ^I\}`$. The vector space $`V`$ of polynomials in the ghosts $`\eta ^A`$ and their time derivatives splits as the direct sum $`V^{2k}`$ of vector spaces with definite pure ghost number $`2k`$. The space $`V^0`$ is one-dimensional and given by the constants. We may choose $`1`$ as basis vector for $`V^0`$, so let us turn to the less trivial spaces $`V^{2k}`$ with $`k0`$. These spaces are themselves the direct sums of finite dimensional vector spaces $`V_r^{2k}`$ containing the polynomials with exactly $`r`$ time derivatives of the $`\eta `$’s (e.g., $`_0\eta ^A_{00}\eta ^B`$ is in $`V_3^4`$). The following lemma provides a basis of $`V^{2k}`$ for $`k0`$:
###### Lemma 7.1
Let $`V^{2k}`$ be the vector space of polynomials in the variables $`\eta ^{A(l)}`$ with fixed pure ghost number $`2k0`$. $`V^{2k}`$ is the direct sum
$$V^{2k}=V_0^{2k}V_1^{2k}\mathrm{},$$
(7.10)
where $`V_m^{2k}`$ is the subspace of $`V^{2k}`$ containing the polynomials with exactly $`m`$ derivatives of $`\eta ^A`$. One has dim$`V_m^{2k}`$ dim$`V_{m+1}^{2k}`$. There exist a basis of $`V_m^{2k}`$
$$\{\omega _{(m)}^{I_m}:I_m=1,\mathrm{},q_m;m=0,\mathrm{}\},$$
(7.11)
which fulfills
$$\omega _{(m)}^{I_m}=_0\omega _{(m1)}^{I_m}(I_m=1,\mathrm{},q_{m1}).$$
(7.12)
In other words, the first $`q_{m1}`$ basis vectors of $`V_m^{2k}`$ are directly constructed from the basis vectors of $`V_{m1}^{2k}`$ by taking their time derivative $`_0`$.
* We will prove the lemma by induction. For $`m=0`$, take an arbitrary basis of $`V_0^{2k}`$ (space of polynomials in the undifferentiated ghosts $`\eta ^A`$ of degree $`k`$). Assume now that a basis with the required properties exists up to order $`m1`$. Let $`\{\omega _{(m1)}^I;I=0,\mathrm{},q_{m1}\}`$ be a basis with those properties for $`V_{m1}^{2k}`$. We want to prove that it is possible to construct a basis of $`V_m^{2k}`$ where the first $`q_{m1}`$ basis vectors are the time derivatives of the basis vectors of $`V_{m1}^{2k}`$. We only have to show that the $`_0\omega _{(m1)}^I`$ are linearly independent (because they can always be completed to form a basis of $`V_m^{2k}`$). In other words, we must prove that
$$\underset{I=1}{\overset{q_{m1}}{}}\lambda _I_0\omega _{(m1)}^I=_0(\underset{I=1}{\overset{q_{m1}}{}}\lambda _I\omega _{(m1)}^I)=0$$
(7.13)
implies $`\lambda _I=0`$. But (7.13) is equivalent to
$$\underset{I=1}{\overset{q_{m1}}{}}\lambda _I\omega _{(m1)}^I=K,$$
(7.14)
where $`K`$ is a constant (algebraic Poincaré lemma in form degree 0). $`K`$ must be equal to zero because we are in pure ghost number $`0`$. By hypothesis, the $`\omega _{(m1)}^I`$ are linearly independant, hence the $`\lambda _I`$ must be all equal to zero, which ends the proof. $`\mathrm{}`$
## 8 Cohomology of $`\gamma `$ modulo $`d`$ at positive antighost number
Let be $`a^p`$ a local $`p`$-form of antighost number $`k0`$ fulfilling
$$\gamma a^p+db^{p1}=0.$$
(8.1)
We want to show that if we add to $`a^p`$ an adequate $`d`$-trivial term, the equation (8.1) reduces to $`\gamma a^p=0`$.
From (8.1), using the algebraic Poincaré lemma and the fact that $`\gamma `$ is nilpotent and anticommute with $`d`$, we can derive the descent equations
$`\gamma a^p`$ $`+`$ $`db^{p1}=0`$ (8.2)
$`\gamma b^{p1}`$ $`+`$ $`dc^{p2}=0`$
$`\mathrm{}`$
$`\gamma e^{q+1}`$ $`+`$ $`df^q=0`$ (8.4)
$`\gamma f^q`$ $`=`$ $`0,`$ (8.5)
Indeed, the fact that the antighost number is strictly positive eliminates the constants. \[E.g., from (8.1), one derives $`d\gamma b^{p1}=0`$ and thus $`\gamma b^{p1}+dc^{p2}=\text{constant}`$, but the constant must vanish since it must have strictly positive antighost number.\] We suppose $`q<p`$, since otherwise $`\gamma a^p=0`$, which is the result we want to prove. The equation (8.5) tells us that $`f^q`$ is a cocycle of $`\gamma `$. It must be non-trivial in $`H^q(\gamma )`$ because if $`f^q=\gamma g^q`$, then (8.4) becomes $`\gamma (e^{q+1}dg^q)=0`$. The redefinition $`e^{{}_{}{}^{}q+1}=e^{q+1}dg^q`$ does not affect the descent equation before (8.4), which means that the descent stops one step earlier, at $`q1`$.
Using theorem 7.1, we deduce from (8.5) that
$$f^q=\underset{m,I_m}{}[\stackrel{~}{P}_{I_m}^{(m)}(\chi )+dx^0\stackrel{~}{Q}_{I_m}^{(m)}(\chi )]\omega _{(m)}^{I_m},$$
(8.6)
where $`\stackrel{~}{P}_{I_m}^{(m)}`$ and $`\stackrel{~}{Q}_{I_m}^{(m)}`$ are local spatial forms of respective degree $`q`$ and $`q1`$. We take the basis elements $`\omega _{(m)}^{I_m}`$ to fulfill the conditions of lemma 7.1. Differentiating (8.6), we find
$`df^q`$ $`=`$ $`{\displaystyle \underset{m,I_m}{}}\{\stackrel{~}{d}\stackrel{~}{P}_{I_m}^{(m)}\omega _{(m)}^{I_m}+\gamma (\stackrel{~}{P}_{I_m}^{(m)}\widehat{\omega }_{(m)}^{I_m})`$ (8.7)
$`+dx^0[(_0\stackrel{~}{P}_{I_m}^{(m)}\stackrel{~}{d}\stackrel{~}{Q}_{I_m}^{(m)})\omega _{(m)}^{I_m}+\stackrel{~}{P}_{I_m}^{(m)}_0\omega _{(m)}^{I_m}]\}.`$
The local function $`\widehat{\omega }_{(m)}^{I_m}`$ is defined by $`\stackrel{~}{d}\omega _{(m)}^{I_m}=\gamma \widehat{\omega }_{(m)}^{I_m}`$ ( and exists thanks to equation (5.13)).
Now, we will show that the component $`\stackrel{~}{P}_{I_m}^{(m)}`$ can be eliminated from $`f^q`$ by a trivial redefinition of $`f^q`$. In order to satisfy (8.4), the term independent of $`dx^0`$ and the coefficient of the term linear in $`dx^0`$ in (8.7) must separately be $`\gamma `$-exact. The second condition gives explicitly
$$\underset{m,I_m}{}[(_0\stackrel{~}{P}_{I_m}^{(m)}\stackrel{~}{d}\stackrel{~}{Q}_{I_m}^{(m)})\omega _{(m)}^{I_m}+\stackrel{~}{P}_{I_m}^{(m)}_0\omega _{(m)}^{I_m}]=\gamma \beta ,$$
(8.8)
To analyze precisely this equation, we define a degree $`T`$ by
$$T(\chi )=0,T(\eta ^{A(m)})=m.$$
(8.9)
In fact, $`T`$ simply counts the number of time derivative of $`\eta ^A`$. We can decompose (8.8) according to the degree $`T`$. Let $`p`$ be the highest degree occuring in $`f^q`$. Then, the highest degree occuring in (8.8) is $`p+1`$ and we must have
$$\underset{I=1}{\overset{q_p}{}}\stackrel{~}{P}_I^{(p)}_0\omega _{(p)}^I=\gamma \beta _{p+1}.$$
(8.10)
From the proof of the lemma 7.1, we find that
$`\stackrel{~}{P}_I^{(p)}=0(I=1,\mathrm{},q_p)`$ (8.11)
because the $`_0\omega _{(p)}^I`$ are linearly independent. In $`T`$-degree $`p`$, (8.8) gives then
$`\gamma \beta _p`$ $`=`$ $`{\displaystyle \underset{I=1}{\overset{q_p}{}}}\stackrel{~}{d}\stackrel{~}{Q}_I^{(p)}\omega _{(p)}^I+{\displaystyle \underset{I=1}{\overset{q_{p1}}{}}}\stackrel{~}{P}_I^{(p1)}_0\omega _{(p1)}^I`$ (8.12)
$`=`$ $`{\displaystyle \underset{I=1}{\overset{q_{p1}}{}}}(\stackrel{~}{P}_I^{(p1)}\stackrel{~}{d}\stackrel{~}{Q}_I^{(p)})\omega _{(p)}^I{\displaystyle \underset{I=q_{p1}+1}{\overset{q_p}{}}}\stackrel{~}{d}\stackrel{~}{Q}_I^{(p)}\omega _{(p)}^I,`$ (8.13)
where we have used the property (7.12) of the basis $`\{\omega ^I\}`$. This implies that
$$\stackrel{~}{P}_I^{(p1)}=\stackrel{~}{d}\stackrel{~}{Q}_I^{(p)}(I=1,\mathrm{},q_{p1})$$
(8.14)
Inserting this equation in (8.6), we find that $`\stackrel{~}{P}_I^{(p1)}`$ can be removed from $`f^q`$ by eliminating a trivial cocycle of $`\gamma `$ modulo $`d`$ and redefining $`\stackrel{~}{Q}_I^{(p1)}`$. It only affects $`e^{q+1}`$ by a $`d`$-exact term. Next, the equation (8.8) at $`T`$-degree $`p1`$ shows that $`\stackrel{~}{P}_I^{(p2)}`$ is also $`\stackrel{~}{d}`$-exact and can thus also be removed. Proceeding in the same way until the order 1 in $`T`$, we have proved that all the $`\stackrel{~}{P}_I^{(m)}`$ can be eliminated from $`f^q`$.
Looking back at (8.8) and taking into account that $`\stackrel{~}{P}_{I_m}^{(m)}`$ can be set equal to zero by the above argument, we find that
$$\stackrel{~}{d}\stackrel{~}{Q}_{I_m}^{(m)}=0.$$
(8.15)
Now, we must use the invariant Poincaré lemma (invariant means in the algebra $``$ of gauge-invariant forms) stating that
###### Theorem 8.1
Let be $`\stackrel{~}{P}(\chi )`$ a local spatial form of degree $`q<5`$, then
$$\stackrel{~}{d}\stackrel{~}{P}(\chi )=0\stackrel{~}{P}(\chi )=\stackrel{~}{R}(F^{A(l)})+\stackrel{~}{d}\stackrel{~}{Q}(\chi ),$$
(8.16)
where $`\stackrel{~}{R}(F^{A(l)})`$ is a polynomial in the curvature forms $`F^A=\frac{1}{6}F_{ijk}^Adx^idx^jdx^k`$ and all their time derivatives (with coefficients that may involve $`dx^k`$, which takes care of the constant forms).
* The set of the generators of the algebra $``$ is
$$\{\chi \}=\{_{i_1\mathrm{}i_k}F_{ijk}^{A(l)},_{i_1\mathrm{}i_k}\varphi _M^{(l)},\eta ^{A(l)},dx^\mu \}$$
(8.17)
The 1-form $`dx^0`$ is not present in our problem since $`\stackrel{~}{P}`$ is a spatial local form (it only involves $`dx^k`$). Considering $`l`$ and $`A`$ as only one label (call it $`\alpha `$) and forgetting about $`dx^0`$, the set (8.17) is the same as the corresponding set of generators of the algebra $``$($`H(\gamma )`$ in pureghost number 0) for a system of spatial two-forms $`\{A_{ij}^\alpha A_{ij}^A,_0A_{ij}^A,_{00}A_{ij}^A,\mathrm{}\}`$ in 5 dimensions. Consequently, we can simply use the results demonstrated in for a system of $`p`$-forms in any dimension. $`\mathrm{}`$
We assumed before that $`f^q`$ is of degree $`q<6`$, hence $`\stackrel{~}{Q}_I`$ is of degree $`<5`$. Thus, (8.15) implies
$$\stackrel{~}{Q}_{I_m}^{(m)}=\stackrel{~}{d}\stackrel{~}{R}_{I_m}^{(m)},$$
(8.18)
where $`\stackrel{~}{R}_{I_m}^{(m)}`$ is a spatial form which only depends on the variables $`\chi `$. There is no exterior polynomial in the curvatures in $`\stackrel{~}{Q}_{I_m}^{(m)}`$ because $`\stackrel{~}{Q}_{I_m}^{(m)}`$ has strictly positive antighost number. We can therefore conclude that $`f^q`$ is trivial in $`H^q(\gamma d)`$ and can be eliminated by redefining $`e^{q+1}`$. The true bottom is then one step higher. We can proceed in the same way until we arrive at $`\gamma a^{}_{}{}^{}p=0`$ with $`a^{}_{}{}^{}p=a^p+dg^{p1}`$. This can be translated into the following theorem
###### Theorem 8.2
Let be a local form $`a`$ of antighost number $`0`$ fulfilling $`\gamma a+db=0`$. There exists a local form $`c`$ such as $`a^{}:=a+dc`$ satisfies $`\gamma a^{}=0`$.
## 9 Cohomology of $`\gamma `$ modulo $`d`$ at zero antighost number
Now, we want to study $`H^{6,0}(\gamma d)`$ in pureghost number 0. Let be $`a^{(6,0)}𝒜`$ of form degree 6, of antighost and pureghost number 0, and fulfilling $`\gamma a^{(6,0)}+da^{(5,1)}=0`$. If $`a^{(5,1)}`$ is trivial $`\gamma `$ modulo $`d`$, this equation reduces to $`\gamma a^{(6,0)}+db^{(5,0)}=0`$, which gives $`a^{(6,0)}=f(_{\mu _1\mathrm{}\mu _k}F_{ijk}^A)d^6x`$ plus a term trivial in the cohomology of $`\gamma `$ modulo $`d`$.
Otherwise, we can derive the non trivial descent equations
$`\gamma a^{(6,0)}`$ $`+`$ $`da^{(5,1)}=0`$ (9.1)
$`\gamma a^{(5,1)}`$ $`+`$ $`da^{(4,2)}=0`$
$`\mathrm{}`$
$`\gamma a^{(7g,g1)}`$ $`+`$ $`da^{(6g,g)}=0`$ (9.3)
$`\gamma a^{(6g,g)}`$ $`=`$ $`0,`$ (9.4)
because $`pureghost(\gamma a^{(6i,i)})>0`$ eliminates the constants. If $`a^{(6g,g)}`$ is trivial $`\gamma `$ modulo $`d`$, the bottom is really one step higher.
Eq. (9.4) implies that
$$a^{(6g,g)}=\underset{I}{}(\stackrel{~}{P}_I^{6g}(\chi )+dx^0\stackrel{~}{Q}_I^{5g}(\chi ))\omega ^I+\gamma b^{(6g,g1)},$$
(9.5)
where $`\stackrel{~}{P}_I^{6g}`$ and $`\stackrel{~}{Q}_I^{5g}`$ are local spatial forms, the superscript giving the form degree. Because the pureghost number of $`\eta `$ is two, $`a^{(6g,g)}`$ is non trivial only for $`g`$ even. So, three cases are of interest: $`g=0,2,4`$.
The case $`g=0`$ corresponds to $`\gamma a^{(6,0)}=0`$ and has been already studied so let us assume $`g>0`$. The equations (9.3) and (9.5) imply together
$$\underset{I}{}(_0\stackrel{~}{P}_I^{6g}\stackrel{~}{d}\stackrel{~}{Q}_I^{5g})\omega ^I+\underset{I}{}\stackrel{~}{P}_I^{6g}_0\omega ^I=\gamma \beta .$$
(9.6)
Repeating the same analysis as for the equation (8.8), we arrive at the conclusion that $`\stackrel{~}{P}_I^{6g}`$ is trivial in the invariant cohomology of $`\stackrel{~}{d}`$ (or vanishes) and can thus be removed from $`a^{(6g,g)}`$ by the addition of trivial terms in the cohomology of $`\gamma `$ modulo $`d`$ and a redefinition of $`\stackrel{~}{Q}_I^{5g}`$. The case $`g=6`$ is then eliminated because in that case $`\stackrel{~}{Q}_I^{5g}`$ is not present at all. Hence, there remains only two cases to examine: $`g=2`$ and $`g=4`$.
Once $`\stackrel{~}{P}_I^{6g}`$ is removed, the equation (9.6) gives $`\stackrel{~}{d}\stackrel{~}{Q}_I^{5g}=0`$. Using the invariant Poincaré lemma, we find $`\stackrel{~}{Q}_I^{5g}=\stackrel{~}{R}_I^{5g}(F^{A(l)})+\stackrel{~}{d}\stackrel{~}{S}_I^{(4g,g)}(\chi )`$. Hence, the form of the bottom is
$$a^{(6g,g)}=dx^0\underset{I}{}\stackrel{~}{R}_I^{5g}(F^{A(l)})\omega ^I+\gamma b^{(6g,g1)}+dc^{(5g,g)}.$$
(9.7)
But $`F^{A(l)}`$ is of form degree 3, thus if $`g=4`$, $`\stackrel{~}{R}_I^{5g}`$ must be a constant spatial $`1`$-form. In that instance, the $`\omega ^I`$ must be quadratic in the ghosts $`\eta ^{A(l)}`$. The lift of such a bottom is obstructed (i.e., leads to no $`a^{6,0}`$) unless it is trivial (see ), so that the case $`g=4`$ need not be considered. \[In the algebra of $`x`$-dependent local forms, the argument is simpler: the bottom is always trivial and removable since it involves a constant 1-form, which is trivial.\]
It only remains to examine the case $`g=2`$. $`\stackrel{~}{R}`$ must then be a 3-form. One can take $`\stackrel{~}{R}`$ linear in $`F^{A(l)}`$. In that case, the lift gives Chern-Simons terms, which are linear combinations of $`dx^0F^{A(l)}A^{B(m)}`$, with $`A^{B(m)}=\frac{1}{2}A_{ij}^{B(m)}dx^idx^j`$. Or one can take $`\stackrel{~}{R}`$ to be a constant 3-form. The corresponding deformation is linear in the 2-form $`A^{A(l)}`$ with coefficients that are constant forms. This second possibility is not $`SO(5)`$ invariant and leads to equations of motion that are not Lorentz invariant. It will not be considered further.
Dropping the latter possibility, all these results can be summarized in the
###### Theorem 9.1
The non trivial elements of $`H_0^{6,0}(\gamma d)`$ are of two types: (i) those that descend trivially; they are of the form $`f(_{\mu _1\mathrm{}\mu _k}B_{ij}^A)d^6x`$; (ii) those that descend non trivially; they are linear combinations of the Chern-Simons terms $`_0^lB^{Aij}_0^mA_{ij}^Bd^6x`$.
Note that the kinetic term in the free action is precisely of the Chern-Simons type (with $`l=0`$ and $`m=1`$).
## 10 Invariant cohomology of $`\delta `$ modulo $`\stackrel{~}{d}`$ in antighost number $`2,4,6,\mathrm{}`$
To pursue the analysis, we need some results on the cohomology of the Koszul-Tate differential $`\delta `$ as well as on its mod-$`d`$ and mod-$`\stackrel{~}{d}`$ cohomologies.
We can rewrite the action of the Koszul-Tate differential in the following way
$`\delta A_{ij}^{A(l)}`$ $`=`$ $`\delta C_i^{A(l)}=\delta \eta ^{A(l)}=0,`$ (10.1)
$`\delta A^{A(l)ij}`$ $`=`$ $`2_kF^{A(l)kij}ϵ^{ijklm}_kA_{lm}^{A(l+1)},`$ (10.2)
$`\delta C^{A(l)i}`$ $`=`$ $`_jA^{A(l)ij},`$ (10.3)
$`\delta \eta ^{A(l)}`$ $`=`$ $`_iC^{A(l)i}.`$ (10.4)
If we regard $`A`$ and $`l`$ as only one label, these equations corresponds to an infinite number of coupled non-chiral 2-forms in 5 dimensions.
It is useful to introduce a degree $`N`$ defined as
$`N(\mathrm{\Phi }_M^{})=1,N(\mathrm{\Phi }^M)=0,`$ (10.5)
$`N(_k)=1,N(_0)=0`$ (10.6)
$`N(dx^\mu )=0.`$ (10.7)
$`N`$ counts the number of spatial derivatives as well as the antifields (with equal weight given to each). According to this degree, $`\delta `$ decomposes as $`\delta _0+\delta _1`$. The differential $`\delta _1`$ acts exactly in the same way as the Koszul-Tate differential for a system of free 2-forms in 5 dimensions.
We are now able to prove the
###### Theorem 10.1
$`H_i(\delta )=0`$ for $`i>0`$, where $`i`$ is the antighost number, i.e, the cohomology of $`\delta `$ is empty in antighost number strictly greater than zero.
* From , we know that $`H_i(\delta _1)=0`$. Let be $`a𝒜`$ a $`\delta `$-closed local function of antighost number $`i>0`$. We decompose $`a`$ according to the degree $`N`$
$$a=a_1+\mathrm{}+a_m.$$
(10.8)
The expansion stops because $`a`$ is polynomial in the antifields and the derivatives. Furthermore, $`a_0=0`$ because $`antigh(a)=i>0`$. The equation $`\delta a=0`$ gives in $`N`$-degree $`m+1`$: $`\delta _1a_m=0`$. But $`H_i(\delta _1)=0`$, hence $`a_m=\delta _1b_{m1}`$. We can define an $`a^{}`$ as being
$$a^{}=a\delta b_{m1}=a_1+\mathrm{}+a_{m2}+a_{m1}^{},$$
(10.9)
with $`a_{m1}^{}=a_{m1}\delta _0b_{m1}`$. We can proceed in the same way as before with $`a^{}`$, whose component of higher $`N`$-degree is of degree less than $`m`$. We will then find a new $`a^{}`$ of highest degree less than $`m1`$, and so on, each time lowering the $`N`$-degree. After a finite number of steps, we arrive at $`a^{^{}}=a_1^{^{}}=a\delta b`$. Then, $`\delta a=0`$ implies $`\delta _1a_1^{}=0`$. Hence, $`a_1^{}=\delta _1b_0=\delta b_0`$ because $`\delta _0\mathrm{\Phi }^M=0`$. In conclusion $`a=\delta b`$, with $`b=b_0+\mathrm{}+b_{m1}`$. $`\mathrm{}`$
Of course, this theorem is really a consequence of general known results on the cohomology of the Koszul-Tate differential. It simply confirms, in a sense, that we have correctly taken into account all gauge symmetries and reducibility identities in constructing the antifield spectrum.
The cohomological space $`H_k^{5,inv}(\delta \stackrel{~}{d})`$ is defined as $`H_k^5(\delta \stackrel{~}{d})`$ in the space of local spatial forms that belongs to $``$, i.e., that are invariant. We want to compute it for $`k`$ even and $`0`$. To do this, we will proceed as in the proof of theorem 10.1. We first prove the requested result for $`\delta _1`$; we then use “cohomological perturbation” techniques to extend the result to $`\delta `$.
###### Lemma 10.1
For $`k=2,4,\mathrm{}`$
$$H_k^{5,inv}(\delta _1\stackrel{~}{d})=0.$$
(10.10)
Again, this result is simply a particular case of more general results, which were previously known, but for completeness, we prove it here.
* Firstly, the theorem 9.1 of says that for a linear gauge theory of reducibility order $`p`$ in $`n`$ dimensions $`H_k^n(\delta d)=0`$ for $`k>p+2`$. A system of abelian spatial 2-forms in 5 dimensions is a linear gauge theory of reducibility order 1 (see section 3), thus, we can state that $`H_k^5(\delta _1\stackrel{~}{d})=0`$ for $`k>3`$.
Secondly, the theorem 7.4 of gives here : $`H_2^5(\delta _1\stackrel{~}{d})=0`$.
Finally, the theorem 10.1 of says that for a system of space-time p-form gauge fields of the same degree $`H_k^n(\delta d)H_k^{n,inv}(\delta d)`$ for $`k>0`$. For the system under consideration here, this can be translated into: $`H_k^5(\delta _1\stackrel{~}{d})H_k^{5,inv}(\delta _1\stackrel{~}{d})`$ for $`k>0`$. Putting all these results together completes the proof. $`\mathrm{}`$
Let be $`a^5(\chi )`$ a local spatial $`5`$-form in $``$ of strictly positive and even antighost number, satisfying
$$\delta a^5(\chi )+\stackrel{~}{d}b^4(\chi )=0.$$
(10.11)
We can decompose $`a^5`$ and $`b^4`$ according to the degree $`N`$
$`a^5=a_1^5+\mathrm{}+a_n^5,`$ (10.12)
$`b^4=b_1^4+\mathrm{}+b_m^4.`$ (10.13)
$`a_0^5=0`$ and $`b_0^4=0`$ because $`a^5`$ and $`b^4`$ are of antighost number $`>0`$. We can always suppose $`mn`$ because if $`m>n`$, (10.11) gives in $`N`$-degree $`m+1`$: $`\stackrel{~}{d}b_m^4=0`$. Using the invariant Poincaré lemma, this yields $`b_m^4=\stackrel{~}{d}c_{m1}^3`$. Hence, $`b_m^4`$ only contributes to $`b^4`$ by a $`\stackrel{~}{d}`$-trivial term which can be eliminated. Proceeding in the same way until $`m=n`$, we arrive at the equation
$$\delta _1a_n^5(\chi )+\stackrel{~}{d}b_n^4(\chi )=0.$$
(10.14)
It has already been noticed above that the algebra $``$ without dependence on $`dx^0`$ is the same as for a system of spatial 2-forms. We can thus use the lemma 10.1 in (10.14) to find that
$`a_n^5(\chi )=\delta _1e_{n1}^5(\chi )+\stackrel{~}{d}f_{n1}^4(\chi ).`$ (10.15)
Therefore, $`a^{}_{}{}^{}5=a^5\delta e_{n1}^5\stackrel{~}{d}f_{n1}^4`$ satisfies the same properties as $`a^5`$, except that its component of highest $`N`$-degree is of degree $`<n`$. We can now apply the same reasoning as before to $`a^{}_{}{}^{}5`$, and so on, until we arrive at
$$a^{}_{}{}^{}5=a_1^{}_{}{}^{}5=a^5\delta (\underset{i=1}{\overset{n1}{}}e_i^5)\stackrel{~}{d}(\underset{i=1}{\overset{n1}{}}f_i^4)$$
(10.16)
This leads to
$$a_1^{}_{}{}^{}5=\delta _1e_0^5(\chi )+\stackrel{~}{d}f_0^4(\chi ).$$
(10.17)
But $`\delta _1e_0^5=\delta e_0^5`$ because $`\delta _0\mathrm{\Phi }^M=0`$. Eventually, we have $`a^5=\delta e^5(\chi )+\stackrel{~}{d}f^4(\chi )`$, with $`e^5=\underset{i=0}{\overset{n1}{}}e_i^5`$ and $`f^4=\underset{i=0}{\overset{n1}{}}f_i^4`$. This gives the awaited theorem:
###### Theorem 10.2
For $`k=2,4,\mathrm{}`$
$$H_k^{5,inv}(\delta \stackrel{~}{d})=0.$$
(10.18)
## 11 Decomposition of the Wess-Zumino equation
We now have all the necessary tools to solve the Wess-Zumino consistency condition that controls the consistent deformations (to first-order) of the action,
$$sa^6+db^5=0,$$
(11.1)
where $`a^6`$ and $`b^5`$ are local forms of respective form degrees 6 and 5, and ghost number 0 and 1. These forms are defined up to the following allowed redefinitions
$`a^6a^6+sf^6+dg^5`$ (11.2)
$`b^5b^5+sg^5+dh^4,`$ (11.3)
which preserve (11.1). We can decompose $`a^6`$ and $`b^5`$ according to antighost number, which gives
$`a^6`$ $`=`$ $`a_0^6+\mathrm{}+a_k^6,`$ (11.4)
$`b^5`$ $`=`$ $`b_0^5+\mathrm{}+b_q^5,`$ (11.5)
with $`a_k^60`$.
We suppose $`k>0`$ and we will show that $`a_k^6`$ can be eliminated if we redefine $`a^6`$ in an appropriate way. In antighost number $`k`$, the equation (11.1) just reads
$$\gamma a_k^6+db_k^5=0.$$
(11.6)
We can always assume $`kq`$ because if $`q>k`$, the equation (11.1) gives in highest antighost number $`db_q^5=0`$. Using the algebraic Poincaré lemma, we find that $`b_q^5=dc_q^4`$. Hence, we can remove the component $`b_q^5`$ up to a $`d`$-trivial redefinition of $`b^5`$.
From the theorems 7.1 and 8.2, we know that Eq. (11.6) implies
$$a_k^6=\underset{I}{}P_I(\chi )\omega ^I+\gamma f_k^6+dg_k^5.$$
(11.7)
The $`\gamma `$ modulo $`d`$ trivial part of $`a_k^6`$ can be eliminated by redefining $`a^6`$ in the following way
$$a^6a^6sf_k^6dg_k^5.$$
(11.8)
We notice that $`H_k^{6,0}(\gamma )`$ is non trivial only in even antighost number $`k`$ (because $`\eta `$ is of pureghost number 2). This implies that we can assume $`k`$ to be even.
The Wess-Zumino consistency condition in antighost number $`k1`$ is
$$\gamma a_{k1}^6+\delta a_k^6+db_{k1}^5=0.$$
(11.9)
The term $`b_{k1}^5`$ is invariant because (11.9) implies $`d(\gamma b_{k1}^5)=0`$. Therefore, the algebraic Poincaré lemma gives $`\gamma b_{k1}^5+dc_{k1}^4=0`$ because $`k>1`$. From ¿From the theorem 8.2 we know that we can suppose $`\gamma b_{k1}^5=0`$ without affecting $`a^6`$. Furthermore, if $`b_{k1}^5=\gamma c_{k1}^5`$ we can eliminate $`b_{k1}^5`$ by redefining $`b^5`$ in the following way: $`b^5b^5sc_{k1}^5`$, which does not modify $`a_k^6`$.
Therefore, we can assume
$`a_k^6`$ $`=`$ $`{\displaystyle \underset{I}{}}dx^0\stackrel{~}{P}_I^5\omega ^I,`$ (11.10)
$`b_{k1}^5`$ $`=`$ $`{\displaystyle \underset{I}{}}(\stackrel{~}{Q}_I^5+dx^0\stackrel{~}{R}_I^4)\omega ^I.`$ (11.11)
The $`\stackrel{~}{P}_I^5`$, $`\stackrel{~}{Q}_I^5`$, and $`\stackrel{~}{R}_I^4`$ are local spatial forms belonging to $``$.
Inserting (11.10) and (11.11) in (11.9), we find
$`\gamma a_{k1}^6`$ $`=`$ $`{\displaystyle \underset{I}{}}\{\stackrel{~}{d}\stackrel{~}{Q}_I^5\omega ^I\gamma [(\stackrel{~}{Q}_I^5+dx^0\stackrel{~}{R}_I^4)\widehat{\omega }^I]`$ (11.13)
$`+dx^0[(\delta \stackrel{~}{P}_I^5+\stackrel{~}{d}\stackrel{~}{R}_I^4_0\stackrel{~}{Q}_I^5)\omega ^I\stackrel{~}{Q}_I^5_0\omega ^I]\},`$
with $`\stackrel{~}{d}\omega ^I=\gamma \widehat{\omega }^I`$. This implies that
$$\underset{I}{}[(\delta \stackrel{~}{P}_I^5+\stackrel{~}{d}\stackrel{~}{R}_I^4_0\stackrel{~}{Q}_I^5)\omega ^I\stackrel{~}{Q}_I^5_0\omega ^I]=\gamma \beta .$$
(11.14)
If we analyse this equation in the same way as the equation (8.8), we can prove that $`\stackrel{~}{Q}_I^5=\delta \stackrel{~}{P}_I^5+\stackrel{~}{d}\stackrel{~}{R}_I^4`$ (or simply vanishes). Inserting these equations in (11.11), we find that $`b_{k1}^5`$ is of the form
$$b_{k1}^5=\delta c_k^5+de_{k1}^4+\gamma f_{k1}^5+dx^0\underset{I}{}\stackrel{~}{R}_I^{}_{}{}^{}4(\chi )\omega ^I,$$
(11.15)
where $`c_k^5`$ and $`e_{k1}^4`$ belong to $`H(\gamma )`$. In conclusion, we can eliminate $`\stackrel{~}{Q}_I^5`$ from $`b_{k1}^5`$ by redefining $`a^6`$ and $`b^5`$ in the following way
$`a^6a^6d(c_k^5+f_{k1}^5),`$ (11.16)
$`b^5b^5s(c_k^5+f_{k1}^5)de_{k1}^4,`$ (11.17)
which does not affect the condition $`\gamma a_k^6=0`$, because $`\gamma c_k^5=0`$.
Therefore, we can finally assume
$$a_k^6=\underset{I}{}dx^0\stackrel{~}{P}_I^5(\chi )\omega ^I,b_{k1}^5=\underset{I}{}dx^0\stackrel{~}{R}_I^4(\chi )\omega ^I.$$
(11.18)
The equation (11.9) becomes
$$\gamma a_{k1}^{^{}}+dx^0\underset{I}{}(\delta \stackrel{~}{P}_I^5(\chi )+\stackrel{~}{d}\stackrel{~}{R}_I^4(\chi ))\omega ^I=0,$$
(11.19)
which implies that $`\delta \stackrel{~}{P}_I^5(\chi )+\stackrel{~}{d}\stackrel{~}{R}_I^4(\chi )=0`$. We know that we are in even antighost number, thus we can use the theorem 10.2 to find that $`\stackrel{~}{P}_I^5=\delta \stackrel{~}{S}_I^5(\chi )+\stackrel{~}{d}\stackrel{~}{T}_I^4(\chi )`$. Hence,
$$a_k^6=sf_{k+1}^6+dg_k^5+\gamma h_k^6,$$
(11.20)
where we have defined
$`f_{k+1}^6=dx^0{\displaystyle \underset{I}{}}\stackrel{~}{S}_I^5\omega ^I,g_k^5=dx^0{\displaystyle \underset{I}{}}\stackrel{~}{T}_I^4\omega ^I,`$ (11.21)
$`h_k^6=dx^0{\displaystyle \underset{I}{}}\stackrel{~}{T}_I^4\widehat{\omega }^I,\stackrel{~}{d}\omega ^I=\gamma \widehat{\omega }^I.`$ (11.22)
Thus $`a_k^6`$ can be completely eliminated by redefing $`a^6`$ as
$$a^{}_{}{}^{}6=a^6s(f_{k+1}^6+h_k^6)dg_k^5,$$
(11.23)
which only affects the components of antighost number $`<k`$. Repeating the argument at lower antighost numbers enables one to remove successively $`a_{k1}`$, $`a_{k2}`$, …, up to $`a_1`$. This completes the proof of the fact that there is no non trivial dependence on the antifields for the elements of $`H^{6,0}(sd)`$.
For antifield-independent local forms, the cocycle condition $`H^{6,0}(sd)`$ reduces to the cocycle condition for $`H^{6,0}(\gamma d)`$. Furthermore, $`\gamma `$-exact (mod-$`d`$) solutions are also $`s`$-exact. Thus, we are led to consider $`H^{6,0}(\gamma d)`$. This cohomology is given by the theorem 9.1. \[The terms in that cohomology that vanish on-shell are trivial in the $`s`$-cohomology.\] Thus, the only consistent deformations of the free action for a system of abelian chiral $`2`$-forms are either functions of the curvatures or of the Chern-Simons type. In both cases, the integrated deformations are off-shell gauge invariant and yield no modification of the gauge transformations.
## 12 Final comments and conclusions
We have shown that the most general first-order consistent deformation of a set of free chiral $`2`$-forms cannot modify (non trivially) the original gauge transformations and a fortiori, their algebra, which remains abelian. Thus, there is no room for a non-abelian, local, generalization of the theory analogous to the Yang-Mills construction.
This result holds in fact to all orders, since the allowed deformations involve the gauge-invariant curvatures or Chern-Simons terms. The addition of such terms to the original action yields a new action which is evidently gauge-invariant under the original gauge transformations to all orders.
One can show along identical lines that the rigidity of the gauge symmetries is actually valid for a set of chiral $`2p`$-forms in $`2p+2`$ dimensions, for any $`p>0`$. If one includes other fields, one may deform the gauge transformations, but the possibilities are severely limited . For instance, in 10 dimensions, the only couplings of a chiral 4-form to 2-forms are those present in type IIB supergravity.
Acknowledgments: X.B. and M.H. are supported in part by the “Actions de Recherche Concertées” of the “Direction de la Recherche Scientifique - Communauté Française de Belgique”, by IISN - Belgium (convention 4.4505.86) and by Proyectos FONDECYT 1970151 and 7960001 (Chile). A.S. is supported in part by the FWO and by the European Commission TMR programme ERBFMRX-CT96-0045 in which he is associated to K. U. Leuven.
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# “Scars” in parametrically excited surface waves
## Abstract
We consider the Faraday surface waves of a fluid in a container with a non-integrable boundary shape. We show that, at sufficiently low frequencies, the wave patterns are “scars” selected by the instability of the corresponding periodic orbits, the dissipation at the container side walls, and interaction effects which reflect the nonlinear nature of the Faraday waves. The results explain the observation of a limited number of scars with anomalous strengths in recent experiments by Kudrolli, Abraham and Gollub.
A prominent feature of wave chaos is the “scar” phenomenon. Namely, wave functions are concentrated along the short periodic orbits of the underlying classical system. In general, the phenomenon appears whenever the Eikonal approximation holds, and the corresponding rays exhibit ballistic chaotic dynamics. Since their discovery by numerical studies, scars have been observed experimentally in microwave cavities, tunneling diodes, and capillary waves.
However, theoretical studies of scars have not addressed several factors which may be present in real systems, such as dissipation, external forcing, nonequilibrium effects, and interactions associated with nonlinear terms of the wave equations. In this Letter we study a system where these factors are significant: the surface waves in a container of non-integrable shape.
This work was motivated by recent experiments of Kudrolli et al. who studied Faraday waves in a stadium shaped container. The surface waves were excited parametrically by oscillating the container in direction perpendicular to the fluid surface at rest. The experimental findings are quite intriguing, see Fig. 1: I. Out of all known scars of the stadium billiard only three scarring patterns were identified; II. These patterns appear in, approximately, 90$`\%`$ of the cases, and their magnitudes, compared to the random background, is unexpectedly large; III. While probability densities associated with eigenstates of the stadium billiard were always symmetric, nonsymmetric wave patterns have been also observed.
The purpose of this Letter is to explain this behavior. In particular, we show that the excitation threshold for observation of a scarring pattern is governed by the Lyapunov exponent of the corresponding periodic orbit, bulk dissipation, and dissipation at the boundary of the vessel. These factors suppress most of the scarring patterns as well as the random background. The amplitudes of those scars which still can be excited are determined by nonlinear effects. These effects (interactions) also break the symmetry between degenerate scars, which results in nonsymmetric wave patterns.
To begin with, consider the surface waves of an infinite system with an infinite fluid depth. Within the linear approximation, the amplitude of a plain wave, $`a_𝐤e^{i𝐤𝐫}`$,
satisfies the equation:
$`{\displaystyle \frac{^2a_𝐤}{t^2}}+2\gamma _b{\displaystyle \frac{a_𝐤}{t}}+\omega _k^2a_𝐤=0,`$ (1)
where $`k`$ is the wave number, $`\gamma _b=\nu k^2`$ is the bulk dissipation rate, and $`\nu `$ is the viscosity. The angular frequency, $`\omega _k`$, satisfies the dispersion relation:
$`\omega _k^2=gk+{\displaystyle \frac{\mathrm{\Sigma }}{\rho }}k^3,`$ (2)
where $`g`$ is the acceleration of gravity, $`\mathrm{\Sigma }`$ is the surface tension, and $`\rho `$ is the fluid density.
In a finite system, the boundary causes two effects. First, it leads to quantization of the wave vector $`kk_\alpha `$ (thus also $`\omega _k\omega _\alpha `$). When the container is rim full, the fluid surface is believed to be pinned to boundary. Thus the eigenmodes of the surface waves, $`\psi _\alpha (𝐫)`$, are solutions of the Helmholtz equation in a domain of the container shape with Dirichlet boundary conditions.
The second effect of the boundary is an additional dissipation due to the viscous flow in the vicinity of the container side walls. This dissipation, in a system with size $`L`$, is of the order $`\sqrt{\nu \omega }/L`$. For large containers it is negligible compared to the bulk dissipation rate. However, in small systems and for specific wave patterns, this dissipation becomes comparable to $`\gamma _b`$.
Parametric excitation of surface waves, due to vertical shaking, is accounted by introducing time dependent acceleration, $`gg+A\mathrm{cos}(2\omega t)`$, into Eq. (1). The parameter space $`(\omega ,A)`$ of the resulting Mathieu equation is characterized by tongues of instability where $`a_𝐤(t)`$ grows exponentially in time. This growth saturates due to nonlinear effects.
To take these interactions into account, it is convenient to employ a formalism analogous to second quantization. For simplicity, we first consider an infinite system. Let us introduce the amplitude variables, $`a_𝐤^{}`$ and $`a_𝐤`$, analogous to boson creation and annihilation operators at momentum state $`𝐤`$. These variables are $`c`$-numbers which can be viewed as boson operators in the limit of a large number of particles. The equations of motion, then, take the form:
$$\frac{}{t}a_𝐤+\gamma _ba_𝐤=i\frac{}{a_𝐤^{}},$$
(3)
where $``$ is the Hamiltonian of the system which consists of three terms: $`=_0+_P+_{\text{int}}`$. The first component, $`_0=_𝐤\omega _ka_𝐤^{}a_𝐤`$, describes the unforced surface waves in the linear approximation, i.e. Eq. (1).
The second component, $`_P`$, is the pumping Hamiltonian responsible for the parametric excitation of the Faraday waves. Confining our attention to excitations near the threshold of the first instability tongue of the Mathieu equation, we can write $`_P`$ as
$$_P=\frac{1}{2}\underset{k}{}ha_𝐤^{}a_𝐤^{}e^{i2\omega t}+C.c.;h=\frac{Ak\omega ^2}{4\omega _k},$$
(4)
where $`A`$ and $`2\omega `$ are the amplitude and the frequency of the vibrations. In the quadratic approximation, $`_0+_P`$, the equations of motion (3) yield
$`a_𝐤,a_𝐤^{}e^{\sigma t},\text{where}\sigma =\gamma _b+\sqrt{h^2(\omega \omega _k)^2}.`$ (5)
Thus the threshold for instability, at resonance $`(\omega =\omega _k)`$, is $`h\gamma _b`$. Under this condition, the rate of energy pumped into the system exceeds the dissipation rate.
The third component of the Hamiltonian, $`_{\text{int}}`$, is an interaction term. In the mean field approximation it takes the form
$`_{\text{int}}={\displaystyle \underset{\mathrm{𝐤𝐤}^{}}{}}I_k^{(1)}(\theta )|a_𝐤|^2|a_𝐤^{}|^2+I_k^{(2)}(\theta )a_𝐤^{}^{}a_𝐤^{}^{}a_𝐤a_𝐤,`$ (6)
where $`I_k^{(j)}(\theta )=T_k^{(j)}(\theta )+i\mathrm{\Gamma }_k^{(j)}(\theta )`$ are complex matrix elements which depend on the length of the vectors, $`k=|𝐤|=|𝐤^{}|`$ as well as the angle, $`\theta `$, between them. The exact form of these matrix elements has been derived by Milner. Here we only note that the interaction is local and $`T`$ is of order $`\omega k^2`$, while $`\mathrm{\Gamma }\nu k^4`$.
Several mechanisms limit the exponential growth of the wave amplitude. They are related to different components of the interaction matrix elements. For instance, $`\mathrm{\Gamma }_k^{(j)}(\theta )`$ are associated with nonlinear damping, while $`T_k^{(2)}(\theta )`$ governs the missphasing mechanism. Missphasing is the phase difference between the pump and the excited wave, which appears due to the nonlinearity and suppresses the energy absorption.
The actual limiting mechanism is determined by the system size and the excitation frequency. In large systems, boundary effects can be neglected, and patterns are selected by the local properties of the wave equation. This is the regime of pattern formation, where the real parts of the interaction terms merely renormalize $`h`$ and $`\omega _k`$. The dominant limiting mechanism, in this case, is nonlinear damping. In small systems, contrarily, the wave patterns are globally determined by the boundary, and missphasing turns out to be more effective.
A system can be considered as small provided that the dissipation time, $`1/\nu k^2`$, exceeds the typical time it takes a surface wave to cross the system. For capillary waves ($`k>\sqrt{g\rho /\mathrm{\Sigma }}`$) this requirement yields the inequality
$`\omega <\omega _c={\displaystyle \frac{\mathrm{\Sigma }}{L\rho \nu }}.`$ (7)
In what follows we assume that this condition is satisfied, the wave patterns are determined by the boundary, and missphasing is the dominant limiting mechanism.
Consider now the Faraday waves in a container of non-integrable boundary shape having no discrete symmetries. The modes of the corresponding linear wave equation, $`\psi _\alpha (𝐫)`$, are eigenstates of a billiard system of the shape of the container. Thus, $`\psi _\alpha (𝐫)`$ are approximately random Gaussian functions scarred by the periodic orbits of the underlying classical dynamics. In principle, one can use $`\psi _\alpha (𝐫)`$ as a basis for the many body approach described above. However, this basis is inconvenient. In the experiment, several modes are excited simultaneously (approximately 10), and the resulting wave patterns are determined by properties of the short time dynamics of the system. Therefore, it is more natural to use scars as the skeleton of the excited wave patterns.
Below we develop a “scar representation”. Instead of dealing with the amplitudes, $`n_\alpha =|a_\alpha |^2`$, of the eigenfunctions, $`\psi _\alpha (𝐫)`$, we consider the amplitudes of scars. Using Wigner representation, one can present $`_0`$ as
$`_0={\displaystyle \underset{\alpha }{}}\omega _\alpha n_\alpha ={\displaystyle ϵ𝑑ϵ𝑑𝐱W(𝐱;ϵ)n(𝐱)}`$ (8)
where $`𝐱=(𝐫,𝐤)`$ is a phase space coordinate, $`n(𝐱)`$ is the wave amplitude at $`𝐱`$, and $`W(𝐱;ϵ)`$ is the Wigner representation of the imaginary part of the Green function:
$`W(𝐱;ϵ)={\displaystyle \underset{\alpha }{}}\delta (ϵ\omega _\alpha ){\displaystyle 𝑑𝐪e^{i𝐤𝐪}\psi _\alpha (𝐫\frac{𝐪}{2})\psi _\alpha ^{}(𝐫+\frac{𝐪}{2})}.`$ (9)
Assuming the semiclassical limit $`kL1`$ as well as the capillary regime, $`k>(g\rho /\mathrm{\Sigma })^{1/2}`$, we can approximate $`W(𝐱;ϵ)`$ as
$`W(𝐱;ϵ)\delta [ϵH_0(𝐱)]\left(1+{\displaystyle \underset{p}{}}W_p(𝐱;ϵ)\right),`$ (10)
where $`H_0(𝐱)=(\mathrm{\Sigma }k^3/\rho )^{1/2}`$. The first term in the parenthesis is the average of the random excitation over the energy shell. The second term is a sum over the primitive periodic orbits of the system, $`p`$. Each term of this sum represents the corresponding scar contribution:
$`W_p(𝐱;ϵ)4\pi \text{Re}{\displaystyle \frac{e^{\frac{u_p}{2}+i(S_p\phi _p)}}{1e^{\frac{u_p}{2}+i(S_p\phi _p)}}}\stackrel{~}{\delta }(𝐗_p).`$ (11)
Here $`S_p(\omega )=l_p(\rho \omega ^2/\mathrm{\Sigma })^{1/3}`$ is the action of the periodic orbit with length $`l_p`$, instability exponent $`u_p`$, and Maslov phase $`\phi _p`$. $`𝐗_p`$ denotes a coordinate on the Poincare section of the periodic orbit. The function $`\stackrel{~}{\delta }(𝐗_p)`$ is localized along the $`p`$-th periodic orbit with the width of order $`\sqrt{L/k}`$. Its integral is unity so that $`\stackrel{~}{\delta }(𝐗_p)`$ becomes a $`\delta `$-function in the classical limit (i.e. at $`k\mathrm{}`$).
Substituting (10) in (8), we obtain
$`_0={\displaystyle 𝑑ϵϵ\overline{d}(ϵ)\overline{n}(ϵ)}+{\displaystyle \underset{p,m}{}}\omega _p^{(m)}n_p,`$ (12)
where $`\overline{d}(ϵ)=𝒜(\rho ^2ϵ/\mathrm{\Sigma }^2)^{\frac{1}{3}}/3\pi `$ is the mean density of states in a container of area $`𝒜`$,
$`\overline{n}={\displaystyle \frac{1}{\overline{d}}}{\displaystyle 𝑑𝐱\delta [ϵH_0(𝐱)]n(𝐱)}`$ (13)
is the average amplitude of the random background, and
$`n_p={\displaystyle \frac{1}{t_p}}{\displaystyle 𝑑tn(𝐱_p(t))}`$ (14)
is the amplitude of the scar $`p`$. Here $`t_p`$ is the period of the corresponding orbit, and $`𝐱_p(t)`$ is its coordinates parameterized by the time. The eigenfrequencies of the scars, $`\omega _p^{(m)}`$, are the poles of (11), i.e. the solutions of the equation $`S_p(\omega _p^{(m)})=2\pi m+\phi _piu_p/2`$. Expanding to leading order in $`u_p`$ we have
$`\omega _p^{(m)}\left({\displaystyle \frac{2\pi m+\phi _p}{l_p(\rho /\mathrm{\Sigma })^{1/3}}}\right)^{3/2}i{\displaystyle \frac{\lambda _p}{2}},`$ (15)
where $`\lambda _p=u_p/t_p`$ is the Lyapunov exponent of the orbit.
Now one can substitute the Hamiltonian (12) into the equations of motion (3) and see that the Lyapunov exponent plays the role of a dissipation. Indeed, $`\lambda _p/2`$ is the rate at which the “particle escapes from the scar”. This rate should be added to the dissipation rate. The Lyapunov exponents of few short periodic orbits of the stadium billiard are tabulated in Table I. Comparison of these values with the experimental bulk dissipation rate, $`\gamma _b2`$ sec<sup>-1</sup>, shows that their effect is significant.
Next, we consider the dissipation at the side walls of the container. This dissipation has been calculated for particular cases, e.g., the square and the circular containers. From these results it is evident that boundary dissipation depends on the particular form of the wave pattern. To obtain the result for a general container shape, and arbitrary scarring pattern, we first note that boundary dissipation takes place very close to the boundary, at distance of order $`l_D=\sqrt{2\nu /\omega }10^2`$cm. Viewing the scar as a plain wave scattered from the boundary, each collision point adds its own contribution to the dissipation rate. These contributions can be calculated by traditional methods. Assuming pinning of the fluid surface at the boundary, the resulting dissipation rate of a scar can be presented as a sum over the reflection points of the orbit from the boundary:
$`\gamma _p{\displaystyle \frac{(2\nu \omega )^{1/2}}{l_p}}{\displaystyle \underset{n}{}}{\displaystyle \frac{1\frac{1}{2}\mathrm{cos}^2(\theta _n)}{\mathrm{cos}(\theta _n)}},`$ (16)
where $`\theta _n`$ is the angle between the orbit and the normal to the boundary. Examples of $`\gamma _p`$ for various scars of the stadium billiard are presented in Table I. Note that the anomalously large dissipation rate of whispering gallery modes (patterns Nos. 8, 9) results from $`\theta _n`$ being close to $`\pi /2`$, see (16).
The boundary dissipation of the random background, $`\overline{\gamma }`$, can be obtained from the result for a long ergodic orbit. However, the result diverges since ergodicity implies that $`\theta _n`$ can become arbitrarily close to $`\pi /2`$. Cutting this divergence by the wave resolution, $`\mathrm{\Delta }\theta _n1/kL`$, we obtain
$`\overline{\gamma }{\displaystyle \frac{(2\omega \nu )^{1/2}}{L}}\mathrm{ln}(kL).`$ (17)
Collecting the various dissipation terms and the contribution of the Lyapunov exponent, the threshold for excitation of scars at resonance $`(\omega =\text{Re}\omega _p^{(m)})`$ is:
$`h>h_p=\gamma _b+\gamma _p+\lambda _p/2`$ (18)
Therefore, large boundary dissipation prevents the excitation of random Gaussian patterns, as well as other modes such as the whispering gallery modes. Moreover, increasing the boundary dissipation, e.g. by lowering the fluid level, should eventually suppress all of the scars except the horizontal one (No. 2 in Table I). This is precisely what is observed in the experiment.
Turning to the interaction effects, we first note that matrix elements off diagonal in scars, $`I_{p,p^{}}^{(j)}`$, are usually much smaller than the diagonal ones, $`I_{p,p}^{(j)}`$. The reason is that scars are localized objects, whereas the interaction between any two of them is proportional to their intersection area.
Consider, therefore, the scattering from a scar to itself. Here, the wave vectors, $`k`$ and $`k^{}`$, are either parallel or antiparallel, except for small regions of self intersections. Thus the diagonal matrix elements of a scar associated with an orbit, $`p`$, which has a time reversal counterpart $`p`$ (such as Nos. 3–9 in Table I) are $`I_{p,p}^{(1)}I_k^{(1)}(0)`$, $`I_{p,p}^{(1)}I_k^{(1)}(\pi )`$ and $`I_{p,p}^{(2)}I_k^{(2)}(0)`$. The interaction matrix elements of scars associated with self retracing orbits (such as No. 1, 2, and 10 in Table I) are $`I_{p,p}^{(1)}=I_{p,p}^{(2)}(I_k^{(1)}(0)+I_k^{(1)}(\pi )+I_k^{(2)}(0))/2`$.
If only one scar is excited, its amplitude can be evaluated exactly from the steady state limit of Eq. (3). For scars associated with self retracing orbits, the result is
$`n_p^{(1)}={\displaystyle \frac{1}{2T}}\sqrt{h^2h_p^2},`$ (19)
where $`T=\text{Re}I_{p,p}^{(1)}`$. We assumed that the system was at resonance, $`\omega =\text{Re}\omega _p`$, and neglected the imaginary part of $`I_{p,p}^{(j)}`$.
Let now two scars be in resonance with the driving frequency, one can calculate the amplitude of the scar $`p`$ in the presence of the second scar $`p^{}`$ by perturbation theory in the off diagonal matrix element between these scars. The result is
$`n_p^{(2)}n_p^{(1)}\left(1\tau _1Q\tau _2Q^2\right)`$ (20)
where $`\tau _j=\text{Re}I_{pp^{}}^{(j)}/2T`$ and $`Q=n_p^{}^{(1)}/n_p^{(1)}`$. According to Eq. (20) even a relatively weak interaction can suppress the scar with the higher threshold amplitude. Thus one should expect to observe only the dominating scar.
Finally, we comment on the role of discrete symmetries in the shape of the container. The exact eigenfunctions in this case can be classified according to the irreducible representation of the corresponding symmetry group. In the case of the stadium billiard this symmetry is not sufficient to produce an exact degeneracy of eigenvalues. However, the eigenfrequencies of scars, associated with an excitation of a set of modes, are degenerate. The scars associated with orbits Nos. 3, 4 and 11 in Table I can be used as examples. The symmetry between the scars is usually weakly broken by perturbations such as inhomogeneous pumping or deviation from an exact horizontal state of the container. Interaction, in turn, drives the system into a nonsymmetric wave pattern.
To summarize, we have shown that Faraday wave patterns in a container of non-integrable boundary shape are mainly scars. The instability threshold for these patterns is determined by bulk dissipation, Lyapunov exponents, and dissipation near the container side walls. The latter is too strong for random patterns and whispering gallery modes to appear. Nonsymmetric wave patterns, in a symmetric container, take place due to the existence of degenerate scars and interaction effects between them.
We are grateful to Arshad Kudrolli and Jerry Gollub for many discussions and communications, and for the experimental data presented in this Letter. This research was supported by Grant No. 9800065 from the USA-Israel Binational Science Foundation (BSF).
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# Precanonical perspective in quantum gravity Based on the talk at the Third Meeting on Constrained Dynamics and Quantum Gravity QG99 (Villasimius, Sardinia, Italy, Sept. 13-17, 1999). To appear in Nucl. Phys. B Proc. Suppl. (2000)
## 1 INTRODUCTION
In spite of the noticeable progress in the quantum theory of gravity during the last decade, mainly owing to the Ashtekar program of non-pertubative canonical quantum gravity and the string/M–theory, it is difficult to escape a feeling that the genuine conceptual foundations of the synthesis of general relativity and the quantum theory are still awaiting of their discovery. In particular, the drastic differences between the way the physics is described general relativistically and quantum theoretically urge us to inquire if the presently known procedures of field quantization, that is the specific way the quantum paradigm is implemented, are adequate to the problem of quantization of gravity. The particular concern is due to the different status of the time variable in quantum theory and in general relativity, in addition to the characteristic to the latter diffeomorphism covariance and a dynamical character of the space-time. The main objective of what we call the precanonical approach to field quantization is to elaborate a procedure which would treat all space-time variables on equal footing, more in accordance with the relativity theory.
The idea of precanonical approach is suggested by a long-known in the calculus of variations fact that the Hamiltonian formulation can be alternatively extended to field theory in the form of the De Donder-Weyl (DW) canonical equations
$$_\mu y^a=\frac{H}{p_a^\mu },_\mu p_a^\mu =\frac{H}{y^a},$$
(1)
where given a Lagrangian density $`L(y^a,y_\mu ^a,x^\nu )`$, a function of field variables $`y^a`$, their space-time derivatives (first jets) $`y_\mu ^a`$ and space-time variables $`x^\mu `$, one introduces new Hamiltonian-like variables $`p_a^\mu :=L/y_\mu ^a`$ (polymomenta) and $`H=H(y^a,p_a^\mu ,x^\nu ):=_\mu y^ap_a^\mu L`$ (the DW Hamiltonian function). Obviously, this formulation is manifestly covariant, it treats space and time variables on equal footing, i.e. requires no usual $`3+1`$ decomposition, and corresponds to what could be viewed as a “multi-time” generalization of the Hamiltonian formulation from mechanics to field theory. No picture of fields as infinite-dimensional mechanical systems evolving in time is implied in (1); instead, fields are described rather as systems varying in space and time. These features make formulation (1) an attractive alternative to the conventional “instantaneous” Hamiltonian formalism as a basis of quantization, especially in the context of general relativity. Besides, the obstacles to the DW Legendre transformation $`y_\mu ^ap_a^\mu `$ in general are different from the usual constraints, thus suggesting a possibility of surmounting the usual constraints analysis when quantizing.
The term “precanonical” refers to the fact that formulation (1) and the related constructions are in a sense intermediate between the covariant Lagrangian description and the “instantaneous” canonical Hamiltonian description; in mechanics precanonical structures coincide with the canonical ones while in field theory they are different.
Field quantization stemming from formulation (1) has been considered in ; its application to general relativity has been discussed in a recent preprint by the author and will be outlined below. Briefly, quantization based on DW Hamiltonian formulation (1) leads to a generalization of quantum theoretic formalism in which the space-time Clifford algebra replaces the algebra of complex numbers (=the Clifford algebra of $`(0+1)`$–dimensional space-time!) in quantum mechanics. The Clifford algebra appears when quantizing the Poisson brackets which in DW theory are defined on differential forms (c. f. ). This results in representing polymonenta by operators
$$\widehat{p}{}_{a}{}^{\mu }=i\mathrm{}\kappa \gamma ^\mu \frac{}{y^a},$$
(2)
which act on spinor (or the Clifford algebra valued) wave functions $`\mathrm{\Psi }=\mathrm{\Psi }(y^a,x^\mu )`$; $`\gamma ^\mu `$’s denote the imaginary units of the space-time Clifford algebra; the constant $`\kappa `$ of the dimension \[length\]<sup>-3</sup> ensures the dimensional consistency of (2) and is interpreted as a quantity of the ultra-violet cutoff or the fundamental length scale . The wave function is supposed to fulfill a generalized Schrödinger equation
$$i\mathrm{}\kappa \gamma ^\mu _\mu \mathrm{\Psi }=\widehat{H}\mathrm{\Psi },$$
(3)
where $`\widehat{H}`$ is the operator form of the DW Hamiltonian function. This equation was found to be consistent with several aspects of the correspondence principle , for example, it leads to an analogue of the Ehrenfest theorem and can be reduced to the field theoretic DW Hamilton-Jacobi equation (with some additional conditions) in the classical limit. Unfortunately, the details of the relationship between the standard quantum field theory and the present formulation so far remain poorly understood. A possible connection with the Schrödinger functional picture has been discussed in .
## 2 PRECANONICAL QUANTUM GENERAL RELATIVITY: AN OUTLINE
When applying the above approach to gravity the configuration space is to be a bundle of symmetric second rank tensors over the space-time and the wave function is to be a function on this space: $`\mathrm{\Psi }=\mathrm{\Psi }(g^{\mu \nu },x^\mu )`$. The Schrödinger equation for this wave function is obtained by first writing a curved space-time version of (3) and then replacing the metric and the connection by the corresponding operators. This leads to the following guess concerning the wave equation for quantum gravity within the precanonical approach:
$$i\mathrm{}\kappa \widehat{e\overline{)}}\mathrm{\Psi }=\widehat{}\mathrm{\Psi },$$
(4)
where $`\widehat{}`$ is the operator form of DW Hamiltonian density of gravity, $`:=\sqrt{g}H`$, where $`\sqrt{g}:=\sqrt{|\mathrm{det}(g_{\mu \nu })|}=:e`$, and $`\widehat{\overline{)}}:=\gamma ^\mu (_\mu +\widehat{\theta }_\mu )`$ denotes the quantized covariant Dirac operator in which $`\gamma `$-matrices are such that $`\gamma ^\mu \gamma ^\nu +\gamma ^\nu \gamma ^\mu =2g^{\mu \nu }`$ and $`\widehat{\theta }_\mu `$ is the spinor connection operator. Recall that classically $`\theta _\mu =\theta ^{\alpha \beta }{}_{\mu }{}^{}\gamma _{[\alpha }^{}\gamma _{\beta ]}`$, $`\gamma ^\mu =e_A^\mu \gamma ^A`$ and
$$\theta ^{\alpha \beta }{}_{\mu }{}^{}=g^{\nu [\beta }(\mathrm{\Gamma }_{\mu \nu }^{\alpha ]}e_A^{\alpha ]}_\mu e_\nu ^A).$$
(5)
Since the quantum gravity possesses an intrinsic fundamental length scale, the Planck length $`\mathrm{}`$, one can expect that $`\kappa \mathrm{}^3`$.
To find an operator realization of the quantities involved in (4) we first have to formulate the Einstein equations in DW Hamiltonian form. It is given by the set of equations
$$_\alpha h^{\beta \gamma }=/Q_{\beta \gamma }^\alpha ,_\alpha Q_{\beta \gamma }^\alpha =/h^{\beta \gamma },$$
(6)
where $`h^{\alpha \beta }:=\sqrt{g}g^{\alpha \beta }`$ are field variables,
$$Q_{\beta \gamma }^\alpha :=\frac{1}{8\pi G}\left(\delta _{(\beta }^\alpha \mathrm{\Gamma }_{\gamma )\delta }^\delta \mathrm{\Gamma }_{\beta \gamma }^\alpha \right)$$
(7)
are corresponding polymomenta and
$$:=8\pi Gh^{\alpha \gamma }(Q_{\alpha \beta }^\delta Q_{\gamma \delta }^\beta \frac{1}{3}Q_{\alpha \beta }^\beta Q_{\gamma \delta }^\delta )$$
(8)
is the DW Hamiltonian function of gravity. The above formulation of the Einstein equations has a deeper foundation in the theory of Lepagean equivalents in the calculus of variations .
Now, polymomenta can be quantized according to the rule (2) adapted to curved space-time
$$\widehat{Q}_{\beta \gamma }^\alpha =i\mathrm{}\kappa \gamma ^\alpha \left\{\sqrt{g}\frac{}{h^{\beta \gamma }}\right\}_{ord},$$
(9)
where the notation $`\{\mathrm{}\}_{ord}`$ refers to the ordering ambiguity of the expression inside the curly brackets. Plugging (9) into (8) we obtain
$$\widehat{}=\frac{16\pi }{3}G\mathrm{}^2\kappa ^2\{\sqrt{g}h^{\alpha \gamma }h^{\beta \delta }\frac{}{h^{\alpha \beta }}\frac{}{h^{\gamma \delta }}\}_{ord}$$
(10)
When formulating the left hand side of eq. (4) we are led to the fundamental difficulties related to the fact that (i) conceptually, the Dirac operator generally refers to a classical space-time background which is ought to be avoided in quantum gravity and (ii) technically, the last term in the spinor connection (5) cannot be expressed in terms of metric variables. We deal with these difficulties by observing that the tetrads do not enter the present DW formulation of General Relativity underlying the quantization and, therefore, can be treated only as non-quantized $`x`$-depended quantities $`\stackrel{~}{e}_A^\mu (x)`$. The correspondence principle then implies that they should be related to the mean value of the metric as follows
$$\stackrel{~}{e}{}_{A}{}^{\mu }(x)\stackrel{~}{e}{}_{B}{}^{\nu }(x)\eta ^{AB}=g^{\mu \nu }(x)$$
(11)
where
$$g^{\mu \nu }(x)=[dg^{\alpha \beta }]\overline{\mathrm{\Psi }}(g,x)g^{\mu \nu }\mathrm{\Psi }(g,x),$$
(12)
and
$$[dg^{\alpha \beta }]=g^{5/2}\underset{\alpha \beta }{}dg^{\alpha \beta }$$
(13)
is the invariant integration measure on the $`10`$-dimensional space of metric components (c.f. ). Note that (12) is well-defined mathematically as a smooth field.
To quantize the connection coefficients let us note that classically (c.f. (7))
$$\mathrm{\Gamma }_{\beta \gamma }^\alpha =8\pi G\left(\frac{2}{3}\delta _{(\beta }^\alpha Q_{\gamma )\delta }^\delta Q_{\beta \gamma }^\alpha \right).$$
(14)
Now, using (5) and (9) we can write
$`\widehat{\theta }^{\alpha \beta }{}_{\mu }{}^{}=8\pi iG\mathrm{}\kappa \{h^{\nu [\beta }({\displaystyle \frac{2}{3}}\delta _{(\mu }^{\alpha ]}\gamma ^\sigma {\displaystyle \frac{}{h^{\nu )\sigma }}}`$
$`\gamma ^{\alpha ]}{\displaystyle \frac{}{h^{\mu \nu }}})\}_{ord}+\stackrel{~}{\theta }^{\alpha \beta }_\mu (x).`$ (15)
This expression involves the ordering dependent operator part $`(\theta ^{\alpha \beta }{}_{\mu }{}^{})^{op}`$ and an auxiliary spinor connection part $`\stackrel{~}{\theta }^{\alpha \beta }{}_{\mu }{}^{}(x)`$ which (i) accounts for the term in (5) which cannot be expressed in metric variables (hence, cannot be quantized) and (ii) ensures the transformation law of $`\widehat{\theta }{}_{}{}^{\alpha \beta }_\mu `$ is that of a spinor connection. Our assumption is that $`\stackrel{~}{\theta }{}_{}{}^{\alpha \beta }{}_{\mu }{}^{}(x)`$ is given by the standard formula
$$\stackrel{~}{\theta }^{\alpha \beta }{}_{\mu }{}^{}(x)=2g^{\gamma [\alpha }\stackrel{~}{e}{}_{B}{}^{\beta ]}_{[\mu }^{}\stackrel{~}{e}{}_{\gamma ]}{}^{B}+g^{\alpha \gamma }g^{\delta \beta }\stackrel{~}{e}{}_{\mu }{}^{B}_{[\delta }^{}\stackrel{~}{e}_{\gamma ]B}$$
(16)
where the tetrad field $`\stackrel{~}{e}{}_{A}{}^{\mu }(x)`$ is given by (11).
Now, we can formulate the diffeomorphism covariant wave equation for quantum gravity:
$$i\mathrm{}\kappa \stackrel{~}{e\overline{)}}\mathrm{\Psi }+i\mathrm{}\kappa (\sqrt{g}\gamma ^\mu \theta _\mu )^{op}\mathrm{\Psi }=\widehat{}\mathrm{\Psi },$$
(17)
where $`\stackrel{~}{\overline{)}}=\stackrel{~}{e}_A^\mu (x)\gamma ^A(_\mu +\stackrel{~}{\theta }_\mu (x))`$ is the Dirac operator constructed using the self-consistent field $`\stackrel{~}{e}_A^\mu (x)`$ and
$$(\sqrt{g}\gamma ^\mu \theta _\mu )^{op}=4\pi iG\mathrm{}\kappa \left\{\sqrt{g}h^{\mu \nu }\frac{}{h^{\mu \nu }}\right\}_{ord}$$
(18)
is the term corresponding to the operator part of the spinor connection. The idea behind eq. (17) is that classical geometric structures needed to formulate the wave equation are introduced as approximate averaged notions in a self-consistent with the underlying quantum dynamics (determined by $`\widehat{}`$) way. As a consequence of condition (11) eq. (17) essentially becomes a non-linear integro-differential equation describing the non-trivial way in which the wave function $`\mathrm{\Psi }`$ specifies, or “lays down”, the space-time geometry it propagates on.
To complete the description, we also need to impose a gauge-type condition in order to distinguish the physically relevant information. For example, the De Donder-Fock harmonic gauge can be imposed in the form
$$_\mu \sqrt{g}g^{\mu \nu }(x)=0.$$
(19)
In the present context this is a gauge condition on the wave function $`\mathrm{\Psi }(g^{\mu \nu },x^\nu )`$ rather than on the metric field.
## 3 CONCLUSION
The De Donder-Weyl Hamiltonian formulation of the field equations leads to the procedure of quantization of fields, which we suggested to call precanonical, treating space and time variables on equal footing. When applied to general relativity in metric variables this framework leads to a Dirac-like wave equation (17), non-linear and integro-differential, with self-consistently incorporated classical geometric structures. No arbitrarily fixed background structure is present in the formulation; background independence is ensured by self-consistency which in its turn is dictated by the correspondence principle. The averaged self-consistent space-time serves its usual role: to order events (here, different possible configurations of the wave function in the metric space where its linear quantum dynamics is given by DW Hamiltonian operator $`\widehat{}`$) and to interpolate between them. This is the only way to describe physics we are sure about at the present: to describe it in space and time. Here we do have a quantum dynamics of the wave function in the metric space but we also do need classical space-time to order and join together the configurations of the wave function in the metric space (the fibre) in different points of the space-time (the base). Our wave equation prescribes how this ordering is achieved and how, as a consequence of this, the space-time is gaining its metric structure. This reference to classical space-time, even though self-consistent, may well be an approximation. “Quantum space-time”, usually thought to be an essential ingredient of quantum gravity, would imply a totally different way of describing physics. Our approach suggests that this could be achieved by attributing a proper sense to the “operator of the Dirac operator” in the left hand side of (4) without referring to classical geometric notions. A. Connes’ non-commutative geometry can be mentioned as an example of a mathematical framework achieving this goal.
Potential advantages of the present approach are (i) the manifest covariance of the foundations and the results and (ii) the luck of serious mathematical problems with underlying mathematical constructions (as opposite to, e.g., the Wheeler- De Witt geometrodynamics). It also offers a framework for discussing the problem of emergence of classical space-time in quantum gravity and has a potential to enlighten the problem of interpretation of quantum formalism in quantum cosmology: the quantum system described by (17) is “self-referential” in the sense that the classical self-consistent tetrad field can be viewed as a model of the observing degrees of freedom explicitly entering into the description of quantum dynamics. We hope that these intriguing features are a sufficient justification of the further analysis and development of the precanonical approach to quantum fields and quantum gravity, in spite of its so far unclarified connections to the standard quantum field theory.
Acknowledgments. I am grateful to the Organizers of the QG99 Conference for their kind invitation, warm hospitality, and the financial support which allowed me to enjoy these unforgettable days in Villasimius (Sardinia).
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# Effect of Thermal Undulations on the Bending Elasticity and Spontaneous Curvature of Fluid Membranes
## Introduction
The fluid membranes of red blood cells and giant vesicles in water undergo microscopically visible thermal shape fluctuations . The thermal undulations lead to orientational decorrelation on a scale characterized by de Gennes’ persistence length . About fifteen years ago, one of us proposed that the decrease of correlation should be accompanied by a softening of the membrane. Specifically, the effective bending rigidity $`\kappa ^{}`$ was predicted to obey
$$\kappa ^{}=\kappa \alpha \frac{kT}{8\pi }\text{log}M,$$
(1)
where $`\kappa `$ is the bare bending rigidity and $`M`$ the number of molecules making up the membrane. The numerical prefactor was found to be $`\alpha =1`$ in the first derivation and two later ones intended to support this result . In all three papers mean curvature was assumed to be the statistical measure. However, serious mistakes were made in the calculations. Peliti and Leibler, whose result was the first to be published , and all other authors rederived eq. (1) for the renormalized bending rigidity, but with the prefactor $`\alpha =3`$. They took normal displacement of the membrane from its equilibrium shape as measure of integration. Only recently the problem was reconsidered when the mistakes were noted in the calculations employing mean curvature. In addition to correcting the previous calculation for the sphere , a new calculation was done for the almost flat fluctuating membrane. It is based on a local bending mode expansion called hat model . In both cases it was found that fluid membranes should be stiffened, not softened, by their thermal undulations. The form of the relationship (1) is maintained, while the numerical prefactor changes sign, becoming $`\alpha =1`$. However, it has be mentioned that a few years ago Gompper and Kroll found a membrane softening in Monte Carlo simulations of fluctuating vesicles .
In the present work we deal once more with the problem of selecting the correct statistical measure, taking advantage of the analogies between polymers in two-dimensional space and fluid membranes. Subsequently, we consider a particularly transparent situation, cylindrical curvature, to rederive the effective bending rigidity in this case. Employing curvilinear coordinates, we then find the stiffening, i.e. the validity of eq. (1) with $`\alpha =1`$, for more general deformations of the initially flat membrane. In the same tensor notation, we confirm that the elastic modulus of Gaussian curvature is not affected by thermal undulations and show for the first time that this also holds for any spontaneous curvature of the membrane.
## Formulation of the problem
The bending energy of a fluid membrane may be expressed by the Hamiltonian
$`H={\displaystyle \text{d}^2\sigma \sqrt{g}\left(\frac{1}{2}\kappa J^2\kappa c_sJ+\overline{\kappa }K\right)},`$ (2)
where $`\kappa `$ and $`\overline{\kappa }`$ are the bending moduli and $`c_s`$ is the spontaneous curvature. The terms in parentheses represent the surface energy density and $`\sqrt{g}\text{d}^2\sigma `$ is the local surface element of the membrane. Let $`\stackrel{}{X}(\overline{\sigma })`$, $`\overline{\sigma }=(\sigma ^1,\sigma ^2)`$, be a parametrization of the membrane surface with the contravariant coordinates $`\sigma ^1`$ and $`\sigma ^2`$. Then $`\text{d}^2\sigma =\text{d}\sigma ^1\text{d}\sigma ^2`$ and $`g=\text{det}g_{ij}`$, where
$$g_{ij}=_i\stackrel{}{X}_j\stackrel{}{X},\text{with}_i\stackrel{}{X}=\stackrel{}{X}/\sigma ^i,i,j=1,2,$$
are the coefficients of the first fundamental form. $`J`$ and $`K`$ denote trace and determinant of the extrinsic curvature tensor $`J_j^i=g^{ik}J_{kj}`$, respectively, where
$$J_{ij}=\stackrel{}{N}_{ij}\stackrel{}{X},i,j=1,2$$
are the coefficients of the second fundamental form, $`\stackrel{}{N}`$ being (in the case of a closed surface) the outward unit normal vector to the surface. In other words,
$$J=c_1+c_2\text{and}K=c_1c_2,$$
are twice the mean curvature and the Gaussian curvature, respectively, with $`c_1`$ and $`c_2`$ being the principal curvatures. Note that with this convention $`J`$ is negative for spheres and cylinders.
The concept of an effective rigidity can be understood as follows: Starting from a basically flat, weakly fluctuating piece of membrane, we are interested in the free energy $`\mathrm{\Gamma }`$ needed to deform the piece into a fluctuating shape of nonvanishing background curvature $`J`$. It turns out that the effect of the short wavelength fluctuations upon the elastic behaviour on large scales can be absorbed into the material parameters of the bending Hamiltonian (2). Accordingly, the contribution to $`\mathrm{\Gamma }`$ from the first term in (2) may be written as
$$\mathrm{\Gamma }=\frac{1}{2}\kappa ^{}J^2\sqrt{g}\text{d}^2\sigma .$$
(3)
Here the curvature and the surface element refer to the nonfluctuating background and $`\kappa ^{}\kappa `$ is the effective bending rigidity.
As is common in shape calculations, the membrane is assumed to be unstretchable. The free energy $`\mathrm{\Gamma }`$ refers to a background surface of fixed area. The associated real membrane area is slightly larger because of the bending fluctuations. Any shrinkage of base area resulting from the fluctuations can be formally avoided by introducing a (fictitious) reservoir of basically flat membrane area. Alternatively, $`\mathrm{\Gamma }`$ can be taken to refer to the shrunken base area. Such a shrinkage does not affect the bending energy of the base if it is done by an isotropic scale transformation, as happens automatically in the case of a freely fluctuating spherical vesicle. Fortunately, these subtleties seem to be irrelevant in practice.
The free energy of bending, $`\mathrm{\Gamma }`$, can be obtained in two different, but equivalent ways. One may ”freeze” the undulations in an instantaneous configuration and then calculate the bending energy of the rippled membrane at fixed undulation amplitudes. The squared fluctuation mode amplitudes are averaged only afterwards.
Alternatively, one may start from the free energy of the fluctuating flat membrane and investigate how it changes when the membrane is subjected to a background curvature $`J0`$. The fluctuation dependent part of the bending energy $`\mathrm{\Gamma }`$ is now entropic, resulting from the effect of $`J`$ on the mean square amplitudes of the undulation modes.
In both cases we are faced with the problem of choosing the right measure of integration. The undulations as expressed by the measure have to be kept constant when the membrane with ”frozen” undulations is deformed in the large. Correspondingly, the entropic part of the free energy of background deformations has to be calculated from the change of the mode mean square amplitudes in terms of the measure. We will expand now on some of the arguments why mean curvature, or any quantity proportional to it, is a correct measure of integration. The scale of the measure is of no interest because we are dealing with differences, not absolute values of entropy.
## Considerations on the measure of integration
A heuristic argument for choosing mean curvature as statistical measure is the fact that it can be regarded as the physical strain of the fluid membrane, much like curvature in the case of a (stiff) polymer. Accordingly, the superposition of two deformations, e.g. uniform background curvature and a (sinusoidal) ripple, begins with adding the strains as they are in the absence of the other deformation. The simultaneous presence of both deformations will produce coupling terms in the effective total bending energy that are quadratic in both the background curvature and the ripple amplitude. Employing the strain as statistical measure has the particular attraction of avoiding spurious coupling terms. They arise from a change of the ripple amplitude in terms of the strain when other measures are used such as normal displacement from the background or reorientation of the layer normal.
From a technical point of view, mean curvature can serve as measure of integration only if its local values determine the shape of the polymer or membrane, respectively. A one–to–one relationship between shape and local curvatures (which may be molecular) clearly exists for the polymer in the plane, a simplified model system frequently considered in the following. For membranes, a similar invertible mapping is particularly transparent in the hat model which expands the fluctuations of the infinite planar membrane into local mean curvature modes. There are only minor restrictions to the local mean curvature fluctuations in other geometries, such as periodic boundary conditions, cylinders or spheres (see below). However, any (small) membrane deformation compatible with the boundary condition can be completely specified either by local normal displacements or mean curvatures. In all these cases Gaussian curvature fluctuates together with mean curvature, but the former is completely slaved by the latter because of the boundary conditions. The general possibility of invertible mappings between shape and mean curvature, including local area conservation and lateral flow, has been discussed elsewhere .
Finally, there is a cogent reason to select (mean) curvature as statistical measure for polymers in $`2`$D or fluid membranes. It stems from the fact that in both systems every local strain fluctuates independently of all the others (apart from the minor restrictions just mentioned). As a result, the total fluctuation entropy is the sum of the equal entropies of all the molecular (or other local) fluctuations. Independence and equality imply that the entropy $`S_q`$ of a sinusoidal (mean) curvature mode on a flat background does not depend on its wave vector $`q`$. This is because for all $`q0`$ the molecules are in the same way homogeneously distributed over the phases of the sine wave. The $`q`$–independence of $`S_q`$ indicates that mean curvature is a suitable measure of integration since the mean square curvature amplitudes also do not depend on $`q`$. Let us add as a remark that the absence of a wave vector dependent Jacobian can perhaps be regarded as characteristic of a correct measure.
The condition of the mode entropy being independent of $`q`$ is violated if the entropy is calculated in terms of other statistical measures. Taking normal displacement, one generally derives
$$S_q=\frac{1}{2}k\mathrm{log}q^4+const.,$$
from the variation with $`1/q^4`$ of the mean square displacement amplitude of a sinusoidal fluctuation mode. The $`q`$–dependent part of $`S_q`$ is half as strong if normal displacement is replaced by the tilt angle of of the polymer or membrane with respect to the background. Evidently, neither normal displacement nor tilt angle is a correct measure of integration. It should be noted that in many statistical–mechanical calculations the choice of the measure is irrelevant. However, it matters whenever the coupling of two bending deformations is considered as will be demonstrated below.
After these general statements let us address particular model systems to deal with the statistical mechanics of both polymers in a plane and fluid membranes. They will confirm the claim that mean curvature is a correct measure of integration for fluid membranes. Examples of the invertible mapping between membrane shape and local mean curvature will also be given, including the restrictions on the latter.
We first have a look at the comparatively simple problem of a polymer chain fluctuating in a plane . The configurational space of the polymer is swept completely by varying the angles between adjacent links. Disregarding fluctuations of link length, one may speak of local curvatures instead of angles. Each joint between subsequent links has its own free energy arising from its curvature fluctuations. Although the change of a particular angle rotates the two sides of the polymer chain with respect to each other, it does not affect their free energies. We assume here that the polymer is almost straight. In a continuum description such a polymer is equivalent to a (thin) rod or wire.
Using local curvature as the measure of integration in dealing with fluctuating polymer chains is not in conflict with the fundamental measure which for mass points is translation. Local curvature is a shorthand notation for relative displacement if, e.g., each joint of the polymer chain is regarded as a mass point. Then a (small) normal displacement of point $`n+1`$ from the tangent to the polymer defined by the two preceding points is proportional to the curvature at point $`n`$. In contrast, the absolute normal displacement of a nearly straight fluctuating polymer from a base line is not a suitable statistical measure. It implies a freedom of motion for each individual mass point which actually is suppressed by cohesion and bending stiffness. The remaining freedom is fully accounted for by the curvature.
The polymer allows a simple demonstration of the effect of a wrong statistical measure on the curvature amplitude of a ripple. We consider a sinusoidal wave on a polymer with a uniform background curvature $`1/r`$, the negative sign indicating a downward bend. A sketch of the situation is given in Fig. 1. Let us keep all normal displacements $`\nu (s)`$ fixed, $`s`$ being the arc length of the background, while we bend the polymer as a whole. This amounts to using the displacement measure. The procedure results in a weakening of the ripple curvature because the total curvature changes from $`\text{d}^2\nu /\text{d}s^2`$ to $`1/r+\text{d}^2\nu /\text{d}s^2+\nu /r^2`$ up to linear order in $`\nu `$, with the two terms constituting the ripple curvature being of opposite sign. Clearly, the change of the ripple curvature amplitude is physically not acceptable so that normal displacement does not appear to be the correct statistical measure.
Continuing the two–dimensional polymer into the third space dimension results in a bent membrane with zero curvature in the new direction. The similarity of such a membrane with the polymer in a plane strongly suggests that some kind of curvature rather than displacement from a fixed base should be the right measure of integration for fluid membranes.
In order to extend the preceding argument to membranes that can be bent not only in a single direction let us consider a fluid membrane made of mass points. A particularly simple situation arises for the almost flat, weakly fluctuating membrane: Like the polymer, the idealized fluid membrane can be build up by adding the constituent mass points one by one. For easy statistics, we may think of a (local) lattice of mass points, quadratic or hexagonal, which is formed row after row. A new molecule continuing a straight sequence completes a plaquette comprising a certain number of preexisting neighbours. If the number is large enough the plaquette can be used to define the mean curvature at its center. Five molecules is the minimum; they may be thought to form a centered square, as is illustrated in Fig. 2. The height of the last molecule added, relative to a suitably chosen tangent plane, finalizes the value of the curvature $`J`$. On a quadratic lattice, $`J`$ may be taken (in lowest order) to be the sum of the (normal) curvatures of the two diagonals of the plaquette. Let us emphasize, that it is not the relative height of the final molecule that gives a parametrization of the phase space of the particle in the middle of the plaquette, but the simultaneously varying mean curvature. This quantity is invariant with respect to the direction along which rows are formed one by one. The tangent plane may be replaced by the base plane, shifted to the local height, if the fluctuations are weak enough to leave the membrane almost flat. Obviously, the last simplification is possible in the example of Fig. 2.
Regardless of the measure of integration, it is convenient to describe the thermal undulations superimposed on a background shape of given $`J`$ and $`K`$ in terms of normal displacement $`\nu (\overline{\sigma })`$. To linear order, $`\nu `$ and the associated extra curvature $`\xi =J^{}J`$ are related through
$$\xi =[J^22K+D^2]\nu ,$$
(4)
$`D^2`$ being the Laplace–Beltrami operator (see eq. (7) below and subsequent definitions). The difference $`\xi `$ is equivalent to $`J`$ as statistical measure.
Undulation modes can be defined as eigenfunctions of the operator represented by the bracketed expression. For surfaces of uniform background curvature, such as planes, cylinders and spheres, this ensures that for each mode the functions $`\xi (\overline{\sigma })`$ and $`\nu (\overline{\sigma })`$ are proportional to each other in a first order approximation, the factor of proportionality depending on the wave number. The transformation (4) can be inverted whenever for a given undulation mode the operator in brackets does not vanish or affect the membrane area linearly in the mode amplitude. Well–known examples of such exceptions are the three $`l=1`$ modes of the fluctuating sphere which actually represent translations and the $`l=0`$ mode which requires area changes of the unstretchable membrane. With periodic boundary conditions the only lost mode is uniform mean curvature ($`q=0`$) as it is incompatible with the boundary conditions. In the case of the cylinder two long–wavelength curvature modes are lost because they represent translations normal to the cylinder axis. Periodic boundary conditions along the axis and area conservation suppress two other mean curvature modes. The loss of modes implies that the local mean curvatures are not totally free to fluctuate independently of each other. However, these restrictions seem irrelevant in the present context for membranes comprising large numbers of molecules.
## Basic formulas
In order to calculate the bending Hamiltonian (2) for a fluctuating membrane, we first present expressions for the mean curvatures and the surface elements of membrane shapes lying close to a given background shape. The latter is also referred to as the projected shape.
If $`\stackrel{}{X}(\overline{\sigma }),\overline{\sigma }=(\sigma ^1,\sigma ^2)`$, is a parametrization of the background configuration, any slightly deviating configuration $`\stackrel{}{X}^{}(\overline{\sigma })`$ may be represented in normal gauge by
$$\stackrel{}{X}^{}(\overline{\sigma })=\stackrel{}{X}(\overline{\sigma })+\nu (\overline{\sigma })\stackrel{}{N}(\overline{\sigma }),$$
(5)
$`\stackrel{}{N}(\overline{\sigma })`$ being the local outward unit normal vector to $`\stackrel{}{X}(\overline{\sigma })`$. From (5) we find that up to quadratic order in $`\nu `$ the surface element of the deformed membrane obeys (cf.)
$$\sqrt{g^{}}=\sqrt{g}(1\nu J+\frac{1}{2}(D\nu )^2+K\nu ^2).$$
(6)
In the same approximation the curvature $`J^{}`$ becomes
$`J^{}=J+(D^2+J^22K)\nu {\displaystyle \frac{1}{2}}J(D\nu )^2`$ (7)
$`+J^{ij}D_i\nu D_j\nu +(J^33JK)\nu ^2+2J^{ij}\nu D_iD_j\nu +`$
$`+\nu D_n\nu g^{ij}(D_jJ_i^n+J\mathrm{\Gamma }_{ij}^n),`$
Here, $`g_{ij}`$ and $`J_{ij}`$ are the coefficients of the first and second fundamental forms of the background surface, respectively. $`\mathrm{\Gamma }_{ij}^l`$ are the Gauss–Christoffel symbols of the second kind derived from $`g_{ij}`$. $`D_i`$ is the covariant derivative being defined as $`D_i\nu =_i\nu =\nu /\sigma ^i`$, and $`D_iD_j\nu =_{ij}\nu \mathrm{\Gamma }_{ij}^k_k\nu `$. $`D^2=g^{ij}D_iD_j`$ is the Laplace–Beltrami operator and we use the abbreviation $`(D\nu )^2=D^i\nu D_i\nu `$. The variation of the Gaussian curvature $`K`$ need not to be considered in (2) thanks to the Gauss–Bonnet theorem .
## Renormalization of $`\kappa `$ for the cylinder via frozen mode amplitudes
Let us now derive the effective bending rigidity for cylindrical background curvature in the ”frozen mode” approach. We first write down reduced versions of (6) and (7) that apply to this special case. Using geodesic coordinates $`x=r\phi `$ and $`z`$ to parametrize the cylinder, where $`\phi `$ and $`z`$ are usual polar coordinates, we find with $`J=1/r`$:
$$\sqrt{g^{}}=1+\frac{1}{r}\nu +\frac{1}{2}(\nu )^2+O(\nu ^3)$$
(8)
and
$`J^{}={\displaystyle \frac{1}{r}}+{\displaystyle \frac{1}{r^2}}\nu +\mathrm{\Delta }\nu +{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{r}}(\nu )^2`$ (9)
$`{\displaystyle \frac{1}{r}}(_x\nu )^22{\displaystyle \frac{1}{r}}\nu _{xx}\nu {\displaystyle \frac{1}{r^3}}\nu ^2+O(\nu ^3).`$
The linear part of $`J^{}`$ in (9) leads to the transformation (4) and its inversion for the cylinder,
$$(1/r^2+\mathrm{\Delta })\nu =\xi \nu =\mathrm{\Delta }^1(1+1/r^2\mathrm{\Delta }^1)^1\xi ,$$
(10)
where $`\mathrm{\Delta }`$ is the usual Laplacian for a flat base. Obviously, the thermal undulations on a cylinder can be expanded, like those on a flat background, into sinusoidal waves which at the same time are curvature and normal displacement modes.
We will consider a quadratic piece of membrane subjected to a weak cylindrical background bend and satisfying periodic boundary conditions. Alternatively, the membrane may be closed to form a cylinder so that the edge of the square equals the circumference in length. In the displacement representation the sine and cosine deformation modes read
$$\nu =a_\stackrel{}{q}\mathrm{sin}(\stackrel{}{q}\stackrel{}{r})\text{and}\nu =b_\stackrel{}{q}\mathrm{cos}(\stackrel{}{q}\stackrel{}{r}),$$
(11)
where $`\stackrel{}{r}=(x,z)`$ and $`\stackrel{}{q}=(q_x,q_z)`$ with $`q_x=2\pi n/L`$ and $`q_z=2\pi m/L`$. The numbers $`n`$ and $`m`$ are running from $`L/(2a)`$ to $`L/(2a)`$, but only over half the $`\stackrel{}{q}`$ plane. The constant mode $`(n,m)=(0,0)`$ is excluded. $`L`$ denotes the size of the background shape and $`a`$ is the molecular distance in the membrane. The operator acting on $`\nu `$ in (10) is $`(1/r^2\stackrel{}{q}^{2})`$ for a given mode.
In order to renormalize the bending rigidity let us first consider the effect of a single sinusoidal mode. A straightforward calculation starting from (8) and (9) and involving a partial integration then results in the Hamiltonian
$`{\displaystyle \frac{1}{2}\kappa J^2\sqrt{g^{}}\text{d}x\text{d}z}=`$ (12)
$`{\displaystyle \frac{1}{2}}\kappa \left[{\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{1}{2}}(1/r^2\stackrel{}{q}^{2})^2a_\stackrel{}{q}^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{r^2}}({\displaystyle \frac{1}{2}}q_x^2+{\displaystyle \frac{3}{2}}q_z^2{\displaystyle \frac{1}{r^2}})a_\stackrel{}{q}^2\right]A_{base}`$
where $`A_{base}=L^2`$ denotes the area of the nonfluctuating background shape. Terms linear in $`\nu `$ do not occur as they vanish on integration. Of the three terms in the brackets, the second one equals the mean value of $`\xi ^2`$, thus giving the regular bending energy of the mode. Note that if the curvature measure applies, the curvature difference $`\xi `$ is kept constant in the ”frozen mode” approach when the membrane is deformed. Accordingly, the third term is due to the coupling of the mode curvatures to cylindrical curvature.
The statistical averages of the squared mode amplitudes are obtained from the equipartition theorem,
$$a_\stackrel{}{q}^2=b_\stackrel{}{q}^2=\frac{2kT}{A_{base}\kappa (1/r^2\stackrel{}{q}^{2})^2}.$$
(13)
Inserting (13) in (12) and absorbing the third term of (12) in the first leads to
$$\frac{1}{2}\kappa J^2\sqrt{g^{}}\text{d}x\text{d}z=\frac{1}{2}\kappa \left(1+\frac{\frac{1}{2}q_x^2+\frac{3}{2}q_z^21/r^2}{(1/r^2\stackrel{}{q}^{2})^2}\frac{kT}{A_{base}\kappa }\right)\frac{1}{r^2}A_{base}+\frac{1}{2}kT.$$
(14)
Adding up the contributions of all fluctuation modes to the renormalization of $`\kappa `$ yields to lowest order in $`kT/\kappa `$
$$\kappa ^{}=\kappa \left(1+\underset{\stackrel{}{q}}{}\frac{q_x^2+3q_z^2}{(\stackrel{}{q}^{2})^2}\frac{kT}{A_{base}\kappa }\right),$$
(15)
where $`1/r^2`$ is omitted because of $`q>1/r`$. The sum is assumed to be small as compared to unity. It does not contain $`(n,m)=(0,0)`$ and some seemingly divergent low wave vector modes. The upper limits of $`n`$ and $`m`$ are large numbers. Replacing the sum by the double integral
$$\underset{\stackrel{}{q}}{}\frac{A_{base}}{8\pi ^2}_{2\pi /L}^{\pi /a}\text{d}^2\stackrel{}{q},$$
we arrive with good accuracy at
$$\kappa ^{}=\kappa +\frac{kT}{8\pi }\mathrm{log}M,$$
(16)
$`M`$ being the number of molecules in the membrane (or half this number for the bilayer). Evidently, we have recovered the result derived previously for the fluctuating spherical vesicle, i.e. eq. (1) with $`\alpha =1`$.
The effect of thermal undulations on the bending rigidity of fluid membranes was attributed to a coupling between the mode curvatures and the background curvature. To confirm the role of coupling, we may for a moment compensate the uniform background curvature by a spontaneous curvature. For the cylinder this means $`c_s=1/r`$ and $`J^{}c_s=\xi `$. Accordingly, the surface energy density is simply $`\frac{1}{2}\kappa \xi ^2=\frac{1}{2}\kappa ((\frac{1}{r^2}+\mathrm{\Delta })\nu )^2`$ up to quadratic order in $`\xi `$. The additional quadratic terms arising in the case of interest, $`c_s=0`$, can therefore be regarded as coupling terms that give rise to a renormalization of the bending rigidity.
For a direct understanding of the coupling let us consider two special types of ripples, those exactly parallel or orthogonal to the cylindrical background curvature.
Parallel ripple: When the ripple is parallel we can use the analogy to a polymer in a plane. This makes it attractive to use a more direct approach than the general calculation based on $`\nu (x,y)`$. Starting from a straight and stiff polymer or a wire, we impose on it a periodic deformation
$$\xi (s)=c_q\mathrm{sin}(qs),$$
where $`\xi `$ is curvature, $`s`$ arc length and $`q`$ wave vector. The undulation diminishes the projected or background length of the wire. Alternatively, it increases the total arc length $`L^{}`$ of the wire (when connected to a reservoir) at fixed projected length $`L`$. In both cases the ratio of the two lengths is
$$L^{}/L=1+\frac{1}{4}\frac{c_q^2}{q^2}$$
(17)
to lowest order in $`\xi `$.
A frozen bending mode of the polymer is similar to a wire with a permanent wave. Any sinusoidal bending deformation allows a given angle difference between the ends of the wire to be spread over $`L^{}`$ instead of $`L`$. This reduces the associated bending energy density by $`(L/L^{})^2`$ and the total bending energy by $`L/L^{}`$ as compared to a polymer or wire without waviness. The analogous ratio for parallel ripples in membranes is $`A/A^{}`$, where $`A^{}`$ and $`A`$ are the real and projected areas, respectively. Absorbing the effect into a renormalized bending rigidity, we have $`\kappa ^{}=\kappa A/A^{}`$. This agrees exactly with what we obtain from (14) for a single mode of type $`(q_x,0)`$ if we make use of $`A^{}=A(1+\frac{1}{4}q_x^2a_{q_x,0}^2)`$. The negative sign in front of the correction term in $`\kappa ^{}`$ indicates the softening in this case.
An amusing illustration of the effect of parallel undulations is a rather stiff polymer closed to form a ring. Its thermal undulations in terms of curvature are not affected by the uniform bend and vice versa. Only the decrease of the ring radius in the presence of undulations signals a lowering of the effective bending stiffness of the polymer.
Orthogonal ripple: Ripples orthogonal to the cylindrical base curvature stiffen the membrane, in contrast to parallel ones. This is due to the fact that at the crests of such a ripple the extra curvature $`\xi `$ is of the same sign as the cylindrical curvature, and at the same time the membrane area is increased as compared to the background shape. In linear approximations the two effects obey
$$J^2=\frac{1}{r^2}2\frac{1}{r}\left(\frac{1}{r^2}+_{zz}\right)\nu +O(\nu ^2)$$
and
$$\sqrt{g^{}}=1+\frac{\nu }{r}+O(\nu ^2).$$
Calculating from these two expressions the total bending energy of a membrane with a single excited mode of wave vector $`(0,q_z)`$, we find the renormalized bending rigidity of the cylinder to be
$$\kappa ^{}=\kappa (1(q_z^2)a_{0,q_z}^2).$$
(18)
According to this formula the orthogonal ripple stiffens the membrane. It is four times as effective as a softening, parallel ripple of equal wavelength.
There is a second correction quadratic in $`\nu `$ to the bending energy due to the coupling between background curvature and orthogonal ripple. It originates from the sloped regions of the ripple where the membrane partially escapes cylindrical curvature. The reduced principal curvature is $`(\mathrm{cos}\phi )/r=(1/r)(1\phi ^2/2+O(\phi ^4))`$, $`\phi `$ being the slope angle. For a single mode of type $`(0,q_z)`$ the second effect and the increase of membrane area due to the slope result in a softening given by
$$\kappa ^{}=\kappa (1\frac{1}{4}q_z^2a_{0,q_z}^2).$$
(19)
Combining the opposite rigidity corrections of (18) and (19) results in a stiffening which agrees with (14).
At this point let us briefly consider what would be different in a derivation of $`\kappa ^{}`$ based on the displacement measure. Keeping $`\nu `$ fixed instead of $`\xi `$ while the background is cylindrically bent implies that the regular part of the bending energy density is $`1/2\kappa (\mathrm{\Delta }\nu )^2`$ instead of $`1/2\kappa (\mathrm{\Delta }\nu \nu /r^2)^2`$. Accordingly, the single mode Hamiltonian (14) would have to be recast in the form
$`{\displaystyle \frac{1}{2}\kappa J^2\sqrt{g^{}}\text{d}x\text{d}z}=`$
$`{\displaystyle \frac{1}{2}\kappa \left(\frac{1}{r^2}+(\stackrel{}{q}^{2})^2\nu ^2+\frac{1}{r^2}(\frac{5}{2}q_x^2\frac{1}{2}q_z^2+O(r^2))\nu ^2\right)\text{d}x\text{d}z},`$
if again the third term is to represent the coupling of the undulations to cylindrical background curvature. The negative signs before $`q_x^2`$ and $`q_z^2`$ in the coupling term suggest that both parallel and perpendicular ripples soften the membrane. Absorbing the third term into the first renormalizes $`\kappa `$. The same manipulations as above, including thermal averaging and summation over all modes, would lead to
$$\kappa ^{}=\kappa 3\frac{kT}{8\pi }\mathrm{log}M.$$
This is the result obtained by all other authors. They generally used the displacement measure and took the approach via undulation mode entropies.
The undulations of unstretchable fluid membranes must be accompanied by lateral displacements. They can be analyzed mode by mode and consist of local and global parts. The local part redistributes membrane material between the crests and troughs of a ripple in a bent surface and in general varies linearly with the mode amplitude. The global part takes account of the area absorbed by the undulation and goes with the square of the amplitude. We have disregarded lateral motion, with one exception: In the case of the wavy polymer it is automatically included by employing arc length as independent variable. Inspection shows that lateral displacement need not be taken into account in calculations of the effective bending rigidity of polymers and membranes in the usual approximation that is quadratic in the mode amplitudes .
## Renormalization of $`\kappa `$, $`c_s`$ and $`\overline{\kappa }`$ via undulation mode entropies
The effective bending rigidity has been derived above for uniform cylindrical curvature and earlier for the spherical membrane, both times in the ”frozen mode” approach. The same stiffening was obtained in both calculations. Let us try now to extend the result to more general geometries, showing simultaneously that the undulations do not affect the spontaneous curvature $`c_s`$ and the modulus of Gaussian curvature, $`\overline{\kappa }`$, of the membrane. Obviously, we have to limit ourselves to ”compact” shapes such as squares, circles and spheres. For the following calculations it is advantageous to follow the approach based on fluctuation mode entropies. The curvature increment $`\xi =J^{}J`$ will again be used as measure of integration.
In order to calculate the free energy of the fluctuating membrane and to renormalize its elastic parameters, we employ a field–theoretic approach as has been introduced by Förster in this frame. Considering some background configuration with parametrization $`\stackrel{}{X}(\overline{\sigma })`$, which describes the effective shape, we expand the bending Hamiltonian in (2) to second order around $`\stackrel{}{X}(\overline{\sigma })`$. We drop linear terms, assuming the background to be an equilibrium shape. The free energy $`F`$ of $`\stackrel{}{X}`$ may then be written as
$$F(\stackrel{}{X})=H(\stackrel{}{X})+\frac{1}{2}kT\mathrm{log}det\frac{\delta ^2H(\xi )}{\delta \xi (\sigma )\delta \xi (\sigma ^{})}+O((kT)^2),$$
(20)
where $`\xi `$ parametrizes the configurations $`\stackrel{}{X}^{}`$ in the vicinity of $`\stackrel{}{X}`$. After expanding the r.h.s. in powers of background curvature one may introduce effective elastic parameters marked by a prime, such as $`\kappa ^{}`$, and an effective Hamiltonian $`H^{}`$ employing the new parameters in conjunction with $`J`$ and $`K`$ of the background shape $`\stackrel{}{X}`$.
To perform this program we first have to calculate the bending energies of the fluctuating membrane. Starting from eq. (7) we obtain
$`J^2=J^2+2J(D^2+J^22K)\nu +((D^2+J^22K)\nu )^2`$ (21)
$`J^2(D\nu )^2+2JJ^{ij}D_i\nu D_j\nu +2J(J^33JK)\nu ^2`$
$`+4JJ^{ij}\nu D_iD_j\nu +2J\nu D_n\nu g^{ij}(D_jJ_i^n+J\mathrm{\Gamma }_{ij}^n)+O(\nu ^3).`$
Combining this with (6) yields
$`J^2\sqrt{g^{}}=[J^2J^3\nu +2J(D^2+J^22K)\nu {\displaystyle \frac{1}{2}}J^2(D\nu )^2`$ (22)
$`2J^2\nu (D^2+J^22K)\nu `$
$`+4JJ^{ij}\nu D_iD_j\nu +2JJ^{ij}D_i\nu D_j\nu +((D^2+J^22K)\nu )^2`$
$`+JK\nu ^2+2J(J^33JK)\nu ^2+2J\nu D_n\nu g^{ij}(D_jJ_i^n+J\mathrm{\Gamma }_{ij}^n)+O(\nu ^3)]\sqrt{g}.`$
We also need
$`J^{}\sqrt{g^{}}=[J+(D^2+J^22K)\nu J^2\nu `$ (23)
$`J\nu (D^2+J^22K)\nu +2J^{ij}\nu D_iD_j\nu +J^{ij}D_i\nu D_j\nu `$
$`+JK\nu ^2+(J^33JK)\nu ^2+\nu D_n\nu g^{ij}(D_jJ_i^n+J\mathrm{\Gamma }_{ij}^n)+O(\nu ^3)]\sqrt{g}.`$
In the following the background curvatures $`J`$ and $`K`$ are assumed to be weak and slowly varying. For all wave vectors except for the smallest ones the predominant terms quadratic in $`\nu `$ are those that contain at least two derivatives acting on $`\nu `$. In calculating the effect of the short wavelength fluctuations upon the long distance elastic behaviour of membranes comprising a large number of molecules, it is sufficient to take into account only the latter. Giving rise to cutoff dependent effects, they are the only ones that can contribute to the renormalization of the elastic parameters. To replace the normal displacement $`\nu `$ by $`\xi `$ we need the inversion of (4)
$$\nu =D^2(1+(J^22K)D^2)^1\xi ,$$
(24)
In the spirit of the above simplification, (24) can be approximated for terms which are quadratic in $`\nu `$ at the stage of $`J^{}`$ and $`g^{1/2}`$ (see eqs. (6) and (7)) by
$$\nu =D^2\xi .$$
Disregarding irrelevant terms, we may thus write the bending energy as
$`{\displaystyle }\text{d}^2\sigma \sqrt{g^{}}({\displaystyle \frac{\kappa }{2}}J^2+c_s\kappa J^{})={\displaystyle }\text{d}^2\sigma \sqrt{g}[{\displaystyle \frac{\kappa }{2}}(J^2`$ (25)
$`+2J\xi J^3D^2(1+(J^22K)D^2)^1\xi `$
$`+\xi ({\displaystyle \frac{1}{2}}J^2D^iD_iD^4+2JJ^{ij}D_iD_jD^42J^2D^2+I)\xi `$
$`+\kappa c_s(J+\xi +\xi (JD^2+J^{ij}D_iD_j/D^4)\xi )+\text{irrel. terms}],`$
where we have performed some partial integrations, assuming either periodic boundary conditions or undulations limited to a patch of membrane.
We are now in a position to set up the functional integral sweeping curvature space that gives the effective free energy $`F`$ of the background shape $`\stackrel{}{X}`$ with fields $`J`$ and $`K`$. In mechanical equilibrium linear terms in (25) drop out so that we may write
$`F=kT\mathrm{log}{\displaystyle }𝒟[\xi ]\mathrm{exp}[{\displaystyle \frac{1}{kT}}{\displaystyle }\text{d}^2\sigma \sqrt{g}({\displaystyle \frac{\kappa }{2}}J^2\kappa c_sJ`$ (26)
$`+{\displaystyle \frac{\kappa }{2}}(\xi (1+{\displaystyle \frac{J^2}{2}}D^2+{\displaystyle \frac{2JJ^{ij}D_iD_j}{D^4}}`$
$`{\displaystyle \frac{2J^2}{D^2}}+2c_s{\displaystyle \frac{J^{ij}D_iD_jJD^2}{D^4}}+\text{irrel. terms})\xi ))].`$
Extracting the background energy and integrating the multiple Gaussian we are led to
$`F={\displaystyle \text{d}^2\sigma \sqrt{g}\left(\frac{\kappa }{2}J^2\kappa c_sJ\right)}`$ (27)
$`+{\displaystyle \frac{kT}{2}}\text{Tr}^{}\mathrm{log}\left[{\displaystyle \frac{\kappa }{2kT}}\left(I+{\displaystyle \frac{J^2}{2D^2}}+{\displaystyle \frac{2JJ^{ij}D_iD_j}{D^4}}{\displaystyle \frac{2J^2}{D^2}}+2c_s{\displaystyle \frac{J^{ij}D_iD_jJD^2}{D^4}}\right)\right],`$
where $`I`$ is the identical operator. The prime on Tr indicates exclusion of some modes at the low wave number cutoff. We assume rather weak background bends, such that
$$J^2,K<1/L^2,$$
(28)
and may thus expand the logarithm in powers of them. A further limitation of $`J^2`$ may arise from $`c_s`$. The traces occuring in (27) are known to be :
$`\text{Tr}^{}{\displaystyle \frac{J^2}{D^2}}`$ $``$ $`{\displaystyle \text{d}^2\sigma \sqrt{g}J^2\frac{2}{4\pi }\mathrm{log}\frac{L}{a}}\text{and}`$ (29)
$`\text{Tr}^{}{\displaystyle \frac{2JJ^{ij}D_iD_j}{D^4}}`$ $``$ $`{\displaystyle \text{d}^2\sigma \sqrt{g}J^2\frac{2}{4\pi }\mathrm{log}\frac{L}{a}},`$ (30)
$`L`$ and $`a`$ being upper and lower cutoffs due to finite membrane size and molecular distance, respectively. They are easily calculated for uniformly bent shapes such as squares, spheres and cylinders. One obtains
$`F={\displaystyle \text{d}^2\sigma \sqrt{g}\frac{1}{2}\kappa J^2\left(1+\frac{kT}{4\pi \kappa }\mathrm{log}\frac{L}{a}\right)}+{\displaystyle \frac{kT}{2}}\text{Tr}^{}\mathrm{log}\left({\displaystyle \frac{\kappa }{2kT}}I\right)`$ (31)
$`{\displaystyle \text{d}^2\sigma \sqrt{g}c_s\kappa J\left(1+\frac{kT}{4\pi \kappa }\mathrm{log}\frac{L}{a}\right)}.`$
The middle term of (31) is equal to the free energy of the corresponding flat reference membrane. After subtracting it we are left with the free energy of bending which may be written in the renormalized form
$$\mathrm{\Gamma }=\text{d}^2\sigma \sqrt{g}(\frac{1}{2}\kappa ^{}J^2\kappa ^{}c_s^{}J),$$
(32)
where
$`\kappa ^{}`$ $`=`$ $`\kappa +{\displaystyle \frac{kT}{4\pi }}\mathrm{log}{\displaystyle \frac{L}{a}}`$ (33)
$`c_s^{}`$ $`=`$ $`c_s.`$ (34)
For comparison with (1), note that $`L^2/a^2=M`$. Evidently, the effective bending rigidity comes out exactly as calculated for the sphere and the cylinder, while the spontaneous curvature is not modified by the thermal undulations.
Formula (32) omits the Gaussian curvature term of the bending energy. We dropped it in the unrenormalized bending energy since, as a consequence of the Gauss–Bonnet theorem, the variation of Gaussian curvature does not affect its integral if the membrane is closed or satisfies periodic boundary conditions. Nevertheless, an additional Gaussian curvature term could have emerged in the calculation of the effective bending rigidity, as it does when normal displacement is the measure of integration . From the fact that such a term is absent in eq. (31) we may infer that $`\overline{\kappa }`$ is not affected by thermal undulations, i.e.
$$\overline{\kappa }^{}=\overline{\kappa },$$
(35)
in agreement with previous considerations .
In the preceding calculation of the effective bending rigidity we did not specify the shape of the membrane, except for stipulating that it be ”compact”. Although the formula for $`\kappa ^{}`$ seems to be rather robust, because of the logarithmic dependence on wave vector cutoffs, it is not expected to hold for filamentous membranes. Specifically, the bending rigidity of the membrane closed to form a cylinder will not stiffen indefinitely as the cylinder is made longer. This is why we chose periodic boundary conditions with an axial period as long as the circumference in our calculation for the cylinder.
The renormalized elastic parameters calculated above are independent of the base curvatures $`J`$ and $`K`$. It is an open question to which extent they remain valid if these curvatures are nonuniform over the base area. In the case of a wavy background one commonly uses as upper limit in (31) its wavelength or a fraction thereof instead of the membrane size $`L`$. This is because ripples of larger wavelength are part of the immutable background curvature.
## Conclusion
We have renormalized the bending elastic parameters of fluctuating fluid membranes for cylindrical and more general membrane shapes. The same logarithmic increase of the bending rigidity with membrane size was obtained as previously for spherical vesicles. In addition, the modulus of Gaussian curvature and, for the first time, spontaneous curvature were shown not to be affected to lowest order by thermal undulations. Apart from the illustrative treatment of the cylindrical membrane, we used arbitrary curvilinear coordinates in the calculations to obtain generally valid results as far as possible.
The predicted stiffening of fluid membranes by their thermal undulations may seem contrary to intuition, although there is no conflict with de Gennes’ derivation of the persistence length of membrane orientation. In the case of polymers the loss of orientational correlation unavoidably results in a breakdown of bending elasticity. This could be different in the case of membranes because of their two-dimensional connectedness and, in particular, the well-known scale invariance of bending energies. Moreover, it remains to be checked if the stiffening survives at very low bending rigidities $`(\kappa kT)`$, where the nearly flat approximation is inapplicable at any scale.
## Acknowledgement
We are grateful for their interest to H. Kleinert, A.M.J. Schakel and M.E.S. Borelli. W.H. thanks D. Nelson for comments and H.A.P. is indebted to G. Foltin for enlightening discussions.
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# Equilibrium sequences of irrotational binary polytropic stars : The case of double polytropic stars
## I Introduction
Coalescing binary neutron stars (BNSs) are considered to be one of the most promising sources of gravitational waves for laser interferometers such as TAMA300, GEO600, VIRGO and LIGO. If we have accurate theoretical templates of inspiraling phase of BNSs, we can determine the mass and the spin of neutron stars from the gravitational wave signals in the inspiraling phase. We may also extract the informations on the equation of state of a neutron star from the signals in the premerging phase. For this purpose it is important to complete theoretical templates of gravitational waves in the premerging phase as well as in the inspiraling phase. Moreover, the study of BNSs in the premerging phase is motivated by the need to provide the realistic initial condition for the merging phase simulations.
In order to obtain accurate theoretical templates of gravitational waves from BNSs around the premerging phase, a general relativistic hydrostatic problem with compressible equation of state must be numerically solved. In solving this problem, we suppose that BNSs reach quasi-equilibrium states because the timescale of the orbital decay driven by the radiation reaction is much longer than the orbital period of BNSs until their innermost stable circular orbit (ISCO). (The formalism for solving quasi-equilibrium figures of irrotational binary neutron stars in general relativity is given in the references .) Furthermore, we regard the internal state of a neutron star as an irrotational or nearly irrotational one because the viscosity of a neutron star is not large enough for synchronization even near the ISCO.
Recently Bonazzola, Gourgoulhon and Marck have numerically studied the irrotational binary neutron stars in general relativity with the conformally flat condition. More recently, Uryū and Eriguchi have also numerically solved the same problem as that of Bonazzola, Gourgoulhon and Marck. In order to check the validity of these numerical calculations, it is necessary to compare them with analytic or semi-analytic ones. However, we have not obtained such analytic solutions in general relativity yet. There are solutions of irrotational BNSs only in the first post-Newtonian (1PN) approximation of general relativity. For example, Lombardi, Rasio and Shapiro have semi-analytically studied irrotational BNSs by using the energy variational method, and one of the authors of the present paper (KT) has analytically calculated equilibrium sequences of irrotational BNSs by using the tensor virial method. The former one restricts the internal motion within the plane orthogonal to the rotational axis assuming the shape of the star to an ellipsoidal one in order to treat compressible equation of state. While in the latter one the velocity component along the orbital axis is included at 1PN order and the shape of the star at 1PN order is not restricted to an ellipsoidal one although the fluid is incompressible.
The situation is also the same in Newtonian gravity, that is, even in Newtonian gravity useful analytic or semi-analytic solutions for equilibrium sequences of irrotational binary systems are not obtained. Let us consider numerically constructed stationary structures of irrotational binary stars computed by Uryū and Eriguchi in Newtonian gravity. One may think that semi-analytic solutions by Lai, Rasio and Shapiro may be used to check the validity of numerical solutions. However, in numerical solutions of Uryū and Eriguchi the velocity component along the orbital axis exists while in those of Lai, Rasio and Shapiro such a component is assumed to be zero from the beginning. Therefore new analytic or semi-analytic solutions are needed to check numerical solutions even in Newtonian gravity. Such a check of numerical solutions is extremely important because in the numerical calculation, there is a possibility to obtain another solution although the binding energy of a binary neutron star is almost the same value, and to lead a different conclusion.
In the previous paper, we showed such a new, almost analytic solution to an equilibrium of irrotational binary polytropic stars for the polytropic index $`n=1`$ in Newtonian gravity by expanding all physical quantities in a power of $`ϵa_0/R`$, where $`R`$ and $`a_0`$ are the orbital separation of the binary system and the radius of each star for $`R=\mathrm{}`$. In that paper, we have extended the method developed by Chandrasekhar more than 65 years ago for corotating fluids to the one for irrotational fluids. In this paper, we show semi-analytic solutions for arbitrary polytropic indices by numerically solving ordinary differential equations.
This paper is organized as follows. In §2, we formulate the method to solve the irrotational binary polytropic stars. In §3, the physical values, i.e., the central density of a star, the orbital angular velocity, the total energy and total angular momentum of the binary system are calculated. The numerical results are presented in §4. Section 5 is devoted to summary and discussions.
Although a binary system consists of two stars, we pay particular attention to one of two stars. We call it star 1 whose mass is $`M_1`$ and the companion one star 2 whose mass is $`M_2`$. In this paper, we adopt two corotating coordinate systems. First one is a Cartesian coordinate system $`𝑿`$ whose origin is located at the center of mass of the binary system. For calculational convenience, we choose the orbital axis as $`X_3`$, and we take the direction of $`X_1`$ from the center of mass of star 2 to that of star 1. The second coordinate system is the spherical one $`𝒓=(r,\theta ,\phi )`$ whose origin is located at the center of mass of star 1. We use units of $`G=1`$.
## II Formulation
Since we treat irrotational fluids in Newtonian gravity, the basic equations for obtaining equilibrium configurations of binary systems are the equation of state, the Euler equation, the equation of continuity and the Poisson equation:
$`P=K\rho ^{1+\frac{1}{n}},`$ (1)
$`\left[{\displaystyle \frac{dP}{\rho }}U+{\displaystyle \frac{1}{2}}v^2𝒗(𝛀\times 𝒓)\right]=0,`$ (2)
$`𝒗=\left(𝒗𝛀\times 𝒓\right){\displaystyle \frac{\rho }{\rho }},`$ (3)
$`\mathrm{\Delta }U=4\pi \rho ,`$ (4)
where $`P`$, $`\rho `$, $`n`$, $`U`$ and $`\mathrm{\Omega }`$ are the pressure, the density, the polytropic index, the gravitational potential and the orbital angular velocity, respectively. $`K`$ is a constant related to entropy and $`𝒗`$ represents the velocity field in the inertial frame. The gravitational potential $`U`$ is separated into two parts, i.e., the contribution from star 1 to itself $`U^{11}`$ and that from star 2 to star 1 $`U^{21}`$. We can also express the gravitational potential contributed from star 1 to star 2 as $`U^{12}`$. These gravitational potentials are written as
$`U^{21}`$ $`=`$ $`{\displaystyle \frac{M_2}{R}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(1)^l\left({\displaystyle \frac{r}{R}}\right)^lP_l(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (6)
$`+{\displaystyle \frac{3I_{11}^{}}{2R^3}}\left[13\left({\displaystyle \frac{r}{R}}\right)P_1(\mathrm{sin}\theta \mathrm{cos}\phi )+O(R^2)\right]+\mathrm{higher}\mathrm{order}\mathrm{terms},`$
$`U^{12}`$ $`=`$ $`{\displaystyle \frac{M_1}{R}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{r^{}}{R}}\right)^lP_l(\mathrm{sin}\theta ^{}\mathrm{cos}\phi ^{})`$ (8)
$`+{\displaystyle \frac{3I_{11}}{2R^3}}\left[1+3\left({\displaystyle \frac{r^{}}{R}}\right)P_1(\mathrm{sin}\theta ^{}\mathrm{cos}\phi ^{})+O(R^2)\right]+\mathrm{higher}\mathrm{order}\mathrm{terms},`$
where the superscript means the term concerned with star 2 and $`P_l`$ denotes the Legendre function of order $`l`$. $`I_{11}`$ is the reduced quadrupole moment defined by
$`I_{11}`$ $`=`$ $`{\displaystyle \frac{1}{3}}\left(2I_{11}I_{22}I_{33}\right),`$ (9)
$`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle _{star1}}d^3x\rho r^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi ).`$ (10)
In the irrotational fluid case, we can express $`𝒗`$ as a gradient of a scalar function $`\mathrm{\Phi }`$, i.e.,
$$𝒗=\mathrm{\Phi }.$$
(11)
Following the Lane-Emden equation we first express the density as
$$\rho =\rho _c\mathrm{\Theta }^n(\xi ,\theta ,\phi )$$
(12)
with $`\xi =r/\alpha `$. Here the definitions of $`\alpha `$ and $`\alpha ^{}`$ are
$`\alpha `$ $``$ $`\left[{\displaystyle \frac{K(1+n)\rho _c^{\frac{1}{n}1}}{4\pi }}\right]^{1/2},`$ (13)
$`\alpha ^{}`$ $``$ $`\left[{\displaystyle \frac{K^{}(1+n^{})\rho _{c}^{}{}_{}{}^{\frac{1}{n^{}}1}}{4\pi }}\right]^{1/2},`$ (14)
with $`\rho _c`$ being the central density. Let us define $`ϵa_0/R`$, where $`R`$ and $`a_0`$ are the orbital separation of the binary system and the radius of star 1 for $`R=\mathrm{}`$. The radius of a spherical star $`a_0`$ is given by $`a_0=\alpha \xi _1`$ as usual. Also we define $`ϵ^{}a_0^{}/R`$ for star 2. Note that in the case of Newtonian gravity, we have two degrees of freedom for describing the differences of mass, i.e., $`\rho _c`$ and $`\alpha `$. Therefore, we can express the differences of mass by changing only $`\rho _c`$ with fixing $`\alpha `$. In this case, we can make the radius of each star and each expansion parameter coincide. Although we have adopted such a situation in Ref. , in the present paper we express equations concerned with star 2 by using $`ϵ^{}`$, $`\alpha ^{}`$ and so on.
We expand $`\mathrm{\Theta }`$ in a power series of a parameter $`ϵ`$ as<sup>*</sup><sup>*</sup>*Note that we express the Lane-Emden function for star 2 as $`\overline{\mathrm{\Theta }}`$ and expand it in a power series of a parameter $`ϵ^{}`$ as $`\overline{\mathrm{\Theta }}=_iϵ_{}^{}{}_{}{}^{i}\overline{\mathrm{\Theta }}_i`$.
$$\mathrm{\Theta }=\underset{i=0}{\overset{\mathrm{}}{}}ϵ^i\mathrm{\Theta }_i.$$
(15)
Since the shape of star 1 is spherical when $`R`$ is large, the lowest order term of $`\mathrm{\Theta }`$ is the solution of the Lane-Emden equation. Then we expand $`\mathrm{\Theta }_i`$ by spherical harmonics as
$$\mathrm{\Theta }_i=\underset{l,m}{}{}_{}{}^{(i)}\psi _{lm}^{}(\xi )Y_l^m(\theta ,\phi ).$$
(16)
Now we consider the orbital motion of star 1. In the spherical coordinate system, it becomes
$$𝛀\times 𝒓=\mathrm{\Omega }R\left(\stackrel{~}{𝛀}\times 𝝃\right)_{orb}+\mathrm{\Omega }\alpha \left(\stackrel{~}{𝛀}\times 𝝃\right)_{fig},$$
(17)
where
$`\left(\stackrel{~}{𝛀}\times 𝝃\right)_{orb}`$ $`=`$ $`{\displaystyle \frac{1}{1+p}}(\mathrm{sin}\theta \mathrm{sin}\phi ,\mathrm{cos}\theta \mathrm{sin}\phi ,\mathrm{cos}\phi ),`$ (18)
$`\left(\stackrel{~}{𝛀}\times 𝝃\right)_{fig}`$ $`=`$ $`(0,0,\xi \mathrm{sin}\theta ),`$ (19)
and $`pM_1/M_2`$. The first term on the right-hand side of Eq. (17) comes from the orbital motion of the center of mass of star 1 and the second term comes from the fluid motion around the center of mass of star 1. The orbital motion of star 2 becomes
$$𝛀\times 𝒓^{}=\mathrm{\Omega }R\left(\stackrel{~}{𝛀}\times 𝝃\right)_{orb}^{}+\mathrm{\Omega }\alpha ^{}\left(\stackrel{~}{𝛀}\times 𝝃\right)_{fig}^{},$$
(20)
where
$`\left(\stackrel{~}{𝛀}\times 𝝃\right)_{orb}^{}`$ $`=`$ $`{\displaystyle \frac{p}{1+p}}(\mathrm{sin}\theta ^{}\mathrm{sin}\phi ^{},\mathrm{cos}\theta ^{}\mathrm{sin}\phi ^{},\mathrm{cos}\phi ^{}),`$ (21)
$`\left(\stackrel{~}{𝛀}\times 𝝃\right)_{fig}^{}`$ $`=`$ $`(0,0,\xi ^{}\mathrm{sin}\theta ^{}),`$ (22)
in the spherical coordinate system whose origin is located at the center of mass of star 2. Here we have expressed the spherical coordinate system for star 2 as $`𝒓^{}=(r^{},\theta ^{},\phi ^{})`$ with $`\xi ^{}=r^{}/\alpha ^{}`$.
Next, we rewrite the equation of continuity (3) as
$$\mathrm{\Delta }\mathrm{\Phi }=n(\mathrm{\Phi }𝛀\times 𝒓)\frac{\mathrm{\Theta }}{\mathrm{\Theta }}.$$
(23)
The condition for $`\mathrm{\Phi }`$ at the stellar surface is
$$(\mathrm{\Phi }𝛀\times 𝒓)(\mathrm{\Theta })|_{surf}=0,$$
(24)
since $`\mathrm{\Theta }=0`$ at the surface. We expand $`\mathrm{\Phi }`$ also as
$$\mathrm{\Phi }=\underset{i=0}{\overset{\mathrm{}}{}}ϵ^i\mathrm{\Phi }_i.$$
(25)
When $`R`$ is large, the shape of star 1 is spherical and star 1 has only the orbital motion of the center of mass in the inertial frame with no intrinsic spin so that the gradient of the lowest order term of $`\mathrm{\Phi }`$ should agree with the orbital motion of the center of mass of star 1, i.e.,
$$\mathrm{\Phi }_0=\mathrm{\Omega }R(\stackrel{~}{𝛀}\times 𝝃)_{orb}.$$
(26)
Considering $`\mathrm{\Phi }_0`$, we normalize $`\mathrm{\Phi }`$ as
$$\stackrel{~}{\mathrm{\Phi }}=\frac{ϵ\mathrm{\Phi }}{\mathrm{\Omega }\alpha a_0}.$$
(27)
We again expand $`\stackrel{~}{\mathrm{\Phi }}_i`$ by spherical harmonics as
$$\stackrel{~}{\mathrm{\Phi }}_i=\underset{l,m}{}{}_{}{}^{(i)}\varphi _{lm}^{}(\xi )\stackrel{~}{Y}_l^m(\theta ,\phi ),$$
(28)
where $`\stackrel{~}{Y}_l^m`$ is defined by using the spherical harmonics $`Y_l^m`$ as
$$\stackrel{~}{Y}_l^m(\theta ,\phi )C_{lm}\frac{}{\phi }Y_l^m(\theta ,\phi )$$
(29)
with constants $`C_{lm}`$. Note that the gradient of the lowest order term of $`\mathrm{\Phi }`$ has order $`ϵ^1`$ except for the implicit dependence in $`\mathrm{\Omega }`$. When we normalize $`\mathrm{\Phi }`$ using Eq. (27), the lowest order term of $`\stackrel{~}{\mathrm{\Phi }}`$ has order $`ϵ^0`$. We must keep this difference in mind.
The orbital angular velocity is derived from the force balance equation, i.e., the first tensor virial relation, defined by
$$_{star1}d^3x\frac{P}{x_1}=0,$$
(30)
where $`x_1=r\mathrm{sin}\theta \mathrm{cos}\phi `$. This equation is expressed by using Eq. (2) as
$$0=_{star1}d^3x\rho \left[\frac{U^{11}}{x_1}+\frac{U^{21}}{x_1}\frac{1}{2}\frac{}{x_1}(\mathrm{\Phi })^2+\frac{}{x_1}\left\{(\mathrm{\Phi })(𝛀\times 𝒓)\right\}\right].$$
(31)
If we substitute Eq. (17), $`\mathrm{\Theta }_0`$ and $`\mathrm{\Phi }_0`$ into Eq. (31), we obtain the orbital angular velocity in the lowest order as
$$\mathrm{\Omega }_0^2=\frac{M_{tot}}{R^3},$$
(32)
where $`M_{tot}=M_1+M_2`$. Note here that we also expand $`\mathrm{\Omega }^2`$ as
$$\mathrm{\Omega }^2=\underset{i=0}{\overset{\mathrm{}}{}}ϵ^i\mathrm{\Omega }_i^2.$$
(33)
Finally, we express the gravitational potential by rewriting Eq. (2) as
$$U=K(1+n)\rho ^{\frac{1}{n}}+\frac{1}{2}v^2𝒗(𝛀\times 𝒓)+U_0,$$
(34)
where $`U_0`$ is constant. Substituting Eq. (34) into the Poisson equation (4), we obtain the equation to determine the equilibrium figure as
$$\alpha ^2\mathrm{\Delta }\mathrm{\Theta }=\mathrm{\Theta }^n\frac{1}{8\pi \rho _c}\mathrm{\Delta }\left[(\mathrm{\Phi })^22(\mathrm{\Phi })(𝛀\times 𝒓)\right].$$
(35)
Now a solution can be obtained iteratively. Firstly, $`\mathrm{\Theta }_i`$ is determined by demanding that the gravitational potential and its normal derivative are continuous at the stellar surface, that is,
$`U_{int}(\xi =\mathrm{\Xi })`$ $`=`$ $`U_{ext}(\xi =\mathrm{\Xi }),`$ (36)
$`{\displaystyle \frac{U_{int}}{\xi }}(\xi =\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{U_{ext}}{\xi }}(\xi =\mathrm{\Xi }),`$ (37)
where $`\mathrm{\Xi }(\theta ,\phi )`$ expresses the surface,
$$\mathrm{\Theta }(\xi =\mathrm{\Xi }(\theta ,\phi ))=0.$$
(38)
Substituting $`\mathrm{\Theta }_i`$ and Eq. (17) into Eqs. (23) and (31), we obtain $`\mathrm{\Phi }_i`$ and $`\mathrm{\Omega }_i^2`$. Secondly we substitute $`\mathrm{\Phi }_i`$ and $`\mathrm{\Omega }_i^2`$ into Eq. (35) and derive $`\mathrm{\Theta }_{i+1}`$. We have continued this procedure up to order $`ϵ^6`$ in this paper.
## III Derivation of the Physical Values
In this section, we calculate physical values such as the baryon mass of star 1, the orbital angular velocity, the total energy and total angular momentum of the binary system. We will just show the main results. All the details will be shown in the Appendix.
### A Mass and Central Density
The baryon mass of star 1 is calculated by integrating the density over the volume of star 1 as
$`M_1`$ $`=`$ $`{\displaystyle _{star1}}d^3x\rho ,`$ (39)
$`=`$ $`\alpha ^3\rho _c{\displaystyle }d\xi d\theta d\phi \xi ^2\mathrm{sin}\theta [\mathrm{\Theta }_0^n+ϵ^3n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_3+ϵ^4n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_4+ϵ^5n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_5`$ (41)
$`+ϵ^6n\{\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_6+{\displaystyle \frac{1}{2}}(n1)\mathrm{\Theta }_0^{n2}\mathrm{\Theta }_3^2\}].`$
Note here that the integration range for $`\xi `$ is $`[0,\mathrm{\Xi }(\theta ,\phi )]`$.
For $`n>0`$ case, the mass of star 1 becomes
$`M_1`$ $`=`$ $`4\pi \rho _c\alpha ^3\xi _1^2\mathrm{\Theta }_{0,\xi }(\xi _1)\left[1+ϵ^6\left({\displaystyle \frac{1}{\mathrm{\Theta }_{0,\xi }(\xi _1)}}\right)\left({\displaystyle \frac{({}_{}{}^{(3)}\psi _{2}^{}(\xi _1))^2}{10\mathrm{\Theta }_{0,\xi }(\xi _1)}}n\mathrm{\Theta }_0^{n1}(\xi _1)+{}_{}{}^{(6)}\psi _{0,\xi }^{}(\xi _1)\right)\right],`$ (42)
$`=`$ $`4\pi \rho _c\alpha ^3\xi _1^2\mathrm{\Theta }_{0,\xi }(\xi _1)\left[1+ϵ^6\left({\displaystyle \frac{1}{\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1)}}\right)\left(\stackrel{~}{c}_6+{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)\right)\right].`$ (43)
where we define
$${}_{,\xi }{}^{}\frac{d}{d\xi }.$$
(44)
We have used Eq. (H72) in deriving Eq. (43) from Eq. (42). Here $`\stackrel{~}{c}_6`$ is a constant, and one can see the expression of $`\stackrel{~}{c}_6`$ in Appendix H. We assume that the baryon mass $`M_1=4\pi \rho _{c0}\alpha _0^3\xi _1^2\mathrm{\Theta }_{0,\xi }(\xi _1)`$ is conserved throughout the equilibrium sequences of the binary system, where $`\rho _{c0}`$ and $`\alpha _0`$ denote the values for a spherical star. Therefore, the central density and $`\alpha `$ are expressed as
$`\rho _c`$ $`=`$ $`\rho _{c0}\left(1ϵ^6{\displaystyle \frac{\delta \rho _c}{\rho _c}}\right),`$ (45)
$`\alpha `$ $`=`$ $`\left[{\displaystyle \frac{K(1+n)\rho _{c0}^{\frac{1}{n}1}}{4\pi }}\right]^{1/2}\left(1ϵ^6{\displaystyle \frac{1n}{2n}}{\displaystyle \frac{\delta \rho _c}{\rho _c}}\right),`$ (46)
where we defineThe expression of $`\delta \rho _c/\rho _c`$ seems to diverge if $`n<1`$ when one see Eq. (42). However physically $`\delta \rho _c/\rho _c`$ should be finite. In reality we determined $`\delta \rho _c/\rho _c`$ by using the virial relation as shown in the Appendix I.
$$\frac{\delta \rho _c}{\rho _c}\frac{2n}{3n}\left(\frac{1}{\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1)}\right)\left(\stackrel{~}{c}_6+{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)\right).$$
(47)
### B Orbital angular velocity
We obtain the orbital angular velocity from Eq. (31) as
$`\mathrm{\Omega }^2`$ $`=`$ $`{\displaystyle \frac{M_{tot}}{R^3}}\left[1+{\displaystyle \frac{9}{2R^2}}\left(ϵ^3{\displaystyle \frac{\overline{I}_{11}}{M_1}}+ϵ_{}^{}{}_{}{}^{3}{\displaystyle \frac{\overline{I}_{11}^{}}{M_2}}\right)+O(ϵ^7)\right],`$ (48)
$`=`$ $`{\displaystyle \frac{M_{tot}}{a_0^3}}ϵ^3\left[1+{\displaystyle \frac{9ϵ^2}{2a_0^2}}\left(ϵ^3{\displaystyle \frac{\overline{I}_{11}}{M_1}}+ϵ_{}^{}{}_{}{}^{3}{\displaystyle \frac{\overline{I}_{11}^{}}{M_2}}\right)+O(ϵ^7)\right],`$ (49)
where $`\overline{I}_{11}`$ denotes $`I_{11}/ϵ^3`$.
### C Total energy
The total energy is written as
$$E=\mathrm{\Pi }_{tot}+(W_{self})_{tot}+(W_{int})_{tot}+T_{tot},$$
(50)
where $`\mathrm{\Pi }_{tot}`$, $`(W_{self})_{tot}`$, $`(W_{int})_{tot}`$ and $`T_{tot}`$ denote the total internal energy, the total self-gravity energy, the total interaction energy and the total kinetic energy, respectively. We calculate each energy in the following.
#### 1 Internal energy
The definition of the internal energy of star 1 is
$$\mathrm{\Pi }_1=n_{star1}d^3xP.$$
(51)
For $`n>0`$ case, we obtain
$$\mathrm{\Pi }_1=\frac{n}{1+n}\frac{M_1^2}{a_0\xi _1^3(\mathrm{\Theta }_{0,\xi }(\xi _1))^2}\left[\left(1ϵ^6\frac{5n}{2n}\frac{\delta \rho _c}{\rho _c}\right)_0^{\xi _1}𝑑\xi \xi ^2\mathrm{\Theta }_0^{1+n}+ϵ^6(1+n)_0^{\xi _1}𝑑\xi \xi ^2\left\{\mathrm{\Theta }_0^n^{(6)}\psi _0+\frac{1}{10}n\mathrm{\Theta }_0^{n1}({}_{}{}^{(3)}\psi _{2}^{})^2\right\}\right].$$
(52)
Therefore, the total internal energy is expressed as
$`\mathrm{\Pi }_{tot}`$ $`=`$ $`{\displaystyle \frac{n}{1+n}}{\displaystyle \frac{M_1^2}{a_0\xi _1^3(\mathrm{\Theta }_{0,\xi }(\xi _1))^2}}\left[\left(1ϵ^6{\displaystyle \frac{5n}{2n}}{\displaystyle \frac{\delta \rho _c}{\rho _c}}\right){\displaystyle _0^{\xi _1}}𝑑\xi \xi ^2\mathrm{\Theta }_0^{1+n}+ϵ^6(1+n){\displaystyle _0^{\xi _1}}𝑑\xi \xi ^2\left\{\mathrm{\Theta }_0^n{}_{}{}^{(6)}\psi _{0}^{}+{\displaystyle \frac{1}{10}}n\mathrm{\Theta }_0^{n1}({}_{}{}^{(3)}\psi _{2}^{})^2\right\}\right]`$ (55)
$`+{\displaystyle \frac{n^{}}{1+n^{}}}{\displaystyle \frac{M_2^2}{a_0^{}\xi _{1}^{}{}_{}{}^{3}(\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}))^2}}[(1ϵ_{}^{}{}_{}{}^{6}{\displaystyle \frac{5n^{}}{2n^{}}}{\displaystyle \frac{\delta \rho _c^{}}{\rho _c^{}}}){\displaystyle _0^{\xi _1^{}}}d\xi ^{}\xi _{}^{}{}_{}{}^{2}\overline{\mathrm{\Theta }}_0^{1+n^{}}`$
$`+ϵ_{}^{}{}_{}{}^{6}(1+n^{}){\displaystyle _0^{\xi _1^{}}}d\xi ^{}\xi _{}^{}{}_{}{}^{2}\{\overline{\mathrm{\Theta }}_0^n^{}{}_{}{}^{(6)}\overline{\psi }_{0}^{}+{\displaystyle \frac{1}{10}}n^{}\overline{\mathrm{\Theta }}_0^{n^{}1}({}_{}{}^{(3)}\overline{\psi }_{2}^{})^2\}],`$
where the superscript means the terms concerned with star 2, and $`\overline{}`$ also means the functions concerned with star 2.
#### 2 Self-gravity energy
The definition of the self-gravity energy of star 1 is
$$W_{self,1}=\frac{1}{2}_{star1}d^3x\rho U^{11}.$$
(57)
For $`n>0`$ case, we obtain
$$W_{self,1}=\frac{1+n}{2n}\mathrm{\Pi }_1+ϵ^6\pi \rho _c\left[\frac{3\mu \overline{I}_{11}}{\xi _1^2(1+p)}2\alpha ^2M_1\stackrel{~}{c}_6\right]+2\pi \rho _c\alpha ^2M_1\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1),$$
(58)
where $`\mu `$ is defined as
$$\mu \frac{M_{tot}}{4\pi \rho _c\alpha ^3\xi _1}=\left(\frac{1+p}{p}\right)\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1).$$
(59)
The self-gravity energy of star 2 is
$$W_{self,2}=\frac{1+n^{}}{2n^{}}\mathrm{\Pi }_2+ϵ_{}^{}{}_{}{}^{6}\pi \rho _c^{}\left[\frac{3\mu ^{}p\overline{I}_{11}^{}}{\xi _{1}^{}{}_{}{}^{2}(1+p)}2\alpha _{}^{}{}_{}{}^{2}M_2\stackrel{~}{c}_6^{}\right]+2\pi \rho _c^{}\alpha _{}^{}{}_{}{}^{2}M_2\xi _1^{}\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}).$$
(60)
Here, $`\stackrel{~}{c}_6^{}`$ is a constant similar to $`\stackrel{~}{c}_6`$, and
$$\mu ^{}\frac{M_{tot}}{4\pi \rho _c^{}\alpha _{}^{}{}_{}{}^{3}\xi _1^{}}=(1+p)\xi _1^{}\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}).$$
(61)
Therefore, the total self-gravity energy becomes
$`(W_{self})_{tot}`$ $`=`$ $`{\displaystyle \frac{1+n}{2n}}\mathrm{\Pi }_1+ϵ^6{\displaystyle \frac{M_1^2}{2a_0}}\left({\displaystyle \frac{3}{2p}}{\displaystyle \frac{\overline{I}_{11}}{M_1a_0^2}}+{\displaystyle \frac{\stackrel{~}{c}_6}{\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1)}}\right){\displaystyle \frac{M_1^2}{2a_0}}\left(1ϵ^6{\displaystyle \frac{1}{n}}{\displaystyle \frac{\delta \rho _c}{\rho _c}}\right)`$ (63)
$`{\displaystyle \frac{1+n^{}}{2n^{}}}\mathrm{\Pi }_2+ϵ_{}^{}{}_{}{}^{6}{\displaystyle \frac{M_2^2}{2a_0^{}}}\left({\displaystyle \frac{3p}{2}}{\displaystyle \frac{\overline{I}_{11}^{}}{M_2a_{0}^{}{}_{}{}^{2}}}+{\displaystyle \frac{\stackrel{~}{c}_6^{}}{\xi _1^{}\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{})}}\right){\displaystyle \frac{M_2^2}{2a_0^{}}}\left(1ϵ_{}^{}{}_{}{}^{6}{\displaystyle \frac{1}{n^{}}}{\displaystyle \frac{\delta \rho _c^{}}{\rho _c^{}}}\right).`$
#### 3 Interaction energy
The definition of the interaction energy of star 1 is
$`W_{int,1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{star1}}d^3x\rho U^{21},`$ (64)
$`=`$ $`{\displaystyle \frac{M_1M_2}{2a_0}}ϵϵ^3{\displaystyle \frac{3M_1M_2}{4a_0^3}}\left(ϵ^3{\displaystyle \frac{\overline{I}_{11}}{M_1}}+ϵ_{}^{}{}_{}{}^{3}{\displaystyle \frac{\overline{I}_{11}^{}}{M_2}}\right).`$ (65)
Therefore, the total interaction energy becomes
$$(W_{int})_{tot}=\frac{M_1M_2}{a_0}ϵϵ^3\frac{3M_1M_2}{2a_0^3}\left(ϵ^3\frac{\overline{I}_{11}}{M_1}+ϵ_{}^{}{}_{}{}^{3}\frac{\overline{I}_{11}^{}}{M_2}\right).$$
(66)
#### 4 Kinetic energy
The definition of the kinetic energy of star 1 is
$`T_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{star1}}d^3x\rho v_1^2,`$ (67)
$`=`$ $`{\displaystyle \frac{M_1M_2}{2a_0(1+p)}}ϵ\left[1+{\displaystyle \frac{9ϵ^2}{2a_0^2}}\left(ϵ^3{\displaystyle \frac{\overline{I}_{11}}{M_1}}+ϵ_{}^{}{}_{}{}^{3}{\displaystyle \frac{\overline{I}_{11}^{}}{M_2}}\right)\right].`$ (68)
The kinetic energy of star 2 is also given as
$$T_2=\frac{M_1M_2p}{2a_0(1+p)}ϵ\left[1+\frac{9ϵ^2}{2a_0^2}\left(ϵ^3\frac{\overline{I}_{11}}{M_1}+ϵ_{}^{}{}_{}{}^{3}\frac{\overline{I}_{11}^{}}{M_2}\right)\right].$$
(69)
Then the total kinetic energy becomes
$$T_{tot}=\frac{M_1M_2}{2a_0}ϵ+ϵ^3\frac{9M_1M_2}{4a_0^3}\left(ϵ^3\frac{\overline{I}_{11}}{M_1}+ϵ_{}^{}{}_{}{}^{3}\frac{\overline{I}_{11}^{}}{M_2}\right).$$
(70)
### D Total angular momentum
The total angular momentum is calculated from the equation;
$$𝑱=_{star1,star2}d^3x\rho (𝑹\times 𝒗),$$
(71)
where $`𝑹`$ for star 1 is written as
$$𝑹=\frac{R}{1+p}(\mathrm{sin}\theta \mathrm{cos}\phi ,\mathrm{cos}\theta \mathrm{cos}\phi ,\mathrm{sin}\phi )+r(1,0,0),$$
(72)
and for star 2 is
$$𝑹^{}=\frac{pR}{1+p}(\mathrm{sin}\theta ^{}\mathrm{cos}\phi ^{},\mathrm{cos}\theta ^{}\mathrm{cos}\phi ^{},\mathrm{sin}\phi ^{})+r^{}(1,0,0)$$
(73)
in the spherical coordinate system. Note that since the velocity field $`𝒗=\mathrm{\Phi }`$ is written in the spherical coordinate system, it is convenient to calculate $`𝑹\times 𝒗`$ in the spherical one. However, we should be careful to integrate it over the volume of each star, because in the spherical coordinate system, the unit vectors depend on the coordinates. For example, the unit vector $`(1,0,0)`$ in the Cartesian coordinate system becomes $`(\mathrm{sin}\theta \mathrm{cos}\phi ,\mathrm{cos}\theta \mathrm{cos}\phi ,\mathrm{sin}\phi )`$ in the spherical one.
As a result, we obtain the total angular momentum;
$$𝑱=(0,0,J)$$
(74)
in the Cartesian coordinate system. Here $`J`$ is calculated as
$$J=\frac{M_1M_2}{M_1+M_2}R^2\mathrm{\Omega }\left[1+\mathrm{higher}\mathrm{order}\mathrm{terms}\mathrm{than}O(ϵ^6)\right].$$
(75)
### E Virial Relation
By using the above energies, we can calculate the virial relation which is written as
$$\frac{3}{n}\mathrm{\Pi }_1+\frac{3}{n^{}}\mathrm{\Pi }_2+(W_{self})_{tot}+(W_{int})_{tot}+2T_{tot}=0.$$
(76)
This equation is used to check the solutions as shown in the previous letter. However in this paper we use it to determine $`\delta \rho _c/\rho _c`$ (see Appendix I).
## IV Numerical results
In the following sections, we set $`M_1=M_2=M`$, $`n=n^{}`$, $`\alpha =\alpha ^{}`$, $`\rho _c=\rho _c^{}`$ and so on, i.e., the identical star binary for simplicity, and investigate the cases of different polytropic indices.
First, we show the Lane-Emden function of star 1 up to $`O(ϵ^6)`$. It is written as
$$\mathrm{\Theta }=\mathrm{\Theta }_0+ϵ^3\mathrm{\Theta }_3+ϵ^4\mathrm{\Theta }_4+ϵ^5\mathrm{\Theta }_5+ϵ^6\mathrm{\Theta }_6,$$
(77)
where
$`\mathrm{\Theta }_3`$ $`=`$ $`{}_{}{}^{(3)}\psi _{2}^{}(\xi )P_2(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (78)
$`\mathrm{\Theta }_4`$ $`=`$ $`{}_{}{}^{(4)}\psi _{3}^{}(\xi )P_3(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (79)
$`\mathrm{\Theta }_5`$ $`=`$ $`{}_{}{}^{(5)}\psi _{4}^{}(\xi )P_4(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (80)
$`\mathrm{\Theta }_6`$ $`=`$ $`{}_{}{}^{(6)}\psi _{0}^{}(\xi )+{}_{}{}^{(6)}\psi _{2}^{}(\xi )P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{}_{}{}^{(6)}\psi _{22}^{}(\xi )P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (82)
$`+{}_{}{}^{(6)}\psi _{4}^{}(\xi )P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{}_{}{}^{(6)}\psi _{5}^{}(\xi )P_5(\mathrm{sin}\theta \mathrm{cos}\phi ).`$
The functions $`\mathrm{\Theta }_0`$, $`^{(i)}`$$`\psi _j`$ and their derivatives for each polytropic index are given in Tables IVIII.
Next, we give the velocity potential up to $`O(ϵ^6)`$. We can write it as
$$\stackrel{~}{\mathrm{\Phi }}=\stackrel{~}{\mathrm{\Phi }}_0+ϵ^4\stackrel{~}{\mathrm{\Phi }}_4+ϵ^5\stackrel{~}{\mathrm{\Phi }}_5+ϵ^6\stackrel{~}{\mathrm{\Phi }}_6,$$
(83)
where
$`\stackrel{~}{\mathrm{\Phi }}_0`$ $`=`$ $`{\displaystyle \frac{1}{1+p}}\xi \mathrm{sin}\theta \mathrm{sin}\phi ,`$ (84)
$`\stackrel{~}{\mathrm{\Phi }}_4`$ $`=`$ $`{}_{}{}^{(4)}\varphi _{2}^{}(\xi )P_2^2(\mathrm{cos}\theta )\mathrm{sin}2\phi ,`$ (85)
$`\stackrel{~}{\mathrm{\Phi }}_5`$ $`=`$ $`{}_{}{}^{(5)}\varphi _{3}^{}(\xi )\left[P_3^1(\mathrm{cos}\theta )\mathrm{sin}\phi {\displaystyle \frac{1}{2}}P_3^3(\mathrm{cos}\theta )\mathrm{sin}3\phi \right],`$ (86)
$`\stackrel{~}{\mathrm{\Phi }}_6`$ $`=`$ $`{}_{}{}^{(6)}\varphi _{4}^{}(\xi )\left[P_4^2(\mathrm{cos}\theta )\mathrm{sin}2\phi {\displaystyle \frac{1}{4}}P_4^4(\mathrm{cos}\theta )\mathrm{sin}4\phi \right].`$ (87)
The functions $`{}_{}{}^{(4)}\varphi _{2}^{}`$, $`{}_{}{}^{(5)}\varphi _{3}^{}`$ and $`{}_{}{}^{(6)}\varphi _{4}^{}`$ and their derivatives are shown in Tables IXXII. The derivatives are related to the velocity field in the inertial frame. Note here that in Eq. (84) we write the dependence of the mass ratio $`p=M_1/M_2`$ although we consider only $`p=1`$ case in this section.
Here we directly show the existence of the velocity component along the orbital axis. In this paper, we take $`X_3`$ as the orbital axis. Since the origin of the corotating coordinate system $`𝑿`$ is located at the center of mass of the binary system, it is convenient to consider another Cartesian coordinate system which origin is located at the center of mass of star 1. We call it $`𝒙=(x_1,x_2,x_3)`$. The relations between $`𝒙`$ and the spherical coordinate system $`𝒓`$ are
$`x_1`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{cos}\phi ,`$ (88)
$`x_2`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{sin}\phi ,`$ (89)
$`x_3`$ $`=`$ $`r\mathrm{cos}\theta ,`$ (90)
as usual. Then, the velocity component along the orbital axis in the inertial frame is written as
$$v_3=\frac{\mathrm{\Phi }}{x_3}.$$
(91)
In the corotating frame, we can express the velocity field $`𝒖`$ as
$`𝒖`$ $`=`$ $`\mathrm{\Phi }𝛀\times 𝒓,`$ (92)
$`=`$ $`\mathrm{\Omega }\alpha (\stackrel{~}{𝛀}\times 𝝃)_{fig}+ϵ^4\mathrm{\Phi }_4+ϵ^5\mathrm{\Phi }_5+ϵ^6\mathrm{\Phi }_6,`$ (93)
$`=`$ $`\mathrm{\Omega }\alpha (\stackrel{~}{𝛀}\times 𝝃)_{fig}+\mathrm{\Omega }a_0\left[ϵ^3\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_4+ϵ^4\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_5+ϵ^5\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_6\right],`$ (94)
where $`\stackrel{~}{}=\alpha `$ is a nabla defined by using $`(\xi ,\theta ,\phi )`$ and given in Eq. (A4). Then, the velocity component along the orbital axis in the corotating frame becomes
$$u_3=\mathrm{\Omega }a_0\left[ϵ^3\frac{\stackrel{~}{\mathrm{\Phi }}_4}{\stackrel{~}{x}_3}+ϵ^4\frac{\stackrel{~}{\mathrm{\Phi }}_5}{\stackrel{~}{x}_3}+ϵ^5\frac{\stackrel{~}{\mathrm{\Phi }}_6}{\stackrel{~}{x}_3}\right],$$
(95)
where we have used the normalized coordinate system,
$`\stackrel{~}{x}_1`$ $`=`$ $`\xi \mathrm{sin}\theta \mathrm{cos}\phi ,`$ (96)
$`\stackrel{~}{x}_2`$ $`=`$ $`\xi \mathrm{sin}\theta \mathrm{sin}\phi ,`$ (97)
$`\stackrel{~}{x}_3`$ $`=`$ $`\xi \mathrm{cos}\theta ,`$ (98)
for the calculational convenience. Note that since $`\mathrm{\Phi }`$ is proportional to $`\mathrm{\Omega }/ϵ`$, the equation (93) does not express the real dependence on $`ϵ`$.
Next, we calculate the term for $`\stackrel{~}{\mathrm{\Phi }}_4`$ in Eq. (95). Then we obtain
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\stackrel{~}{x}_3}}`$ $`=`$ $`\mathrm{cos}\theta {\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\xi }}{\displaystyle \frac{\mathrm{sin}\theta }{\xi }}{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\theta }},`$ (99)
$`=`$ $`3\left({\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}{\displaystyle \frac{2{}_{}{}^{(4)}\varphi _{2}^{}}{\xi }}\right)\mathrm{sin}^2\theta \mathrm{cos}\theta \mathrm{sin}2\phi .`$ (100)
If $`{}_{}{}^{(4)}\varphi _{2}^{}(\xi )`$ is proportional to $`\xi ^2`$, the velocity component along the orbital axis disappearsFor the $`n=0`$ case, i.e., the incompressible fluid case, $`{}_{}{}^{(4)}\varphi _{2}^{}`$ is proportional to $`\xi ^2`$.. However, $`{}_{}{}^{(4)}\varphi _{2}^{}(\xi )`$ is not proportional to $`\xi ^2`$ at all for $`n>0`$. Moreover, we can show the velocity components along the orbital axis for $`\mathrm{\Phi }_5`$ and $`\mathrm{\Phi }_6`$ as
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_5}{\stackrel{~}{x}_3}}`$ $`=`$ $`6\left[\left({\displaystyle \frac{d{}_{}{}^{(5)}\varphi _{3}^{}}{d\xi }}{\displaystyle \frac{3{}_{}{}^{(5)}\varphi _{3}^{}}{\xi }}\right)(15\mathrm{sin}^2\theta \mathrm{cos}^2\phi )+{\displaystyle \frac{2{}_{}{}^{(5)}\varphi _{3}^{}}{\xi }}\right]\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{sin}\phi ,`$ (101)
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_6}{\stackrel{~}{x}_3}}`$ $`=`$ $`30\left[\left({\displaystyle \frac{d{}_{}{}^{(6)}\varphi _{4}^{}}{d\xi }}{\displaystyle \frac{4{}_{}{}^{(6)}\varphi _{4}^{}}{\xi }}\right)(37\mathrm{sin}^2\theta \mathrm{cos}^2\phi )+{\displaystyle \frac{6{}_{}{}^{(6)}\varphi _{4}^{}}{\xi }}\right]\mathrm{sin}^2\theta \mathrm{cos}\theta \mathrm{sin}\phi \mathrm{cos}\phi .`$ (102)
In these equations, we can see that even if $`{}_{}{}^{(5)}\varphi _{3}^{}\xi ^3`$ and $`{}_{}{}^{(6)}\varphi _{4}^{}\xi ^4`$, there remain the velocity components along the orbital axis. The velocity component along the orbital axis is order $`ϵ^3`$ higher. However, if one would like to pay attention to the internal state of neutron stars in the binary system, it should be taken into account.
We give also the other components of the velocity field as
$`u_1`$ $`=`$ $`\mathrm{\Omega }x_2+\mathrm{\Omega }a_0\left[ϵ^3{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\stackrel{~}{x}_1}}+ϵ^4{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_5}{\stackrel{~}{x}_1}}+ϵ^5{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_6}{\stackrel{~}{x}_1}}\right],`$ (103)
$`u_2`$ $`=`$ $`\mathrm{\Omega }x_1+\mathrm{\Omega }a_0\left[ϵ^3{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\stackrel{~}{x}_2}}+ϵ^4{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_5}{\stackrel{~}{x}_2}}+ϵ^5{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_6}{\stackrel{~}{x}_2}}\right],`$ (104)
where
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\stackrel{~}{x}_1}}`$ $`=`$ $`\mathrm{sin}\theta \mathrm{cos}\phi {\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\xi }}+{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\phi }{\xi }}{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\theta }}{\displaystyle \frac{\mathrm{sin}\phi }{\xi \mathrm{sin}\theta }}{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\phi }},`$ (105)
$`=`$ $`6\left[\left({\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}{\displaystyle \frac{2{}_{}{}^{(4)}\varphi _{2}^{}}{\xi }}\right)\mathrm{sin}^2\theta \mathrm{cos}^2\phi +{\displaystyle \frac{{}_{}{}^{(4)}\varphi _{2}^{}}{\xi }}\right]\mathrm{sin}\theta \mathrm{sin}\phi ,`$ (106)
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\stackrel{~}{x}_2}}`$ $`=`$ $`\mathrm{sin}\theta \mathrm{sin}\phi {\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\xi }}+{\displaystyle \frac{\mathrm{cos}\theta \mathrm{sin}\phi }{\xi }}{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\theta }}+{\displaystyle \frac{\mathrm{cos}\phi }{\xi \mathrm{sin}\theta }}{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\phi }},`$ (107)
$`=`$ $`6\left[\left({\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}{\displaystyle \frac{2{}_{}{}^{(4)}\varphi _{2}^{}}{\xi }}\right)\mathrm{sin}^2\theta \mathrm{sin}^2\phi +{\displaystyle \frac{{}_{}{}^{(4)}\varphi _{2}^{}}{\xi }}\right]\mathrm{sin}\theta \mathrm{cos}\phi ,`$ (108)
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_5}{\stackrel{~}{x}_1}}`$ $`=`$ $`6\left[\left({\displaystyle \frac{d{}_{}{}^{(5)}\varphi _{3}^{}}{d\xi }}{\displaystyle \frac{3{}_{}{}^{(5)}\varphi _{3}^{}}{\xi }}\right)(15\mathrm{sin}^2\theta \mathrm{cos}^2\phi ){\displaystyle \frac{8{}_{}{}^{(5)}\varphi _{3}^{}}{\xi }}\right]\mathrm{sin}^2\theta \mathrm{sin}\phi \mathrm{cos}\phi ,`$ (109)
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_5}{\stackrel{~}{x}_2}}`$ $`=`$ $`6\left[\left({\displaystyle \frac{d{}_{}{}^{(5)}\varphi _{3}^{}}{d\xi }}{\displaystyle \frac{3{}_{}{}^{(5)}\varphi _{3}^{}}{\xi }}\right)(15\mathrm{sin}^2\theta \mathrm{cos}^2\phi )\mathrm{sin}^2\theta \mathrm{sin}^2\phi +{\displaystyle \frac{{}_{}{}^{(5)}\varphi _{3}^{}}{\xi }}\left\{1+\mathrm{sin}^2\theta (2\mathrm{sin}^2\phi 5\mathrm{cos}^2\phi )\right\}\right],`$ (110)
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_6}{\stackrel{~}{x}_1}}`$ $`=`$ $`30\left[\left({\displaystyle \frac{d{}_{}{}^{(6)}\varphi _{4}^{}}{d\xi }}{\displaystyle \frac{4{}_{}{}^{(6)}\varphi _{4}^{}}{\xi }}\right)(37\mathrm{sin}^2\theta \mathrm{cos}^2\phi )\mathrm{sin}^3\theta \mathrm{sin}\phi \mathrm{cos}^2\phi +{\displaystyle \frac{3{}_{}{}^{(6)}\varphi _{4}^{}}{\xi }}(15\mathrm{sin}^2\theta \mathrm{cos}^2\phi )\mathrm{sin}\theta \mathrm{sin}\phi \right],`$ (111)
$`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_6}{\stackrel{~}{x}_2}}`$ $`=`$ $`30[({\displaystyle \frac{d{}_{}{}^{(6)}\varphi _{4}^{}}{d\xi }}{\displaystyle \frac{4{}_{}{}^{(6)}\varphi _{4}^{}}{\xi }})(37\mathrm{sin}^2\theta \mathrm{cos}^2\phi )\mathrm{sin}^3\theta \mathrm{sin}^2\phi \mathrm{cos}\phi `$ (113)
$`+{\displaystyle \frac{{}_{}{}^{(6)}\varphi _{4}^{}}{\xi }}\{3+\mathrm{sin}^2\theta (7\mathrm{cos}^2\phi 6\mathrm{sin}^2\phi )\}\mathrm{sin}\theta \mathrm{cos}\phi ].`$
In Table XIII, we show the first zero point of $`\mathrm{\Theta }_0`$, the orbital separation at the contact point, the energies, the quadrupole moment, and the change in the central density. In this table, we have used
$$\stackrel{~}{E}\frac{E}{M^2/a_0}=\stackrel{~}{E}_{self}+ϵ\stackrel{~}{E}_{point}+ϵ^6\stackrel{~}{E}_{quad}.$$
(114)
Each energy is given in Table XIV, where we have defined as
$`\stackrel{~}{\mathrm{\Pi }}_{tot}`$ $`=`$ $`\stackrel{~}{\mathrm{\Pi }}_{self}+ϵ^6\stackrel{~}{\mathrm{\Pi }}_{quad},`$ (115)
$`(\stackrel{~}{W}_{self})_{tot}`$ $`=`$ $`(\stackrel{~}{W}_{self})_{self}+ϵ^6(\stackrel{~}{W}_{self})_{quad},`$ (116)
$`(\stackrel{~}{W}_{int})_{tot}`$ $`=`$ $`ϵ(\stackrel{~}{W}_{int})_{point}+ϵ^6(\stackrel{~}{W}_{int})_{quad},`$ (117)
$`\stackrel{~}{T}_{tot}`$ $`=`$ $`ϵ\stackrel{~}{T}_{point}+ϵ^6\stackrel{~}{T}_{quad}.`$ (118)
Accordingly, we can express
$`\stackrel{~}{E}_{self}`$ $`=`$ $`\stackrel{~}{\mathrm{\Pi }}_{self}+(\stackrel{~}{W}_{self})_{self},`$ (119)
$`\stackrel{~}{E}_{point}`$ $`=`$ $`(\stackrel{~}{W}_{int})_{point}+\stackrel{~}{T}_{point},`$ (120)
$`\stackrel{~}{E}_{quad}`$ $`=`$ $`\stackrel{~}{\mathrm{\Pi }}_{quad}+(\stackrel{~}{W}_{self})_{quad}+(\stackrel{~}{W}_{int})_{quad}+\stackrel{~}{T}_{quad}.`$ (121)
The total angular momentum for the identical star binary becomes
$$J=\frac{M}{2}R^2\mathrm{\Omega },$$
(122)
and the orbital angular velocity is calculated as
$$\mathrm{\Omega }^2=\frac{2M}{a_0^3}ϵ^3\left[1+\frac{9ϵ^5}{a_0^2}\frac{\overline{I}_{11}}{M}+O(ϵ^7)\right].$$
(123)
We show the total energy, the total angular momentum and the orbital angular velocity along the equilibrium sequences of the binary system in Table XV.
In Figs. 14, we show (a) the total energy and (b) the total angular momentum as functions of the orbital separation, and (c) the orbital angular velocity as functions of the total angular momentum for each polytropic index, where we have used
$$\rho _0\frac{3M}{4\pi a_0^3}$$
(124)
In these figures, we compare our results (solid lines) with those of numerical calculations computed by Uryū and Eriguchi (open triangles) and those of semi-analytic ones proposed by Lai, Rasio and Shapiro (filled circles). For the smaller orbital separation, we must expand the physical values up to higher order than $`O(ϵ^6)`$ in order to include the effect of the spin of each star. In the case of the irrotational binary system, there is no intrinsic spin for the large distant stars. However, when the stars reach close range, the small spins are produced by the deformations of the stars. For example, the effect of these spins appears in the total energy at $`O(ϵ^9)`$<sup>§</sup><sup>§</sup>§This is because the spin kinetic energy is given by the volume integral of $`(ϵ^4\mathrm{\Phi }_4)^2\mathrm{\Omega }^2ϵ^6(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_4)^2`$, where $`\mathrm{\Omega }^2`$ is $`O(ϵ^3)`$. The definition of $`\stackrel{~}{}`$ is in Eq. (A4).. The effect of the spin in the case of the smaller polytropic index is larger than that in the case of the larger polytropic index. The reason is that the matter is not concentrated in the case of the smaller polytropic index, and the quadrupole moment becomes larger (see Table XIII). Therefore, at the smaller orbital separation, the deviation between our results and those of Uryū and Eriguchi for $`n=0.5`$ (Fig. 1) is larger than that for $`n=1.5`$ (Fig. 3).
Finally, we represent the accuracy of our numerical calculations by comparing numerical and analytic solutions for $`n=1`$ cases. In Fig. 5, we show the absolute values of the relative error for $`{}_{}{}^{(3)}\psi _{2}^{}`$, $`{}_{}{}^{(4)}\psi _{3}^{}`$, $`{}_{}{}^{(5)}\psi _{4}^{}`$, $`{}_{}{}^{(6)}\psi _{2}^{}`$, $`{}_{}{}^{(6)}\psi _{4}^{}`$ and $`{}_{}{}^{(6)}\psi _{5}^{}`$, where we define the relative error by
$$\left|\frac{\mathrm{Numerical}\mathrm{solution}}{\mathrm{Analytic}\mathrm{solution}}1\right|.$$
(125)
We can see from Fig. 5 that the relative errors are within $`10^8`$. Therefore, we regard that our results have enough accuracy.
## V Summary and Discussions
### A Summary
In this paper, we have calculated the equilibrium solutions of irrotational binary polytropic stars in Newtonian gravity by expanding all physical quantities in a power of $`ϵ`$. We have presented the results of the cases of several polytropic indices ($`n=0.5,1,1.5`$, and $`2`$). In particular, we have shown the velocity fields by solving the equation of continuity. It is found that there exists the small velocity component along the orbital axis (see §IV). It agrees with the numerical calculations performed by Uryū and Eriguchi and Bonazzola, Gourgoulhon and Marck.
Furthermore, we have given the figures and tables of the total energy, total angular momentum and orbital angular velocity along the equilibrium sequences for each polytropic indices. We can see from these figures that our results agree with those of Lai, Rasio and Shapiro and Uryū and Eriguchi for $`R/a_0>3`$.
Since our solutions are correct if $`ϵ1`$, they can be used to check the validity of numerical solutions. For any numerical codes, one can ask to solve an equilibrium for large $`R`$, and compare numerically derived velocity distribution and so on with our semi-analytic solutions.
However, since we expanded physical quantities up to $`O(ϵ^6)`$, it may not be enough to discuss about the behavior of solutions for small $`R`$. In order to apply our solutions in the case of small $`R`$ and to check the validity of numerical codes in this case, further higher order calculations are needed.
### B Discussions
It is important to compare the velocity field which we obtain by solving the continuity equation with that given by Lai, Rasio and Shapiro. We can see the velocity field they give in their papers or in the Chandrasekhar’s textbook. In the irrotational case, it becomes
$`(u_{LRS})_1`$ $`=`$ $`{\displaystyle \frac{2a_1^2}{a_1^2+a_2^2}}\mathrm{\Omega }x_2,`$ (126)
$`(u_{LRS})_2`$ $`=`$ $`{\displaystyle \frac{2a_2^2}{a_1^2+a_2^2}}\mathrm{\Omega }x_1,`$ (127)
$`(u_{LRS})_3`$ $`=`$ $`0,`$ (128)
in the corotating frame. Here $`a_i`$ denotes the length of the principal axis parallel to the $`x_i`$-axis. We can rewrite these component of the velocity field as
$`(u_{LRS})_1`$ $`=`$ $`\mathrm{\Omega }x_2+\left({\displaystyle \frac{a_1^2a_2^2}{a_1^2+a_2^2}}\right)\mathrm{\Omega }x_2,`$ (129)
$`(u_{LRS})_2`$ $`=`$ $`\mathrm{\Omega }x_1+\left({\displaystyle \frac{a_1^2a_2^2}{a_1^2+a_2^2}}\right)\mathrm{\Omega }x_1.`$ (130)
The second terms in the above equations are order $`ϵ^3`$ because the deviation between $`a_1`$ and $`a_2`$ is produced by the tidal force and the effect of the tidal force is order $`ϵ^3`$.
On the other hand, if we restrict the form of the function in our velocity field as $`{}_{}{}^{(4)}\varphi _{2}^{}=(b_2/\xi _1)\xi ^2`$, $`{}_{}{}^{(5)}\varphi _{3}^{}=(\alpha b_3/\xi _1)\xi ^3`$ and $`{}_{}{}^{(6)}\varphi _{4}^{}=(\alpha ^2b_4/\xi _1)\xi ^4`$ as in the case of the incompressible fluid, we can express our velocity field by using Eqs. (103), (104) and (95) as
$`(u_{present})_1`$ $`=`$ $`\mathrm{\Omega }x_2+\mathrm{\Omega }\left[6ϵ^3b_2x_248ϵ^4b_3x_1x_2+90ϵ^5b_4x_2(r^25x_1^2)\right],`$ (131)
$`(u_{present})_2`$ $`=`$ $`\mathrm{\Omega }x_1+\mathrm{\Omega }\left[6ϵ^3b_2x_1+6ϵ^4b_3(r^25x_1^2+2x_2^2)+30ϵ^5b_4r(3r^2+7x_1^26x_2^2)\right],`$ (132)
$`(u_{present})_3`$ $`=`$ $`\mathrm{\Omega }\left[12ϵ^4b_3x_2x_3+180ϵ^5b_4x_1x_2x_3\right],`$ (133)
where $`r^2=x_1^2+x_2^2+x_3^2`$. Therefore, we find that the form of the velocity field we obtain in the case of the incompressible fluid coincides with that given by Lai, Rasio and Shapiro up to $`O(ϵ^3)`$ including the value of $`b_2`$. This means that the velocity field of Lai, Rasio and Shapiro is correct only in the case of the “incompressible” equation of state and “ellipsoidal” figures, because the velocity field at order $`ϵ^3`$ is produced by the ellipsoidal deformation of star 1 (see Appendix F).
Finally, we discuss about the configuration of each star. When we pay attention to star 1, the equation for the stellar surface is written as
$`\mathrm{\Xi }(\theta ,\phi )`$ $`=`$ $`\xi _1+ϵ^3S_3(\theta ,\phi )+ϵ^4S_4(\theta ,\phi )+ϵ^5S_5(\theta ,\phi )+ϵ^6S_6(\theta ,\phi ),`$ (134)
$`=`$ $`\xi _1+ϵ^3{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{{}_{}{}^{(4)}\psi _{3}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^5{\displaystyle \frac{{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (138)
$`+ϵ^6{\displaystyle \frac{1}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}[{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}({\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\xi _1}}+{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1))\{{\displaystyle \frac{18}{35}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{2}{7}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{1}{5}}\}`$
$`+{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)+{}_{}{}^{(6)}\psi _{2}^{}(\xi _1)P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{}_{}{}^{(6)}\psi _{22}^{}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$
$`+{}_{}{}^{(6)}\psi _{4}^{}(\xi _1)P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{}_{}{}^{(6)}\psi _{5}^{}(\xi _1)P_5(\mathrm{sin}\theta \mathrm{cos}\phi )],`$
where $`S_i`$ are defined in Appendix A and we have used the relation (C27). Accordingly, we can express the length of the principal axis. Although the real length of the axis is written as $`\alpha \mathrm{\Xi }`$, we show the results divided by $`\alpha `$.
(1) The opposite direction to star 2:
$`\mathrm{\Xi }({\displaystyle \frac{\pi }{2}},0)`$ $`=`$ $`\xi _1+ϵ^3{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}+ϵ^4{\displaystyle \frac{{}_{}{}^{(4)}\psi _{3}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}+ϵ^5{\displaystyle \frac{{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}`$ (141)
$`+ϵ^6{\displaystyle \frac{1}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}[{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}({\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\xi _1}}+{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1))+{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)+{}_{}{}^{(6)}\psi _{2}^{}(\xi _1)+3{}_{}{}^{(6)}\psi _{22}^{}(\xi _1)`$
$`+{}_{}{}^{(6)}\psi _{4}^{}(\xi _1)+{}_{}{}^{(6)}\psi _{5}^{}(\xi _1)],`$
(2) The direction to star 2:
$`\mathrm{\Xi }({\displaystyle \frac{\pi }{2}},\pi )`$ $`=`$ $`\xi _1+ϵ^3{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}ϵ^4{\displaystyle \frac{{}_{}{}^{(4)}\psi _{3}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}+ϵ^5{\displaystyle \frac{{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}`$ (144)
$`+ϵ^6{\displaystyle \frac{1}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}[{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}({\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\xi _1}}+{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1))+{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)+{}_{}{}^{(6)}\psi _{2}^{}(\xi _1)+3{}_{}{}^{(6)}\psi _{22}^{}(\xi _1)`$
$`+{}_{}{}^{(6)}\psi _{4}^{}(\xi _1){}_{}{}^{(6)}\psi _{5}^{}(\xi _1)],`$
(3) The (opposite) direction to the orbital motion:
$`\mathrm{\Xi }({\displaystyle \frac{\pi }{2}},{\displaystyle \frac{\pi }{2}})`$ $`=`$ $`\mathrm{\Xi }({\displaystyle \frac{\pi }{2}},{\displaystyle \frac{3\pi }{2}}),`$ (145)
$`=`$ $`\xi _1ϵ^3{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{2|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}+ϵ^5{\displaystyle \frac{3{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}{8|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}`$ (148)
$`+ϵ^6{\displaystyle \frac{1}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}[{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{4|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}({\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\xi _1}}+{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1))+{}_{}{}^{(6)}\psi _{0}^{}(\xi _1){\displaystyle \frac{1}{2}}{}_{}{}^{(6)}\psi _{2}^{}(\xi _1)3{}_{}{}^{(6)}\psi _{22}^{}(\xi _1)`$
$`+{\displaystyle \frac{3}{8}}{}_{}{}^{(6)}\psi _{4}^{}(\xi _1)],`$
(4) The direction parallel to the rotational axis:
$`\mathrm{\Xi }(0,0)`$ $`=`$ $`\mathrm{\Xi }(\pi ,0),`$ (149)
$`=`$ $`\xi _1ϵ^3{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{2|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}+ϵ^5{\displaystyle \frac{3{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}{8|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}`$ (151)
$`+ϵ^6{\displaystyle \frac{1}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}\left[{\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{4|\mathrm{\Theta }_{0,\xi }(\xi _1)|}}\left({\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\xi _1}}+{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1)\right)+{}_{}{}^{(6)}\psi _{0}^{}(\xi _1){\displaystyle \frac{1}{2}}{}_{}{}^{(6)}\psi _{2}^{}(\xi _1)+{\displaystyle \frac{3}{8}}{}_{}{}^{(6)}\psi _{4}^{}(\xi _1)\right].`$
We can see from these equations and Tables IVIII that the axis to star 2 is the longest, and the deviation between the axis to star 2 and that opposite to star 2 appears at order $`ϵ^4`$. On the contrary, the deviation between the axis to the orbital motion and that parallel to the rotational axis appears at order $`ϵ^6`$, and the difference is the effect of the deformation induced by the spin of the figure ($`{}_{}{}^{(6)}\psi _{22}^{}`$).
When we see the quadrupole moments in Eq. (138), we find that the coeffient of the higher order term is not large as
$$ϵ^3\frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}\left[1+ϵ^3\left\{\frac{2}{7|\mathrm{\Theta }_{0,\xi }(\xi _1)|}\left(\frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\xi _1}+\frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }(\xi _1)\right)+\frac{{}_{}{}^{(6)}\psi _{2}^{}(\xi _1)}{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}\right\}\right]P_2(\mathrm{sin}\theta \mathrm{cos}\phi ).$$
(152)
Therefore we can expect the convergence of these terms. However, there is another quadrupole moment induced by the spin of the figure. The term $`ϵ^6({}_{}{}^{(6)}\psi _{22}^{}(\xi _1)/|\mathrm{\Theta }_{0,\xi }(\xi _1)|)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ seems to be as effective as the leading quadrupole term for the smaller orbital separation. This means that the terms concerned with the spin of the figure which appear at order $`ϵ^9`$ in the total energy may change the behavior of the total energy. Furthermore, the hexadecapole moments in Eq. (138),
$$ϵ^5\frac{{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}\left[1+ϵ\left\{\frac{18}{35{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}\frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{|\mathrm{\Theta }_{0,\xi }(\xi _1)|}\left(\frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\xi _1}+\frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }(\xi _1)\right)+\frac{{}_{}{}^{(6)}\psi _{4}^{}(\xi _1)}{{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)}\right\}\right]P_4(\mathrm{sin}\theta \mathrm{cos}\phi ),$$
(153)
does not seem to converge at order $`ϵ^6`$. However, since these terms will appear at order $`ϵ^{10}`$ in the total energy, they do not have so much effect. Anyway, if we discuss about the behavior of the total energy, the total angular momentum and so on for the smaller orbital separation such as $`R/a_0<3`$, we must calculate at least up to order $`ϵ^9`$.
###### Acknowledgements.
We would like to thank K. Ioka and K. Nakao for useful discussions. This work was partly supported by a Grant-in-Aid for Scientific Research Fellowship (No.9402: KT) and Grant-in-Aid of Scientific Research (No.11640274, 09NP0801: TN) of the Japanese Ministry of Education, Science, Sports and Culture.
## A Summary of the equations and their boundary conditions
In the following appendices, we derive equations order by order of $`ϵ`$, and describe in detail. In this appendix, we summarize the equations which should be solved numerically with their boundary conditions in each order.
First of all, we give the total equations which include all terms up to $`O(ϵ^6)`$. The equations for determination of the velocity potential and stellar configuration are written as
$`\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}`$ $`=`$ $`n\left[\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{𝛀}\times 𝝃)_{orb}{\displaystyle \frac{ϵ}{\xi _1}}(\stackrel{~}{𝛀}\times 𝝃)_{fig}\right]{\displaystyle \frac{\stackrel{~}{}\mathrm{\Theta }}{\mathrm{\Theta }}},`$ (A1)
$`\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }`$ $`=`$ $`\mathrm{\Theta }^n{\displaystyle \frac{\mathrm{\Omega }^2\xi _1^2}{8\pi \rho _cϵ^2}}\stackrel{~}{\mathrm{\Delta }}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }})^22(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }})\left\{(\stackrel{~}{𝛀}\times 𝝃)_{orb}+{\displaystyle \frac{ϵ}{\xi _1}}(\stackrel{~}{𝛀}\times 𝝃)_{fig}\right\}\right],`$ (A2)
where
$`\stackrel{~}{\mathrm{\Delta }}`$ $`=`$ $`\alpha ^2\mathrm{\Delta }={\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{}{\xi }}\left(\xi ^2{\displaystyle \frac{}{\xi }}\right)+{\displaystyle \frac{1}{\xi ^2\mathrm{sin}\theta }}{\displaystyle \frac{}{\theta }}\left(\mathrm{sin}\theta {\displaystyle \frac{}{\theta }}\right)+{\displaystyle \frac{1}{\xi ^2\mathrm{sin}^2\theta }}{\displaystyle \frac{^2}{\phi ^2}},`$ (A3)
$`\stackrel{~}{}`$ $`=`$ $`\alpha ={\displaystyle \frac{}{\xi }}\widehat{𝝃}+{\displaystyle \frac{1}{\xi }}{\displaystyle \frac{}{\theta }}\widehat{𝜽}+{\displaystyle \frac{1}{\xi \mathrm{sin}\theta }}{\displaystyle \frac{}{\phi }}\widehat{𝝋}.`$ (A4)
The boundary condition for the velocity field is
$$[\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}(\stackrel{~}{𝛀}\times 𝝃)_{orb}\frac{ϵ}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}](\stackrel{~}{}\mathrm{\Theta })|_{surf}=0.$$
(A5)
The internal and external gravitational potentials which should be matched at the stellar surface are
$`\stackrel{~}{U}_{int}`$ $``$ $`{\displaystyle \frac{U^{11}}{4\pi \rho _c\alpha ^2}}`$ (A6)
$`=`$ $`\mathrm{\Theta }+{\displaystyle \frac{\mathrm{\Omega }^2\xi _1^2}{8\pi \rho _cϵ^2}}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }})^22(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }})\left\{(\stackrel{~}{𝛀}\times 𝝃)_{orb}+{\displaystyle \frac{ϵ}{\xi _1}}(\stackrel{~}{𝛀}\times 𝝃)_{fig}\right\}\right]+\stackrel{~}{U}_0`$ (A8)
$`{\displaystyle \frac{\mu }{1+p}}ϵ{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(1)^lϵ^l\left({\displaystyle \frac{\xi }{\xi _1}}\right)^lP_l(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{3\mu }{2M_{tot}}}\left({\displaystyle \frac{I_{11}^{}}{a_0^2}}\right)ϵ^3\left[13ϵ{\displaystyle \frac{\xi }{\xi _1}}P_1(\mathrm{sin}\theta \mathrm{cos}\phi )+\mathrm{}\right],`$
$`\stackrel{~}{U}_{ext}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }}+ϵ{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{1:lm}}{\xi ^{l+1}}}Y_l^m+ϵ^2{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{2:lm}}{\xi ^{l+1}}}Y_l^m+ϵ^3{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{3:lm}}{\xi ^{l+1}}}Y_l^m+ϵ^4{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{4:lm}}{\xi ^{l+1}}}Y_l^m+ϵ^5{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{5:lm}}{\xi ^{l+1}}}Y_l^m`$ (A10)
$`+ϵ^6{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{6:lm}}{\xi ^{l+1}}}Y_l^m+\mathrm{},`$
where $`\kappa _0`$ and $`\kappa _{i:lm}`$ are multipole moments. $`\stackrel{~}{U}_0`$ is constant and expanded as
$$\stackrel{~}{U}_0=c_0+ϵc_1+ϵ^2c_2+ϵ^3c_3+ϵ^4c_4+ϵ^5c_5+ϵ^6c_6.$$
(A11)
The velocity potential and configuration function are expanded up to $`O(ϵ^6)`$ as
$`\stackrel{~}{\mathrm{\Phi }}`$ $`=`$ $`\stackrel{~}{\mathrm{\Phi }}_0+ϵ\stackrel{~}{\mathrm{\Phi }}_1+ϵ^2\stackrel{~}{\mathrm{\Phi }}_2+ϵ^3\stackrel{~}{\mathrm{\Phi }}_3+ϵ^4\stackrel{~}{\mathrm{\Phi }}_4+ϵ^5\stackrel{~}{\mathrm{\Phi }}_5+ϵ^6\stackrel{~}{\mathrm{\Phi }}_6,`$ (A12)
$`\mathrm{\Theta }`$ $`=`$ $`\mathrm{\Theta }_0+ϵ\mathrm{\Theta }_1+ϵ^2\mathrm{\Theta }_2+ϵ^3\mathrm{\Theta }_3+ϵ^4\mathrm{\Theta }_4+ϵ^5\mathrm{\Theta }_5+ϵ^6\mathrm{\Theta }_6.`$ (A13)
The boundary conditions for $`\mathrm{\Theta }`$ are
$`\stackrel{~}{U}_{int}(\mathrm{\Xi })`$ $`=`$ $`\stackrel{~}{U}_{ext}(\mathrm{\Xi }),`$ (A14)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\mathrm{\Xi }),`$ (A15)
where $`\mathrm{\Xi }`$ denotes the first zero point of the function $`\mathrm{\Theta }`$, i.e., the stellar surface and is formally expressed as
$$\mathrm{\Xi }(\theta ,\phi )=\xi _1+ϵS_1(\theta ,\phi )+ϵ^2S_2(\theta ,\phi )+ϵ^3S_3(\theta ,\phi )+ϵ^4S_4(\theta ,\phi )+ϵ^5S_5(\theta ,\phi )+ϵ^6S_6(\theta ,\phi ).$$
(A16)
The regularity conditions at the center of each star is
$$\frac{\mathrm{\Theta }}{\xi }(\xi =0)=0,$$
(A17)
and also we normalize the central value of $`\mathrm{\Theta }`$ as
$$\mathrm{\Theta }(\xi =0)=1.$$
(A18)
In the following, we show the equation for the velocity potential, its boundary condition, the equation for determination of the figure and its boundary conditions (A14) and (A15).
### 1 0th order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_0=n\left[\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_0(\stackrel{~}{𝛀}\times 𝝃)_{orb}\right]\frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0},$$
(A19)
and its boundary condition is
$$[\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_0(\stackrel{~}{𝛀}\times 𝝃)_{orb}](\stackrel{~}{}\mathrm{\Theta }_0)|_{\xi _1}=0.$$
(A20)
The equation for determination of the figure is
$$\frac{1}{\xi ^2}\frac{d}{d\xi }\left(\xi ^2\frac{d\mathrm{\Theta }_0}{d\xi }\right)=\mathrm{\Theta }_0^n,$$
(A21)
and its boundary conditions at the stellar surface are
$`c_0={\displaystyle \frac{\kappa _0}{\xi _1}},`$ (A22)
$`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1)={\displaystyle \frac{\kappa _0}{\xi _1^2}}.`$ (A23)
The regularity condition at the center of the star is
$$\frac{d\mathrm{\Theta }_0}{d\xi }(\xi =0)=0,$$
(A24)
and the normalization of $`\mathrm{\Theta }_0`$ at the center of the star is
$$\mathrm{\Theta }_0(\xi =0)=1.$$
(A25)
### 2 1st order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_1=n\left[\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_1\frac{1}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}\right]\frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0},$$
(A26)
and its boundary condition is
$$(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_1)(\stackrel{~}{}\mathrm{\Theta }_0)|_{\xi _1}=0.$$
(A27)
The equation for determination of the figure is
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_1=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_1,$$
(A28)
and its boundary conditions at the stellar surface are
$`c_1{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}}={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_1+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (A29)
$`S_1{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_1}{\xi }}(\xi _1)={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_1{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (A30)
The regularity condition at the center of the star is
$$\frac{\mathrm{\Theta }_1}{\xi }(\xi =0)=0,$$
(A31)
and the normalization of $`\mathrm{\Theta }_1`$ at the center of the star is
$$\mathrm{\Theta }_1(\xi =0)=0,$$
(A32)
because we take $`\mathrm{\Theta }_0(\xi =0)=1`$.
### 3 2nd order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_2=n(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_2)\frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0},$$
(A33)
and its boundary condition is
$$(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_2)(\stackrel{~}{}\mathrm{\Theta }_0)|_{\xi _1}=0.$$
(A34)
The equation for determination of the figure is
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_2=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_2,$$
(A35)
and its boundary conditions at the stellar surface are
$`c_2={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_2+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (A36)
$`S_2{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_2}{\xi }}(\xi _1)={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_2{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (A37)
The regularity condition at the center of the star is
$$\frac{\mathrm{\Theta }_2}{\xi }(\xi =0)=0,$$
(A38)
and the normalization of $`\mathrm{\Theta }_2`$ at the center of the star is
$$\mathrm{\Theta }_2(\xi =0)=0,$$
(A39)
because we take $`\mathrm{\Theta }_0(\xi =0)=1`$.
### 4 3rd order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_3=n(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_3)\frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0},$$
(A40)
and its boundary condition is
$$(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_3)(\stackrel{~}{}\mathrm{\Theta }_0)|_{\xi _1}=0.$$
(A41)
The equation for determination of the figure is
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_3=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_3,$$
(A42)
and its boundary conditions at the stellar surface are
$`c_3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_3+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (A43)
$`S_3{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_3}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_3{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (A44)
The regularity condition at the center of the star is
$$\frac{\mathrm{\Theta }_3}{\xi }(\xi =0)=0,$$
(A45)
and the normalization of $`\mathrm{\Theta }_3`$ at the center of the star is
$$\mathrm{\Theta }_3(\xi =0)=0,$$
(A46)
because we take $`\mathrm{\Theta }_0(\xi =0)=1`$.
### 5 4th order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_4=\frac{n}{\mathrm{\Theta }_0}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_4)(\stackrel{~}{}\mathrm{\Theta }_0)\frac{1}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_3)\right],$$
(A47)
and its boundary condition is
$$[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_4)(\stackrel{~}{}\mathrm{\Theta }_0)\frac{1}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_3)]|_{\xi _1}=0.$$
(A48)
The equation for determination of the figure is
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_4=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_4,$$
(A49)
and its boundary conditions at the stellar surface are
$`c_4+{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_4+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (A50)
$`S_4{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_4}{\xi }}(\xi _1)+{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_4{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (A51)
The regularity condition at the center of the star is
$$\frac{\mathrm{\Theta }_4}{\xi }(\xi =0)=0,$$
(A52)
and the normalization of $`\mathrm{\Theta }_4`$ at the center of the star is
$$\mathrm{\Theta }_4(\xi =0)=0,$$
(A53)
because we take $`\mathrm{\Theta }_0(\xi =0)=1`$.
### 6 5th order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_5=\frac{n}{\mathrm{\Theta }_0}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_5)(\stackrel{~}{}\mathrm{\Theta }_0)\frac{1}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_4)\right],$$
(A54)
and its boundary condition is
$$[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_5)(\stackrel{~}{}\mathrm{\Theta }_0)\frac{1}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_4)]|_{\xi _1}=0.$$
(A55)
The equation for determination of the figure is
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_5=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_5,$$
(A56)
and its boundary conditions at the stellar surface are
$`c_5{\displaystyle \frac{\mu }{1+p}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_5+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (A57)
$`S_5{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_5}{\xi }}(\xi _1){\displaystyle \frac{4\mu }{\xi _1(1+p)}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_5{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (A58)
The regularity condition at the center of the star is
$$\frac{\mathrm{\Theta }_5}{\xi }(\xi =0)=0,$$
(A59)
and the normalization of $`\mathrm{\Theta }_5`$ at the center of the star is
$$\mathrm{\Theta }_5(\xi =0)=0,$$
(A60)
because we take $`\mathrm{\Theta }_0(\xi =0)=1`$.
### 7 6th order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_6=\frac{n}{\mathrm{\Theta }_0}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_6)(\stackrel{~}{}\mathrm{\Theta }_0)\frac{1}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_5)\right],$$
(A61)
and its boundary condition is
$$[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_6)(\stackrel{~}{}\mathrm{\Theta }_0)\frac{1}{\xi _1}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_5)]|_{\xi _1}=0.$$
(A62)
The equation for determination of the figure is
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_6=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_6\frac{1}{2}n(n1)\mathrm{\Theta }_0^{n2}\mathrm{\Theta }_3^2+\frac{\mu }{\xi _1}\stackrel{~}{\mathrm{\Delta }}\frac{\stackrel{~}{\mathrm{\Phi }}_4}{\phi },$$
(A63)
and its boundary conditions at the stellar surface are
$`c_6{\displaystyle \frac{3\mu \overline{I}_{11}^{}}{2M_{tot}a_0^2}}{\displaystyle \frac{\mu \delta }{2(1+p)^2}}{\displaystyle \frac{2\mu }{\xi _1(1+p)}}S_3P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{\mu }{1+p}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{\mu }{\xi _1}}{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_4}{\phi }}(\xi _1)`$ (A64)
$`={\displaystyle \frac{\kappa _0}{\xi _1^2}}\left({\displaystyle \frac{S_3^2}{\xi _1}}S_6\right){\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+2}}}S_3Y_l^m+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (A65)
$`S_6{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{1}{2}}S_3^2{\displaystyle \frac{d^3\mathrm{\Theta }_0}{d\xi ^3}}(\xi _1)+S_3{\displaystyle \frac{^2\mathrm{\Theta }_3}{\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_6}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}S_3P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{5\mu }{\xi _1(1+p)}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (A66)
$`{\displaystyle \frac{\mu }{\xi _1}}{\displaystyle \frac{^2\stackrel{~}{\mathrm{\Phi }}_4}{\xi \phi }}(\xi _1)={\displaystyle \frac{\kappa _0}{\xi _1^3}}\left({\displaystyle \frac{3S_3^2}{\xi _1}}2S_6\right)+{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(l+1)(l+2){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+3}}}S_3Y_l^m{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}(l+1){\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+2}}}Y_l^m,`$ (A67)
where
$$\delta \frac{9}{2a_0^2}\left[\frac{\overline{I}_{11}}{M_1}+\left(\frac{a_0^{}}{a_0}\right)^3\frac{\overline{I}_{11}^{}}{M_2}\right].$$
(A68)
The regularity condition at the center of the star is
$$\frac{\mathrm{\Theta }_6}{\xi }(\xi =0)=0,$$
(A69)
and the normalization of $`\mathrm{\Theta }_6`$ at the center of the star is
$$\mathrm{\Theta }_6(\xi =0)=0,$$
(A70)
because we take $`\mathrm{\Theta }_0(\xi =0)=1`$.
## B Solutions at zeroth order
Since the configuration of star 1 is spherical at 0th order, the equation for determination of the figure becomes Lane-Emden equation,
$$\frac{1}{\xi ^2}\frac{d}{d\xi }\left(\xi ^2\frac{d\mathrm{\Theta }_0}{d\xi }\right)=\mathrm{\Theta }_0^n.$$
(B1)
For example, we obtain analytic solutions,
$`\mathrm{\Theta }_0`$ $`=`$ $`1{\displaystyle \frac{\xi ^2}{6}}(\mathrm{for}n=0),`$ (B2)
$`\mathrm{\Theta }_0`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\xi }{\xi }}(\mathrm{for}n=1),`$ (B3)
or numerical solutions by solving the ordinary differential equation (B1).
The velocity field at 0th order is the same as the orbital motion so that the velocity potential is determined by
$$\mathrm{\Phi }_0=\mathrm{\Omega }R(\stackrel{~}{𝛀}\times 𝝃)_{orb}.$$
(B4)
Then, the solution should be
$$\mathrm{\Phi }_0=\frac{\mathrm{\Omega }R}{1+p}r\mathrm{sin}\theta \mathrm{sin}\phi .$$
(B5)
From this equation, we find that it is convenient to normalize $`\mathrm{\Phi }_0`$ as
$$\stackrel{~}{\mathrm{\Phi }}_0=\frac{1}{1+p}\xi \mathrm{sin}\theta \mathrm{sin}\phi .$$
(B6)
In the following, we use the normalized velocity potential defined by
$$\stackrel{~}{\mathrm{\Phi }}=\frac{ϵ\mathrm{\Phi }}{\mathrm{\Omega }\alpha a_0}.$$
(B7)
Here we show the useful relations:
$`(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_0)^2={\displaystyle \frac{1}{(1+p)^2}},`$ (B8)
$`(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_0)(\stackrel{~}{𝛀}\times 𝝃)_{orb}={\displaystyle \frac{1}{(1+p)^2}},`$ (B9)
$`(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_0)(\stackrel{~}{𝛀}\times 𝝃)_{fig}={\displaystyle \frac{1}{1+p}}\xi \mathrm{sin}\theta \mathrm{cos}\phi .`$ (B10)
The orbital angular velocity can be calculated from Eq. (31) as
$$0=_{star1}d^3x\rho \left(\frac{M_2}{R^2}+\frac{\mathrm{\Omega }_0^2R}{1+p}\right).$$
(B11)
Therefore, we have
$$\mathrm{\Omega }_0^2=\frac{M_{tot}}{R^3},$$
(B12)
where $`M_{tot}=M_1+M_2`$.
For 0th order, we have the internal and external gravitational potentials and their derivatives as
$`\stackrel{~}{U}_{int}(\xi )`$ $`=`$ $`\mathrm{\Theta }_0+c_0,`$ (B13)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\xi )`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }},`$ (B14)
$`\stackrel{~}{U}_{ext}(\xi )`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }},`$ (B15)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\xi )`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi ^2}}.`$ (B16)
They become at the stellar surface as
$`\stackrel{~}{U}_{int}(\xi _1)`$ $`=`$ $`c_0,`$ (B17)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\xi _1)`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1),`$ (B18)
$`\stackrel{~}{U}_{ext}(\xi _1)`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1}},`$ (B19)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\xi _1)`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1^2}}.`$ (B20)
Therefore, we can decide two constants at 0th order from the boundary conditions at the stellar surface as
$`\kappa _0`$ $`=`$ $`\xi _1^2{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1),`$ (B21)
$`c_0`$ $`=`$ $`\xi _1{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1).`$ (B22)
## C Solutions at first order
The equation for the velocity potential is
$`\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_1`$ $`=`$ $`n\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_1){\displaystyle \frac{1}{\xi _1}}(\stackrel{~}{𝛀}\times 𝝃)_{fig}\right]{\displaystyle \frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0}},`$ (C1)
$`=`$ $`n(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_1){\displaystyle \frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0}},`$ (C2)
where we have used the fact that $`\stackrel{~}{}\mathrm{\Theta }_0`$ has only $`\xi `$ component and $`(\stackrel{~}{𝛀}\times 𝝃)_{fig}`$ has only $`\phi `$ component. The boundary condition is
$$(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_1)(\stackrel{~}{}\mathrm{\Theta }_0)|_{\xi _1}=0.$$
(C3)
When we expand $`\stackrel{~}{\mathrm{\Phi }}_1`$ as
$$\stackrel{~}{\mathrm{\Phi }}_1=\underset{l,m}{}{}_{}{}^{(1)}\varphi _{lm}^{}(\xi )\stackrel{~}{Y}_l^m(\theta ,\phi ),$$
(C4)
the boundary conditions are that $`{}_{}{}^{(1)}\varphi _{lm}^{}`$ is regular at $`\xi =0`$ and
$$\frac{d{}_{}{}^{(1)}\varphi _{lm}^{}}{d\xi }(\xi _1)=0.$$
(C5)
We cannot determine the absolute value of $`\stackrel{~}{\mathrm{\Phi }}_1`$ from the equation and the boundary condition. Therefore, $`\stackrel{~}{\mathrm{\Phi }}_1`$ should be zero. Note that $`\stackrel{~}{\mathrm{\Phi }}_1=\mathrm{constant}`$ is a solution. However, since the physical value is a velocity field $`\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_1`$, the constant solution is meaningless.
Next, we consider the deformation of the figure at $`O(ϵ)`$. The equation is written as
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_1=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_1.$$
(C6)
We expand $`\mathrm{\Theta }_1`$ as
$$\mathrm{\Theta }_1=\underset{l,m}{}{}_{}{}^{(1)}\psi _{lm}^{}(\xi )Y_l^m(\theta ,\phi ).$$
(C7)
The surface equation up to $`O(ϵ)`$ is written as
$$\mathrm{\Xi }=\xi _1+ϵS_1(\theta ,\phi ),$$
(C8)
where
$$S_1=\frac{\mathrm{\Theta }_1(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}.$$
(C9)
The gravitational potentials and their derivatives up to $`O(ϵ)`$ are
$`\stackrel{~}{U}_{int}`$ $`=`$ $`\mathrm{\Theta }ϵ{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}}+c_0+ϵc_1,`$ (C10)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }}{\xi }},`$ (C11)
$`\stackrel{~}{U}_{ext}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }}+ϵ{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{1:lm}}{\xi ^{l+1}}}Y_l^m,`$ (C12)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi ^2}}ϵ{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{1:lm}}{\xi ^{l+2}}}Y_l^m.`$ (C13)
They become at the stellar surface as
$`\stackrel{~}{U}_{int}(\mathrm{\Xi })`$ $`=`$ $`c_0+ϵ\left[c_1{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}}\right],`$ (C14)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1)+ϵ\left[S_1{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_1}{\xi }}(\xi _1)\right],`$ (C15)
$`\stackrel{~}{U}_{ext}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1}}\left(1ϵ{\displaystyle \frac{S_1}{\xi _1}}\right)+ϵ{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (C16)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1^2}}\left(1ϵ{\displaystyle \frac{2S_1}{\xi _1}}\right)ϵ{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (C17)
Then, the boundary conditions for $`O(ϵ)`$ become
$`c_1{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}}={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_1+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (C18)
$`S_1{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_1}{\xi }}(\xi _1)={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_1{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+2}}}Y_l^m,`$ (C19)
and also $`\mathrm{\Theta }_1`$ should be satisfied with the conditions at the center of the star,
$`{\displaystyle \frac{\mathrm{\Theta }_1}{\xi }}(0)`$ $`=`$ $`0,`$ (C20)
$`\mathrm{\Theta }_1(0)`$ $`=`$ $`0.`$ (C21)
### 1 The case of $`n=0`$
In this case, since we can express $`{}_{}{}^{(1)}\psi _{lm}^{}={}_{}{}^{(1)}A_{lm}^{}\xi ^l`$ by using a constant $`{}_{}{}^{(1)}A_{lm}^{}`$, the boundary conditions can be written as
$`c_1{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}}={\displaystyle \underset{l,m}{}}\left({}_{}{}^{(1)}A_{lm}^{}\xi _1^l+{\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+1}}}\right)Y_l^m,`$ (C22)
$`0={\displaystyle \underset{l,m}{}}\left[(3l){}_{}{}^{(1)}A_{lm}^{}\xi _1^{l1}(l+1){\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (C23)
Therefore, we obtain
$`c_1`$ $`=`$ $`{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}},`$ (C24)
$`\kappa _{1:lm}`$ $`=`$ $`0,`$ (C25)
$`\mathrm{\Theta }_1`$ $`=`$ $`0.`$ (C26)
### 2 The case of $`n>0`$
In this case, we have the relation;
$$\frac{d^2\mathrm{\Theta }_0}{d\xi ^2}(\xi _1)=\frac{2}{\xi _1}\frac{d\mathrm{\Theta }_0}{d\xi }(\xi _1),$$
(C27)
from the Lane-Emden equation at the stellar surface. Then, the boundary conditions can be written as
$`c_1{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}}={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(1)}\psi _{lm}^{}(\xi _1)+{\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (C28)
$`0={\displaystyle \underset{l,m}{}}\left[{\displaystyle \frac{d{}_{}{}^{(1)}\psi _{lm}^{}}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{1:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (C29)
Therefore, we obtain
$`c_1`$ $`=`$ $`{\displaystyle \frac{\mu (3+2p)}{2(1+p)^2}},`$ (C30)
$`\kappa _{1:lm}`$ $`=`$ $`0,`$ (C31)
$`\mathrm{\Theta }_1`$ $`=`$ $`0.`$ (C32)
## D Solutions at second order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_2=n(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_2)\frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0}.$$
(D1)
The boundary condition is
$$(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_2)(\stackrel{~}{}\mathrm{\Theta }_0)|_{\xi _1}=0.$$
(D2)
Then, $`\stackrel{~}{\mathrm{\Phi }}_2`$ should be zero as in the case of $`\stackrel{~}{\mathrm{\Phi }}_1`$.
Next, we consider the deformation of the figure at $`O(ϵ^2)`$. The equation is written as
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_2=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_2.$$
(D3)
We expand $`\mathrm{\Theta }_2`$ as
$$\mathrm{\Theta }_2=\underset{l,m}{}{}_{}{}^{(2)}\psi _{lm}^{}(\xi )Y_l^m(\theta ,\phi ).$$
(D4)
The surface equation up to $`O(ϵ^2)`$ is written as
$$\mathrm{\Xi }=\xi _1+ϵ^2S_2(\theta ,\phi ),$$
(D5)
where
$$S_2=\frac{\mathrm{\Theta }_2(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}.$$
(D6)
The gravitational potentials and their derivatives up to $`O(ϵ^2)`$ are
$`\stackrel{~}{U}_{int}`$ $`=`$ $`\mathrm{\Theta }+c_0+ϵ^2c_2,`$ (D7)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }}{\xi }},`$ (D8)
$`\stackrel{~}{U}_{ext}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }}+ϵ^2{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{2:lm}}{\xi ^{l+1}}}Y_l^m,`$ (D9)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi ^2}}ϵ^2{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{2:lm}}{\xi ^{l+2}}}Y_l^m.`$ (D10)
They become at the stellar surface as
$`\stackrel{~}{U}_{int}(\mathrm{\Xi })`$ $`=`$ $`c_0+ϵ^2c_2,`$ (D11)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1)+ϵ^2\left[S_2{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_2}{\xi }}(\xi _1)\right],`$ (D12)
$`\stackrel{~}{U}_{ext}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1}}\left(1ϵ^2{\displaystyle \frac{S_2}{\xi _1}}\right)+ϵ^2{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (D13)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1^2}}\left(1ϵ^2{\displaystyle \frac{2S_2}{\xi _1}}\right)ϵ^2{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (D14)
Then, the boundary conditions for $`O(ϵ^2)`$ become
$`c_2={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_2+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (D15)
$`S_2{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_2}{\xi }}(\xi _1)={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_2{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+2}}}Y_l^m,`$ (D16)
and also $`\mathrm{\Theta }_2`$ should be satisfied with the conditions at the center of the star,
$`{\displaystyle \frac{\mathrm{\Theta }_2}{\xi }}(0)`$ $`=`$ $`0,`$ (D17)
$`\mathrm{\Theta }_2(0)`$ $`=`$ $`0.`$ (D18)
### 1 The case of $`n=0`$
In this case, since we can express $`{}_{}{}^{(2)}\psi _{lm}^{}={}_{}{}^{(2)}A_{lm}^{}\xi ^l`$ by using a constant $`{}_{}{}^{(2)}A_{lm}^{}`$, the boundary conditions can be written as
$`c_2={\displaystyle \underset{l,m}{}}\left({}_{}{}^{(2)}A_{lm}^{}\xi _1^l+{\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+1}}}\right)Y_l^m,`$ (D19)
$`0={\displaystyle \underset{l,m}{}}\left[(3l){}_{}{}^{(2)}A_{lm}^{}\xi _1^{l1}(l+1){\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (D20)
Therefore, we obtain
$`c_2`$ $`=`$ $`0,`$ (D21)
$`\kappa _{2:lm}`$ $`=`$ $`0,`$ (D22)
$`\mathrm{\Theta }_2`$ $`=`$ $`0.`$ (D23)
### 2 The case of $`n>0`$
In this case, the boundary conditions can be written as
$`c_2={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(2)}\psi _{lm}^{}(\xi _1)+{\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (D24)
$`0={\displaystyle \underset{l,m}{}}\left[{\displaystyle \frac{d{}_{}{}^{(2)}\psi _{lm}^{}}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{2:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (D25)
Therefore, we obtain
$`c_2`$ $`=`$ $`0,`$ (D26)
$`\kappa _{2:lm}`$ $`=`$ $`0,`$ (D27)
$`\mathrm{\Theta }_2`$ $`=`$ $`0.`$ (D28)
## E Solutions at third order
The equation for the velocity potential is
$$\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_3=n(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_3)\frac{\stackrel{~}{}\mathrm{\Theta }_0}{\mathrm{\Theta }_0}.$$
(E1)
The boundary condition is
$$(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_3)(\stackrel{~}{}\mathrm{\Theta }_0)|_{\xi _1}=0.$$
(E2)
Then, $`\stackrel{~}{\mathrm{\Phi }}_3`$ should be zero as in the cases of $`\stackrel{~}{\mathrm{\Phi }}_1`$ and $`\stackrel{~}{\mathrm{\Phi }}_2`$.
Next, we consider the deformation of the figure at $`O(ϵ^3)`$. The equation is written as
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_3=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_3.$$
(E3)
We expand $`\mathrm{\Theta }_3`$ as
$$\mathrm{\Theta }_3=\underset{l,m}{}{}_{}{}^{(3)}\psi _{lm}^{}(\xi )Y_l^m(\theta ,\phi ).$$
(E4)
The surface equation up to $`O(ϵ^3)`$ is written as
$$\mathrm{\Xi }=\xi _1+ϵ^3S_3(\theta ,\phi ),$$
(E5)
where
$$S_3=\frac{\mathrm{\Theta }_3(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}.$$
(E6)
The gravitational potentials and their derivatives up to $`O(ϵ^3)`$ are
$`\stackrel{~}{U}_{int}`$ $`=`$ $`\mathrm{\Theta }+c_0+ϵ^3\left[c_3{\displaystyle \frac{\mu }{\xi _1^2(1+p)}}\xi ^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$ (E7)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }}{\xi }}ϵ^3{\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}\xi P_2(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (E8)
$`\stackrel{~}{U}_{ext}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }}+ϵ^3{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{3:lm}}{\xi ^{l+1}}}Y_l^m,`$ (E9)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi ^2}}ϵ^3{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi ^{l+2}}}Y_l^m.`$ (E10)
They become at the stellar surface as
$`\stackrel{~}{U}_{int}(\mathrm{\Xi })`$ $`=`$ $`c_0+ϵ^3\left[c_3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$ (E11)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1)+ϵ^3\left[S_3{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_3}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$ (E12)
$`\stackrel{~}{U}_{ext}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1}}\left(1ϵ^3{\displaystyle \frac{S_3}{\xi _1}}\right)+ϵ^3{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (E13)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1^2}}\left(1ϵ^3{\displaystyle \frac{2S_3}{\xi _1}}\right)ϵ^3{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (E14)
Then, the boundary conditions for $`O(ϵ^3)`$ become
$`c_3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_3+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (E15)
$`S_3{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_3}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_3{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+2}}}Y_l^m,`$ (E16)
and also $`\mathrm{\Theta }_3`$ should be satisfied with the conditions at the center of the star,
$`{\displaystyle \frac{\mathrm{\Theta }_3}{\xi }}(0)`$ $`=`$ $`0,`$ (E17)
$`\mathrm{\Theta }_3(0)`$ $`=`$ $`0.`$ (E18)
### 1 The case of $`n=0`$
In this case, since we can express $`{}_{}{}^{(3)}\psi _{lm}^{}={}_{}{}^{(3)}A_{lm}^{}\xi ^l`$ by using a constant $`{}_{}{}^{(3)}A_{lm}^{}`$, the boundary conditions can be written as
$`c_3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left({}_{}{}^{(3)}A_{lm}^{}\xi _1^l+{\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+1}}}\right)Y_l^m,`$ (E19)
$`{\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[(3l){}_{}{}^{(3)}A_{lm}^{}\xi _1^{l1}(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (E20)
Therefore, we obtain
$`c_3`$ $`=`$ $`0,`$ (E21)
$`\mathrm{\Theta }_3`$ $`=`$ $`{}_{}{}^{(3)}A_{2}^{}\xi ^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (E22)
where
$`{}_{}{}^{(3)}A_{2}^{}`$ $`=`$ $`{\displaystyle \frac{5\mu }{2(1+p)\xi _1^2}},`$ (E23)
$`\kappa _{3:2}`$ $`=`$ $`{\displaystyle \frac{3\mu \xi _1^3}{2(1+p)}}.`$ (E24)
### 2 The case of $`n=1`$
In this case, since we can express $`{}_{}{}^{(3)}\psi _{lm}^{}={}_{}{}^{(3)}B_{lm}^{}j_l`$ by using a constant $`{}_{}{}^{(3)}B_{lm}^{}`$ and a spherical Bessel function $`j_l`$, the boundary conditions can be written as
$`c_3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(3)}B_{lm}^{}j_l(\xi _1)+{\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (E25)
$`{\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(3)}B_{lm}^{}{\displaystyle \frac{dj_l}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (E26)
Therefore, we obtain
$`c_3`$ $`=`$ $`0,`$ (E27)
$`\mathrm{\Theta }_3`$ $`=`$ $`{}_{}{}^{(3)}B_{2}^{}j_2(\xi )P_2(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (E28)
where
$`{}_{}{}^{(3)}B_{2}^{}`$ $`=`$ $`{\displaystyle \frac{5\mu }{1+p}},`$ (E29)
$`\kappa _{3:2}`$ $`=`$ $`{\displaystyle \frac{\mu \xi _1}{1+p}}(15\xi _1^2),`$ (E30)
$`j_2(\xi )`$ $`=`$ $`{\displaystyle \frac{1}{\xi ^3}}\left[(3\xi ^2)\mathrm{sin}\xi 3\xi \mathrm{cos}\xi \right].`$ (E31)
### 3 The case of $`n>0`$
In this case, the boundary conditions can be written as
$`c_3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(3)}\psi _{lm}^{}(\xi _1)+{\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (E32)
$`{\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{lm}^{}}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{3:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (E33)
Therefore,
$`c_3`$ $`=`$ $`0,`$ (E34)
$`\mathrm{\Theta }_3`$ $`=`$ $`{}_{}{}^{(3)}\psi _{2}^{}(\xi )P_2(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (E35)
and the ordinary differential equation for $`{}_{}{}^{(3)}\psi _{2}^{}`$ is written as
$$\left[\frac{1}{\xi ^2}\frac{d}{d\xi }\left(\xi ^2\frac{d}{d\xi }\right)\frac{6}{\xi ^2}\right]{}_{}{}^{(3)}\psi _{2}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(3)}\psi _{2}^{}.$$
(E36)
The coefficient of $`{}_{}{}^{(3)}\psi _{2}^{}`$ and the constant $`\kappa _{3:2}`$ are determined from two equations;
$`3{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1)={\displaystyle \frac{5\mu }{1+p}},`$ (E37)
$`\kappa _{3:2}=\xi _1^3\left[{}_{}{}^{(3)}\psi _{2}^{}(\xi _1){\displaystyle \frac{\mu }{1+p}}\right]={\displaystyle \frac{\xi _1^3}{3}}\left[{\displaystyle \frac{2\mu }{1+p}}\xi _1{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1)\right].`$ (E38)
### 4 Reduced quadrupole moments
At this order, we obtain the reduced quadrupole moments as
$`I_{11}`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle _{star1}}d^3x\rho r^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (E39)
$`=`$ $`\{\begin{array}{cc}ϵ^3{\displaystyle \frac{8\pi }{5}}\rho _c\alpha ^5\xi _1^5{}_{}{}^{(3)}A_{2}^{},(n=0)\hfill & \\ ϵ^3{\displaystyle \frac{8\pi }{15}}\rho _c\alpha ^5\xi _1^3\left({\displaystyle \frac{15}{\xi _1^2}}1\right){}_{}{}^{(3)}B_{2}^{},(n=1)\hfill & \\ ϵ^3{\displaystyle \frac{8\pi }{15}}\rho _c\alpha ^5{\displaystyle _0^{\xi _1}}d\xi n\mathrm{\Theta }_0^{n1}(\xi ){}_{}{}^{(3)}\psi _{2}^{}(\xi )\xi ^4.(n>0)\hfill & \end{array}`$ (E43)
## F Solutions at fourth order
The equation for the velocity potential is
$`\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_4`$ $`=`$ $`{\displaystyle \frac{n}{\mathrm{\Theta }_0}}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_4)(\stackrel{~}{}\mathrm{\Theta }_0){\displaystyle \frac{1}{\xi _1}}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_3)\right],`$ (F1)
$`=`$ $`{\displaystyle \frac{n}{\mathrm{\Theta }_0}}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_4)(\stackrel{~}{}\mathrm{\Theta }_0)+{\displaystyle \frac{1}{2\xi _1}}{}_{}{}^{(3)}\psi _{2}^{}P_2^2(\mathrm{cos}\theta )\mathrm{sin}2\phi \right].`$ (F2)
Therefore, we have the form of $`\stackrel{~}{\mathrm{\Phi }}_4`$ and the equation for it as
$`\stackrel{~}{\mathrm{\Phi }}_4={}_{}{}^{(4)}\varphi _{2}^{}(\xi )P_2^2(\mathrm{cos}\theta )\mathrm{sin}2\phi ,`$ (F3)
$`\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{6}{\xi ^2}}+{\displaystyle \frac{n}{\mathrm{\Theta }_0}}{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}{\displaystyle \frac{d}{d\xi }}\right]{}_{}{}^{(4)}\varphi _{2}^{}={\displaystyle \frac{n}{2\xi _1\mathrm{\Theta }_0}}{}_{}{}^{(3)}\psi _{2}^{}.`$ (F4)
The boundary condition is
$$\frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }(\xi _1)\frac{d\mathrm{\Theta }_0}{d\xi }(\xi _1)+\frac{1}{2\xi _1}{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)=0.$$
(F5)
Next, we consider the deformation of the figure at $`O(ϵ^4)`$. The equation is written as
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_4=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_4.$$
(F6)
We expand $`\mathrm{\Theta }_4`$ as
$$\mathrm{\Theta }_4=\underset{l,m}{}{}_{}{}^{(4)}\psi _{lm}^{}(\xi )Y_l^m(\theta ,\phi ).$$
(F7)
The surface equation up to $`O(ϵ^4)`$ is written as
$$\mathrm{\Xi }=\xi _1+ϵ^3S_3(\theta ,\phi )+ϵ^4S_4(\theta ,\phi ),$$
(F8)
where
$$S_4=\frac{\mathrm{\Theta }_4(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}.$$
(F9)
The gravitational potentials and their derivatives up to $`O(ϵ^4)`$ are
$`\stackrel{~}{U}_{int}`$ $`=`$ $`\mathrm{\Theta }+c_0ϵ^3{\displaystyle \frac{\mu }{\xi _1^2(1+p)}}\xi ^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4\left[c_4+{\displaystyle \frac{\mu }{\xi _1^3(1+p)}}\xi ^3P_3(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$ (F10)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }}{\xi }}ϵ^3{\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}\xi P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{3\mu }{\xi _1^3(1+p)}}\xi ^2P_3(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (F11)
$`\stackrel{~}{U}_{ext}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }}+ϵ^3{\displaystyle \frac{\kappa _{3:2}}{\xi ^3}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{4:lm}}{\xi ^{l+1}}}Y_l^m,`$ (F12)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi ^2}}ϵ^3{\displaystyle \frac{3\kappa _{3:2}}{\xi ^4}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^4{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{4:lm}}{\xi ^{l+2}}}Y_l^m.`$ (F13)
They become at the stellar surface as
$`\stackrel{~}{U}_{int}(\mathrm{\Xi })`$ $`=`$ $`c_0ϵ^3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4\left[c_4+{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$ (F14)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1)+ϵ^3\left[S_3{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_3}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )\right]`$ (F16)
$`+ϵ^4\left[S_4{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_4}{\xi }}(\xi _1)+{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$
$`\stackrel{~}{U}_{ext}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1}}\left(1ϵ^3{\displaystyle \frac{S_3}{\xi _1}}ϵ^4{\displaystyle \frac{S_4}{\xi _1}}\right)+ϵ^3{\displaystyle \frac{\kappa _{3:2}}{\xi _1^3}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (F17)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1^2}}\left(1ϵ^3{\displaystyle \frac{2S_3}{\xi _1}}ϵ^4{\displaystyle \frac{2S_4}{\xi _1}}\right)ϵ^3{\displaystyle \frac{3\kappa _{3:2}}{\xi _1^4}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^4{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+2}}}Y_l^m.`$ (F18)
Then, the boundary conditions for $`O(ϵ^4)`$ become
$`c_4+{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_4+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (F19)
$`S_4{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_4}{\xi }}(\xi _1)+{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_4{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+2}}}Y_l^m,`$ (F20)
and also $`\mathrm{\Theta }_4`$ should be satisfied with the conditions at the center of the star,
$`{\displaystyle \frac{\mathrm{\Theta }_4}{\xi }}(0)`$ $`=`$ $`0,`$ (F21)
$`\mathrm{\Theta }_4(0)`$ $`=`$ $`0.`$ (F22)
### 1 The case of $`n=0`$
In this case, since we can express $`{}_{}{}^{(4)}\psi _{lm}^{}={}_{}{}^{(4)}A_{lm}^{}\xi ^l`$ by using a constant $`{}_{}{}^{(4)}A_{lm}^{}`$, the boundary conditions can be written as
$`c_4+{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left({}_{}{}^{(4)}A_{lm}^{}\xi _1^l+{\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+1}}}\right)Y_l^m,`$ (F23)
$`{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[(3l){}_{}{}^{(4)}A_{lm}^{}\xi _1^{l1}(l+1){\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (F24)
Therefore, we obtain
$`c_4`$ $`=`$ $`0,`$ (F25)
$`\mathrm{\Theta }_4`$ $`=`$ $`{}_{}{}^{(4)}A_{3}^{}\xi ^3P_3(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (F26)
where
$`{}_{}{}^{(4)}A_{3}^{}`$ $`=`$ $`{\displaystyle \frac{7\mu }{4(1+p)\xi _1^3}},`$ (F27)
$`\kappa _{4:3}`$ $`=`$ $`{\displaystyle \frac{3\mu \xi _1^4}{4(1+p)}}.`$ (F28)
### 2 The case of $`n=1`$
In this case, since we can express $`{}_{}{}^{(4)}\psi _{lm}^{}={}_{}{}^{(4)}B_{lm}^{}j_l`$ by using a constant $`{}_{}{}^{(4)}B_{lm}^{}`$, the boundary conditions can be written as
$`c_4+{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(4)}B_{lm}^{}j_l(\xi _1)+{\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (F29)
$`{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(4)}B_{lm}^{}{\displaystyle \frac{dj_l}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (F30)
Therefore, we obtain
$`c_4`$ $`=`$ $`0,`$ (F31)
$`\mathrm{\Theta }_4`$ $`=`$ $`{}_{}{}^{(4)}B_{3}^{}j_3(\xi )P_3(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (F32)
where
$`{}_{}{}^{(4)}B_{3}^{}`$ $`=`$ $`{\displaystyle \frac{7\mu \xi _1}{3(1+p)}},`$ (F33)
$`\kappa _{4:3}`$ $`=`$ $`{\displaystyle \frac{5\mu \xi _1^2}{3(1+p)}}(212\xi _1^2),`$ (F34)
$`j_3(\xi )`$ $`=`$ $`{\displaystyle \frac{1}{\xi ^4}}\left[(156\xi ^2)\mathrm{sin}\xi \xi (15\xi ^2)\mathrm{cos}\xi \right].`$ (F35)
### 3 The case of $`n>0`$
In this case, the boundary conditions can be written as
$`c_4+{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(4)}\psi _{lm}^{}(\xi _1)+{\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (F36)
$`{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{\displaystyle \frac{d{}_{}{}^{(4)}\psi _{lm}^{}}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{4:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (F37)
Therefore, we obtain
$`c_4`$ $`=`$ $`0,`$ (F38)
$`\mathrm{\Theta }_4`$ $`=`$ $`{}_{}{}^{(4)}\psi _{3}^{}(\xi )P_3(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (F39)
and the ordinary differential equation for $`{}_{}{}^{(4)}\psi _{3}^{}`$ is written as
$$\left[\frac{1}{\xi ^2}\frac{d}{d\xi }\left(\xi ^2\frac{d}{d\xi }\right)\frac{12}{\xi ^2}\right]{}_{}{}^{(4)}\psi _{3}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(4)}\psi _{3}^{}.$$
(F40)
The coefficient of $`{}_{}{}^{(4)}\psi _{3}^{}`$ and the constant $`\kappa _{4:3}`$ are determined from two equations;
$`4{}_{}{}^{(4)}\psi _{3}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(4)}\psi _{3}^{}}{d\xi }}(\xi _1)={\displaystyle \frac{7\mu }{1+p}},`$ (F41)
$`\kappa _{4:3}=\xi _1^4\left[{}_{}{}^{(4)}\psi _{3}^{}(\xi _1)+{\displaystyle \frac{\mu }{1+p}}\right]={\displaystyle \frac{\xi _1^4}{4}}\left[{\displaystyle \frac{3\mu }{1+p}}+\xi _1{\displaystyle \frac{d{}_{}{}^{(4)}\psi _{3}^{}}{d\xi }}(\xi _1)\right].`$ (F42)
### 4 Octupole moments
At this order, we can obtain the octupole moments, but we do not represent them here.
## G Solutions at fifth order
The equation for the velocity potential is
$`\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_5`$ $`=`$ $`{\displaystyle \frac{n}{\mathrm{\Theta }_0}}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_5)(\stackrel{~}{}\mathrm{\Theta }_0){\displaystyle \frac{1}{\xi _1}}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_4)\right],`$ (G1)
$`=`$ $`{\displaystyle \frac{n}{\mathrm{\Theta }_0}}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_5)(\stackrel{~}{}\mathrm{\Theta }_0){\displaystyle \frac{1}{4\xi _1}}{}_{}{}^{(4)}\psi _{3}^{}\left\{P_3^1(\mathrm{cos}\theta )\mathrm{sin}\phi {\displaystyle \frac{1}{2}}P_3^3(\mathrm{cos}\theta )\mathrm{sin}3\phi \right\}\right].`$ (G2)
Therefore, we have the form of $`\stackrel{~}{\mathrm{\Phi }}_5`$ and the equation for it as
$`\stackrel{~}{\mathrm{\Phi }}_5={}_{}{}^{(5)}\varphi _{3}^{}(\xi )\left[P_3^1(\mathrm{cos}\theta )\mathrm{sin}\phi {\displaystyle \frac{1}{2}}P_3^3(\mathrm{cos}\theta )\mathrm{sin}3\phi \right],`$ (G3)
$`\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{12}{\xi ^2}}+{\displaystyle \frac{n}{\mathrm{\Theta }_0}}{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}{\displaystyle \frac{d}{d\xi }}\right]{}_{}{}^{(5)}\varphi _{3}^{}={\displaystyle \frac{n}{4\xi _1\mathrm{\Theta }_0}}{}_{}{}^{(4)}\psi _{3}^{}.`$ (G4)
The boundary condition is
$$\frac{d{}_{}{}^{(5)}\varphi _{3}^{}}{d\xi }(\xi _1)\frac{d\mathrm{\Theta }_0}{d\xi }(\xi _1)\frac{1}{4\xi _1}{}_{}{}^{(4)}\psi _{3}^{}(\xi _1)=0.$$
(G5)
Next, we consider the deformation of the figure at $`O(ϵ^5)`$. The equation is written as
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_5=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_5.$$
(G6)
We expand $`\mathrm{\Theta }_5`$ as
$$\mathrm{\Theta }_5=\underset{l,m}{}{}_{}{}^{(5)}\psi _{lm}^{}(\xi )Y_l^m(\theta ,\phi ).$$
(G7)
The surface equation up to $`O(ϵ^5)`$ is written as
$$\mathrm{\Xi }=\xi _1+ϵ^3S_3(\theta ,\phi )+ϵ^4S_4(\theta ,\phi )+ϵ^5S_5(\theta ,\phi ),$$
(G8)
where
$$S_5=\frac{\mathrm{\Theta }_5(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}.$$
(G9)
The gravitational potentials and their derivatives up to $`O(ϵ^5)`$ are
$`\stackrel{~}{U}_{int}`$ $`=`$ $`\mathrm{\Theta }+c_0ϵ^3{\displaystyle \frac{\mu }{\xi _1^2(1+p)}}\xi ^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{\mu }{\xi _1^3(1+p)}}\xi ^3P_3(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (G11)
$`+ϵ^5\left[c_5{\displaystyle \frac{\mu }{\xi _1^4(1+p)}}\xi ^4P_4(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }}{\xi }}ϵ^3{\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}\xi P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{3\mu }{\xi _1^3(1+p)}}\xi ^2P_3(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (G13)
$`ϵ^5{\displaystyle \frac{4\mu }{\xi _1^4(1+p)}}\xi ^3P_4(\mathrm{sin}\theta \mathrm{cos}\phi ),`$
$`\stackrel{~}{U}_{ext}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }}+ϵ^3{\displaystyle \frac{\kappa _{3:2}}{\xi ^3}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{\kappa _{4:3}}{\xi ^4}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^5{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{5:lm}}{\xi ^{l+1}}}Y_l^m,`$ (G14)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi ^2}}ϵ^3{\displaystyle \frac{3\kappa _{3:2}}{\xi ^4}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^4{\displaystyle \frac{4\kappa _{4:3}}{\xi ^5}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^5{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{5:lm}}{\xi ^{l+2}}}Y_l^m.`$ (G15)
They become at the stellar surface as
$`\stackrel{~}{U}_{int}(\mathrm{\Xi })`$ $`=`$ $`c_0ϵ^3{\displaystyle \frac{\mu }{1+p}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^5\left[c_5{\displaystyle \frac{\mu }{1+p}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$ (G16)
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1)+ϵ^3\left[S_3{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_3}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1(1+p)}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )\right]`$ (G19)
$`+ϵ^4\left[S_4{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_4}{\xi }}(\xi _1)+{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )\right]`$
$`+ϵ^5\left[S_5{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_5}{\xi }}(\xi _1){\displaystyle \frac{4\mu }{\xi _1(1+p)}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )\right],`$
$`\stackrel{~}{U}_{ext}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1}}\left(1ϵ^3{\displaystyle \frac{S_3}{\xi _1}}ϵ^4{\displaystyle \frac{S_4}{\xi _1}}ϵ^5{\displaystyle \frac{S_5}{\xi _1}}\right)+ϵ^3{\displaystyle \frac{\kappa _{3:2}}{\xi _1^3}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{\kappa _{4:3}}{\xi _1^4}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^5{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (G20)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1^2}}\left(1ϵ^3{\displaystyle \frac{2S_3}{\xi _1}}ϵ^4{\displaystyle \frac{2S_4}{\xi _1}}ϵ^5{\displaystyle \frac{2S_5}{\xi _1}}\right)ϵ^3{\displaystyle \frac{3\kappa _{3:2}}{\xi _1^4}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^4{\displaystyle \frac{4\kappa _{4:3}}{\xi _1^5}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (G22)
$`ϵ^5{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+2}}}Y_l^m.`$
Then, the boundary conditions for $`O(ϵ^5)`$ become
$`c_5{\displaystyle \frac{\mu }{1+p}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{\kappa _0}{\xi _1^2}}S_5+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (G23)
$`S_5{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_5}{\xi }}(\xi _1){\displaystyle \frac{4\mu }{\xi _1(1+p)}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \frac{2\kappa _0}{\xi _1^3}}S_5{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+2}}}Y_l^m,`$ (G24)
and also $`\mathrm{\Theta }_5`$ should be satisfied with the conditions at the center of the star,
$`{\displaystyle \frac{\mathrm{\Theta }_5}{\xi }}(0)`$ $`=`$ $`0,`$ (G25)
$`\mathrm{\Theta }_5(0)`$ $`=`$ $`0.`$ (G26)
### 1 The case of $`n=0`$
In this case, since we can express $`{}_{}{}^{(5)}\psi _{lm}^{}={}_{}{}^{(5)}A_{lm}^{}\xi ^l`$ by using a constant $`{}_{}{}^{(5)}A_{lm}^{}`$, the boundary conditions can be written as
$`c_5{\displaystyle \frac{\mu }{1+p}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left({}_{}{}^{(5)}A_{lm}^{}\xi _1^l+{\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+1}}}\right)Y_l^m,`$ (G27)
$`{\displaystyle \frac{4\mu }{\xi _1(1+p)}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[(3l){}_{}{}^{(5)}A_{lm}^{}\xi _1^{l1}(l+1){\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (G28)
Therefore, we obtain
$`c_5`$ $`=`$ $`0,`$ (G29)
$`\mathrm{\Theta }_5`$ $`=`$ $`{}_{}{}^{(5)}A_{4}^{}\xi ^4P_4(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (G30)
where
$`{}_{}{}^{(5)}A_{4}^{}`$ $`=`$ $`{\displaystyle \frac{3\mu }{2(1+p)\xi _1^4}},`$ (G31)
$`\kappa _{5:4}`$ $`=`$ $`{\displaystyle \frac{\mu \xi _1^5}{2(1+p)}}.`$ (G32)
### 2 The case of $`n=1`$
In this case, since we can express $`{}_{}{}^{(5)}\psi _{lm}^{}={}_{}{}^{(5)}B_{lm}^{}j_l`$ by using a constant $`{}_{}{}^{(5)}B_{lm}^{}`$, the boundary conditions can be written as
$`c_5{\displaystyle \frac{\mu }{1+p}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(5)}B_{lm}^{}j_l(\xi _1)+{\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (G33)
$`{\displaystyle \frac{4\mu }{\xi _1(1+p)}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(5)}B_{lm}^{}{\displaystyle \frac{dj_l}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (G34)
Therefore, we obtain
$`c_5`$ $`=`$ $`0,`$ (G35)
$`\mathrm{\Theta }_5`$ $`=`$ $`{}_{}{}^{(5)}B_{4}^{}j_4(\xi )P_4(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (G36)
where
$`{}_{}{}^{(5)}B_{4}^{}`$ $`=`$ $`{\displaystyle \frac{9\mu \xi _1^2}{(1+p)(15\xi _1^2)}},`$ (G37)
$`\kappa _{5:4}`$ $`=`$ $`{\displaystyle \frac{\mu }{1+p}}\left({\displaystyle \frac{\xi _1^4105\xi _1^2+945}{\xi _1^2(15\xi _1^2)}}\right),`$ (G38)
$`j_4(\xi )`$ $`=`$ $`{\displaystyle \frac{1}{\xi ^5}}\left[(10545\xi ^2+\xi ^4)\mathrm{sin}\xi \xi (10510\xi ^2)\mathrm{cos}\xi \right].`$ (G39)
### 3 The case of $`n>0`$
In this case, the boundary conditions can be written as
$`c_5{\displaystyle \frac{\mu }{1+p}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(5)}\psi _{lm}^{}(\xi _1)+{\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (G40)
$`{\displaystyle \frac{4\mu }{\xi _1(1+p)}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )={\displaystyle \underset{l,m}{}}\left[{\displaystyle \frac{d{}_{}{}^{(5)}\psi _{lm}^{}}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{5:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (G41)
Therefore, we obtain
$`c_5`$ $`=`$ $`0,`$ (G42)
$`\mathrm{\Theta }_5`$ $`=`$ $`{}_{}{}^{(5)}\psi _{4}^{}(\xi )P_4(\mathrm{sin}\theta \mathrm{cos}\phi ),`$ (G43)
and the ordinary differential equation for $`{}_{}{}^{(5)}\psi _{4}^{}`$ is written as
$$\left[\frac{1}{\xi ^2}\frac{d}{d\xi }\left(\xi ^2\frac{d}{d\xi }\right)\frac{20}{\xi ^2}\right]{}_{}{}^{(5)}\psi _{4}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(5)}\psi _{4}^{}.$$
(G44)
The coefficient of $`{}_{}{}^{(5)}\psi _{4}^{}`$ and the constant $`\kappa _{5:4}`$ are determined from two equations;
$`5{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(5)}\psi _{4}^{}}{d\xi }}(\xi _1)={\displaystyle \frac{9\mu }{1+p}},`$ (G45)
$`\kappa _{5:4}=\xi _1^5\left[{}_{}{}^{(5)}\psi _{4}^{}(\xi _1){\displaystyle \frac{\mu }{1+p}}\right]={\displaystyle \frac{\xi _1^5}{5}}\left[{\displaystyle \frac{4\mu }{1+p}}\xi _1{\displaystyle \frac{d{}_{}{}^{(5)}\psi _{4}^{}}{d\xi }}(\xi _1)\right].`$ (G46)
### 4 Orbital angular velocity
At this order, the orbital angular velocity has the 5th higher order term than the monopole term. The total angular velocity is
$$\mathrm{\Omega }^2=\frac{M_{tot}}{R^3}\left[1+\frac{9}{2R^2}\left(ϵ^3\frac{\overline{I}_{11}}{M_1}+ϵ_{}^{}{}_{}{}^{3}\frac{\overline{I}_{11}^{}}{M_2}\right)\right].$$
(G47)
For the calculational convenience, we rewrite the above equation as
$$\mathrm{\Omega }^2=\frac{M_{tot}}{R^3}\left(1+ϵ^5\delta \right),$$
(G48)
where
$$\delta \frac{9}{2a_0^2}\left[\frac{\overline{I}_{11}}{M_1}+\left(\frac{a_0^{}}{a_0}\right)^3\frac{\overline{I}_{11}^{}}{M_2}\right].$$
(G49)
## H Solutions at sixth order
The equation for the velocity potential is
$`\stackrel{~}{\mathrm{\Delta }}\stackrel{~}{\mathrm{\Phi }}_6`$ $`=`$ $`{\displaystyle \frac{n}{\mathrm{\Theta }_0}}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_6)(\stackrel{~}{}\mathrm{\Theta }_0){\displaystyle \frac{1}{\xi _1}}(\stackrel{~}{𝛀}\times 𝝃)_{fig}(\stackrel{~}{}\mathrm{\Theta }_5)\right],`$ (H1)
$`=`$ $`{\displaystyle \frac{n}{\mathrm{\Theta }_0}}\left[(\stackrel{~}{}\stackrel{~}{\mathrm{\Phi }}_6)(\stackrel{~}{}\mathrm{\Theta }_0){\displaystyle \frac{1}{12\xi _1}}{}_{}{}^{(5)}\psi _{4}^{}\left\{P_4^2(\mathrm{cos}\theta )\mathrm{sin}2\phi {\displaystyle \frac{1}{4}}P_4^4(\mathrm{cos}\theta )\mathrm{sin}4\phi \right\}\right].`$ (H2)
Therefore, we have the form of $`\stackrel{~}{\mathrm{\Phi }}_6`$ and the equation for it as
$`\stackrel{~}{\mathrm{\Phi }}_6={}_{}{}^{(6)}\varphi _{4}^{}(\xi )\left[P_4^2(\mathrm{cos}\theta )\mathrm{sin}2\phi {\displaystyle \frac{1}{4}}P_4^4(\mathrm{cos}\theta )\mathrm{sin}4\phi \right],`$ (H3)
$`\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{20}{\xi ^2}}+{\displaystyle \frac{n}{\mathrm{\Theta }_0}}{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}{\displaystyle \frac{d}{d\xi }}\right]{}_{}{}^{(6)}\varphi _{4}^{}={\displaystyle \frac{n}{12\xi _1\mathrm{\Theta }_0}}{}_{}{}^{(5)}\psi _{4}^{}.`$ (H4)
The boundary condition is
$$\frac{d{}_{}{}^{(6)}\varphi _{4}^{}}{d\xi }(\xi _1)\frac{d\mathrm{\Theta }_0}{d\xi }(\xi _1)\frac{1}{12\xi _1}{}_{}{}^{(5)}\psi _{4}^{}(\xi _1)=0.$$
(H5)
Next, we consider the deformation of the figure at $`O(ϵ^6)`$. The equation is written as
$$\stackrel{~}{\mathrm{\Delta }}\mathrm{\Theta }_6=n\mathrm{\Theta }_0^{n1}\mathrm{\Theta }_6\frac{1}{2}n(n1)\mathrm{\Theta }_0^{n2}\mathrm{\Theta }_3^2+\frac{2\mu }{\xi _1}\stackrel{~}{\mathrm{\Delta }}{}_{}{}^{(4)}\varphi _{2}^{}(\xi )P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi .$$
(H6)
We expand $`\mathrm{\Theta }_6`$ as
$$\mathrm{\Theta }_6=\underset{l,m}{}{}_{}{}^{(6)}\psi _{lm}^{}(\xi )Y_l^m(\theta ,\phi ).$$
(H7)
The surface equation up to $`O(ϵ^6)`$ is written as
$$\mathrm{\Xi }=\xi _1+ϵ^3S_3(\theta ,\phi )+ϵ^4S_4(\theta ,\phi )+ϵ^5S_5(\theta ,\phi )+ϵ^6S_6(\theta ,\phi ),$$
(H8)
where
$$S_6=\frac{1}{\mathrm{\Theta }_{0,\xi }(\xi _1)}\left[\frac{S_3^2}{2}\frac{d^2\mathrm{\Theta }_0}{d\xi ^2}(\xi _1)+S_3\frac{\mathrm{\Theta }_3}{\xi }(\xi _1)+\mathrm{\Theta }_6(\xi _1)\right].$$
(H9)
The gravitational potentials and their derivatives up to $`O(ϵ^6)`$ are
$`\stackrel{~}{U}_{int}`$ $`=`$ $`\mathrm{\Theta }+c_0ϵ^3{\displaystyle \frac{\mu }{\xi _1^2(1+p)}}\xi ^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{\mu }{\xi _1^3(1+p)}}\xi ^3P_3(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H12)
$`ϵ^5{\displaystyle \frac{\mu }{\xi _1^4(1+p)}}\xi ^4P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$
$`+ϵ^6\left[c_6{\displaystyle \frac{3\mu \overline{I}_{11}^{}}{2M_{tot}a_0^2}}\left({\displaystyle \frac{a_0^{}}{a_0}}\right)^3{\displaystyle \frac{\mu \delta }{2(1+p)^2}}+{\displaystyle \frac{\mu }{\xi _1^5(1+p)}}\xi ^5P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{}_{}{}^{(4)}\varphi _{2}^{}P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi \right],`$
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Theta }}{\xi }}ϵ^3{\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}\xi P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{3\mu }{\xi _1^3(1+p)}}\xi ^2P_3(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^5{\displaystyle \frac{4\mu }{\xi _1^4(1+p)}}\xi ^3P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H14)
$`+ϵ^6\left[{\displaystyle \frac{5\mu }{\xi _1^5(1+p)}}\xi ^4P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi \right],`$
$`\stackrel{~}{U}_{ext}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi }}+ϵ^3{\displaystyle \frac{\kappa _{3:2}}{\xi ^3}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{\kappa _{4:3}}{\xi ^4}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^5{\displaystyle \frac{\kappa _{5:4}}{\xi ^5}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^6{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{6:lm}}{\xi ^{l+1}}}Y_l^m,`$ (H15)
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi ^2}}ϵ^3{\displaystyle \frac{3\kappa _{3:2}}{\xi ^4}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^4{\displaystyle \frac{4\kappa _{4:3}}{\xi ^5}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^5{\displaystyle \frac{5\kappa _{5:4}}{\xi ^6}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H17)
$`ϵ^6{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{6:lm}}{\xi ^{l+2}}}Y_l^m.`$
They become at the stellar surface as
$`\stackrel{~}{U}_{int}(\mathrm{\Xi })`$ $`=`$ $`c_0ϵ^3{\displaystyle \frac{\mu }{1+p}}\left(1+{\displaystyle \frac{2S_3}{\xi _1}}ϵ^3\right)P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^4{\displaystyle \frac{\mu }{1+p}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^5{\displaystyle \frac{\mu }{1+p}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H19)
$`+ϵ^6\left[c_6{\displaystyle \frac{3\mu \overline{I}_{11}^{}}{2M_{tot}a_0^2}}\left({\displaystyle \frac{a_0^{}}{a_0}}\right)^3{\displaystyle \frac{\mu \delta }{2(1+p)^2}}+{\displaystyle \frac{\mu }{1+p}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{}_{}{}^{(4)}\varphi _{2}^{}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi \right],`$
$`{\displaystyle \frac{\stackrel{~}{U}_{int}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{d\mathrm{\Theta }_0}{d\xi }}(\xi _1)+ϵ^3\left[S_3{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_3}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1(1+p)}}\left(1+{\displaystyle \frac{S_3}{\xi _1}}ϵ^3\right)P_2(\mathrm{sin}\theta \mathrm{cos}\phi )\right]`$ (H24)
$`+ϵ^4\left[S_4{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_4}{\xi }}(\xi _1)+{\displaystyle \frac{3\mu }{\xi _1(1+p)}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )\right]`$
$`+ϵ^5\left[S_5{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_5}{\xi }}(\xi _1){\displaystyle \frac{4\mu }{\xi _1(1+p)}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )\right]`$
$`+ϵ^6[S_6{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{1}{2}}S_3^2{\displaystyle \frac{d^3\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+S_3{\displaystyle \frac{^2\mathrm{\Theta }_3}{\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_6}{\xi }}(\xi _1)+{\displaystyle \frac{5\mu }{\xi _1(1+p)}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi )`$
$`{\displaystyle \frac{2\mu }{\xi _1}}{\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi ],`$
$`\stackrel{~}{U}_{ext}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1}}\left[1ϵ^3{\displaystyle \frac{S_3}{\xi _1}}ϵ^4{\displaystyle \frac{S_4}{\xi _1}}ϵ^5{\displaystyle \frac{S_5}{\xi _1}}+ϵ^6\left\{\left({\displaystyle \frac{S_3}{\xi _1}}\right)^2{\displaystyle \frac{S_6}{\xi _1}}\right\}\right]+ϵ^3\left(1ϵ^3{\displaystyle \frac{3S_3}{\xi _1}}\right){\displaystyle \frac{\kappa _{3:2}}{\xi _1^3}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H26)
$`+ϵ^4{\displaystyle \frac{\kappa _{4:3}}{\xi _1^4}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^5{\displaystyle \frac{\kappa _{5:4}}{\xi _1^5}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+ϵ^6{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+1}}}Y_l^m,`$
$`{\displaystyle \frac{\stackrel{~}{U}_{ext}}{\xi }}(\mathrm{\Xi })`$ $`=`$ $`{\displaystyle \frac{\kappa _0}{\xi _1^2}}\left[1ϵ^3{\displaystyle \frac{2S_3}{\xi _1}}ϵ^4{\displaystyle \frac{2S_4}{\xi _1}}ϵ^5{\displaystyle \frac{2S_5}{\xi _1}}+ϵ^6\left\{3\left({\displaystyle \frac{S_3}{\xi _1}}\right)^22{\displaystyle \frac{S_6}{\xi _1}}\right\}\right]ϵ^3\left(1ϵ^3{\displaystyle \frac{4S_3}{\xi _1}}\right){\displaystyle \frac{3\kappa _{3:2}}{\xi _1^4}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H28)
$`ϵ^4{\displaystyle \frac{4\kappa _{4:3}}{\xi _1^5}}P_3(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^5{\displaystyle \frac{5\kappa _{5:4}}{\xi _1^6}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )ϵ^6{\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+2}}}Y_l^m.`$
Then, the boundary conditions for $`O(ϵ^6)`$ become
$`\stackrel{~}{c}_6{\displaystyle \frac{2\mu }{\xi _1(1+p)}}S_3P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{\mu }{1+p}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{}_{}{}^{(4)}\varphi _{2}^{}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (H29)
$`={\displaystyle \frac{\kappa _0}{\xi _1^2}}\left({\displaystyle \frac{S_3^2}{\xi _1}}S_6\right){\displaystyle \frac{3\kappa _{3:2}}{\xi _1^4}}S_3P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \underset{l,m}{}}{\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+1}}}Y_l^m,`$ (H30)
$`S_6{\displaystyle \frac{d^2\mathrm{\Theta }_0}{d\xi ^2}}(\xi _1)+{\displaystyle \frac{1}{2}}S_3^2{\displaystyle \frac{d^3\mathrm{\Theta }_0}{d\xi ^3}}(\xi _1)+S_3{\displaystyle \frac{^2\mathrm{\Theta }_3}{\xi ^2}}(\xi _1)+{\displaystyle \frac{\mathrm{\Theta }_6}{\xi }}(\xi _1){\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}S_3P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{5\mu }{\xi _1(1+p)}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H31)
$`{\displaystyle \frac{2\mu }{\xi _1}}{\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi ={\displaystyle \frac{\kappa _0}{\xi _1^3}}\left({\displaystyle \frac{3S_3^2}{\xi _1}}2S_6\right)+{\displaystyle \frac{12\kappa _{3:2}}{\xi _1^5}}S_3P_2(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \underset{l,m}{}}(l+1){\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+2}}}Y_l^m,`$ (H32)
where
$$\stackrel{~}{c}_6c_6\frac{3\mu \overline{I}_{11}^{}}{2M_{tot}a_0^2}\left(\frac{a_0^{}}{a_0}\right)^3\frac{\mu \delta }{2(1+p)^2}.$$
(H33)
Also $`\mathrm{\Theta }_6`$ should be satisfied with the conditions at the center of the star,
$`{\displaystyle \frac{\mathrm{\Theta }_6}{\xi }}(0)`$ $`=`$ $`0,`$ (H34)
$`\mathrm{\Theta }_6(0)`$ $`=`$ $`0.`$ (H35)
In the calculation of this section, we need the expression of $`[P_2(\mathrm{sin}\theta \mathrm{cos}\phi )]^2`$. It is written as
$$\left[P_2(\mathrm{sin}\theta \mathrm{cos}\phi )\right]^2=\frac{18}{35}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+\frac{2}{7}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+\frac{1}{5}.$$
(H36)
### 1 The case of $`n=0`$
In this case, since we can express $`{}_{}{}^{(3)}\psi _{2}^{}={}_{}{}^{(3)}A_{2}^{}\xi ^2`$ and $`{}_{}{}^{(6)}\psi _{lm}^{}={}_{}{}^{(6)}A_{lm}^{}\xi ^l`$, by using the constants $`{}_{}{}^{(3)}A_{2}^{}`$ and $`{}_{}{}^{(6)}A_{lm}^{}`$, the boundary conditions can be written as
$`\stackrel{~}{c}_6+{\displaystyle \frac{\mu }{1+p}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{}_{}{}^{(4)}\varphi _{2}^{}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (H37)
$`={\displaystyle \frac{{}_{}{}^{(3)}A_{2}^{}\xi _1^2}{\mathrm{\Theta }_{0,\xi }(\xi _1)}}\left[{\displaystyle \frac{2\mu }{\xi _1(1+p)}}+{\displaystyle \frac{1}{2}}{}_{}{}^{(3)}A_{2}^{}\xi _1+{\displaystyle \frac{3\kappa _{3:2}}{\xi _1^4}}\right]\left({\displaystyle \frac{18}{35}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{2}{7}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{1}{5}}\right)`$ (H38)
$`+{\displaystyle \underset{l,m}{}}\left({}_{}{}^{(6)}A_{lm}^{}\xi _1^l+{\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+1}}}\right)Y_l^m,`$ (H39)
$`{\displaystyle \frac{5\mu }{\xi _1(1+p)}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (H40)
$`={\displaystyle \frac{{}_{}{}^{(3)}A_{2}^{}\xi _1^2}{\mathrm{\Theta }_{0,\xi }(\xi _1)}}\left[{\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}+{\displaystyle \frac{1}{2}}{}_{}{}^{(3)}A_{2}^{}{\displaystyle \frac{12\kappa _{3:2}}{\xi _1^5}}\right]\left({\displaystyle \frac{18}{35}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{2}{7}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{1}{5}}\right)`$ (H41)
$`+{\displaystyle \underset{l,m}{}}\left[(3l){}_{}{}^{(6)}A_{lm}^{}\xi _1^{l1}(l+1){\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (H42)
Therefore, we obtain
$`\stackrel{~}{c}_6`$ $`=`$ $`{\displaystyle \frac{45\mu ^2}{2(1+p)^2\xi _1^2}},`$ (H43)
$`\mathrm{\Theta }_6`$ $`=`$ $`{}_{}{}^{(6)}A_{2}^{}\xi ^2P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{}_{}{}^{(6)}A_{22}^{}\xi ^2P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi +{}_{}{}^{(6)}A_{4}^{}\xi ^4P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H45)
$`+{}_{}{}^{(6)}A_{5}^{}\xi ^5P_5(\mathrm{sin}\theta \mathrm{cos}\phi ),`$
where
$`{}_{}{}^{(6)}A_{2}^{}`$ $`=`$ $`{\displaystyle \frac{225\mu ^2}{28(1+p)^2\xi _1^4}},`$ (H46)
$`{}_{}{}^{(6)}A_{22}^{}`$ $`=`$ $`{\displaystyle \frac{75\mu ^2}{8(1+p)\xi _1^4}},`$ (H47)
$`{}_{}{}^{(6)}A_{4}^{}`$ $`=`$ $`{\displaystyle \frac{45\mu ^2}{14(1+p)^2\xi _1^6}},`$ (H48)
$`{}_{}{}^{(6)}A_{5}^{}`$ $`=`$ $`{\displaystyle \frac{11\mu }{8(1+p)\xi _1^5}}.`$ (H49)
### 2 The case of $`n=1`$
In this case, since we can express $`{}_{}{}^{(3)}\psi _{2}^{}={}_{}{}^{(3)}B_{2}^{}j_2`$ and $`{}_{}{}^{(6)}\psi _{lm}^{}={}_{}{}^{(6)}B_{lm}^{}`$ except for $`{}_{}{}^{(6)}\psi _{22}^{}`$ by using the constants $`{}_{}{}^{(3)}B_{2}^{}`$ and $`{}_{}{}^{(6)}B_{lm}^{}`$, the boundary conditions can be written as
$`\stackrel{~}{c}_6+{\displaystyle \frac{\mu }{1+p}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{}_{}{}^{(4)}\varphi _{2}^{}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (H50)
$`={\displaystyle \frac{{}_{}{}^{(3)}B_{2}^{}j_2(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}}\left[{\displaystyle \frac{2\mu }{\xi _1(1+p)}}+{}_{}{}^{(3)}B_{2}^{}{\displaystyle \frac{dj_2}{d\xi }}(\xi _1)+{\displaystyle \frac{3\kappa _{3:2}}{\xi _1^4}}\right]\left({\displaystyle \frac{18}{35}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{2}{7}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{1}{5}}\right)`$ (H51)
$`+{\displaystyle \underset{l,m}{}}\left({}_{}{}^{(6)}B_{lm}^{}j_l(\xi _1)+{\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+1}}}\right)Y_l^m,`$ (H52)
$`{\displaystyle \frac{5\mu }{\xi _1(1+p)}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (H53)
$`={\displaystyle \frac{{}_{}{}^{(3)}B_{2}^{}j_2(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}}\left[{\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}+\left({\displaystyle \frac{3}{\xi _1^2}}+{\displaystyle \frac{\xi _1}{2}}{\displaystyle \frac{d^3\mathrm{\Theta }_0}{d\xi ^3}}(\xi _1)\right){}_{}{}^{(3)}B_{2}^{}j_2(\xi _1)+{}_{}{}^{(3)}B_{2}^{}{\displaystyle \frac{d^2j_2}{d\xi ^2}}(\xi _1){\displaystyle \frac{12\kappa _{3:2}}{\xi _1^5}}\right]`$ (H54)
$`\times \left({\displaystyle \frac{18}{35}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{2}{7}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{1}{5}}\right)`$ (H55)
$`+{\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(6)}B_{lm}^{}{\displaystyle \frac{dj_l}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (H56)
Therefore, we obtain
$`\stackrel{~}{c}_6`$ $`=`$ $`{\displaystyle \frac{45\mu ^2}{2(1+p)^2\xi _1^2}},`$ (H57)
$`\mathrm{\Theta }_6`$ $`=`$ $`{}_{}{}^{(6)}B_{2}^{}j_2(\xi )P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{}_{}{}^{(6)}\psi _{22}^{}(\xi )P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi +{}_{}{}^{(6)}B_{4}^{}j_4(\xi )P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H59)
$`+{}_{}{}^{(6)}B_{5}^{}j_5(\xi )P_5(\mathrm{sin}\theta \mathrm{cos}\phi ),`$
where
$`{}_{}{}^{(6)}B_{2}^{}`$ $`=`$ $`{\displaystyle \frac{225\mu ^2}{7(1+p)^2\xi _1^2}},`$ (H60)
$`{}_{}{}^{(6)}B_{4}^{}`$ $`=`$ $`{\displaystyle \frac{405\mu ^2}{7(1+p)^2(15\xi _1^2)}},`$ (H61)
$`{}_{}{}^{(6)}B_{5}^{}`$ $`=`$ $`{\displaystyle \frac{11\mu \xi _1^3}{(1+p)(10510\xi _1^2)}},`$ (H62)
$`j_5(\xi )`$ $`=`$ $`{\displaystyle \frac{1}{\xi ^6}}\left[(945420\xi ^2+15\xi ^4)\mathrm{sin}\xi \xi (945105\xi ^2+\xi ^4)\mathrm{cos}\xi \right],`$ (H63)
and the coefficient of $`{}_{}{}^{(6)}\psi _{22}^{}`$ is determined from the equation;
$$3{}_{}{}^{(6)}\psi _{22}^{}(\xi _1)+\xi _1\frac{d{}_{}{}^{(6)}\psi _{22}^{}}{d\xi }(\xi _1)=\frac{2\mu }{\xi _1}\left[3{}_{}{}^{(4)}\varphi _{2}^{}(\xi _1)+\xi _1\frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }(\xi _1)\right].$$
(H64)
### 3 The case of $`n>0`$
In this case, the boundary conditions can be written as
$`\stackrel{~}{c}_6+{\displaystyle \frac{\mu }{1+p}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{}_{}{}^{(4)}\varphi _{2}^{}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (H65)
$`={\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}}\left[{\displaystyle \frac{2\mu }{\xi _1(1+p)}}+{\displaystyle \frac{d{}_{}{}^{(3)}\psi _{2}^{}}{d\xi }}(\xi _1)+{\displaystyle \frac{3\kappa _{3:2}}{\xi _1^4}}\right]\left({\displaystyle \frac{18}{35}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{2}{7}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{1}{5}}\right)`$ (H66)
$`+{\displaystyle \underset{l,m}{}}\left[{}_{}{}^{(6)}\psi _{lm}^{}(\xi _1)+{\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+1}}}\right]Y_l^m,`$ (H67)
$`{\displaystyle \frac{5\mu }{\xi _1(1+p)}}P_5(\mathrm{sin}\theta \mathrm{cos}\phi ){\displaystyle \frac{2\mu }{\xi _1}}{\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}(\xi _1)P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi `$ (H68)
$`={\displaystyle \frac{{}_{}{}^{(3)}\psi _{2}^{}(\xi _1)}{\mathrm{\Theta }_{0,\xi }(\xi _1)}}\left[{\displaystyle \frac{2\mu }{\xi _1^2(1+p)}}+\left({\displaystyle \frac{3}{\xi _1^2}}{\displaystyle \frac{1}{2\mathrm{\Theta }_{0,\xi }(\xi _1)}}{\displaystyle \frac{d^3\mathrm{\Theta }_0}{d\xi ^3}}(\xi _1)\right){}_{}{}^{(3)}\psi _{2}^{}(\xi _1)+{\displaystyle \frac{d^2{}_{}{}^{(3)}\psi _{2}^{}}{d\xi ^2}}(\xi _1){\displaystyle \frac{12\kappa _{3:2}}{\xi _1^5}}\right]`$ (H69)
$`\times \left({\displaystyle \frac{18}{35}}P_4(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{2}{7}}P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{\displaystyle \frac{1}{5}}\right)`$ (H70)
$`+{\displaystyle \underset{l,m}{}}\left[{\displaystyle \frac{d{}_{}{}^{(6)}\psi _{lm}^{}}{d\xi }}(\xi _1)(l+1){\displaystyle \frac{\kappa _{6:lm}}{\xi _1^{l+2}}}\right]Y_l^m.`$ (H71)
Therefore, we obtain
$`\stackrel{~}{c}_6`$ $`=`$ $`{\displaystyle \frac{\xi _1}{10\mathrm{\Theta }_{0,\xi }(\xi _1)}}n\mathrm{\Theta }_0^{n1}(\xi _1)\left({}_{}{}^{(3)}\psi _{2}^{}(\xi _1)\right)^2{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)\xi _1{\displaystyle \frac{d{}_{}{}^{(6)}\psi _{0}^{}}{d\xi }}(\xi _1),`$ (H72)
$`\mathrm{\Theta }_6`$ $`=`$ $`{}_{}{}^{(6)}\psi _{0}^{}(\xi )+{}_{}{}^{(6)}\psi _{2}^{}(\xi )P_2(\mathrm{sin}\theta \mathrm{cos}\phi )+{}_{}{}^{(6)}\psi _{22}^{}(\xi )P_2^2(\mathrm{cos}\theta )\mathrm{cos}2\phi +{}_{}{}^{(6)}\psi _{4}^{}(\xi )P_4(\mathrm{sin}\theta \mathrm{cos}\phi )`$ (H74)
$`+{}_{}{}^{(6)}\psi _{5}^{}(\xi )P_5(\mathrm{sin}\theta \mathrm{cos}\phi ),`$
where the forms of $`{}_{}{}^{(6)}\psi _{i}^{}`$ are determined by the equations
$`{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){}_{}{}^{(6)}\psi _{0}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(6)}\psi _{0}^{}{\displaystyle \frac{1}{10}}n(n1)\mathrm{\Theta }_0^{n2}({}_{}{}^{(3)}\psi _{2}^{})^2,`$ (H75)
$`\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{6}{\xi ^2}}\right]{}_{}{}^{(6)}\psi _{2}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(6)}\psi _{2}^{}{\displaystyle \frac{1}{7}}n(n1)\mathrm{\Theta }_0^{n2}({}_{}{}^{(3)}\psi _{2}^{})^2,`$ (H76)
$`\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{6}{\xi ^2}}\right]{}_{}{}^{(6)}\psi _{22}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(6)}\psi _{22}^{}+{\displaystyle \frac{2\mu }{\xi _1}}\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{6}{\xi ^2}}\right]{}_{}{}^{(4)}\varphi _{2}^{},`$ (H77)
$`\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{20}{\xi ^2}}\right]{}_{}{}^{(6)}\psi _{4}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(6)}\psi _{4}^{}{\displaystyle \frac{9}{35}}n(n1)\mathrm{\Theta }_0^{n2}({}_{}{}^{(3)}\psi _{2}^{})^2,`$ (H78)
$`\left[{\displaystyle \frac{1}{\xi ^2}}{\displaystyle \frac{d}{d\xi }}\left(\xi ^2{\displaystyle \frac{d}{d\xi }}\right){\displaystyle \frac{30}{\xi ^2}}\right]{}_{}{}^{(6)}\psi _{5}^{}=n\mathrm{\Theta }_0^{n1}{}_{}{}^{(6)}\psi _{5}^{},`$ (H79)
and the coefficients of $`{}_{}{}^{(6)}\psi _{i}^{}`$ are determined by the boundary conditions;
$`3{}_{}{}^{(6)}\psi _{2}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(6)}\psi _{2}^{}}{d\xi }}(\xi _1)={\displaystyle \frac{\xi _1}{7\mathrm{\Theta }_{0,\xi }(\xi _1)}}n\mathrm{\Theta }_0^{n1}(\xi _1)\left({}_{}{}^{(3)}\psi _{2}^{}(\xi _1)\right)^2,`$ (H80)
$`3{}_{}{}^{(6)}\psi _{22}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(6)}\psi _{22}^{}}{d\xi }}(\xi _1)={\displaystyle \frac{2\mu }{\xi _1}}\left[3{}_{}{}^{(4)}\varphi _{2}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(4)}\varphi _{2}^{}}{d\xi }}(\xi _1)\right],`$ (H81)
$`5{}_{}{}^{(6)}\psi _{4}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(6)}\psi _{4}^{}}{d\xi }}(\xi _1)={\displaystyle \frac{9\xi _1}{35\mathrm{\Theta }_{0,\xi }(\xi _1)}}n\mathrm{\Theta }_0^{n1}(\xi _1)\left({}_{}{}^{(3)}\psi _{2}^{}(\xi _1)\right)^2,`$ (H82)
$`6{}_{}{}^{(6)}\psi _{5}^{}(\xi _1)+\xi _1{\displaystyle \frac{d{}_{}{}^{(6)}\psi _{5}^{}}{d\xi }}(\xi _1)={\displaystyle \frac{11\mu }{1+p}}.`$ (H83)
## I The change in the central density
In this section, we present the expression of the change in the central density at $`O(ϵ^6)`$ by using the first tensor virial relation.
First, we divide the first tensor virial relation (31),
$$\frac{3}{n}\mathrm{\Pi }_1+\frac{3}{n^{}}\mathrm{\Pi }_2+(W_{self})_{tot}+(W_{int})_{tot}+2T_{tot}=0,$$
(I1)
into three parts of $`O(ϵ^0)`$, $`O(ϵ)`$ and $`O(ϵ^6)`$. After that we demand that these three parts should be zero independently.
(a) The 0th order of $`ϵ`$:
$$\frac{M_1^2}{2a_0}\left[\frac{5n}{1+n}\frac{1}{\xi _1^3(\mathrm{\Theta }_{0,\xi }(\xi _1))^2}_0^{\xi _1}𝑑\xi \xi ^2\mathrm{\Theta }_0^{1+n}1\right]+\frac{M_2^2}{2a_0^{}}\left[\frac{5n^{}}{1+n^{}}\frac{1}{\xi _{1}^{}{}_{}{}^{3}(\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}))^2}_0^{\xi _1^{}}𝑑\xi ^{}\xi _{}^{}{}_{}{}^{2}\overline{\mathrm{\Theta }}_0^{1+n^{}}1\right]=0.$$
(I2)
If we assume that both of the terms concerned with star 1 and star 2 become zero independently, we have two relations;
$`{\displaystyle \frac{5n}{1+n}}{\displaystyle \frac{1}{\xi _1^3(\mathrm{\Theta }_{0,\xi }(\xi _1))^2}}{\displaystyle _0^{\xi _1}}𝑑\xi \xi ^2\mathrm{\Theta }_0^{1+n}=1,`$ (I3)
$`{\displaystyle \frac{5n^{}}{1+n^{}}}{\displaystyle \frac{1}{\xi _{1}^{}{}_{}{}^{3}(\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}))^2}}{\displaystyle _0^{\xi _1^{}}}𝑑\xi ^{}\xi _{}^{}{}_{}{}^{2}\overline{\mathrm{\Theta }}_0^{1+n^{}}=1.`$ (I4)
The above relations are proved using the Lane-Emden equations.
(b) The 1st order of $`ϵ`$:
It is clear that this term becomes zero.
(c) The 6th order of $`ϵ`$:
By using the relations;
$`\stackrel{~}{c}_6`$ $`=`$ $`\left({\displaystyle \frac{3n}{2n}}\right)\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1){\displaystyle \frac{\delta \rho _c}{\rho _c}}{}_{}{}^{(6)}\psi _{0}^{}(\xi _1),`$ (I5)
$`\stackrel{~}{c}_6^{}`$ $`=`$ $`\left({\displaystyle \frac{3n^{}}{2n^{}}}\right)\xi _1^{}\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}){\displaystyle \frac{\delta \rho _c^{}}{\rho _c^{}}}{}_{}{}^{(6)}\overline{\psi }_{0}^{}(\xi _1^{}),`$ (I6)
$`\mu `$ $`=`$ $`\left({\displaystyle \frac{1+p}{p}}\right)\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1),`$ (I7)
$`\mu ^{}`$ $`=`$ $`(1+p)\xi _1^{}\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}),`$ (I8)
we can express the equation at order $`ϵ^6`$ as
$`{\displaystyle \frac{M_1^2}{a_0}}[{\displaystyle \frac{5n}{2(1+n)}}{\displaystyle \frac{1}{\xi _1^3(\mathrm{\Theta }_{0,\xi }(\xi _1))^2}}[{\displaystyle \frac{5n}{2n}}{\displaystyle \frac{\delta \rho _c}{\rho _c}}{\displaystyle _0^{\xi _1}}d\xi \xi ^2\mathrm{\Theta }_0^{1+n}(1+n){\displaystyle _0^{\xi _1}}d\xi \xi ^2\{\mathrm{\Theta }_0^n{}_{}{}^{(6)}\psi _{0}^{}+{\displaystyle \frac{1}{10}}n\mathrm{\Theta }_0^{n1}({}_{}{}^{(3)}\psi _{2}^{})^2\}]`$ (I9)
$`+{\displaystyle \frac{15}{4p}}{\displaystyle \frac{\overline{I}_{11}}{M_1a_0^2}}{\displaystyle \frac{{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)}{2\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1)}}+{\displaystyle \frac{n1}{4n}}{\displaystyle \frac{\delta \rho _c}{\rho _c}}]`$ (I10)
$`+{\displaystyle \frac{M_2^2}{a_0^{}}}\left({\displaystyle \frac{a_0^{}}{a_0}}\right)^6[{\displaystyle \frac{5n^{}}{2(1+n^{})}}{\displaystyle \frac{1}{\xi _{1}^{}{}_{}{}^{3}(\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}))^2}}`$ (I11)
$`\times \left[{\displaystyle \frac{5n^{}}{2n^{}}}{\displaystyle \frac{\delta \rho _c^{}}{\rho _c^{}}}{\displaystyle _0^{\xi _1^{}}}𝑑\xi ^{}\xi _{}^{}{}_{}{}^{2}\overline{\mathrm{\Theta }}_0^{1+n^{}}(1+n^{}){\displaystyle _0^{\xi _1^{}}}𝑑\xi ^{}\xi _{}^{}{}_{}{}^{2}\left\{\overline{\mathrm{\Theta }}_0^n^{}{}_{}{}^{(6)}\overline{\psi }_{0}^{}+{\displaystyle \frac{1}{10}}n^{}\overline{\mathrm{\Theta }}_0^{n^{}1}({}_{}{}^{(3)}\overline{\psi }_{2}^{})^2\right\}\right]`$ (I12)
$`+{\displaystyle \frac{15p}{4}}{\displaystyle \frac{\overline{I}_{11}^{}}{M_2a_{0}^{}{}_{}{}^{2}}}{\displaystyle \frac{{}_{}{}^{(6)}\overline{\psi }_{0}^{}(\xi _1^{})}{2\xi _1^{}\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{})}}+{\displaystyle \frac{n^{}1}{4n^{}}}{\displaystyle \frac{\delta \rho _c^{}}{\rho _c^{}}}]=0.`$ (I13)
If we assume that both of the terms concerned with star 1 and star 2 become zero independently, we obtain two relations;
$`{\displaystyle \frac{\delta \rho _c}{\rho _c}}`$ $`=`$ $`{\displaystyle \frac{2n}{3n}}[{\displaystyle \frac{15}{4p}}{\displaystyle \frac{\overline{I}_{11}}{M_1a_0^2}}{\displaystyle \frac{{}_{}{}^{(6)}\psi _{0}^{}(\xi _1)}{2\xi _1\mathrm{\Theta }_{0,\xi }(\xi _1)}}+{\displaystyle \frac{5n}{2\xi _1^3(\mathrm{\Theta }_{0,\xi }(\xi _1))^2}}{\displaystyle _0^{\xi _1}}d\xi \xi ^2\{\mathrm{\Theta }_0^n{}_{}{}^{(6)}\psi _{0}^{}+{\displaystyle \frac{1}{10}}n\mathrm{\Theta }_0^{n1}({}_{}{}^{(3)}\psi _{2}^{})^2\}],`$ (I14)
$`{\displaystyle \frac{\delta \rho _c^{}}{\rho _c^{}}}`$ $`=`$ $`{\displaystyle \frac{2n^{}}{3n^{}}}\left[{\displaystyle \frac{15p}{4}}{\displaystyle \frac{\overline{I}_{11}^{}}{M_2a_{0}^{}{}_{}{}^{2}}}{\displaystyle \frac{{}_{}{}^{(6)}\overline{\psi }_{0}^{}(\xi _1^{})}{2\xi _1^{}\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{})}}+{\displaystyle \frac{5n^{}}{2\xi _{1}^{}{}_{}{}^{3}(\overline{\mathrm{\Theta }}_{0,\xi ^{}}(\xi _1^{}))^2}}{\displaystyle _0^{\xi _1^{}}}𝑑\xi ^{}\xi _{}^{}{}_{}{}^{2}\left\{\overline{\mathrm{\Theta }}_0^n^{}{}_{}{}^{(6)}\overline{\psi }_{0}^{}+{\displaystyle \frac{1}{10}}n^{}\overline{\mathrm{\Theta }}_0^{n^{}1}({}_{}{}^{(3)}\overline{\psi }_{2}^{})^2\right\}\right],`$ (I15)
where we have used Eqs. (I3) and (I4).
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# A Uniform Analysis of the Ly-𝛼 forest at 𝑧=0-5: II. Measuring the mean intensity of the extragalactic ionizing background using the proximity effect
## 1. Introduction
The spectra of quasi-stellar objects (QSOs) blueward of Ly$`\alpha `$ emission show a prodigious number of absorption lines primarily due to Ly$`\alpha `$ absorption by intervening neutral hydrogen along the line of sight to the QSO (Lynds, 1971; Sargent et al. 1980; Weymann, Carswell, & Smith, 1981). In models of these Ly$`\alpha `$ systems, they are in photoionization equilibrium with a background ultraviolet radiation field. This radiation field is modeled as the integrated emission from QSOs and young galaxies (Bechtold et al. 1987; Miralda-Escudé & Ostriker 1990; Madau 1992, Meiksin & Madau 1993, Madau & Shull 1996, Haardt & Madau 1996, Fardal et al. 1998).
The basic trend in the number density evolution of Ly$`\alpha `$ absorbers (12 $``$ log(N<sub>HI</sub>) $``$ 16 cm<sup>-2</sup>) at high redshift in a given QSO spectrum is a steep power law increase with redshift (Sargent et al. 1980). However, a decrease in the number of lines near the emission redshift of an individual QSO relative to the number of lines expected from the power law distribution has been observed (Carswell et al. 1982, Murdoch et al. 1986; Tytler 1987; Bajtlik, Duncan, & Ostriker 1988, hereafter BDO; Lu, Wolfe, & Turnshek 1991, hereafter LWT; Bechtold 1994, hereafter B94). The simplest explanation for this $`\mathrm{`}\mathrm{`}`$inverse” effect is enhanced ionization of HI in the vicinity of the QSO by ultraviolet photons from the QSO itself. Thus, the name $`\mathrm{`}\mathrm{`}`$proximity effect” is also used. This interpretation, along with the assumptions about the spectrum of the background and the photoionization of the nearby intergalactic medium (IGM) by the QSOs, allows for a measurement of the mean intensity of the ionizing background at the Lyman limit of hydrogen (BDO), which can be compared to estimates of the integrated emission from QSOs.
BDO find that $`J(\nu _0)`$ is approximately constant over a redshift range $`1.7<z<3.8`$. Expressing $`J(\nu _0)`$ as $`J_{21}`$ $`\times `$ 10<sup>-21</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>, these authors find $`J_{21}1`$. Their best fit form for the dependence of $`J(\nu _0)`$ on redshift rules out a luminosity dependent cutoff in the QSO luminosity function (BWLM; Schmidt, Schneider, & Gunn 1986). Several other authors have carried out this analysis on other data sets (LWT, Kulkarni & Fall 1993, hereafter KF93, B94, Williger et al. 1994, Cristiani et al, 1995, Giallongo et al. 1996.) LWT found $`J_{21}1`$ for $`1.7<z<3.8`$; and B94 found a value consistent with this, $`J_{21}3`$ for $`1.6<z<4.1`$. Srianand & Khare (1996) compile a sample of 69 QSOs from the literature (54 from B94) and obtain $`J_{21}6`$ for $`1.7<z<4.1`$. Williger et al. (1994) find a lower value, $`J_{21}0.10.3`$ for z $``$4.2 from a single z=4.5 QSO; and Giallongo et al. (1996) find $`J_{21}0.5`$ for $`2.8<z<4.1`$, based on high resolution spectra of six QSOs. KF93 re-express the BDO formalism using maximum likelihood techniques and use this to derive $`J(\nu _0)`$ 6 $`\times `$ 10<sup>-24</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup> for z $``$ 0.5 from the sample of Bahcall et al. (1993).
In this paper, a homogeneous sample of QSO spectra is used to measure the mean intensity of the ionizing background via the standard proximity effect analysis and the maximum likelihood analysis of KF93. We have presented new Lyman $`\alpha `$ forest data for 39 objects with $`1.9<z_{em}<2.5`$ in Paper I; and supplemented this sample with 59 objects from the literature, and one QSO from the Hamburg/CfA Bright Quasar Survey (Dobrzycki, Engels, & Hagen 1999). The spectra comprising our dataset are of moderate resolution, $``$1 $`\AA `$ FWHM; but this is not a disadvantage, as this analysis requires good absorption line statistics and therefore many QSO sight lines. This is difficult to achieve at high resolution, the major reason for using a large set of moderate resolution spectra such as this one.
A major uncertainty in the proximity effect analysis is in the systemic redshifts of the QSOs. Redshifts measured from low ionization permitted lines (e.g. Balmer lines or Mg II) or forbidden lines (e.g. \[OIII\] $`\lambda \lambda `$4959,5007) lines have been shown to be redshifted with respect to Ly$`\alpha `$ and C IV emission by up to $``$250 km s<sup>-1</sup> (Boroson & Green 1992, Laor et al. 1995). B94 found that increasing the values of the QSO redshifts by 1000 km s<sup>-1</sup> caused the best fit value of $`J_{21}`$ to be decreased by a factor of 3. We therefore obtained emission line spectra for several objects in our sample in order to examine redshift differences between Ly$`\alpha `$ and \[OIII\] $`\lambda \lambda `$4959,5007, Mg II or Balmer emission. We investigate the effect of these shifts on the value of $`J(\nu _0)`$ derived.
In Section 2 the spectra used to measure Lyman limit fluxes for some of our sample objects and spectra used to perform systemic redshift measurements of several QSOs in our sample are presented. In Section 3 the techniques used to measure the mean intensity of the HI ionizing background at the Lyman limit are discussed, which includes discussions of the QSO systemic redshifts in Section 3.5 and the HI ionization rate measurements in Section 3.6. In Section 4 we present Lyman $`\alpha `$ forest simulations and investigate curve-of-growth effects in the proximity effect analysis. Section 5 is a summary and discussion of our results and the possible systematic effects entering into our analysis.
Throughout this paper, we assume values of 75 km s<sup>-1</sup> Mpc<sup>-1</sup> for H<sub>0</sub> and 0.5 for $`q_0`$.
## 2. Data
### 2.1. Spectrophotometry
Spectrophotometry of 12 sample objects in the spectral region between Ly$`\alpha `$ and C IV emission was obtained at the Steward Observatory (SO) Bok Telescope with the Boller and Chivens (B&C) Spectrograph and the 1200 $`\times `$ 800 CCD on the nights of September 22, 1992, November 29, 1994, and March 28, 1995. Observations were made with a 400 l mm<sup>-1</sup> grating with $`\lambda _b`$=4889 $`\AA `$ in the first order and a 4.5$`\mathrm{}`$ slit. Spectrophotometry of the object 1422+231 was obtained at the SO B&C using a 600 l mm<sup>-1</sup> grating with $`\lambda _b`$=6681 $`\AA `$ in the first order and a 1.5$`\mathrm{}`$ slit on April 22, 1996; and the object 1603+383 was observed by A.D. as part of the Hamburg/CfA Bright Quasar Survey on July 4, 1995 with the Fred Lawrence Whipple Observatory 1.5-meter Tillinghast telescope and FAST spectrograph, using a 300 l mm<sup>-1</sup> grating with $`\lambda _b`$=4750 in the first order and a 3$`\mathrm{}`$ slit. See Table 1 for a summary.
All observations except those of 1422+231 and 1603+383 were made with the slit set at the parallactic angle. This should not seriously effect the spectrophotometry of 1603+383 as it was observed at a small airmass. Additionally, however, the observation of 1422+231 was made with a slit width that is somewhat small for optimal spectrophotometry. In any case, as discussed further below, both 1422+231 and 1603+383 are excluded from the proximity effect analysis due to the fact that 1422+231 is a gravitational lens and the presence of associated absorption in the spectrum of 1603+383. Any small errors in the spectrophotometry of the 74 objects used in the proximity effect analysis should not significantly bias the results of this work.
Object spectra were bias corrected and extracted using standard IRAF packages using He-Ne-Ar and quartz calibration exposures taken at each telescope position to perform the wavelength calibration and to correct for pixel-to-pixel variations, respectively. The data were then flux calibrated using standard star exposures. The column density of Galactic neutral hydrogen along the line of sight to each object was found using the program COLDEN, made available by J. M<sup>c</sup>Dowell; and the spectra were corrected for the Galactic reddening calculated from the relation $`N_{HI}/E(BV)`$ = 4.8 $`\times `$ 10<sup>21</sup> atoms cm<sup>-2</sup> magnitude<sup>-1</sup> (Bohlin et al. 1978). The spectra and the power law continuum fits are shown in Figure 1.
### 2.2. QSO Systemic Redshifts
For the present absorption line sample, the QSO narrow emission lines discussed above all lie redward of $``$7600 $`\AA `$, and into the near infrared. Spectra of four objects in this sample were obtained at the MMT with the infrared spectrometer FSpec (Williams et al. 1993) on May 20, 1994 (1207+399 and 1422+231) and April 1, 1996 (1408+009, and 1435+638) using a 75 l mm<sup>-1</sup> grating and a 1.2$`\mathrm{}`$ slit giving a resolution of $``$34 $`\AA `$ in the K band. A series of exposures of each object was taken. Between each exposure, the object was moved along the slit. The total integration time is listed in Table 2. One object, 0836+710, was observed on March 28, 1995 with the B&C, the 1200x800 CCD, a 300 l mm<sup>-1</sup> grating with $`\lambda _b`$=6693 $`\AA `$ in the first order, and a 4.5$`\mathrm{}`$ slit. Infrared spectra of eight objects in this sample, 0000-263, 0014+813, 0636+680, 0956+122, 1159+124, 1208+101, 2126-158, were obtained using FSpec, OSIRIS on the Cerro Tololo Inter-American Observatory 4 m telescope, and CRSP on the Kitt Peak National Observatory 4 m telescope (Kuhn 1996). A summary of these observations is given in Table 2 and the spectra are displayed in Figure 2.
## 3. Ly$`\alpha `$ Forest Statistics for $`z_{abs}z_{em}`$: The Proximity Effect
### 3.1. Spectrophotometry
In order to perform the proximity effect analysis, the flux of each QSO at the Lyman limit is needed. The spectrophotometry data discussed above was used for this purpose. A power law of the form $`f_\nu \nu ^\alpha `$ was fit to the continua of these objects. The straight line fit to log($`f_\nu `$) vs. log($`\nu `$) was done using a robust estimation technique; and emission lines found by visually inspecting the spectrum were excluded from the points used in the fit. The measured flux at 1450 $`\AA `$ and the value of $`\alpha `$ derived from this fit were used to determine the flux at 912 $`\AA `$. For the objects we did not observe, we proceed as follows. If a flux measurement at a rest UV wavelength other than 912 $`\AA `$ exists along with a published spectral index, we use these to extrapolate to the Lyman limit. If no spectral index is available, we use the value of $`0.46`$ (Francis 1996). The object 2134+004 has a variable continuum (Perez et al. 1989, Corbin 1992). Therefore, although we have spectrophotometry from our own observations of this object, we take the flux measurement of these authors from their averaged spectrum produced from observations made over several months. We use this with the spectral index we derive to extrapolate to 912 $`\AA `$.
If no rest UV spectrophotometry of an object exists, we estimate $`f_\nu `$ at 5500 $`\AA `$ (observed) from the V magnitude given in Table 1 of Paper I with an extinction correction applied. The extinction correction was calculated using the column density of neutral hydrogen from COLDEN and the Seaton (1979) re-normalization of the composite UV-optical reddening curve of Nandy et al. (1975, and references therein). A rest-frame composite QSO spectrum (Zheng et al. 1997) with an arbitrary flux scale was redshifted by the appropriate amount for each object. The flux in the V filter was calculated by convolving this spectrum with the V filter transmission as a function of wavelength. A scaling factor was calculated so that when the redshifted QSO composite spectrum was multiplied by this factor, the resulting magnitude matched the magnitude listed in Table 1 of Paper I. The flux at 1450 $`\AA `$ was then taken from this scaled spectrum and this flux was extrapolated to the Lyman limit using the spectral index given in Table 3. A zero point flux density for the V filter of 3.81 $`\times `$ 10<sup>-20</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> (Johnson 1966) was used.
The asterisks in Table 3 mark QSOs which are known lenses or which show associated absorption in their spectra. Associated absorption is defined to be any Lyman $`\alpha `$ absorption within $``$5000 km s<sup>-1</sup> of the QSO redshift which also shows metal lines. (See Paper I for a description of the metal line systems identified in each QSO spectrum.) These objects were excluded from the proximity effect analysis on the grounds that gas associated with the QSO or QSO host galaxy is not part of the general intergalactic medium and bulk motions within this gas may skew the results. The spectrophotometric properties adopted for the 59 QSOs from the literature are listed in Table 5 of B94.
### 3.2. Number of Lines with $`z_{abs}z_{em}`$
The first method we use to demonstrate the proximity effect is to compare the number of lines predicted if there was no effect from the equation
$$\frac{d𝒩}{dz}=𝒜_0(1+z)^\gamma .$$
(1)
with the number of lines counted in the spectrum as a function of distance from the QSO,
$$\mathrm{\Delta }𝒩=𝒩_{pred}𝒩_{obs}.$$
(2)
The number of lines predicted is found by integrating Equ. 1,
$$𝒩_{pred}=\frac{𝒜_0}{\gamma +1}((1+z_{max})^{\gamma +1}(1+z_{min})^{\gamma +1}).$$
(3)
The bins in luminosity distance from the QSO are defined according to the relation,
$$\mathrm{\Delta }R=1687.5\frac{\mathrm{\Delta }z}{(1+z_{em})^{5/2}}h^1\mathrm{Mpc}.$$
(4)
We use $`h=0.75`$. Figure 3 plots the distribution in z and Lyman limit luminosity of the QSOs in our sample.
The dataset was divided into low luminosity and high luminosity subsamples at log\[$`L(\nu _0)`$\]=31.1, such that there were equal numbers of objects in each subsample. The Lyman limit luminosity of each object was calculated according to the expression
$$L(\nu _0)=4\pi d_L^2\frac{f(\nu _0)}{(1+z_{em})}$$
(5)
where the luminosity distance to the QSO, $`d_L`$ is given by
$$d_L=\frac{c\{q_0z+(q_01)[(1+2q_0z)^{1/2}1]\}}{q_0^2H_0}$$
(6)
for $`q_0>0`$. In this paper, we use a value of 0.5 for $`q_0`$. Figure 4 plots the fractional deficit of lines, $`(𝒩_{pred}𝒩_{obs})/𝒩_{pred}`$, for the total sample and the high and low luminosity subsamples.
For the total sample, a 5.5$`\sigma `$ deficit of lines is found in the 0-1.5 $`h^1`$ Mpc bin. The low luminosity subsample shows a deficit of lower significance (3.6$`\sigma `$) than the high luminosity subsample (4.6$`\sigma `$). These deficits are expected for a proximity effect caused by enhanced ionization of HI from the QSO flux; and the marginally higher significance for high luminosity objects further suggests that this picture is legitimate.
### 3.3. Photoionization Model
We follow the formalism outlined in BDO to calculate a value of the mean intensity of the ionizing background in the redshift range $`1.7<z<3.4`$. The column density of a Ly$`\alpha `$ absorber in the immediate vicinity of a QSO will be modified from the value that it would have if the QSO were not present. The amount by which the column density of HI will be reduced due to ionization by UV photons from the QSO is given by
$$N=N_0(1+\omega )^1$$
(7)
where $`N`$ is the observed column density of the absorber, and $`N_0`$ is the column density that the absorber would have if the QSO were absent. The column density distribution of the general Ly$`\alpha `$ absorber population was been shown to follow a power law over several orders of magnitude in column density,
$$𝒩N^\beta ,$$
(8)
which, for a fixed limiting column density, $`N_{thr}`$, (corresponding to the limiting rest equivalent width) can be integrated to give the total number of lines with column densities equal to or larger than the limiting value, $`𝒩(NN_{thr})=N_{thr}^{(\beta 1)}`$. Thus, a proximity effect-corrected redshift distribution for a fixed rest equivalent width threshold can be derived:
$$\frac{d𝒩}{dz}=𝒜_0(1+z)^\gamma [1+\omega (z)]^{(\beta 1)}$$
(9)
where $`\omega `$ represents a flux-scaled distance of each cloud from the QSO
$$\omega =\frac{F^Q(\nu _0)}{4\pi J(\nu _0)}.$$
(10)
Here, $`F^Q(\nu _0)`$ is the Lyman limit flux density due to the QSO at the position of a given absorber,
$$F^Q(\nu _0)=\frac{L(\nu _0)}{4\pi r_L^2}$$
(11)
where $`r_L`$ is the luminosity distance between the QSO and the absorber. We remove the dominant dependence of the line density on redshift by introducing a coevolving coordinate, $`X_\gamma `$, given by
$$X_\gamma =(1+z)^\gamma 𝑑z.$$
(12)
If no proximity effect existed, the number of lines per coevolving coordinate would be expressed as
$$d𝒩/dX_\gamma =𝒜_0.$$
(13)
In this analysis, we use a value for $`\beta `$ of 1.46 from of Hu et al. (1995) based upon high S/N, high resolution spectra of four QSOs at z $``$ 3, consistent with the value of 1.4 found by Dobrzycki & Bechtold (1996), hereafter DB96, from simulations of Ly$`\alpha `$ forest spectra in QSOs at z $``$ 3. The value of this parameter is an important factor in the ionization model. B94 found that changing the adopted value of $`\beta `$ from 1.7 to 1.4 caused the derived value of $`J_{21}`$ to decrease by a factor of $``$3. Giallongo et al. (1996) find that a double power law provides a better fit to the observed column density distribution in their high resolution spectra than a single power law. The form of their double power law consists of a break at $`N_{HI}=10^{14}`$ cm<sup>-2</sup> and values of $`\beta `$ above and below this break of 1.8 and 1.4 respectively. For this analysis, however, we will use a single power law, as the data of Hu et al. (1995) do not require the double power law form.
The procedure consists of assuming a form for $`J(\nu _0)`$ as a function of z, dividing the lines into the appropriate $`\omega `$ bins, and finding the parameters of the assumed form of $`J(\nu _0)`$ that gives the lowest $`\chi ^2`$ between the binned data and the ionization model. Since no work to date has shown that $`J(\nu _0)`$ evolves significantly with redshift over the range of our sample objects, we will treat the case that $`J(\nu _0)`$ is constant over the redshift range of the data.
Figure 5 plots $`\chi ^2`$ with respect to the constant $`J(\nu _0)`$ photoionization model versus log\[$`J(\nu _0)`$\] and Figure 6 plots the coevolving number density versus $`\omega `$ for the lowest $`\chi ^2`$ value of $`J(\nu _0)`$ for each subsample. The results of this analysis are summarized in Table 4 and are discussed in more detail in Section 5.
### 3.4. Maximum Likelihood Analysis
In addition to the standard BDO analysis, we also used a maximum likelihood method outlined by KF93 to measure the extragalactic ionizing background in a manner that avoids binning of the data. One constructs a likelihood function of the form
$$L=\underset{a}{}f(N_a,z_a)\underset{Q}{}exp[_{z_{min}^Q}^{z_{max}^Q}𝑑z_{N_{min}^Q}^{\mathrm{}}f(N,z)𝑑N],$$
(14)
where the subscripts $`a`$ and $`Q`$ refer to absorbers and QSOs and where $`f(N,z)`$ is the standard equation for the distribution of Lyman $`\alpha `$ absorbers in column density and redshift,
$$f(N,z)=AN^\beta (1+z)^\gamma [1+\omega (z)]^{(\beta 1)}.$$
(15)
The parameter $`\omega `$ is defined as above, but here, the normalization in terms of $`𝒜_0`$ in Equation 1 is given by $`𝒜_0(N_{lim}/N_0)^{\beta 1}(\beta 1)^1`$. With the exception of the case in which a variable threshold is used, $`N_{min}`$ for each QSO is the column density which, according to the curve-of-growth adopted (see KF93), corresponds to an equivalent width of 0.32 $`\AA `$, 2.62 $`\times `$ 10<sup>14</sup> cm<sup>-2</sup>.
Instead of using the method outlined by KF93 whereby the parameters $`A`$,$`\beta `$,$`\gamma `$, and $`J(\nu _0)`$ are all found by minimizing -ln(L) where L is given by the likelihood function above, we chose to take the parameter $`\gamma `$ from a separate maximum likelihood solution to Equation 1 (see Paper I.) Since our spectra are more highly blended than the low redshift data used by KF93, we choose not to determine $`\beta `$ directly from our data using line equivalent widths and the curve-of-growth and instead adopt a value found from high resolution spectra. As described in the previous section, we take $`\beta `$ to be 1.46 (Hu et al. (1995) and solve for $`A`$ by requiring $`f`$(N,z) to give the observed number of lines in the regions of the QSO spectra unaffected by the proximity effect.
We ran two tests on this set of algorithms. The first of these was to attempt to reproduce the results of KF93 with the dataset they used from Bahcall et al. (1993) . Next, we used a high redshift subsample of our complete dataset, the DB96 sample, to compare the results of the maximum likelihood analysis and the BDO analysis to each other and to independent checks on these values (B94, Giallongo et al. 1996).
We were able to reproduce the results of KF93. Using their Sample 2, the Bahcall et al. (1993) sample minus one BAL QSO, PG 0043+039, we obtain ($`\gamma `$, $`\beta `$, log(A)) = (0.23, 1.47, 7.74) and log\[$`J(\nu _0)`$\]=-23.0$`{}_{0.6}{}^{}{}_{}{}^{+0.7}`$ for $`b=35`$ km s<sup>-1</sup>. These agree with the values they find, ($`\gamma `$, $`\beta `$, log$`(A)`$)= (0.21, 1.48, 7.74), and the errors in these values, $`\sigma _\gamma 0.06`$ $`\sigma _\beta 0.05`$, and $`\sigma _{log(A)}0.1`$. Their result for log\[$`J(\nu _0)`$\] for this sample is -23.3$`{}_{0.5}{}^{}{}_{}{}^{+0.7}`$.
The high redshift subsample we created consisted of 518 lines from the 15 objects from DB96 that do not show associated absorption. The QSOs have redshifts between 2.52 and 3.38. Using our maximum likelihood program to solve for the Ly$`\alpha `$ forest statistics, we find $`\gamma =1.926\pm 0.656`$, and log$`(A)=7.03`$ for $`N_{min}`$=2.6 $`\times `$ 10<sup>14</sup> cm<sup>-2</sup> and $`\beta `$= 1.46. This subsample does give similar results in the BDO and the maximum likelihood cases, log\[$`J(\nu _0)`$\]=$`21.40_{0.69}^{+1.10}`$ and log\[$`J(\nu _0)`$\]=$`21.58_{0.23}^{+0.30}`$, respectively. (See rows 1 and 2 of Table 4.) These values agree well with the Giallongo et al. (1996) result of log\[$`J(\nu _0)`$\]=$`21.30\pm 0.7`$ for z=1.7-4.1.
The software we used for the maximum likelihood analysis uses all regions of the QSO spectra between $`z_{min}`$, specified by the spectral coverage or by Ly$`\beta `$ emission, and $`z_{max}`$, specified by Ly$`\alpha `$ emission. Though it does not count lines associated with identified metal line systems, it does not exclude the regions of the spectrum where these lines lie. To ensure that this does not have a significant effect on our resultant solution for the background, we tested a program that does exclude regions of the spectra in the same way that our BDO-style software does. The change in the result was indeed insignificant; but taking these excluded spectral regions into account and binning the data in the same way the BDO-style software does brings the maximum likelihood and the BDO method results into excellent agreement.
Figure 7 plots the log of the ratio of the likelihood function to the maximum value versus log\[$`J(\nu _0)`$\]; and Figure 8 plots the coevolving number distribution of Ly$`\alpha `$ lines with respect to $`\omega `$ just as in Figure 6. The results of this analysis are also summarized in Table 4 and discussed further in Section 5.
### 3.5. Systemic QSO Redshifts
One of the major uncertainties in the proximity effect analysis is in the systemic redshifts of the QSOs. If the true redshift of a QSO is higher than the value used in the analysis, any given cloud is further away from the QSO than assumed. Hence, the influence of the QSO at this cloud is less than inferred and the value of $`J(\nu _0)`$ in reality is lower than the one derived.
For the data presented in Figure 2, an average of several cursor settings at the peak of the emission line was used to determine the line centers. More detailed fits were not done as our purpose lies mainly in determining if any gross shifts between Ly$`\alpha `$ and the Balmer lines/\[OIII\]/Mg II exist for our data; but we found no significant difference between this method and making Gaussian fits to the upper 50% of the emission line profiles.
Ly$`\alpha `$ redshifts were measured from the absorption line spectra when the entire Ly$`\alpha `$ profile was observed, in the same way as was done for the Balmer, \[OIII\], and Mg II lines. Table 5 lists the adopted best redshift value for each emission line for each object supplementing our measurements with measurements from the literature.
Laor et al. (1994) and Laor et al. (1995) found, from a sample of 13 QSO spectra from the Faint Object Spectrograph on Hubble Space Telescope between redshifts of $`z0.16`$ and $`z2.0`$, average velocity shifts between \[OIII\] $`\lambda `$5007 and Ly$`\alpha `$, Mg II, and H$`\beta `$ of $`200\pm 150`$ km s<sup>-1</sup>, $`85\pm 130`$ km s<sup>-1</sup>, and $`75\pm 110`$ km s<sup>-1</sup>, respectively. This agrees with the Corbin & Boroson (1996) result for 48 objects with $`0.03<z<0.77`$. They found mean \[OIII\]-Ly$`\alpha `$ and \[OIII\]-H$`\beta `$ shifts of $`191\pm 101`$ km s<sup>-1</sup> and $`75\pm 57`$ km s<sup>-1</sup>. Thus, Ly$`\alpha `$ is blueshifted with respect to \[OIII\] by $``$200 km s<sup>-1</sup>, while Mg II and H$`\beta `$ are marginally redshifted with respect to \[OIII\]. Tytler & Fan (1992) find a mean \[OIII\]-H$`\beta `$ shift of $`15\pm 37`$ km s<sup>-1</sup> from 8 QSOs with redshifts between $``$0.3 and $``$0.6 and conclude that both Balmer lines and narrow forbidden lines give redshifts within 100 km s<sup>-1</sup> or less of the QSO systemic redshift. They then find a blueshift of Mg II with respect to \[OIII\]/H$`\beta `$ for 100 QSOs of $`101\pm 47`$ km s<sup>-1</sup> which they use as a secondary systemic redshift zero point in their analysis of a large QSO sample. The magnitude of the blueshift of Ly$`\alpha `$ with respect to \[OIII\]/H$`\beta `$ that they derive is $`172\pm 17`$ km s<sup>-1</sup>. The data of Nishihara et al. (1997) for five QSOs at $`z1.5`$ show a negligible redshift of Mg II with respect to \[OIII\], $`31\pm 411`$ km s<sup>-1</sup>. However these five objects show a somewhat larger redshift of H$`\beta `$ with respect to \[OIII\] $`\lambda `$5007, equalling $`260\pm 522`$ km s<sup>-1</sup>, consistent with the fact that these objects have high luminosities. M<sup>c</sup>Intosh et al. (1999b) use the near-infrared spectra of QSOs at 2.0 $``$ z $``$ 2.5 presented in M<sup>c</sup>Intosh et al. (1999a) to examine the redshift differences between \[OIII\] and H$`\beta `$. They supplement their data with data from the literature to measure the redshift differences between \[OIII\] and Mg II. They find that on average, H$`\beta `$ is redshifted relative to \[OIII\] by 520 $`\pm `$ 80 km s<sup>-1</sup> for 21 of their sample objects, while Mg II lies within 50 km s<sup>-1</sup> of the redshift of \[OIII\] for 12 sample objects.
For our sample, we find that Ly$`\alpha `$ is blueshifted with respect to \[OIII\] $`\lambda `$5007 by $`382\pm 1160`$ km s<sup>-1</sup> for 19 QSOs. Mg II emission is blueshifted by an average of $`338\pm 901`$ km s<sup>-1</sup> with respect to \[OIII\] on the basis of seven measurements. We find that H$`\beta `$ is redshifted by $`642\pm 740`$ km s<sup>-1</sup> with respect to \[OIII\] on the basis of five measurements. Including three H$`\alpha `$ redshifts listed in Table 5 with these H$`\beta `$ redshifts leads to a $`507\pm 615`$ km s<sup>-1</sup> redshift of Balmer lines with respect to \[OIII\]. This shift is larger than that discussed above for low reshift QSOs. However, it is consistent with the Nishihara et al. (1997) H$`\beta `$ shift for high luminosity QSOs. Combining our data with that of these authors, we find that Mg II is blueshifted with respect to \[OIII\] by $`184\pm 735`$ km s<sup>-1</sup>. Including the data of M<sup>c</sup>Intosh et al. (1999b) that is not already in our sample gives a blueshift of $`95\pm 603`$ km s<sup>-1</sup>. Similarly, combining our data with that of Nishihara et al. (1997), we find that H$`\beta `$ is redshifted with respect to \[OIII\] by $`451\pm 636`$ km s<sup>-1</sup>. After supplementing this combined data set with the data of M<sup>c</sup>Intosh et al. (1999b), the redshift becomes $`379\pm 516`$ km s<sup>-1</sup>. Finally, combining the data of M<sup>c</sup>Intosh et al. (1999b) with ours gives a Ly$`\alpha `$ blueshift of $`418\pm 920`$ km s<sup>-1</sup> with respect to \[OIII\].
As has been noted in previous work, the standard error in the mean velocity shifts is quite large, on the order of or exceeding the value of the shift itself. We estimate that the wavelength calibration errors in our data contribute a $``$10-30 km s<sup>-1</sup> error in the derived redshifts; and the spread in different redshift measurements of the same species (e.g. Balmer lines or \[OIII\] $`\lambda `$4959 and $`\lambda `$5007) for the same object is typically 100-200 km s<sup>-1</sup>. The observed spreads in the velocity differences of the Ly$`\alpha `$, Mg II, and Balmer emission lines with respect to the QSO systemic redshifts are much larger than this, indicating that it is intrinsic to the QSO population. Figure 9 shows histograms of the emission line redshift differences between \[OIII\] and Ly$`\alpha `$, \[OIII\] and Mg II, and \[OIII\] and Balmer lines. Our results are plotted with those of Laor et al. (1995) and of Nishihara et al. (1997). Our sample shows no well-defined mean \[OIII\]-Balmer shift, just a large scatter in the measurements included. Our sample also shows a large range of \[OIII\]-Ly$`\alpha `$ and \[OIII\]-Mg II shifts with no well-defined mean value. Nonetheless, the mean trend is that the \[OIII\]-Ly$`\alpha `$ shift is different from zero by 1.4$`\sigma `$ for our data, less than the 3.5$`\sigma `$ significance found by Laor et al. (1995). The \[OIII\]-Balmer line shifts for both our data set and for our data combined with that of Nishihara et al. (1997) are more significant, 2.7$`\sigma `$ and 2.8$`\sigma `$ respectively. The \[OIII\]-Mg II shift is consistent with zero in a mean sense, but with large scatter. Thus, though better statistics are desirable, it seems that for these high redshift and relatively high luminosity objects, Balmer lines are not good indicators of the QSO systemic redshift. For the purposes of this study therefore, we treat only the redshifts found from \[OIII\] $`\lambda `$5007 for 19 objects in our sample and Mg II for 16 objects in our sample as systemic QSO redshifts.
### 3.6. The HI Ionization Rate
The HI ionization rate due to a source of UV flux is formally given by the equation:
$$\mathrm{\Gamma }=_{\nu _0}^{\mathrm{}}\frac{4\pi J(\nu )\sigma _{HI}(\nu )}{h\nu }𝑑\nu \mathrm{s}^1.$$
(16)
The calculations of the mean intensity of the ionizing background to date have made a critical assumption, namely that the spectrum of the background and the spectra of the individual QSOs are identical. This allows the expression $`\omega =\mathrm{\Gamma }^Q/\mathrm{\Gamma }^{bg}`$ to reduce to the ratio of the Lyman limit flux density of the QSO, $`J^Q(\nu _0)`$, to that of the background, $`J^{bg}(\nu _0)`$, for each line (BDO). Since the IGM reprocesses the radiation emitted from QSOs, this is not strictly true (Miralda-Escudé & Ostriker 1990, Madau 1991,1992, Meiksin & Madau 1993, Haardt & Madau 1996, Fardal et al. 1998). Furthermore, the value of $`\mathrm{\Gamma }^{bg}`$ is of particular interest as it can be used to infer the value of $`\mathrm{\Omega }_b`$ by comparing the distribution of flux decrements in high resolution QSO spectra to Lyman $`\alpha `$ forest simulations (Rauch et al. 1997). Therefore, we repeat the standard BDO analysis without making this assumption, ie. using $`\omega =\mathrm{\Gamma }^Q/\mathrm{\Gamma }^{bg}`$ and solving for the HI ionization rate from the metagalactic background radiation. The ionization rate for each QSO was calculated using Equation 16, where $`\sigma _{HI}(\nu )=6.3\times 10^{18}(\frac{\nu _0}{\nu })^3`$ cm<sup>2</sup> and where $`J^Q(\nu )=J^Q(\nu _0)(\frac{\nu }{\nu _0})^\alpha `$. For each QSO, $`J^Q(\nu _0)`$ is the same value used in the standard analysis used to solve for $`J^{bg}(\nu _0)`$, and $`\alpha `$ is given in Table 3. For some objects, no $`\alpha `$ listed in this table and a value of 0.46 was used, as described in Section 3.1. As before, the best value will be the one that gives the lowest $`\chi ^2`$ between the model with $`\beta `$=1.46 and the binned data. We use the narrow line redshifts for each QSO discussed above and add 400 km s<sup>-1</sup> to each QSO redshift measured from the Lyman $`\alpha `$ emission line.
Haardt & Madau (1996) present a Gaussian fit to their model for the evolution of $`\mathrm{\Gamma }`$ with redshift,
$$\mathrm{\Gamma }=A(1+z)^Bexp[(zz_c)^2/S]$$
(17)
that agrees with their detailed model for the background to within 10% over the range $`0<z<5`$. The best fit parameters they derive for $`q_0`$=0.5 are $`A`$=6.7 $`\times `$ 10<sup>-13</sup> s<sup>-1</sup>, $`B`$=0.43, $`z_c`$=2.30, and $`S`$=1.95. Fardal et al. (1998) fit their model for the background with the parameter sets $`A`$=5.6 $`\times `$ 10<sup>-13</sup> s<sup>-1</sup>, $`B`$=0.60, $`z_c`$=2.22 , and $`S`$=1.90 and $`A`$=1.26 $`\times `$ 10<sup>-12</sup> s<sup>-1</sup>, $`B`$=0.58, $`z_c`$=2.77 , and $`S`$=2.38 for the Q1 and Q2 luminosity functions, of Pei (1995) respectively. Incorporating this expression for $`\mathrm{\Gamma }(z)`$ with these three different sets of parameters into the BDO style analysis allows us to determine which of these models fits our data best. The results are listed in Table 6 and are discussed in greater depth below in Section 5.
## 4. Simulations and the Curve of Growth
Simulated Lyman $`\alpha `$ forest spectra for the DB96 sample only were produced using the software described in that paper. The simulation input $`\gamma `$ was changed slightly to reflect the maximum likelihood value found by the software used in the analysis described in Paper I. The normalization was chosen to give matching amounts of total absorption in the real and simulated spectra. The parameters used were $`\gamma =2.069`$, $`𝒜_0=4.835`$, $`\beta =1.46`$, log(N$`_{\mathrm{HI}_{\mathrm{min}}}`$) = 13.0, log(N$`_{\mathrm{HI}_{\mathrm{max}}}`$) = 16.0, $`<b>`$= 28.0 km s<sup>-1</sup>, $`\sigma _b`$= 10.0 km s<sup>-1</sup>, and $`b_{\mathrm{cut}}`$= 20.0 km s<sup>-1</sup>.
The proximity effect was included in these simulations by simply modifying each cloud’s column density according to equations 7 and 10. The value of log\[$`J(\nu _0)`$\] from the BDO type analysis on the DB96 sample is -21.40$`{}_{0.69}{}^{}{}_{}{}^{+1.1}`$. Values of -19.0, -20.0, -21.3, -22.0, and -23.0 for log\[$`J(\nu _0)`$\] were input and the analyses described above were used to recover that $`J(\nu _0)`$. Two examples of the simulated spectra are shown in Figure 10.
The analysis considers all lines above a fixed equivalent width threshold of 0.32 $`\AA `$. Thus, as the column densities of lines are modified by the QSO flux from their expected values in the absence of the proximity effect, the equivalent widths of the lines will change according to the curve-of-growth. If a line is saturated, changing its column density will have little effect on its equivalent width, since it lies on the flat part of the curve-of-growth where $`W\sqrt{log(N)}`$. This will mean that for a given equivalent width cutoff in the data, this line will not drop out of the sample as the proximity effect is turned on in the simulations. Since the line deficit will be less than expected for a given input value of $`J(\nu _0)`$, the proximity effect will appear less pronounced and the true $`J(\nu _0)`$ will be overestimated. We found this to be the case from our simulations. As Figure 11 illustrates and Table 7 summarizes, though the values of $`J(\nu _0)`$ recovered from the simulated data were usually consistent with the input values within the 1$`\sigma `$ confidence limits, they were systematically larger than the input values by up to a factor of 3. The largest input values of log\[$`J(\nu _0)`$\], -19.0 and -20.0, give the largest discrepancy between this input value and the log\[$`J(\nu _0)`$\] recovered from the BDO analysis performed on the simulated spectra. The smallest input value of log\[$`J(\nu _0)`$\], -23.0, gives the smallest discrepancy between the input and recovered values. However, the 1$`\sigma `$ confidence limits on this fit are also relatively small, making it the only trial which does not recover the input log\[$`J(\nu _0)`$\] to within those limits.
To demonstrate the effect, Figure 12 compares the simulated line equivalent widths with and without the proximity effect included. The column density of each line from the simulated spectra line lists with no proximity effect were modified according to equations 7 and 10. Figures 12(a-e) plot the nonproximity effect rest equivalent width $`W_{noPE}`$ versus the ratio of the proximity effect and nonproximity effect equivalent widths, $`W_{PE}/W_{noPE}`$. The solid line delineates the detection threshold for the lines in the list for which the proximity effect is included, $`W_{PE}`$= 0.32 $`\AA `$. Absorption lines that fall above this line were not removed from the sample when the proximity effect was turned on, while those below it disappeared. For a given set of QSOs with fixed Lyman limit lumosities, such as this one, the proximity effect signature in their spectra will become less pronounced as the ambient UV background increases. Therefore, as log\[$`J(\nu _0)`$\] increases from -23.0 to -19.0, the magnitude of the proximity effect decreases, and the proximity effect line list differs less and less from the nonproximity effect line list.
## 5. Results and Discussion
Table 4 lists the best fit values of $`J(\nu _0)`$ found for various subsamples of this dataset using both the canonical BDO and the maximum likelihood methods. For the BDO method, the 1$`\sigma `$ confidence limits are found from a $`\mathrm{\Delta }\chi ^2`$ of 8.18 for 7 degrees of freedom. The maximum likelihood method 1$`\sigma `$ confidence limits derive from the fact that ln$`(L/L_{max})`$ is distributed as $`\chi ^2/2`$. The total sample consisting of 74 QSOs with all QSO redshifts based on the Ly$`\alpha `$ emission line gives a best fit value of log\[$`J(\nu _0)`$\] of -20.90$`{}_{0.48}{}^{}{}_{}{}^{+0.61}`$ for the BDO analysis and -20.83$`{}_{0.20}{}^{}{}_{}{}^{+0.23}`$ for the maximum likelihood analysis.
As the results in Table 4 demonstrate, using narrow line redshifts for 35 of the 74 QSOs for which they have been directly measured and Ly$`\alpha `$ redshifts for the rest does not change the result. However, when 400 km s<sup>-1</sup> is added to the Ly$`\alpha `$ redshifts of the objects with no measured narrow line redshift, a value for log\[$`J(\nu _0)`$\] of -21.15$`{}_{0.43}{}^{}{}_{}{}^{+0.17}`$ is derived using the BDO method and log\[$`J(\nu _0)`$\]=-21.17$`{}_{0.15}{}^{}{}_{}{}^{+0.19}`$ is found using the maximum likelihood method. Recall that the mean blueshift of Ly$`\alpha `$ with respect to \[OIII\] for the 19 objects in this paper with \[OIII\] $`\lambda `$5007 measurements was found to be $``$400 km s<sup>-1</sup>. This decrease in the mean intensity of the background derived when larger QSO redshifts are used is to be expected. (cf. Section 3.5) Because this measurement of the background accounts for the systematic blueshift of the Ly$`\alpha `$ emission line with respect to the systemic redshift of each QSO, we consider it to be our best estimate for the mean intensity of the background at the Lyman limit.
These measurements have been made, however, using a photoionization model with somewhat unrealistic assumptions, particularly that Ly$`\alpha `$ absorbers are isothermal and are composed of pure hydrogen. For clouds with a primordial He abundance and which are in thermal and ionization equilibrium, Using CLOUDY to model the ionization state of absorbers with a metal adundance of 10<sup>-2</sup> solar (Cowie et al. 1995, Tytler & Fan 1994) as a function of $`\omega `$, we find that the neutral fraction, $`\chi `$, is proportional to $`(1+\omega )^{1.21}`$. This implies that
$$\frac{d𝒩}{dz}=𝒜_0(1+z)^\gamma [1+\omega (z)]^{1.21(\beta 1)}$$
(18)
In this scenario, the optimal value found for log\[$`J(\nu _0)`$\] is -21.10$`{}_{0.28}{}^{}{}_{}{}^{+0.53}`$. This value is marginally larger than the value discussed above, found under the assumption of absorbers composed of pure hydrogen; but it is not significantly different, so we conclude that the absence of metals in the BDO model has not drastically affected our measurement of the background.
It is worth noting that 16 objects in our sample of objects with no associated absorption show evidence for damped Ly$`\alpha `$ absorption: 0058+019, 0100+130, 0334-204, 0913+072, 0938+119, 0952+338, 0955+472, 1009+299, 1017+280, 1215+333, 1247+267, 1548+092, 1946+770, 2126-158, 2233+131, and 2320+079. The dust in these systems could cause the intrinsic QSO fluxes to be underestimated. This in turn can cause log\[$`J(\nu _0)`$\] to be underestimated by up to a factor of 3, in addition to the sources of error discussed above (Srianand & Khare 1996). Only six of these objects, 0334-204, 0938+119, 0955+472, 1215+333, 2126-158, and 2233+131, appear in our low luminosity subsample, suggesting that this subsample is not preferentially heavily dust-obscured. Nevertheless, the BDO analysis was performed on all 16 objects exhibiting damped Ly$`\alpha `$ systems; and found the best fit value for log\[$`J(\nu _0)`$\] to be -21.45$`{}_{0.53}{}^{}{}_{}{}^{+0.40}`$, a factor of 1.9$`{}_{1.6}{}^{}{}_{}{}^{+8.1}`$ lower than the value obtained for the sample as a whole. This does not allow us to say anything significant about the presence or absence of dust, so we will neglect its influence.
Dividing our line sample into subsamples of high ($`z>2.5`$) and low ($`z<2.5`$) redshift lines, we find marginal evidence for evolution in the intensity of the background, namely that the maximum likelihood background intensity is lower by a factor of about 1.9$`{}_{1.4}{}^{}{}_{}{}^{+3.9}`$ at lower redshift. The BDO results corroborate this, but with larger uncertainties. The factor by which $`J(\nu _0)`$ is found to be lower at lower redshifts is 2.5$`{}_{2.2}{}^{}{}_{}{}^{+27.7}`$. Gravitational lensing could mimic a trend with redshift with about the same order of magnitude, if the high redshift subsample contains a significant number of unknown lenses. However, Figure 3 suggests little if any trend for high luminosity objects to exist at high redshifts in our sample; and the results of Section 3.2 indicate that the high luminosity objects do show a somewhat stronger proximity effect despite the fact that the measured background at high redshift appears to be higher. No other studies have found this evidence of redshift evolution in the background, so we regard it as tentative; and note that it will be interesting to see in future work if this trend can be shown to be real and if it extends smoothly to the low values of $`J(\nu _0)`$ found at redshifts less than 1.5.
Since we find high luminosity objects do not exist preferentially at high redshift in our sample, a simple test can be done to determine whether or not there is a significant number of lensed objects in our sample. If the high luminosity QSOs are indeed intrinsically more luminous, and the proximity effect is a purely photoionization-driven phenomenon, these objects should show a more prominent proximity effect. The results of Section 3.2 suggest this is the case. However, in the analysis, this larger line deficit is normalized to the higher Lyman limit luminosities of this subsample. Therefore, one expects these objects, when analyzed as a separate subsample, to yield a value of $`J(\nu _0)`$ that is consistent with that found for low luminosity objects if the values of the QSO fluxes are not in error due to lensing. If the high luminosity QSOs, or a subset of them, are lensed objects, then they are not necessarily intrinsically more luminous than the low luminosity QSOs. In this case, the influence of the lensed objects on the surrounding IGM will be overestimated and given the observed line deficit, the background will also be overestimated. Table 4 lists the results obtained for the high and low luminosity subsamples of our data set. The values obtained for these subsamples are equal within the uncertainties. This is consistent with there being no significant effects from gravitational lensing in our sample.
### 5.1. HI Ionization Rate
We tested a range of values for $`\mathrm{\Gamma }`$, the HI ionization rate, using our data. The constant value found to fit the data the best is 1.9$`{}_{1.0}{}^{}{}_{}{}^{+1.2}`$ $`\times `$ 10<sup>-12</sup> s<sup>-1</sup>. This value is in reasonable agreement with that predicted by the QSO-dominated model of Haardt & Madau (1996) at this redshift, 1.0 $`\times `$ 10<sup>-12</sup> s<sup>-1</sup> ($`q_0`$=0.5). Using Equation 16 and $`J^Q(\nu )=J^Q(\nu _0)(\frac{\nu }{\nu _0})^\alpha `$, and assuming global QSO spectral indicies of 0, 1.5, and 2, the ionization rate found from our data corresponds to log\[$`J(\nu _0)`$\]= -21.34, -21.17, and -21.12, respectively.
The parameter set ($`A,B,z_c,S`$) found to give the best fit to the data is that of Fardal et al. (1998) for the Q2 luminosity function (1.2 $`\times `$ 10<sup>-12</sup> s<sup>-1</sup>, 0.58, 2.38, 2.77) which, for a redshift of 2.9 yields an ionization rate of 2.7 $`\times `$ 10<sup>-12</sup> s<sup>-1</sup>, in good agreement with our solution, and within a factor of $``$3 of the Haardt & Madau result. Thus, we conclude that a significant contribution to the ionizing background from stellar UV emission is not required at this redshift.
### 5.2. Curve-of-Growth and Other Systematics
On the basis of a curve-of-growth argument, one might expect that weak lines would show a more prominent proximity effect than strong lines. We have compared the results obtained for a constant equivalent width threshold of 0.32 $`\AA `$ with that obtained for lines with 0.16 $`\AA <W<`$ 0.32 $`\AA `$. Instead of finding a more pronounced proximity effect for the weak lines, we find a less significant deficit of lines within 1.5 $`h^1`$ Mpc of the QSOs. This deficit is 4.0$`\sigma `$, versus 5.5$`\sigma `$ for lines with $`W>`$ 0.32 $`\AA `$. As Table 4 lists, the value of log\[$`J(\nu _0)`$\] recovered from these weak lines is correspondingly higher than that found using strong lines, -20.45$`{}_{0.90}{}^{}{}_{}{}^{+0.37}`$ versus -21.15$`{}_{0.43}{}^{}{}_{}{}^{+0.17}`$. Cooke et al. (1997) point out that this could be the result of a higher degree of blending of weaker lines compared to strong ones in crowded spectral regions. The background flux measurement will be an overestimate because blending will cause fewer individual lines to be resolved further from the QSO. Because the reduction in line density near the QSO will work to reduce line blending, the overall effect of line blending will be to suppress the true magnitude of the proximity effect causing $`J(\nu _0)`$ to be overestimated, by a factor of 4.5 in this case. It is difficult to ascertain whether this effect is as strong for lines with $`W>`$ 0.32 $`\AA `$ or whether the curve-of-growth effect discussed in Section 4 which also causes $`J(\nu _0)`$ to be overestimated, is more important. We expect that for lines with $`W>`$ 0.32 $`\AA `$, the effects of blending are reduced somewhat, while the curve-of-growth effects will remain a factor.
We have addressed many of the systematics which could possibly have affected our analysis. A treatment of the QSO systemic redshifts was integrated directly into our analysis and was found to influence the $`J(\nu _0)`$ found by up to a factor of $``$2. Other effects, such as the influences of metals and dust, which can cause $`J(\nu _0)`$ to be underestimated, and the influences of lensing, line blending, and curve-of-growth effects, which can cause $`J(\nu _0)`$ to be overestimated, were treated after the fact in an attempt to understand the magnitude of their effects on the value of $`J(\nu _0)`$ derived. The CLOUDY simulations discussed above indicate that allowing for an absorber metal abundance of 10<sup>-2</sup> solar has little effect on the value of $`J(\nu _0)`$ found from the data. Dust in intervening absorption systems may have affected our result. Though we were unable to quantify this effect with high confidence, it could be on the order of a factor of 2. We assert that QSO flux amplification due to lensing has not significantly biassed our result; and we attempt to minimize the effect of blending discussed above by using only lines with $`W>`$ 0.32 $`\AA `$. Our result may be susceptible to the curve-of-growth effect we addressed through the simulations in Section 4. In those simulations, we found that the discrepancy between in the input and recovered values of $`J(\nu _0)`$ depended upong the input value of $`J(\nu _0)`$ itself. The magnitude of the discrepancy corresponding to the $`J(\nu _0)`$ we found from the data was a factor of $``$2. We therefore suspect that if our result, log\[$`J(\nu _0)`$\]=-21.15$`{}_{0.43}{}^{}{}_{}{}^{+0.17}`$, is systematically biased in any way, it is an overestimate of the true background and could be in error by up to a factor of 2; though this could be balanced somewhat by systematic error due to dust, which works in the opposite direction.
### 5.3. Comparison with Previous Measurements
Our value for $`J(\nu _0)`$ agrees well with other measurements at similar redshift, with the exception of those of B94 and Fernández-Soto et al. (1995) who both derive values four times larger than our best value for $`J(\nu _0)`$, $``$3 $`\times `$ 10<sup>-21</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>. The measurement of B94 does not take into account QSO systemic redshifts, but she notes that if they are blueshifted with respect to Ly$`\alpha `$ by 1000 km s<sup>-1</sup>, this would lower the derived value of $`J(\nu _0)`$ by a factor of 3, bringing it into reasonable agreement with our result. The Fernández-Soto et al. (1995) value is derived from 3 QSO spectra showing a proximity effect due to foreground QSOs. These authors are not able to place an upper limit on their measurement, but our value of 7.0 $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup> for $`J(\nu _0)`$ is consistent with their lower limit of 1.6 $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>. In fact, when these authors examine the proximity effect in a single QSO spectrum due to the background z$``$2 QSO itself, they derive a value for $`J(\nu _0)`$ of 7.9$`{}_{6.0}{}^{}{}_{}{}^{+23.}`$ $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>, which brings their estimate into better agreement with our values for our total sample and for our low redshift subsample within their large errors. Direct measurements of the background at redshifts $``$3-3.5 have been made using long-slit spectroscopy of fields containing optically thick Ly$`\alpha `$ absorbers in efforts to detect fluorescent emission the absorbers produce from the ionizing radiation field incident upon them (Lowenthal et al. 1990, Martínez-González et al. 1995). Recent Keck telescope observations by Bunker et al. (1998) at 2.5 $`<`$ z $`<`$ 4.1 have achieved a factor of 2-10 higher sensitivity and place a firmer direct limit on the background than previous work. Their null signal in a 90-minute integration with a 3$`\mathrm{}`$ slit sets an upper limit on $`J(\nu _0)`$ of 2 $`\times `$ 10<sup>-21</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>.
Cooke et al. (1997) claim that the value for the background at z$``$4 is between their value of 8.0$`{}_{4.0}{}^{}{}_{}{}^{+8.0}`$ $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup> and that of Williger et al. (1994), 1.0-3.0 $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>. Our best value of $`J(\nu _0)`$ at z$``$3, 7.0 $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>, is in agreement with this, although within the uncertainty there is an allowance for the background to decrease as z approaches 4.
Table 8 lists these various measurements of $`J(\nu _0)`$ in the literature as well as the Kulkarni & Fall (1993) measurement at z$``$0.5. Figure 13 also summarizes the literature measurements of $`J(\nu _0)`$ from z$``$0.5 to z=4.5.
The solid curves in Figure 13 delineate the evolution of the mean background intensity as a function of redshift for global background source spectral indicies between 0 and 2, derived from the Haardt & Madau (1996) model for the HI photoionization rate as a function of redshift discussed in Section 3.6. Over 90% of our sample QSO redshifts lie within the FWHM of the Gaussian in the Haardt & Madau (1996) expression using their best fit parameters. At these redshifts, the Haardt & Madau (1996) curves in Figure 13 are turning over. Nonetheless, for comparison with previous work (B94 and references therein), we investigate a power law redshift dependence of the background intensity:
$$J(\nu _0,z)=J(\nu _0,0)(1+z)^j.$$
(19)
Using the BDO method, we executed a crude grid search in an attempt to constrain the power law index and normalization of this power law. The lowest $`\chi ^2`$ (3.86) between the binned data and the BDO photoionization model for a power law background was achieved by ($`j`$, log\[$`J(\nu _0,0)`$\])= (5.12, -23.97), shown by a dashed line in Figure 14. Extending this solution to low redshift gives log\[$`J(\nu _0),0.5)`$\]=-23.0, in good agreement with the measurement of Kulkarni & Fall (1993). The solution ($`j`$, log\[$`J(\nu _0,0)`$\])= (-4.16, -18.76) gives the next lowest $`\chi ^2`$ (4.91); and though it also implies mean background intensities over four orders of magnitude too high at low redshift, it traces the Haardt & Madau model at high redshift, giving log\[$`J(\nu _0,4.5)`$\]=-21.8, in agreement with the Willigher et al. (1994) measurement. It is also shown by a dashed line in Figure 14. Fitting parabolas to the regions near the $`\chi ^2`$ minima in both $`j`$ and log\[$`J(\nu _0,0)`$\] gives the error in each parameter for both of these solutions, (5.12 $`\pm `$ 1.96, -23.97 $`\pm `$ 1.07) and (-4.16 $`\pm `$ 2.36, -18.76 $`\pm `$ 1.31). B94 found a similarly large range of acceptable solutions: $`7<j<4`$ and $`16.5<`$ log\[$`J(\nu _0,0)`$\] $`<23.0`$. The large error bars on these fits indicate that the power law fit to the data is not well-constrained, due possibly to the fact that the mean intensity of the background is turning over at the redshifts of our sample objects, as the Haardt & Madau (1996) model predicts.
### 5.4. Comparison with Models for the Background
Recent models of the ionizing background include not only the integrated emission from QSOs but also a variety of other physical processes such as star formation in young, high redshift galaxies and attenuation of UV photons by Ly$`\alpha `$ absorbers and Lyman limit systems (Miralda-Escudé & Ostriker 1990, Madau 1991, 1992, Meiksin & Madau 1993, Haardt & Madau 1996,Fardal et al. 1998). Madau & Shull (1996) find that the production of metals in Ly$`\alpha `$ absorbers may also be a significant contributor to the UV background at z $``$ 3. Their contribution may be up to 5 $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>, assuming that the bulk of the metals in the Lyman $`\alpha `$ forest did not form at z $`>>`$ 3, and assuming a Lyman continuum escape fraction, $`f_{esc}`$, from a galaxy of $``$0.25. They note, however, that $`f_{esc}`$ is essentially unconstrained.
Past debate about how the space density of QSOs evolves at high redshift (Koo & Kron 1988; Boyle et al. 1991; Irwin et al. 1991; Schmidt et al. 1991; Warren et al. 1994; Kennefick et al. 1995) has been clarified by recent radio surveys (Hook et al. 1995, 1998; Shaver et al. 1996). This work has demonstrated that the space density of radio-loud QSOs decreases rapidly with redshift beyond z$``$3. Since these surveys are unaffected by any presence of dust in the intervening IGM; and since they confirm the behavior seen in optically selected surveys, they indicate that the QSO population is truly declining at high redshift. Nevertheless, the discovery of QSOs with redshifts greater than 4 has brought better agreement between the values of $`J(\nu _0)`$ found via the proximity effect and the values predicted by the models with QSOs primarily contributing to the background (Madau 1992, Meiksin & Madau 1993, Haardt & Madau 1996).
Madau (1992) and Meiksin & Madau (1993) estimate the QSO UV background by integrating the QSO luminosity function (Boyle 1991) and including the effects of attenuation by hydrogen in the IGM. Their estimates however, 1-3 $`\times `$ 10<sup>-22</sup> ergs s<sup>-1</sup> cm<sup>-2</sup> Hz<sup>-1</sup> sr<sup>-1</sup>, are still somewhat lower than the values derived in this paper. The analysis of Haardt & Madau (1996) takes into account the effects of various atomic processes leading to the production of hydrogen-ionizing photons within Ly$`\alpha `$ absorbers and Lyman limits systems themselves. They conclude that observed QSOs can account for number of ionizing photons required by the proximity effect at z$``$4. These authors find a value of log\[$`J(\nu _0)`$\] equal to $``$-21.4 at z=3, in good agreement with the value found in this paper at similar redshifts. The solid lines in Figure 13 show the results from the Haardt & Madau (1996) model for two different values of the global background source spectral index. The lower and upper curves show the evolution of the background for indicies of 0 and 2 respectively. The literature measurements at redshifts between 1.7 and 3.6 agree well with the model predictions. The z$``$0.5 measurement of Kulkarni & Fall (1993) falls below both model curves and the z=4.5 measurement of Williger et al. (1994) falls above them.
Madau, Haardt, & Rees (1999) revisit the issue of the contribution of high redshift, star-forming galaxies to the ionizing background in light of recent work identifying such objects at 2 $`<`$ z $`<`$ 4. (Steidel et al. 1996a,b; Madau et al. 1996; Lowenthal et al. 1997) They calculate the critical photoionization rate necessary to reionize a non-uniform intergalactic medium as a function of redshift. This is compared to the expected contributions from QSOs and young, star-forming galaxies. There are uncertainties in estimating both of these. The QSO luminosity function at z $`>`$ 4 must be extrapolated from that at lower redshifts. There is also still some debate between theory and observations, eg. of the Hubble Deep Field, on the subject of a population of low-luminosity QSOs (see Madau et al. 1999 and references therein) which could cause the QSO luminosity function to steepen with lookback time, making up for the dearth of observed objects at z $`>`$ 4. The estimation of the galaxy contribution of ionizing photons is limited by poor knowledge of luminosity function of Lyman-break galaxies at z $`>`$ 4 as well as by the lack of constraints upon $`f_{esc}`$. Nevertheless, the results are intriguing. Assuming that $`f_{esc}`$=0.5, Madau et al. (1999) find that the contribution of hydrogen-ionizing photons from star-forming galaxies z $``$ 3 could exceed that from QSOs by a factor of more than 3. However, the QSO contribution at this redshift is sufficient, according to these estimates, to ionize the IGM at this redshift. Deharveng et al. (1997) estimate a much lower $`f_{esc}`$ at z=0, less than 1%, based on the local galaxy H$`\alpha `$ luminosity density. Furthermore, Devriendt et al. (1998) make an independent estimation of the galaxy contribution to $`J(\nu _0)`$ assuming damped Ly$`\alpha `$ systems to be the progenitors of present day galaxies. Their semi-analytic models include a treatment of not only HI absorption of Lyman limit photons in the intervening IGM, but also of HI and dust absorption in the interstellar medium of the photon-producing galaxies. Their results show that constraining $`f_{esc}`$ in this way yields a much lower contribution to the UV background from galaxies at z $`>`$ 2. At z $``$ 2.5, their estimated QSO contribution to $`J(\nu _0)`$ is 3 orders of magnitude greater than that expected from galaxies. Our measurement of $`J(\nu _0)`$ is consistent with the UV background being QSO dominated in the models of both these authors and Haardt & Madau (1996).
In the models of Madau et al. (1999), the scenario changes at z $``$ 3.5. At this redshift, the QSO contribution of ionizing photons falls below the critical limit needed to photoionize the IGM; and by z=5, it will fall short of the critical value by a factor of $``$4. This implies that at high redshift, the contribution from young stars may become the dominant contributor to the background, with the caveat that the space density of star-forming galaxies would have to be maintained at the level observed at z $``$ 3, and that most of their UV photons would have to be free to escape into the IGM. The Devriendt et al. (1998) models lead to the conclusion, however, that the galaxy contribution to the UV background is negligible at high redshifts.
In conclusion, the proximity effect data at present reflect that the UV background at 2 $`<`$ z $`<`$ 4 is QSO dominated. The discrepancies between this model at low and high redshifts (Kulkarni & Fall 1993, Williger et al. 1994) indicate that the contribution to the background from galaxies may be of larger relative importance. We plan to undertake an analysis of the proximity effect at low redshifts from a large sample of QSO spectra taken with the Faint Object Spectrograph on the Hubble Space Telescope to place better constraints on the background at 0.5 $`<`$ z $`<`$ 2. Further observations of objects at z $`>`$ 4 are also of particular interest to this subject.
We extend thanks to the staff of the Steward Observatory Bok Telescope for their assistance with the observations, to T. Aldcroft and J. Shields for providing data, to C. Foltz and D. M<sup>c</sup>Intosh for helpful discussions, to J. M<sup>c</sup>Dowell for use of his program COLDEN, and to G. Ferland and associates for making the program CLOUDY available for general use. We also thank S. Morris for a helpful referee report. J. S. acknowledges the support of the National Science Foundation Graduate Research Fellowship and the Zonta Foundation Amelia Earhart Fellowship. J. B. acknowledges support from AST-9058510 and AST-9617060 of the National Science Foundation. A. D. acknowledges support from NASA Contract No. NAS8-39073 (ASC). V. P. K. acknowledges partial support from an award from the William F. Lucas Foundation and the San Diego Astronomers’ Association. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration.
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# Triplet Dispersion in CuGeO3: Perturbative Analysis
## I Introduction
The dispersion of the magnetic excitations is an important source of information on experimental low-dimensional spin systems. Knowledge of the dispersion relation $`\omega (\stackrel{}{k})`$ helps essentially to identify the model appropriate to describe the compound under study. The dispersion relation provides also important insight in the nature of the ground state. Very common in low dimensional systems is the scenario of a singlet $`S=0`$ ground state without magnetic long range order (a “spin liquid”) of which the elementary excitations are triplets $`S=1`$. These systems are generically gapped. Examples are isolated or weakly coupled dimerized spin chains and spin ladders such as (VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub>, the spin-Peierls phase of CuGeO<sub>3</sub>, and SrCu<sub>2</sub>O<sub>3</sub>. A true 2D example is SrCu<sub>2</sub>(BO)<sub>2</sub> which is characterized by frustrated dimers.
In these gapped $`S=1/2`$ systems where the gap is related to some “strong” bond (which can be also the rung of a 2-leg ladder) the elementary triplet excitations are in principle accessible by a perturbative expansion about the limit of isolated dimers. This approach, however, becomes tedious for the description of realistic materials since the expansion parameter is often not really small. Thus one has to compute high orders to achieve quantitative agreement. For this reason various automated approaches have been conceived which leave the tedious part to computers .
In the present article, we will apply the previously introduced perturbation by flow equation to the two-dimensional, though anisotropic, system of CuGeO<sub>3</sub> in its dimerized low-temperature phase . Thereby we extend the previous analysis (Ref. , henceforth cited as \[I\]) considerably. Our starting point remains the same as before. The strongest coupling is given by $`J`$; the other couplings are given relative to $`J`$ as indicated in Fig. 1 (for details see Fig. 1 in \[I\]).
## II Method
The problem to be solved reads
$$H=H_0+\lambda H_S.$$
(1)
As in the chain in Ref. the isolated dimer limit ($`\lambda =0`$ at finite $`\mu /\lambda `$) has an equidistant energy spectrum and the perturbation can alter the number of energy quanta (here: triplets on the dimers) by 2 at maximum. Hence $`H_S`$ can be represented as $`H_S=T_2+T_1+T_0+T_1+T_2`$ where $`T_i`$ stands for the perturbing part changing the number of elementary triplets by $`i`$. The same formalism as in Ref. can be used. This formalism maps the perturbed Hamiltonian by a continuous unitary transformation, the so-called flow equation method , to an effective Hamiltonian $`H_{\mathrm{eff}}`$ which conserves the number of energy quanta, i.e. $`0=[H_{\mathrm{eff}},H_0]`$. The effective Hamiltonian has the form
$$H_{\mathrm{eff}}=H_0+\underset{k=1}{\overset{\mathrm{}}{}}\lambda ^k\underset{|\underset{¯}{m}|=k,M(\underset{¯}{m})=0}{}C(\underset{¯}{m})T(\underset{¯}{m}),$$
(2)
where $`\underset{¯}{m}`$ is a vector of dimension $`k`$ of which the components are in $`\{\pm 2,\pm 1,0\}`$; $`M(\underset{¯}{m})=0`$ signifies that the sum of the components vanishes which reflects the conservation of the number of energy quanta (triplets). The coefficients $`C(\underset{¯}{m})`$ are generally valid fractions computed in Ref. where also further details on the flow equation method can be found.
Since $`H_{\mathrm{eff}}`$ conserves the triplet number the one-triplet sector is particularly easy to solve. Acting on one triplet the action of $`H_{\mathrm{eff}}`$ may only consist in shifting the triplet. This means that the triplet hops on an effective lattice where one site stands for one dimer on the original lattice, see Fig. 2 in \[I\] or Fig. 2.
The full dispersion $`\omega (\stackrel{}{k})`$ is obtained by Fourier transform
$$\omega (\stackrel{}{k})=J\underset{j,n}{}h_{j,n}\mathrm{exp}(i(k_1j+k_2n)).$$
(3)
The hopping amplitudes $`h_{j,n}`$ can be calculated on finite clusters of the (in principle infinite) effective lattice: From the linked cluster theorem we know that the finite order contribution of a short-ranged perturbation does not depend on the cluster size for sufficient large clusters. Carrying out our perturbation within order $`l`$ implies that one allows dimer to dimer hopping processes of length $`l`$ . The minimum cluster for a given amplitude $`h_{j,n}`$ in a given order $`l`$ contains all dimers and links that are involved in a hopping of length $`l`$ starting at dimer $`(0,0)`$ and ending at $`(j,n)`$. The minimum cluster is determined by considering all paths from $`(0,0)`$ to $`(j,n)`$. All dimers and links covered by one of these paths are part of the minimum cluster. In Fig. 3, the computer generated minimum cluster for calculating $`h_{3,1}`$ in order 8 is shown.
Due to the strong anisotropy of the quasi-1D system CuGeO<sub>3</sub> it is reasonable to use higher order terms only along the chains. This simplifies the computational task considerably since the calculation of a hopping process along the chain is much simpler. The cluster to be considered can be chosen smaller. The same is true for hopping processes close to the chain direction. Here we restrict the hopping processes to be at maximum quadratic in the interchain hopping $`\mu `$, which reduces the cluster sizes significantly so that the perturbation order can be enlarged.
## III Analysis of Experimental Data
The results for the $`h_{j,n}`$ are too lengthy to be published in written form. We will provide them in electronic form on our home pages on appearance of this article. In \[I\] the $`h_{j,n}`$ in third order in $`\lambda `$ and $`\mu `$ were presented. A few of these are erroneous. They are corrected herewith . The corrections, however, have no influence on the conclusions in \[I\] (see also discussion below).
Once all amplitudes $`h_{j,n}`$ are calculated the dispersion relation is given by Eq. (3). After rewriting Eq. (3) in terms of $`k_b`$ and $`k_c`$ (the reciprocal basis to $`e_b`$ and $`e_c`$) we add the term $`4t_a\mathrm{cos}(k_a)\mathrm{cos}(k_c)`$ with $`4t_a=0.22`$ meV to account for the dispersion in a-direction (cf. \[I\]). To fix the parameters $`J`$, $`\alpha `$, $`\beta `$, $`\mu `$ and $`\lambda `$ (cf. Fig. 1) we use the one-magnon dispersion data for CuGeO<sub>3</sub> experimentally determined by inelastic neutron scattering . Note that the hopping amplitudes are computed as polynomials over $``$ in the parameters.
As noticed in \[I\] the parameter $`\beta `$ has almost no influence on the shape of the dispersion. Hence we refrain from determining $`\beta `$ from the dispersion but set it beforehand to some reasonable values in the interval $`[0.3,0.3]`$. This choice is motivated by comparing the microscopic direct super-exchange path $`\mu `$ and the shifted super-exchange path $`\mu \beta `$ shown schematically in Fig. 4 ( cf. Fig. 1). There is only one path per Cu<sup>2+</sup>-site for the shifted coupling whereas there are two paths for the direct coupling. Thus we expect $`|\mu \beta |1/2|\mu |`$, i.e. $`|\beta |0.5`$. There are also results from ab-initio calculation for the interchain couplings which indicate the existence of interchain frustration . Further evidence is provided below by the analysis of the susceptibility. Furthermore, we find that for $`|\beta |>0.4`$ fits to the dispersion data become worse.
To determine the remaining parameters we equate four different experimental points with the corresponding parameter dependent dispersion values given by Eq. (3). The parameters are fixed by solving the resulting system of equations. For $`\beta =0.3`$ and $`0.3`$ Figs. 56 show the resulting dispersion curves in c\- and in b-direction, respectively, using all $`h_{j,n}`$ calculated.
As can be seen from Figs. 5 and 6 the plain series up to $`10^{\mathrm{th}}`$ order provides excellent fits. Yet one realizes that the parameter values still change on passing from order to order. So it appears that even at $`10^{\mathrm{th}}`$ order the results are not quantitative. In order to obtain quantitative reliable results we adopt a systematic extrapolation in the order. In each order $`l\{3,4,\mathrm{}10\}`$ we determine the optimum fit parameters. For illustration, Fig. 7 shows results for $`\beta =0.3`$.
Assuming exponential convergence we use
$$f(l)=Xbe^{2cl},$$
(4)
where $`X`$ is the asymptotic value of the parameter considered and $`b`$ and $`c`$ are constants. The choice (4) is motivated on one hand by its obvious applicability (see Fig. 7). On the other it stems from the fact that CuGeO<sub>3</sub> is a quasi one-dimensional gapped spin system. So one expects the magnetic correlations to drop exponentially with distance. Furthermore the order $`l`$ determines the maximum distance over which correlations occur (cf. Ref , namely $`l`$ counted in dimers or $`2l`$ counted in spin sites. Hence the constant $`c`$ in Eq. (4) can be understood as the inverse of a correlation length $`\xi `$. With the usual relation $`\xi v_S/\mathrm{\Delta }`$ for one dimensional systems we obtain $`c1/6`$ based on the rough estimates $`v_S=\pi /2J(11.12\alpha _0)`$ and $`\alpha _00.3;J12\text{meV};\mathrm{\Delta }2\text{meV}`$. This is indeed what is found (cf. Tab. I, II) so that we judge our extrapolations as being well justified.
In Fig. 7 the extrapolations are depicted by lines. The solid lines were obtained by using all calculated parameter values. The dashed lines are obtained from on the last four points, i.e. the results in order 7, 8, 9, and 10. The deviation between these two extrapolations are used as a measure for the extrapolation error. This procedure is carried out for $`\alpha ,\lambda `$, and $`J`$.
There is no systematic dependence of $`\mu `$ on the order $`l`$. The parameter $`\mu `$ oscillates between the two thin horizontal lines in Fig. 7. So we take the average of these two bounds as our estimate for $`\mu `$ and their difference as the error in the determination of $`\mu `$.
Tabs. I and II summarize the results of the fits for the parameters $`\alpha `$, $`\lambda `$ and $`J`$ for four preset values of $`\beta `$. The values for $`\mu `$ are listed in Tab. III.
| | All points considered | | | last four points considered | | |
| --- | --- | --- | --- | --- | --- | --- |
| | $`X`$ | $`b`$ | $`c`$ | $`X`$ | $`b`$ | $`c`$ |
| $`\beta =0.3`$ | | | | | | |
| $`\alpha `$ | 0.245 | 2.61 | 0.249 | 0.236 | 8.24 | 0.338 |
| $`\lambda `$ | 0.867 | 0.501 | 0.155 | 0.840 | 8.66 | 0.418 |
| $`J/`$meV | 13.6 | 11.0 | 0.164 | 13.2 | 105 | 0.357 |
| $`\beta =0.22`$ | | | | | | |
| $`\alpha `$ | 0.232 | 3.27 | 0.268 | 0.228 | 9.50 | 0.343 |
| $`\lambda `$ | 0.863 | 0.681 | 0.184 | 0.842 | 13.6 | 0.444 |
| $`J/`$meV | 13.1 | 11.3 | 0.171 | 12.8 | 68.9 | 0.324 |
| $`\beta =0`$ | | | | | | |
| $`\alpha `$ | 0.218 | 4.16 | 0.294 | 0.226 | 9.20 | 0.334 |
| $`\lambda `$ | 0.862 | 0.902 | 0.213 | 0.848 | 14.94 | 0.442 |
| $`J/`$meV | 12.8 | 11.1 | 0.174 | 12.6 | 44.48 | 0.287 |
| $`\beta =0.3`$ | | | | | | |
| $`\alpha `$ | 0.212 | 4.48 | 0.300 | 0.222 | 9.94 | 0.337 |
| $`\lambda `$ | 0.863 | 0.974 | 0.218 | 0.849 | 16.5 | 0.448 |
| $`J/`$meV | 12.7 | 10.9 | 0.173 | 12.5 | 42.2 | 0.282 |
TABLE I. Extrapolated parameter values $`X`$ according to Eq. (4). The experimental points we used in the fit process for this table are (cf. Figs 56) \[$`(k_b,k_c);\omega (𝐤)/\text{meV}`$\]:
\[$`(0,0);2.1`$\], \[$`(0,0.05);4.55`$\], \[$`(0,0.25);15.7`$\], \[$`(1,0);5.78`$\].
| | All points considered | | | last four points considered | | |
| --- | --- | --- | --- | --- | --- | --- |
| | $`X`$ | $`b`$ | $`c`$ | $`X`$ | $`b`$ | $`c`$ |
| $`\beta =0.3`$ | | | | | | |
| $`\alpha `$ | 0.297 | 1.41 | 0.215 | 0.309 | 1.18 | 0.187 |
| $`\lambda `$ | 0.868 | 0.423 | 0.197 | 0.877 | 0.404 | 0.165 |
| $`J/`$meV | 14.3 | 8.88 | 0.142 | 13.9 | 37.4 | 0.269 |
| $`\beta =0.22`$ | | | | | | |
| $`\alpha `$ | 0.301 | 1.41 | 0.208 | 0.318 | 1.15 | 0.175 |
| $`\lambda `$ | 0.867 | 0.850 | 0.274 | 0.886 | 0.526 | 0.175 |
| $`J/`$meV | 13.6 | 12.0 | 0.191 | 13.8 | 15.2 | 0.191 |
| $`\beta =0`$ | | | | | | |
| $`\alpha `$ | 0.308 | 1.34 | 0.191 | 0.323 | 1.30 | 0.173 |
| $`\lambda `$ | 0.900 | 0.38 | 0.140 | 0.896 | 0.754 | 0.191 |
| $`J/`$meV | 13.6 | 8.00 | 0.133 | 13.7 | 9.31 | 0.139 |
| $`\beta =0.3`$ | | | | | | |
| $`\alpha `$ | 0.314 | 1.28 | 0.180 | 0.326 | 1.37 | 0.173 |
| $`\lambda `$ | 0.913 | 0.369 | 0.127 | 0.903 | 0.832 | 0.192 |
| $`J/`$meV | 13.6 | 7.69 | 0.127 | 13.6 | 9.12 | 0.136 |
TABLE II. Same as in Tab. I based on different experimental points: \[$`(0,0);2.1`$\], \[$`(0,0.05);4.35`$\], \[$`(0,0.25);15.7`$\], \[$`(1,0);5.78`$\].
A closer inspection of Figs. 5 and 6 reveals that we are confronted with a certain arbitrariness of which fit we should favor. The experimental errors enhance this problem. The filled circles in the range of small wave vectors in Fig. 5 represent experimental points which have been measured with a high degree of precision. Thus it is reasonable to fit the theoretical curve as well as possible to these points. Fig. 8 shows an enlargement of this region. The solid line is the 10<sup>th</sup> order fit result for $`\beta =0.3`$ as the solid line in Fig. 5. The depicted arrow indicates the experimental point ($`k_c=0.05,\omega =4.55`$meV) used to obtain Tab. I.
A likewise well suited curve, however, is produced if one uses the point ($`k_c=0.05,\omega =4.35`$meV) for the fit keeping the other points (Tab. II). It is not possible to prefer one of the two lines in Fig. 8 to the other on the basis of their agreement to the experimental data. Hence we choose these two fits as the bounds within which all fits are acceptable. The corresponding parameter values $`X_1`$ (fit 1) and $`X_2`$ (fit 2) provide an interval \[$`X_1,X_2`$\] which we expect to contain the true model parameter $`\overline{X}`$. Hence the latter is estimated by
$`\overline{X}=(\overline{X}_1+\overline{X}_2)/2\pm \mathrm{\Delta }\overline{X}`$ , with $`\overline{X}_i=1/2(X_i^{\mathrm{all}\mathrm{points}}+X_i^{\mathrm{last}4\mathrm{points}})`$, $`\mathrm{\Delta }\overline{X}_i=|\overline{X}_iX_i^{\mathrm{all}\mathrm{points}}|`$ and $`\mathrm{\Delta }\overline{X}=\mathrm{max}\{|\overline{X}\overline{X}_1|,\mathrm{\Delta }\overline{X}_1,\mathrm{\Delta }\overline{X}_2\}`$. The results are summarized in Tab. III.
| parameter | interval | parameter | interval |
| --- | --- | --- | --- |
| $`\beta =0.3`$ | | | |
| $`\alpha `$ | 0.27(4) | $`\alpha _0`$ | 0.25(3) |
| $`\lambda `$ | 0.86(2) | $`\delta `$ | 0.08(1) |
| $`\mu `$ | 0.27(1) | $`\mu _0`$ | 0.29(1) |
| $`J`$/meV | 13.8(5) | $`J_0`$/meV | 12.8(6) |
| $`\beta =0.22`$ | | | |
| $`\alpha `$ | 0.27(4) | $`\alpha _0`$ | 0.25(3) |
| $`\lambda `$ | 0.86(2) | $`\delta `$ | 0.08(1) |
| $`\mu `$ | 0.21(1) | $`\mu _0`$ | 0.23(1) |
| $`J`$/meV | 13.4(4) | $`J_0`$/meV | 12.5(5) |
| $`\beta =0.0`$ | | | |
| $`\alpha `$ | 0.27(5) | $`\alpha _0`$ | 0.25(5) |
| $`\lambda `$ | 0.88(3) | $`\delta `$ | 0.07(2) |
| $`\mu `$ | 0.13(1) | $`\mu _0`$ | 0.14(1) |
| $`J`$/meV | 13.2(6) | $`J_0`$/meV | 12.4(7) |
| $`\beta =0.3`$ | | | |
| $`\alpha `$ | 0.27(5) | $`\alpha _0`$ | 0.25(5) |
| $`\lambda `$ | 0.88(3) | $`\delta `$ | 0.06(2) |
| $`\mu `$ | 0.08(2) | $`\mu _0`$ | 0.09(1) |
| $`J`$/meV | 13.1(5) | $`J_0`$/meV | 12.3(7) |
TABLE III.Final parameter intervals resulting from Tabs. I and II for three different values of $`\beta `$.
For the readers’ convenience Tab. III also gives the results in the more commonly used parameters $`\delta `$, $`\alpha _0`$, $`\mu _0`$, $`J_0`$ and $`\beta `$. This notation is connected to the one used so far in this article by
$`J=J_0(1+\delta )`$ , $`\lambda ={\displaystyle \frac{1\delta }{1+\delta }}`$ (5)
$`\alpha ={\displaystyle \frac{\alpha _0}{1\delta }}`$ , $`\mu ={\displaystyle \frac{\mu _0}{1+\delta }}.`$ (6)
It corresponds to the Hamiltonian depicted in Fig. 9.
## IV Implications for the Susceptibility
The temperature dependence of the homogeneous susceptibility $`\chi (T)`$ is often used to determine the parameters of CuGeO<sub>3</sub> . Already the Curie-Weiss temperature $`\mathrm{\Theta }`$ provides valuable information on the sum of the coupling constants. This is particularly useful to detect frustration. The dispersions are governed by the difference of the direct and the frustrating coupling whereas $`\chi (T)`$ at larger temperatures is more influenced by the sum of direct and frustrating coupling.
The analysis of the Curie-Weiss temperature alone bears some risks. It is easy to calculate but difficult to determine experimentally since it has to be deduced from values at high temperatures where $`\chi (T)`$ is fairly small and hence strongly influenced by background effects (van Vleck, diamagnetism) or by slight structural changes.
A convincing fit for temperatures above 50K is given by Fabricius et al. in Ref. on the basis of frustrated chains. The inclusion of interchain couplings, however, would spoil the excellent agreement and a re-determination of the constant would be necessary. A description of $`\chi (T)`$ on the basis of a two-dimensional model has not been done except for a consideration of the two leading coefficients in an expansion in $`1/T`$ in \[I\]. In Fig. 10 we show the same high quality experimental data as in Ref. and compare it to theoretical curves at four values of $`\beta `$. The theoretical curves are obtained by computing a Dlog Padé approximant $`\chi _0(T)`$ based on the high temperature series provided in Ref. for the frustrated chains. This procedure provides excellent results down to $`TJ/5`$ . The asymptotic behavior of the approximant is chosen such that $`\chi _0(T)`$ vanishes linearly on $`T0`$ as is to be expected for a two-dimensional massless antiferromagnet. Besides this feature the two-dimensionality is incorporated on a chain-mean-field level
$$\chi (T)=\frac{\chi _0(T)}{1+2\mu _0(1+2\beta )\chi _0(T)}.$$
(7)
This relation is exact in linear order in $`\mu `$. Estimates of corrections quadratic in $`\mu `$ indicate that they are negligible for the values of $`\mu `$ and $`\beta `$ in which we are interested.
From the results in Fig. 10 it is evident that the interchain frustration cannot be neglected. Only for a finite value of about $`0.22`$ a very good agreement can be obtained. The agreement to the frustrated chain model is still better since the position of the maximum is also reproduced. But on the basis of the neutron scattering results it is undoubtful that CuGeO<sub>3</sub> is a two-dimensional substance. Furthermore, it must be considered that the previous fit was a two-parameter ($`J`$ and $`\alpha _0`$) fit whereas only one parameter ($`\beta `$) is fitted to obtain Fig. 10. The other parameters ($`J_0,\alpha _0,\mu _0`$) were determined from an entirely different experiment. Hence the agreement for $`\chi (T)`$ corroborates also the validity of the parameters determined in the preceding section.
## V Discussion
We will first discuss our results and propose a set of parameters. Then we will put these results into the context of other results in the literature.
Let us consider the remaining difference between experiment and theory concerning $`\chi (T)`$ in Fig. 10. There are four conceivable sources for it. The first are experimental inaccurcies. We are not in the position to judge this aspect. We just like to comment that from the results and error bars in Fig. 8 it is obvious that the experimental data is not completely consistent so that this explanation is possible.
Second, it is conceivable that the couplings change across the transition, i.e. the intrachain frustration at $`T0`$ (where the dispersions are measured) is different from the one above $`T_{\mathrm{SP}}`$ (where $`\chi (T)`$ is measured). Since we are considering a spin-Peierls transition there is definitely a structural change. So far, however, the assumption that only the dimerization changes worked well. The structural changes in the transition are very small whereas the changes needed to explain the discrepancy are of the order of 20 to 30 % (assuming a change in the intrachain frustration). Estimates point into the direction that the changes on the couplings are unimportant . Yet the estimates concern in the first place the nearest-neighbor coupling only. Quantitative ab initio calculations of the frustration are very difficult , even more so for changes of the frustration. So, again, this explanation is perhaps not the most promising but cannot be excluded either.
Third, the influence of the phonon dynamics is to be considered. It is shown that spin-Peierls systems can be unitarily transformed in such a way that an effective spin model remains at low energies uncoupled to the phonons . The effective couplings in a single chain model become temperature dependent so that this may account for the discrepancy. But it turns out for non-resonant phonons ($`\omega >J`$) that these effects leave the susceptibility fairly unchanged. This is so since these effects become significant for relatively large temperatures where $`\chi (T)`$ depends only on the sum of all couplings which is unchanged by the unitary transformation (this is observed in (VO)<sub>2</sub>P<sub>2</sub>O<sub>7</sub> ). So this reason appears rather unlikely even though it looked plausible at first sight.
Fourth, one has to think about any kind of precursors of the spin-Peierls transition. By this we mean on one side the critical fluctuations which appear in a narrow region ($`3K`$) around the spin-Peierls temperature . On the other side, we mean any precursor which goes beyond a purely static spin model. Experimentally, a finite lattice correlation length can be detected already at $`T40`$K far above the actual transition . From there on deviations from the behavior of a static spin model should be observable. In the adiabatic limit, for instance, the fluctuation yield already a reduction of the susceptibility . What happens in the antiadiabatic limit has not yet been investigated for a two or higher dimensional model. The mapping in Ref. leads to four-point interchain couplings the influence of which is unclear so far.
In view of the above mentioned possible pitfalls of the static model the agreement in Fig. 10 is already very convincing. Summarizing our results we propose the parameters given for $`\beta =0.22`$ in Tab. III to be the ones deduced from the dispersion data. Assessing the reliability of our estimates, we redo the analysis of the susceptibility for $`\alpha _0=0.28`$ (the upper bound of our estimate for $`\alpha _0`$) with the corresponding value of $`J_0=12.8`$meV. Then the optimum $`\chi (T)`$ is obtained for $`\beta =0.15`$; the $`\chi (T)`$ curve is almost identical to the one shown in Fig. 10. So the value of $`\beta `$ cannot be determined very precisely, but it should be in the range $`\beta =0.2(1)`$. A certain dicrepancy between the optimum parameters for the $`T=0`$ dispersion data and for the $`\chi (T>T_{\mathrm{SP}})`$ data remains.
We split the comparison of our findings to previous works into three groups. The first comprises the analyses on the basis of a one-dimensional model . The most striking difference to the results for static spin models is that the dimerization $`\delta `$ is not of the order of 1% but significantly larger. This is not astounding since it has been noted already in \[I\] that the gap is lowered by the interchain hopping. Hence the neglect of the latter requires to lower the gap otherwise, i.e. by assuming a lower dimerization.
Our intrachain frustration is slightly larger than the one of Castilla et al. ($`0.24`$), but significantly smaller than the value of Riera and Dobry ($`0.36`$) or the value of Fabricius et al. ($`0.35`$). Fabricius et al. showed that the value $`\alpha _0=0.24`$ is too small for a single chain model. The difference between the larger frustration value in the single chain model to our value results directly from the interchain coupling. As can be nicely seen in Eq. (7) the interchain coupling lowers the susceptibility without changing (in the chain mean-field approximation) the position of its maximum. The one dimensional models favor a larger intrachain frustration and a concomitant larger coupling $`J`$ in order to reduce the magnitude of the susceptibility.
The claim by Wellein and coworkers that the dimerization experimentally found to be larger than would fit to a static 1D model is due to the phonon dynamics is not compelling. They use a root-mean-square definition of the dimerization which naturally provides larger values for the dimerization since it includes all the fluctuations. The dispersion perpendicular to the chains, however, is an unambiguous experimental fact. Furthermore, Trebst and coworkers do not find a substantial gap renormalization for parameters relevant for CuGeO<sub>3</sub> even though one should take care of different schemes to couple the phonons.
Let us turn to the second group of papers considering the essentially two-dimensional character of CuGeO<sub>3</sub>. The first work used a bond-operator technique . No frustration was considered, hence rather small values of $`J=10.2`$meV and a rather small interchain hopping $`\mu 0.06`$ resulted. The same technique was also applied later again by Brenig including frustration. It turned out, however, that only $`\lambda (12\alpha )`$ and $`\mu (12\beta )`$ matter on the free-boson level. Hence an independent determination of the frustration is not possible. Using additional input ($`\delta =0.012`$) the values $`\alpha _0=0.059`$ and $`\mu (12\beta )=0.054`$ were obtained. In view of the extensive comparisons to numerical results made in Ref. it appears that the bond-operator method overestimates the influence of additional couplings such as dimerization or frustration. Generally, the values for dimerization or frustration tend to be too low. This is confirmed by our findings in the present work.
Compared to \[I\] ($`\alpha _0=0,J=9.8`$meV, $`\delta =0.12,\mu _0=0.34,\beta =0.3`$) the extended series on which our present analysis is based gives a much better handle on intrachain frustration, see Fig. 7. Only in the present work, we are able to assess its value reliably. With respect to the interchain frustration, the present results agree qualitatively with those in \[I\] where such a frustration was proposed for CuGeO<sub>3</sub> first. The use of susceptibility information has been improved in the present work since the whole $`\chi (T)`$ curve is used, not only the leading coefficients.
Bouzerar et al. have carried out an estimate leading to results not too far from ours: $`\delta =0.065,\alpha _0=0.2,J=12.2`$meV, $`\mu =0.15`$. They used just linear order in the interchain hopping without interchain frustration and some square-root averaging with numerical results for chains to describe the dispersion. The intrachain frustration ($`\alpha _0=0.2`$) could only be taken from the Curie-Weiss constant. The resulting $`\chi (T)`$ has similarities with the experimental one.
The third group comprises ab-initio calculations of the exchange couplings and of the spin-phonon couplings. Microscopic calculations find relatively large values of the dimerization between $`0.07`$ and $`0.2`$ in agreement with our findings. (Even though there is also a different result ) Very important for our work are recent results by Drechsler and coworkers supporting the existence of interchain frustration . Werner and coworkers estimate a large dimerization from the spin-phonon couplings and the shifts of the ions ($`\delta =0.11`$). From the balance of elastic and magnetic energy in the D phase they obtain without frustration a lower bound of $`\delta >0.044`$. Assuming critical frustration $`\alpha _0=0.2412`$ they find even $`\delta >0.078`$ which fits very nicely to our findings.
In summary, we provide by the present work a determination in great detail of the coupling parameters ($`\beta =0.2(1)`$ and right column in Tab. III under $`\beta =0.22`$) of CuGeO<sub>3</sub> based on a static dimerized spin model at $`T=0`$. The experimental input comes from inelastic neutron scattering. The implications of the parameters found for the susceptibility are also studied. Very good agreement could be obtained fitting the interchain frustration appropriately. A small discrepancy at low temperatures around $`50`$K indicates that the static spin model is probably insufficient to describe CuGeO<sub>3</sub> completely. By this work, we proved the outstanding possibilities of high-order series expansions (around the dimer limit or around the limit $`T=\mathrm{}`$) in the analysis of experimental data.
## Acknowledgements
The authors like to thank B. Büchner, A. Bühler, U. Löw, B. Marić, E. Müller-Hartmann, and F. Schönfeld for fruitful discussions. The provision of the experimental data by B. Büchner, T. Lorenz and by L. P. Regnault is gratefully acknowledged. Furthermore, we thank the Regional Computing Center of the University of Cologne for its kind and efficient support. This work was supported by the Deutsche Forschungsgemeinschaft in the Schwerpunkt 1073 and in the Sonderforschungsbereich 341.
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# The Type-problem on the Average for random walks on graphs
## I Introduction
In statistical mechanics and field theory Euclidean lattices describe the geometrical structure of crystals and of more abstract geometrical objects, such as discretized flat space-time. However, most real systems have irregular geometry: this is the case for glasses, polymers, amorphous materials, biological structures and fractals in condensed matter, as well as discretized curved space-times in field theory. The geometrical model for these inhomogeneous systems is a general discrete network made of sites and links, i.e. a graph. From this point of view usual lattices are a particular kind of graphs. Statistical models defined on graphs (e.g. harmonic oscillations, random walks, spin system) are the natural way to describe the physical properties of inhomogeneous real structures.
The study of the relation between geometry and physics is one of the most complex and interesting problem of statistical mechanics and field theory on graphs. The main link between this two aspects is provided by random walks. The latter are usually introduced to describe the diffusion of a classical particle and they are related to Markov chains, potential theory and algebraic graph theory on one side , and to many problems of equilibrium and non equilibrium statistical mechanics, disordered systems and field theory on the other .
In particular, the large times asymptotics of random walks provides the most effective method to describe the influence of large scale topology on the physical properties of discrete structures. The definition of the spectral dimension for inhomogeneous networks, generalizing the Euclidean dimension of lattices in field theory and phase transitions, is indeed based on long time behavior of random walks . More generally, this asymptotic regime allows to classify every graph either as locally recursive or transient, according to the probability of ever returning to the starting site: the probability is 1 in the former case and less than 1 in the latter, independently of the site. This classification, first introduced by Polya for regular lattices , is known as the Type-Problem.
Local transience and recurrence describe local properties of physical models on graphs. However, in the study of statistical models on graphs we are in general interested in average (extensive) thermodynamic quantities. Indeed, while on lattices, due to translation invariance, local quantities are the same on all sites and therefore they are equal to their average, on inhomogeneous structures they depend in general on the site and the average behavior can not be reduced to the local one. In the last few years it has become clear that bulk properties are affected by the average values of random walks return probabilities over all starting sites: this is the case for spontaneous breaking of continuous symmetries , critical exponents of the spherical model , harmonic vibrational spectra . Therefore the classification of discrete structure in terms of recurrence on the average and transience on the average appears to be the most suitable. Unfortunately, while for regular lattices the two classifications are equivalent, on more general networks they can be different and one has to study a Type-Problem on the Average .
Recently this problem has acquired particular relevance in the study of spin models on graphs. Indeed it has been shown that spontaneous breaking of continuous symmetries occurs at $`T>0`$ if and only if the underlying network is transient on the average . Moreover this analysis has shown that relevant and new topological properties of infinite graphs are associated to the on the average classification.
In this paper we deal with the Type-Problem on the Average and with the topological and thermodynamic properties arising from it, proving some basic theorems and discussing their relevance for present and future development in statistical physics.
The paper is organized as follows. In the next section we introduce the basic concepts and notations concerning random walks on infinite graphs. In section III we analyze local recurrence and transience properties of random walks, defined by the asymptotic behavior of return probabilities generating function. In section IV we consider thermodynamic averages on infinite graphs and in section V we introduce the topological classification in terms of average properties of random walks, based on transience and recurrence on the average. In section VI we analyze the general relations between average generating functions of return probabilities and consider a further classification holding for transient on the average graphs, which completes the topological description of infinite graphs in terms of random walk behavior. In section VII and VIII we show the relevance of this topological classification in the study of thermodynamic properties of statistical models on inhomogeneous structures. A summary and a discussion of our results are presented in section IX.
## II Random walks on infinite graphs
Let us begin by recalling the basic definitions and results concerning graph theory and random walks on infinite graphs, which will be used in the following. A more detailed and complete treatment can be found in the mathematical reviews by Woess .
A graph $`𝒢`$ is a countable set $`V`$ of vertices (or sites) $`(i)`$ connected pairwise by a set $`E`$ of unoriented links (or bonds) $`(i,j)=(j,i)`$. If the set $`V`$ is finite, $`𝒢`$ is called a finite graph and we will denote by $`N`$ the number of vertices of $`𝒢`$. A subgraph $`𝒮`$ of $`𝒢`$ is a graph whose set of vertices $`SV`$ and whose set of links $`E^{}E`$.
A path in $`𝒢`$ is a sequence of consecutive links $`\{(i,k)(k,h)\mathrm{}(n,m)(m,j)\}`$ and a graph is said to be connected, if for any two points $`i,jV`$ there is always a path joining them. In the following we will consider only connected graphs.
The graph topology can be algebraically represented introducing its adjacency matrix $`A_{ij}`$ given by:
$$A_{ij}=\{\begin{array}{cc}1& \mathrm{if}(i,j)E\hfill \\ 0& \mathrm{if}(i,j)E\hfill \end{array}$$
(1)
The Laplacian matrix $`\mathrm{\Delta }_{ij}`$ is defined by:
$$\mathrm{\Delta }_{ij}=z_i\delta _{ij}A_{ij}$$
(2)
where $`z_i=_jA_{ij}`$, the number of nearest neighbors of $`i`$, is called the coordination number of site $`i`$. Here we will consider graphs with $`z_{max}=sup_iz_i<\mathrm{}`$.
In order to describe disordered structures we introduce a generalization of the adjacency matrix given by the ferromagnetic coupling matrix $`J_{ij}`$, with $`J_{ij}0A_{ij}=1`$ and $`sup_{(i,j)}J_{ij}<\mathrm{}`$, $`inf_{(i,j)}J_{ij}>0`$. One can then define the generalized Laplacian:
$$L_{ij}=I_i\delta _{ij}J_{ij}$$
(3)
where $`I_i=_jJ_{ij}`$.
Every connected graph $`𝒢`$ is endowed with an intrinsic metric generated by the chemical distance $`r_{i,j}`$ which is defined as the number of links in the shortest path connecting vertices $`i`$ and $`j`$.
Let us now introduce the random walk on a graph $`𝒢`$ defining the jumping probability $`p_{ij}`$ between nearest neighbors sites $`i`$ and $`j`$:
$$p_{ij}=\frac{A_{ij}}{z_i}=(Z^1A)_{ij}$$
(4)
where $`Z_{ij}=z_i\delta _{ij}`$. From (4) the probability of reaching in $`t`$ steps site $`j`$ starting from $`i`$ is given by:
$$P_{ij}(t)=(p^t)_{ij}.$$
(5)
Recurrence properties of random walks are studied introducing the probability $`F_{ij}(t)`$ for a walker starting from $`i`$ of reaching for the first time in $`t`$ steps the site $`ji`$, while $`F_{ii}(t)`$ is the probability of returning to the starting point $`i`$ for the first time after $`t`$ steps and $`F_{ii}(0)=0`$. The basic relationship between $`P_{ij}(t)`$ and $`F_{ij}(t)`$ is given by:
$$P_{ij}(t)=\underset{k=0}{\overset{t}{}}F_{ij}(k)P_{jj}(tk)$$
(6)
$`(t>0)`$. From the previous definitions $`F_{ij}_{t=0}^{\mathrm{}}F_{ij}(t)`$ turns out to be the probability of ever reaching the site $`j`$ starting from $`i`$ (or of ever returning to $`i`$ if $`j=i`$). Therefore $`0<F_{ij}1`$. The generating functions $`\stackrel{~}{P}_{ij}(\lambda )`$ and $`\stackrel{~}{F}_{ij}(\lambda )`$ are given by:
$$\stackrel{~}{P}_{ij}(\lambda )=\underset{t=0}{\overset{\mathrm{}}{}}\lambda ^tP_{ij}(t)\stackrel{~}{F}_{ij}(\lambda )=\underset{t=0}{\overset{\mathrm{}}{}}\lambda ^tF_{ij}(t)$$
(7)
where $`\lambda `$ is a complex number. From definition (7) and from the property $`0<F_{ij}1`$ by Abel theorem we have that $`\stackrel{~}{F}_{ij}(\lambda )`$ is a uniformly continuous function for $`\lambda [0,1]`$ and $`0<\stackrel{~}{F}_{ij}(\lambda )1`$, while $`\stackrel{~}{P}_{ij}(\lambda )`$ is continuous for $`\lambda [0,1[`$ but it can diverge for $`\lambda 1^{}`$.
Multiply equations (6) by $`\lambda ^t`$ and then summing over all possible $`t`$ with the initial condition $`P_{ij}(0)=\delta _{ij}`$ we get the basic relations between $`\stackrel{~}{P}_{ij}(\lambda )`$ and $`\stackrel{~}{F}_{ij}(\lambda )`$
$$\stackrel{~}{P}_{ij}(\lambda )=\stackrel{~}{F}_{ij}(\lambda )\stackrel{~}{P}_{jj}(\lambda )+\delta _{ij}$$
(8)
In the following we will call $`\stackrel{~}{P}_i(\lambda )\stackrel{~}{P}_{ii}(\lambda )`$ and $`\stackrel{~}{F}_i(\lambda )\stackrel{~}{F}_{ii}(\lambda )`$.
Before discussing recurrence and transience properties we briefly recall the definition of the Gaussian model on a graph, whose deep relation with random walks will be exploited in the next section.
The Gaussian model on the graph $`𝒢`$ can be defined introducing the field $`\varphi _i`$ which are the functions of $`l^{\mathrm{}}(V)=\{(\varphi _i)_{iV}:sup_i|\varphi _i|<\mathrm{}\}`$. It exist a unique Gaussian probability measure $`d\mu _g\varphi `$ on $`l^{\mathrm{}}(V)`$ with mean zero and covariance $`(L+\mu )^1`$ ($`\mu _{ij}`$ is the diagonal matrix $`\mu _{ij}=\mu \delta _{ij}`$, $`\mu >0`$); $`d\mu _g(\varphi )`$ characterize the Gaussian model and we will write:
$$F(\varphi )=F(\varphi )𝑑\mu _g(\varphi )$$
(9)
and in particular:
$$\varphi _i\varphi _j=(L+\mu )_{ij}^1$$
(10)
Alternately, the Gaussian model can be introduced using standard approach of statistical mechanics via the Hamiltonian:
$$=\underset{i,j𝒢}{}L_{ij}\varphi _i\varphi _j+\underset{i𝒢}{}\mu \varphi _i^2$$
(11)
together with the Boltzmann weight $`\mathrm{exp}()`$, also leading to (10).
## III Local recurrence and local transience
The long time asymptotic behavior of random walks on infinite graphs are determined by the large scale topology of the graph and the quantities $`\stackrel{~}{F}_i(1)`$ and $`lim_{\lambda 1}\stackrel{~}{P}_i(\lambda )`$ can be used to characterize the geometry of the graph itself. In particular a graph is called locally recurrent if
$$\stackrel{~}{F}_i(1)=1\mathrm{or}\mathrm{equivalently}\underset{\lambda 1}{lim}\stackrel{~}{P}_i(\lambda )=\mathrm{}i$$
(12)
On the other hand if:
$$\stackrel{~}{F}_i(1)<1\mathrm{or}\mathrm{equivalently}\underset{\lambda 1}{lim}\stackrel{~}{P}_i(\lambda )<\mathrm{}i$$
(13)
the graph is called locally transient. By standard Markov chains properties (12) and (13) are independent from the site $`i`$ and then they can be consider as properties of the graphs.
Let us prove the independence from $`i`$ of (12). If $`lim_{\lambda 1}\stackrel{~}{P}_i(\lambda )=\mathrm{}`$ then by equation (8) we get $`lim_{\lambda 1}\stackrel{~}{P}_{ji}(\lambda )=\mathrm{}`$ for all $`j`$ ($`0<\stackrel{~}{F}_{ji}(1)1`$); now from (4) and (5) we have that $`z_iP_{ij}(t)=z_jP_{ji}(t)`$ and $`z_i\stackrel{~}{P}_{ij}(\lambda )=z_j\stackrel{~}{P}_{ji}(\lambda )`$. Then also $`lim_{\lambda 1}\stackrel{~}{P}_{ij}(\lambda )=\mathrm{}`$ and from (8) we obtain $`lim_{\lambda 1}\stackrel{~}{P}_j(\lambda )=\mathrm{}`$, $`jV`$. In an analogous way it can be shown that property (13) is independent from the choice of $`i`$.
Local transience and local recurrence satisfy important universality properties . Indeed these properties are not modified if we substitute the jumping probabilities of the random walker (4) with the generalized jumping probability:
$$p_{ij}=\frac{J_{ij}}{I_i}.$$
(14)
In the invariance of the local recurrence properties under a wide class of transformations of the graph itself is also proven. Local recurrence and transience are not modified by the addition a finite number of links or the introduction of second neighbor links on the graph. Notice that these basic invariance properties prove that local recurrence and transience are determined only by the large scale topology of the graph.
## IV Averages on infinite graphs
Let us now consider thermodynamic averages on infinite graphs. The generalized Van Hove sphere $`S_{o,r}𝒢`$ of center $`o`$ and radius $`r`$ is the subgraph of $`𝒢`$ containing all $`i𝒢`$ whose chemical distance from $`o`$ $`r_{i,o}`$ is $`r`$ and all the links of $`𝒢`$ joining them. We will call $`N_{o,r}`$ the number of vertices contained in $`S_{o,r}`$.
The average in the thermodynamic limit $`\overline{\varphi }`$ of a function $`\varphi _i`$ defined on each site $`i`$ of the infinite graph $`𝒢`$ is:
$$\overline{\varphi }\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o,r}}{}}\varphi _i}{N_{o,r}}.$$
(15)
The measure $`|S|`$ of a subset $`S`$ of $`V`$ is the average value $`\overline{\chi (S)}`$ of its characteristic function $`\chi _i(S)`$ defined by $`\chi _i(S)=1`$ if $`iS`$ and $`\chi _i(S)=0`$ if $`iS`$. The measure of a subset of links $`E^{}E`$ is given by:
$$|E^{}|\underset{r\mathrm{}}{lim}\frac{E_r^{}}{N_{o,r}}.$$
(16)
where $`E_r^{}`$ is the number of links of $`E^{}`$ contained in the sphere $`S_{o,r}`$. The normalized trace $`\overline{\mathrm{Tr}}B`$ of a matrix $`B_{ij}`$ is:
$$\overline{\mathrm{Tr}}B\overline{b}$$
(17)
where $`b_iB_{ii}`$. Let us require that:
$$\underset{r\mathrm{}}{lim}\frac{|S_{o,r}|}{N_{o,r}}=0$$
(18)
where $`|S_{o,r}|`$ is the number of the vertices of the sphere $`S_{o,r}`$ connected with the rest of the graph.
Under this hypothesis we now prove that the averages of a bounded from below function $`\varphi _i`$ are independent from the center $`o`$ of the spheres sequence, using the fact that $`\chi _i(S)`$ is bounded and that measures of subsets are always well defined. From the boundedness of the coordination number we get for any couple of vertices $`o`$ and $`o^{}`$:
$$N_{o,r}(z_{max})^{r_{o,o^{}}}|S_{o^{},r}|N_{o^{},r}N_{o,r}+(z_{max})^{r_{o,o^{}}}|S_{o,r}|$$
(19)
and
$$(z_{max})^{r_{o,o^{}}}|S_{o,r}||S_{o^{},r}|(z_{max})^{r_{o,o^{}}}|S_{o,r}|$$
(20)
Let us consider a bounded from below function $`\varphi _i`$. Given two vertices $`o`$ and $`o^{}`$, we have:
$$\frac{{\displaystyle \underset{iS_{o^{},rr_{o,o^{}}}}{}}\varphi _i+{\displaystyle \underset{iS_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}}{}}\varphi _i}{N_{o,r}}=\frac{{\displaystyle \underset{iS_{o,r}}{}}\varphi _i}{N_{o,r}}=\frac{{\displaystyle \underset{iS_{o^{},r+r_{o,o^{}}}}{}}\varphi _i{\displaystyle \underset{iS_{o^{},r+r_{o,o^{}}}\mathrm{\Delta }S_{o,r}}{}}\varphi _i}{N_{o,r}}$$
(21)
where $`S_{o,r}S_{o^{},r+r_{o,o^{}}}`$, $`S_{o^{},rr_{o,o^{}}}S_{o,r}`$ and $`S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}`$ is the symmetric difference between $`S_{o,r}`$ and $`S_{o,rr_{o,o^{}}}`$. $`|S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}|`$ denotes the number of vertices of $`S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}`$ and from (19) we get:
$$|S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}|(z_{max})^{r_{o,o^{}}}|S_{o,r}||S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}|(z_{max})^{r_{o,o^{}}}|S_{o,rr_{o.o^{}}}|$$
(22)
with an analogous equation holding for $`S_{o^{},r+r_{o,o^{}}}\mathrm{\Delta }S_{o,r}`$. Defining $`\overline{\varphi }=0`$ if $`\varphi _i>0`$ for all $`i`$ and $`\overline{\varphi }=|min_i\varphi _i|`$ otherwise, from (21) we have:
$$\frac{{\displaystyle \underset{iS_{o^{},rr_{o,o^{}}}}{}}\varphi _i\overline{\varphi }|S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}|}{N_{o,r}}\frac{{\displaystyle \underset{iS_{o,r}}{}}\varphi _i}{N_{o,r}}\frac{{\displaystyle \underset{iS_{o^{},r+r_{o,o^{}}}}{}}\varphi _i+\overline{\varphi }|S_{o^{},r+r_{o,o^{}}}\mathrm{\Delta }S_{o,r}|}{N_{o,r}}$$
(23)
with property (18) of $`𝒢`$ and inequalities (22) we get:
$$\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o^{},rr_{o,o^{}}}}{}}\varphi _i}{N_{o,r}}\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o,r}}{}}\varphi _i}{N_{o,r}}\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o^{},r+r_{o,o^{}}}}{}}\varphi _i}{N_{o,r}}$$
(24)
since $`N_{o,r}=N_{o^{},rr_{o,o^{}}}+|S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}|=N_{o^{},r+r_{o,o^{}}}|S_{o,r}\mathrm{\Delta }S_{o^{},rr_{o,o^{}}}|`$ using again (18) and (22) we get:
$$\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o^{},rr_{o,o^{}}}}{}}\varphi _i}{N_{o^{},rr_{o,o^{}}}}\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o,r}}{}}\varphi _i}{N_{o,r}}\underset{r\mathrm{}}{lim}\frac{{\displaystyle \underset{iS_{o^{},r+r_{o,o^{}}}}{}}\varphi _i}{N_{o^{},r+r_{o,o^{}}}}$$
(25)
Therefore, if the limit with the spheres centered in $`o^{}`$ exists, it gives the same result using as center any vertex $`o`$.
## V Recurrence and transience on the average
The study of thermodynamic properties of statistical models on infinite graphs requires the introduction of averages of local quantities. The latter are related to random walks by the return probabilities on the average $`\overline{P}`$ and $`\overline{F}`$, which are defined by:
$$\overline{P}=\underset{\lambda 1}{lim}\overline{\stackrel{~}{P}(\lambda )}$$
(26)
$$\overline{F}=\underset{\lambda 1}{lim}\overline{\stackrel{~}{F}(\lambda )}$$
(27)
A graph $`𝒢`$ is called recurrent on the average (ROA) if $`\overline{F}=1`$, while it is transient on the average (TOA) when $`\overline{F}<1`$.
Recurrence and transience on the average are in general independent of the corresponding local properties. The first example of this phenomenon occurring on inhomogeneous structures was found in a class of infinite trees called NTD (Fig. 1) which are locally transient but recurrent on the average .
Moreover, while for local probabilities (8) gives:
$$\stackrel{~}{P}_i(\lambda )=\stackrel{~}{F}_i(\lambda )\stackrel{~}{P}_i(\lambda )+1$$
(28)
an analogous relation for (27) and (26) does not hold since averaging (28) over all sites $`i`$ would involves the average of a product, which due to correlations is in general different from the product of the average. Therefore the double implication $`\stackrel{~}{F}_i(1)=1lim_{\lambda 1}\stackrel{~}{P}_i(\lambda )=\mathrm{}`$ is not true. Indeed there are graphs for which $`\overline{F}<1`$ but $`\overline{P}=\mathrm{}`$ (an example in shown in Fig. 2) and the study of the relation between $`\overline{P}`$ and $`\overline{F}`$ is a non trivial problem.
## VI Pure and mixed transience on the average
In this section we study the relation between $`\overline{P}`$ and $`\overline{F}`$ and we show that a complete picture of the behavior of random walks on graphs can be given by dividing transient on the average graphs into two further classes, which will be called pure and mixed transient on the average (TOA).
First, considering a ROA graph, we prove that if $`\overline{F}=1`$ then $`\overline{P}=\mathrm{}`$. In this case for each $`\delta >0`$ it exists $`ϵ`$ such that if $`1ϵ\lambda <1`$, we have: $`1\delta \overline{\stackrel{~}{F}(\lambda )}1`$. Let us consider the subset $`SV`$ of the sites $`i`$ such that $`\stackrel{~}{F}_i(1ϵ)<1\sqrt{\delta }`$ and we call $`\overline{S}`$ its complement. We obtain:
$$1\delta \overline{\stackrel{~}{F}(1ϵ)}=\overline{\chi (S)\stackrel{~}{F}(1ϵ)}+\overline{\chi (\overline{S})\stackrel{~}{F}(1ϵ)}(1\sqrt{\delta })|S|+|\overline{S}|=1\sqrt{\delta }|S|$$
(29)
From (29) we get $`|S|\sqrt{\delta }`$ and then $`|\overline{S}|1\sqrt{\delta }`$. Exploiting the property that $`\overline{\stackrel{~}{P}(\lambda )}`$ is an increasing function of $`\lambda `$, for each $`\lambda 1ϵ`$ we get:
$$\overline{\stackrel{~}{P}(\lambda )}\overline{\stackrel{~}{P}(1ϵ)}\overline{\chi (\overline{S})(1\stackrel{~}{F}(1ϵ))^1}|\overline{S}|\delta ^{1/2}(1\sqrt{\delta })\delta ^{1/2}$$
(30)
In this way we proved that for arbitrary large value of $`(1\sqrt{\delta })\delta ^{1/2}`$ ($`\delta 0`$), it exists $`ϵ`$ such that for each $`\lambda `$, $`1ϵ\lambda <1`$, we have $`\overline{\stackrel{~}{P}(\lambda )}(1\sqrt{\delta })\delta ^{1/2}`$, and therefore $`\overline{P}=lim_{\lambda 1}\overline{\stackrel{~}{P}(\lambda )}=\mathrm{}`$.
Notice that this proof can be easily generalized to graphs in which there is a positive measure subset $`S`$ such that: $`lim_{\lambda 1}\overline{\chi (S)\stackrel{~}{F}(\lambda )}=|S|`$. Indeed in an analogous way it can be proven that:
$$\overline{P}\underset{\lambda 1}{lim}\overline{\chi (S^{})\stackrel{~}{P}(\lambda )}=\mathrm{}S^{}S,|S^{}|>0$$
(31)
We will call mixed transient on the average a TOA graphs having a positive measure subset $`S`$ such that:
$$\underset{\lambda 1}{lim}\overline{\chi (S)\stackrel{~}{F}(\lambda )}=|S|.$$
(32)
while a graph will be called pure TOA, if:
$$\underset{\lambda 1}{lim}\overline{\chi (S)\stackrel{~}{F}(\lambda )}<|S|SV,|S|>0$$
(33)
Examples of mixed and pure TOA graphs are shown respectively in Fig. 2 and Fig. 3. From the previous proof for mixed TOA graphs we have $`\overline{P}=\mathrm{}`$; let us now study the behavior of $`\overline{P}`$ on pure TOA graphs. We define $`k`$ as
$$k=\underset{SV,|S|>0}{sup}\underset{\lambda 1}{lim}\overline{\chi (S)\stackrel{~}{F}(\lambda )}|S|^1$$
(34)
and since the graphs is pure TOA, $`k<1`$. For each $`0<\lambda ^{}<1`$ we introduce $`S_\lambda ^{}V`$ as the set of the vertices $`i`$ such that $`\stackrel{~}{F}_i(\lambda ^{})>k`$. Exploiting the property that $`\stackrel{~}{F}_i(\lambda )`$ is an increasing function of $`\lambda `$ we have $`\overline{\chi (S_\lambda ^{})\stackrel{~}{F}_i(\lambda )}>k|S_\lambda ^{}|`$ and then $`lim_{\lambda 1}\overline{\chi (S_\lambda ^{})\stackrel{~}{F}_i(\lambda )}>k|S_\lambda ^{}|`$. From (34) we obtain that $`S_\lambda ^{}`$ has zero measure, i. e. it must be $`|S_\lambda ^{}|=0`$. Exploiting the definition (7) we have, for all $`iV`$, $`\stackrel{~}{P}_i(\lambda )(1\lambda )^1`$ and we obtain for $`\overline{\stackrel{~}{P}(\lambda ^{})}`$:
$$\overline{\stackrel{~}{P}(\lambda ^{})}=\overline{\chi (\overline{S}_\lambda ^{})\stackrel{~}{P}(\lambda ^{})}+\overline{\chi (S_\lambda ^{})\stackrel{~}{P}(\lambda ^{})}\overline{\chi (\overline{S}_\lambda ^{})(1\stackrel{~}{F}(\lambda ^{}))^1}+|S_\lambda ^{}|(1\lambda ^{})^1|\overline{S}_\lambda ^{}|(1k)^1(1k)^1$$
(35)
Taking the limit $`\lambda ^{}1`$, we have that for pure TOA graphs $`\overline{P}`$ is finite.
This prove can be generalized to graphs in which there is a positive measure subset $`S`$ such that for all $`S^{}S`$, $`|S^{}|>0`$, $`lim_{\lambda 1}\overline{\chi (S^{})\stackrel{~}{F}(\lambda )}|S^{}|`$ obtaining
$$\underset{\lambda 1}{lim}\overline{\chi (S^{})\stackrel{~}{P}(\lambda )}<\mathrm{}.S^{}S,|S^{}|>0$$
(36)
## VII Random walks and infrared properties of the Gaussian model
The generating function $`\stackrel{~}{P}_i(\lambda )`$ is strictly connected with the correlation functions of the Gaussian model (10) by the following equation:
$$\stackrel{~}{P}_i(\lambda )=\underset{t=0}{\overset{\mathrm{}}{}}\lambda ^t(Z^1A)_{ii}^t=(1\lambda Z^1A)_{ii}^1=[Z\lambda ^1(\mathrm{\Delta }+(1\lambda )\lambda ^1Z)^1]_{ii}$$
(37)
and taking the limit $`\lambda 1`$ we have:
$$\underset{\lambda 1}{lim}\stackrel{~}{P}_i(\lambda )=z_i\underset{\mu 0}{lim}(\mathrm{\Delta }+\mu Z)_{ii}^1$$
(38)
In the invariance of the limit: $`lim_{\mu 0}(\mathrm{\Delta }+\mu )_{ii}^1`$ under a local rescaling of the masses is proven. In particular we have that if $`lim_{\mu 0}(\mathrm{\Delta }+\mu )_{ii}^1`$ is finite than $`lim_{\mu 0}(\mathrm{\Delta }+\mu Z)_{ii}^1<\mathrm{}`$ while when the first limit diverges the latter also diverges. Therefore on locally recurrent graphs we have $`lim_{\mu 0}(\mathrm{\Delta }+\mu )_{ii}^1=\mathrm{}`$ $`i`$ while for locally transient graphs $`lim_{\mu 0}(\mathrm{\Delta }+\mu )_{ii}^1<\mathrm{}`$ $`i`$.
Let us now consider the average generating function. From (37) we get:
$$\underset{\lambda 1}{lim}\overline{\stackrel{~}{P}(\lambda )}=\underset{\mu 0}{lim}\overline{\mathrm{Tr}}[Z(\mathrm{\Delta }+\mu Z)^1]$$
(39)
and since the connectivity of $`𝒢`$ is bounded we get:
$$\underset{\mu 0}{lim}\overline{\mathrm{Tr}}(\mathrm{\Delta }+\mu Z)^1\underset{\lambda 1}{lim}\overline{\stackrel{~}{P}(\lambda )}z_{max}\underset{\mu 0}{lim}\overline{\mathrm{Tr}}(\mathrm{\Delta }+\mu Z)^1$$
(40)
Exploiting the universality properties of the Gaussian model , we have that $`lim_{\mu 0}\overline{\mathrm{Tr}}(\mathrm{\Delta }+\mu Z)^1`$ is finite if and only if $`lim_{\mu 0}\overline{\mathrm{Tr}}(L+\mu )^1<\mathrm{}`$, where $`L`$ is a generalized Laplacian given by (3). Finally from inequalities (40) we get that $`lim_{\mu 0}\overline{\mathrm{Tr}}(L+\mu )^1=lim_{\mu 0}\overline{\varphi _i\varphi _i}`$ diverges if $`\overline{P}=\mathrm{}`$ i.e. on ROA and mixed TOA graphs, while it is finite on pure TOA graphs, where $`\overline{P}<\mathrm{}`$.
## VIII Separability and statistical independence
In this section we prove and discuss an important property characterizing mixed TOA graphs which allows to simplify the study of statistical models on these very inhomogeneous structures. We will show that in this case the graph $`𝒢`$ can be always decomposed in a pure TOA subgraph $`𝒮`$ and a ROA subgraph $`\overline{𝒮}`$ with independent jumping probabilities by cutting a zero measure set of links $`𝒮\{(i,j)E|i𝒮j\overline{𝒮}\}`$. The separability property implies that the two subgraphs are statistically independent and that their thermodynamic properties can be studied separately. Indeed, the partition functions referring to the two subgraphs factorize .
As a first step, from definition (32) the set of vertices $`V`$ of a mixed TOA graph $`𝒢`$ can always be decomposed in two complementary subsets $`S`$ and $`\overline{S}`$ such that
$$\frac{\overline{\chi (S^{})\stackrel{~}{F}(1)}}{|S^{}|}<1$$
(41)
for all $`S^{}S`$ with $`|S^{}|>0`$ and
$$\frac{\overline{\chi (S^{\prime \prime })\stackrel{~}{F}(1)}}{|S^{\prime \prime }|}=1$$
(42)
for all $`S^{\prime \prime }\overline{S}`$ with $`|S^{\prime \prime }|>0`$.
To this decomposition we can associate the two subgraphs $`𝒮`$ and $`\overline{𝒮}`$ defined as follows: $`𝒮`$ has $`S`$ as set of vertices and its links are all the links $`(i,j)𝒢`$ such that $`i,jS`$; in the same way $`\overline{𝒮}`$ has $`\overline{S}`$ as set of vertices and its links are all the links $`(i,j)𝒢`$ such that $`i,j\overline{S}`$. Let us now prove that the measure of the boundary $`|𝒮|`$ (16) is zero.
We introduce $`B_S`$, the border set of $`𝒮`$, defined as the set of the vertices $`iS`$ with $`(i,j)𝒮`$ for some $`j`$ while we will call $`B_{\overline{S}}`$ the border set of $`\overline{𝒮}`$. Proving that $`|𝒮|=0`$ is equivalent to show that the measure of $`B_S`$ and $`B_{\overline{S}}`$ is zero. Indeed we have $`|𝒮|_r|B_S|_rz_{max}|𝒮|_r`$ and $`|𝒮|_r|B_{\overline{S}}|_rz_{max}|𝒮|_r`$, where $`|B_S|_r`$ and $`|B_{\overline{S}}|_r`$ are the number of sites in $`B_S`$ and $`B_{\overline{S}}`$ contained in the sphere $`S_{o,r}`$.
Let us suppose that $`𝒮0`$ and that $`|B_S|0`$, $`|B_{\overline{S}}|0`$. From (31) and (36) we have:
$$\underset{\lambda 1}{lim}\overline{\chi (B_S)\stackrel{~}{P}(\lambda )}\mathrm{}$$
(43)
and
$$\underset{\lambda 1}{lim}\overline{\chi (B_{\overline{S}})\stackrel{~}{P}(\lambda )}=\mathrm{}$$
(44)
We will now derive a relation between $`\overline{\chi (B_S)\stackrel{~}{P}(\lambda )}`$ and $`\overline{\chi (B_{\overline{S}})\stackrel{~}{P}(\lambda )}`$ which can not be satisfied if (43) and (44) hold, leading to a contradiction. This implies that $`|𝒮|=0`$.
Let us evaluate $`\stackrel{~}{P}_i(\lambda )`$ in a site $`iB_S`$
$$\stackrel{~}{P}_i(\lambda )=\underset{t}{}\lambda ^tp_{ii}^t=\underset{t}{}\lambda ^t\underset{jk}{}p_{ik}p_{kj}^{t2}p_{ji}\underset{t}{}\lambda ^t\underset{jB_{\overline{S}}}{}p_{ij}p_{jj}^{t2}p_{ji}$$
(45)
where in the inequality we do not consider the terms in which $`jk`$ and $`jB_{\overline{S}}`$ Exploiting the fact that $`p_{ij}1/z_{max}`$ we get:
$$\stackrel{~}{P}_i(\lambda )\frac{\lambda ^2}{z_{max}^2}\underset{t}{}\lambda ^{t2}\underset{jB_{\overline{S},i}}{}p_{jj}^{t2}=\frac{\lambda ^2}{z_{max}^2}\underset{jB_{\overline{S},i}}{}\stackrel{~}{P}_{jj}(\lambda )$$
(46)
where $`B_{\overline{S},i}`$ is the set of the nearest neighbors sites of $`i`$ which belong to $`B_{\overline{S}}`$. If we take the average over the sites $`iB_S`$ we have:
$$\overline{\chi (B_S)\stackrel{~}{P}(\lambda )}\frac{\lambda ^2}{z_{max}^2}\underset{r\mathrm{}}{lim}\frac{\chi _i(B_S)}{N_{o,r}}\underset{iS_{o,r}}{}\underset{jB_{\overline{S},i}}{}\stackrel{~}{P}_{jj}(\lambda )\frac{\lambda ^2}{z_{max}^2}\underset{r\mathrm{}}{lim}\frac{\chi _j(B_{\overline{S}})}{N_{o,r}}\underset{jS_{o,r}}{}\stackrel{~}{P}_{jj}(\lambda )=\frac{\lambda ^2}{z_{max}^2}\overline{\chi (B_{\overline{S}})\stackrel{~}{P}(\lambda )}$$
(47)
If we take the limit $`\lambda 1`$ we have:
$$\underset{\lambda 1}{lim}\overline{\chi (B_S)\stackrel{~}{P}(\lambda )}\frac{1}{z_{max}^2}\underset{\lambda 1}{lim}\overline{\chi (B_{\overline{S}})\stackrel{~}{P}(\lambda )}$$
(48)
Expressions (48), (43) and (44) can not be satisfied at the same time and therefore one must have $`|𝒮|=0`$
Finally we have to prove that $`𝒮`$ is a pure TOA graph and $`\overline{𝒮}`$ is a ROA graph, i.e. we introduce the restricted jumping probability on $`S`$ and $`\overline{S}`$ $`p_{ij}^S`$ and $`p_{ij}^{\overline{S}}`$, given by $`p_{ij}^S=p_{ij}`$ if $`i,jS`$, $`p_{ij}^S=0`$ otherwise and an analogous definition for $`p_{ij}^{\overline{S}}`$. Then we show that $`S`$ and $`\overline{S}`$ with the new jumping probabilities $`p_{ij}^S`$ and $`p_{ij}^{\overline{S}}`$ are respectively pure TOA and ROA.
More generally for a walker on $`𝒮`$ starting from $`i`$, we call $`P_{ij}^S(t)`$ the probability of reaching site $`j`$ in $`t`$ steps and $`F_{ij}^S(t)`$ the probability of reaching $`j`$ for the first time in $`t`$ steps. We will prove that:
$$\overline{\stackrel{~}{P}^S(\lambda )}=\overline{\underset{t=0}{\overset{\mathrm{}}{}}\lambda ^tP_i^S(t)}=\overline{\chi (S)\stackrel{~}{P}(\lambda )}|S|^1\overline{\stackrel{~}{F}^S(\lambda )}=\overline{\underset{t=0}{\overset{\mathrm{}}{}}\lambda ^tF_i^S(t)}=\overline{\chi (S)\stackrel{~}{F}(\lambda )}|S|^1$$
(49)
where the average of $`\stackrel{~}{P}^S(\lambda )`$ and of $`\stackrel{~}{F}^S(\lambda )`$ is taken considering $`𝒮`$ as the whole graph. Analogous equations hold also for $`\overline{𝒮}`$. From (49), (41) and (36) we easily obtain that $`\overline{P}^S<\mathrm{}`$ i.e. $`𝒮`$ is pure TOA, while if we call $`P_{ij}^{\overline{S}}(t)`$ and $`F_{ij}^{\overline{S}}(t)`$ the probabilities for a random walk on $`\overline{𝒮}`$, we get $`\overline{F}^{\overline{S}}=1`$, i.e. $`\overline{𝒮}`$ is ROA.
To prove equations (49) first we have to show that:
$$\overline{\stackrel{~}{P}^S(\lambda )}=\underset{t=0}{\overset{\mathrm{}}{}}\lambda ^t\overline{P^S(t)}\lambda <1$$
(50)
Equation (50) implies that the thermodynamic average and the sum over the discretized times $`t`$ commute when $`\lambda <1`$. To prove (50) notice that for all $`\lambda <1`$ we have:
$$\overline{\stackrel{~}{P}^S(\lambda )}=\underset{r\mathrm{}}{lim}\underset{iS_{o,r}}{}N_{o,r}^1\left(\underset{t=0}{\overset{\overline{t}}{}}\lambda ^tP_i^S(t)+\underset{t=\overline{t}}{\overset{\mathrm{}}{}}\lambda ^{\overline{t}}P_i^S(t)\right)=\underset{t=0}{\overset{\overline{t}}{}}\lambda ^t\overline{P^S(t)}+\underset{r\mathrm{}}{lim}\underset{iS_{o,r}}{}N_{o,r}^1\underset{t=\overline{t}}{\overset{\mathrm{}}{}}\lambda ^tP_i^S(t)$$
(51)
Now $`_{iS_{o,r}}N_{o,r}^1_{t=\overline{t}}^{\mathrm{}}\lambda ^{\overline{t}}P_i^S(t)\lambda ^{\overline{t}}(1\lambda )^1`$ and letting in (51) $`\overline{t}\mathrm{}`$ we get (50). Obviously an analogous equation holds also for $`F_i^S(t)`$, $`P_i(t)`$ and $`F_i(t)`$. Then we can prove (49) showing that:
$$\overline{P^S(t)}=\overline{\chi (S)P(t)}|S|^1\overline{F^S(t)}=\overline{\chi (S)\stackrel{~}{F}(t)}|S|^1$$
(52)
We define $`d(i,B_S)`$ as the distance between $`i`$ and the cutset $`B_S`$: $`d(i,B_S)=inf_{kB_S}r_{i,k}`$ and will call $`S_t`$ the subset of $`S`$ such that: $`S_t=\{iS|d(i,B_S)t\}`$, exploiting the boundedness of the coordination number we get:
$$|S_t|<(z_{max})^t|B_S|=0$$
(53)
since $`|B_S|=0`$. Taking the average of $`P_i^S(t)`$ we have:
$$\overline{P^S(t)}=\overline{\chi (\overline{S}_t)P^S(t)}+\overline{\chi (S_t)P^S(t)}$$
(54)
Now $`\overline{\chi (S_t)P^S(t)}|S_t|=0`$, and then $`\overline{P^S(t)}=\overline{\chi (\overline{S}_t)P^S(t)}`$. Finally exploiting the fact that on $`\overline{S}_t`$ we have $`P_i^S(t)=P_{ii}(t)`$, we obtain (52). Following analogous steps we obtain the equality for $`\overline{F^S(t)}`$ and for the averages $`\overline{P^{\overline{S}}(t)}`$ and $`\overline{F^{\overline{S}}(t)}`$ defined on $`\overline{S}`$.
## IX Discussion and conclusions
In this paper we have presented a systematic mathematical analysis of the Type Problem for random walks on infinite graphs by considering return probabilities averaged over all sites. After showing that recurrence and transience on the average (ROA and TOA) do not in general coincide with the corresponding local properties, we prove that TOA has to be splitted in two complementary subcases, the pure and the mixed one. Then we show that a mixed TOA graph can always be decomposed in a ROA and a pure TOA subgraphs by cutting a zero measure set of links. This property has deep physical implications, since it allows to decompose a statistical model defined on a mixed TOA graph in two thermodynamically independent models defined respectively on the ROA and pure TOA subgraphs.
In conclusion, we introduced an exhaustive classification of infinite networks in terms on their average recurrence and transience properties, stating the Type Problem on the Average. This classification is the relevant one in the study of thermodynamic properties of statistical models on inhomogeneous structures.
Fig. 1: The NTD tree: the distances between the ramifications increase exponentially. This graph is locally transient and recurrent on the average.
Fig 2: A mixed TOA graph: the cubic lattice is a pure TOA graph while the hairs are ROA.
Fig 3: A pure TOA graph, i.e. an inhomogeneous Bethe lattice in which the distance between ramifications can be 1 or 2.
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# Universal Metric Spaces According to W. Holsztynski
## Abstract
In this note we show, following W. Holsztynski, that there is a continuous metric $`d`$ on $``$ such that any finite metric space is isometrically embeddable into $`(,d)`$.
Let $``$ be a family of metric spaces. A metric space $`U`$ is said to be a *universal space* for $``$ if any space from $``$ is (isometrically) embeddable in $`U`$.
Fréchet proved that $`\mathrm{}^{\mathrm{}}`$ (the space of all bounded sequences of real numbers endowed with the $`sup`$ norm) is a universal space for the family $``$ of all separable metric spaces. Later, Uryson constructed an example of a separable universal space for this $``$ ($`\mathrm{}^{\mathrm{}}`$ is not separable).
In this note we establish the following result which is Theorem 5 in .
###### Theorem 1.
There exists a metric $`d`$ in $``$, inducing the usual topology, such that every finite metric space embeds in $`(,d)`$.
Our proof essentially follows the original Holsztynski’s approach .
We say that a metric space $`(X,d)`$ is $`\epsilon `$*dispersed* if $`d(x,y)\epsilon `$ for all $`xy`$ in $`X`$ $`(\epsilon >0)`$. Clearly, any $`\epsilon `$– metric space is also $`\epsilon ^{}`$–dispersed for any positive $`\epsilon ^{}<\epsilon `$. The following proposition will be used to construct universal spaces for particular families $``$.
###### Proposition 1.
Let $`f:XY`$ be a continuous surjection from a metric space $`(X,d)`$ onto a metric space $`(Y,D)`$. Then $`(X,d_\epsilon )`$ where
$$d_\epsilon (x,y)=\mathrm{max}\{\mathrm{min}\{d(x,y),\epsilon \},D(f(x),f(y))\}$$
(1)
is a universal space for the family of $`\epsilon `$–dispersed subspaces of $`(Y,D)`$ and metrics $`d`$ and $`d_\epsilon `$ are equivalent on $`X`$.
###### Proof.
$`d_\epsilon `$ is a distance function. Indeed, $`d_\epsilon `$ is symmetric and $`d_\epsilon (x,y)=0`$ if and only if $`x=y`$. We have
$`\mathrm{max}\{`$ $`\mathrm{min}\{d(x,y),\epsilon \},D(f(x),f(y))\}+\mathrm{max}\{\mathrm{min}\{d(y,z),\epsilon \},D(f(y),f(z))\}`$
$`D(f(x),f(y))+D(f(y),f(z))D(f(x),f(z))`$
and
$`\mathrm{max}\{`$ $`\mathrm{min}\{d(x,y),\epsilon \},D(f(x),f(y))\}+\mathrm{max}\{\mathrm{min}\{d(y,z),\epsilon \},D(f(y),f(z))\}`$
$`\mathrm{min}\{d(x,y),\epsilon \}+\mathrm{min}\{d(y,z),\epsilon \}=\mathrm{min}\{d(x,y)+d(y,z),\epsilon \}`$
$`\mathrm{min}\{d(x,z),\epsilon \}`$
Hence, $`d_\epsilon (x,y)+d_\epsilon (y,z)d_\epsilon (x,z)`$.
Let $`Z`$ be an $`\epsilon `$–dispersed subspace of $`Y`$. Since $`f`$ is surjective, for any $`zZ`$, there is $`x_zX`$ such that $`f(x_z)=z`$. Let $`X^{}=\{x_z:zZ\}`$. By (1), $`d_\epsilon (x,y)=D(f(x),f(y))`$ for all $`x,yX^{}`$. Thus $`f`$ establishes an isometry between $`(Z,D)`$ and $`(X^{},d_\epsilon )`$.
In what follows, $``$ is the family of all finite metric spaces.
We define
$$I^n=\{\overline{x}=(x_1,\mathrm{},x_n)^n:0x_in,\mathrm{\hspace{0.33em}1}in\}$$
and $`J_n=[n1,n]`$ for $`n1`$. $`I^n`$ is a metric space with the distance function
$$D(\overline{x},\overline{y})=\mathrm{max}\{|x_iy_i|:1in\},$$
and $`J_n`$ is a metric space with the usual distance.
Let $`(X,d)`$ be a finite metric space. We define
$$p=|X|,q=\text{Diam}(X),r=\epsilon ^1,$$
where $`\epsilon =\mathrm{min}\{d(x,y):xy\}`$, and $`n=\mathrm{max}\{p,q,r\}`$. Clearly, $`(X,d)`$ is $`\frac{1}{n}`$–dispersed.
###### Proposition 2.
*(The Kuratowski embedding)* $`(X,d)`$ is embeddable into $`I^n`$.
###### Proof.
Let $`X=\{x_1,\mathrm{},x_p\}`$. We define $`f:XI^n`$ by
$$f(x_i)=(d(x_i,x_1),\mathrm{},d(x_i,x_p),\underset{np\mathrm{\hspace{0.33em}0}\text{’s}}{\underset{}{0,\mathrm{},0}}),$$
for $`1ip`$. (Since $`n\text{Diam}(X)`$, $`f(x_i)I^n`$.) We have, by the triangle inequality,
$$D(f(x_k),f(x_m))=\underset{j}{\mathrm{max}}\{|d(x_k,x_j)d(x_m,x_j)|\}d(x_k,x_m).$$
On the other hand, $`|d(x_k,x_j)d(x_m,x_j)|=d(x_k,x_m)`$ for $`j=m`$. Therefore, $`D(f(x_k),f(x_m))=d(x_k,x_m)`$ for all $`1k,mp`$.
Let $`f_n`$ be a continuous surjection from $`J_n`$ onto $`I^n`$ (a “Peano curve” \[1, IV(4)\]) and $`d_n(x,y)`$ be the distance function on $`J_n`$ defined by (1) for $`\epsilon =\frac{1}{n}`$. Note, that $`d_n`$ is equivalent to the usual distance on $`J_n`$. By Proposition 1, $`J_n`$ is a universal space for any $`\frac{1}{n}`$–dispersed subspace of $`I^n`$. By Proposition 2, any finite metric space is embeddable in $`(J_n,d_n)`$ for some $`n`$.
It is easy to show that there is a continuous distance function on $``$ that coincides with $`d_n(x,y)`$ on $`J_n`$ for all $`n`$. Indeed, let $`d_1(x,y)`$ and $`d_2(x,y)`$ be two continuous distance functions on intervals $`[a,b]`$ and $`[b,c]`$, respectively. Then $`d(x,y)`$ defined by
$$d(x,y)=\{\begin{array}{cc}d_1(x,y),\hfill & \text{if }x,y[a,b]\text{,}\hfill \\ d_2(x,y),\hfill & \text{if }x,y[b,c]\text{,}\hfill \\ d_1(x,b)+d_2(b,y),\hfill & \text{if }x[a,b]\text{ and }y[c,d]\text{.}\hfill \end{array}$$
is a continuous distance function on $`[a,c]`$. In fact, thus defined $`d`$ is equivalent to the usual metric on $`[a,c]`$.
By applying this process consecutively, we obtain a required distance function on $``$.
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# 1 Information cascade: an analogy with hydrodynamic turbulence
## 1 Information cascade: an analogy with hydrodynamic turbulence
The idea of a heterogeneous market that consists of traders acting on different time horizons was first advanced by Müller et al.<sup>?</sup> who investigated the absolute values of FX returns on different time scales. They observed that the price changes over longer time intervals have a stronger influence on those over shorter time intervals than conversely. This has been interpreted as an information flow from long-term to short-term traders which directly influences the volatility on different time scales.
Though being very different in its physical nature, the motion of a turbulent fluid medium is also governed by a hierarchical process in which energy flows from large to small spatial scales<sup>?</sup>. Interestingly enough, this formal analogy between FX dynamics and turbulence leads to similar statistical behavior of both phenomena<sup>?</sup>. For the FX dynamics, the hierarchical process can be implemented in form of a volatility cascade, which has recently been visualized by means of wavelet techniques<sup>?</sup>.
In the present contribution we take this hypothesis as granted and construct a stochastic multiplicative volatility cascade driving the FX price process. In the stochastic cascade model (SCM), that we present here, the volatility at an instant of time $`t`$ is assumed to be created by processes on different time scales reflecting the heterogeneity of the market dynamics. The information flow between traders on different time horizons is modeled by a directional interaction scheme between the levels of the cascade corresponding to the different time scales. More precisely, the flow of information is accounted for by an update of stochastic “transfer” factors relating the subsequent levels of the cascade.
## 2 Stochastic cascade model
We denote the price at time $`t`$ by $`p_t`$ and define the return as $`r_t\mathrm{log}p_t\mathrm{log}p_{t1}`$. In the proposed SCM the returns are described by a stochastic volatility model, i.e.
$$r_t=\sigma _t\xi _t,$$
(1)
where $`\xi _t`$ are i.i.d. Student random numbers with 3 degrees of freedom. The volatility $`\sigma _t`$ is governed by a hierarchical process that reflects the heterogeneity of the market on its different time horizons. Each level $`k`$ of the cascade corresponds to a time horizon of duration $`\tau ^{(k)}`$. The largest horizon $`\tau ^{(0)}`$ representing the first level at the top of the cascade is typically of the order of a few months while the level $`m`$ at the bottom of the cascade has the smallest horizon $`\tau ^{(m)}`$ on which dealers may act and which is of the order of minutes. For the sake of simplicity we assume that the horizon ratio between neighboring levels is constant, i.e. $`\tau ^{(k)}/\tau ^{(k1)}=p<1`$ is independent of $`k`$.
We now assign to each level $`k`$ of the hierarchy a volatility $`\sigma _t^{(k)}`$ which is determined by the respective volatility on the level $`k1`$ and a time-dependent random factor $`a_t^{(k)}`$ that will be specified below. Hence the level volatilities are recursively defined by:
$$\sigma _t^{(k)}=a_t^{(k)}\sigma _t^{(k1)}.$$
(2)
As a toy model we assume that the market dynamics is completely determined by the behavior of the dealers and not subject to exogenous time-dependent influences. Then, the volatility $`\sigma ^{(0)}`$ on the largest horizon is time independent. The volatility on the shortest horizon determines the returns, as given by eq. (1) with $`\sigma _t=\sigma _t^{(m)}`$. Using eq. (2) we find:
$$\sigma _t=\sigma ^{(0)}\underset{k=1}{\overset{m}{}}a_t^{(k)}.$$
(3)
## 3 Updating rule for the volatility factors $`a_t^{(k)}`$
The time-dependence of the factors $`a_t^{(k)}`$ results from the following renewal scheme: At the initial time $`t_0`$ the factors $`a_{t_0}^{(k)}`$ are drawn from independent lognormal distributions $`LN(a_k,\lambda _k^2)`$ with mean values $`a_k`$ and variances $`\lambda _k^2`$, $`k=1,\mathrm{},m`$. At a later time $`t_{n+1}`$ the factor $`a_{t_{n+1}}^{(1)}`$ at the top of the hierarchy maintains the corresponding value $`a_{t_n}^{(1)}`$ at the preceding time $`t_n`$ with a probability $`1w^{(1)}`$ and else is drawn from $`LN(a_1,\lambda _1^2)`$. In the latter case all the subsequent factors $`a_{t_{n+1}}^{(k)}`$, $`k>1`$ are also renewed, i.e. they aquire independent random values from the respective distributions $`LN(a_k,\lambda _k^2)`$. If $`a_{t_{n+1}}^{(1)}`$ coincides with $`a_{t_n}^{(1)}`$ the factor $`a_{t_{n+1}}^{(2)}`$ will be renewed with a probability $`w^{(2)}`$. These rules apply through the whole cascade down to $`k=m`$: A renewal at some level $`k`$ entails one at all higher levels $`k^{}>k`$; if no renewal has taken place up to level $`k`$, the coefficient at level $`k+1`$ will be renewed with probability $`w^{(k+1)}`$ (drawn from the distribution $`LN(a_{k+1}\lambda _{k+1}^2)`$) or, with probability $`1w^{(k+1)}`$, its value remains constant in time, $`a_{t_{n+1}}^{(k+1)}=a_{t_n}^{(k+1)}`$.
The renewal probabilities $`w^{(k)}`$ are given by:
$$w^{(k)}=1(1p^{mk})/(1p^{mk+1}),k=1,\mathrm{},m.$$
(4)
According to this renewal scheme the mean life-times of a factor at level k (measured in no. of time steps of length $`\tau (m)`$) is given by $`p^{km}`$ where $`p`$ is the scaling factor of the time horizons as defined above.
## 4 Results
The time series defined by eqs. (1 \- 4) have been simulated and the parameters have been adjusted in such a way that the model simultaneously reproduces the empirical return distributions (Fig. 3), the scaling laws for the moments<sup>?</sup>, and the autocorrelation function of the volatility (Fig. 4) of USD/CHF FX spot rates. The distributions of returns (cf. Fig. 3) exhibit the well-known heavy tails, which lose weight with increasing time interval. The autocorrelation function of the absolute returns (Fig. 4) shows the characteristic slowly decaying tail. In both cases the SCM simulation (full line) reproduces the observed data (dotted curves) very well. The scaling behavior of the moments and the cross correlation function of the volatility are presented in an extended version of the paper<sup>?</sup>.
To summarize, the SCM is a hierarchical time series model where a net information flow from long to short time horizons is implemented in terms of a random unidirectional action of the volatility at a given time horizon on that at the next shorter one. With only three adjustable parameters this model is able to reproduce different characteristic properties of intra-day FX price series.
Note that the existence of different groups of traders and the volatility updating mechanism are elements similar to those found in market microstructure models. Therefore the SCM may provide a link between market microstructure models and more conventional models.
One of the authors (S.G.) gratefully acknowledges suppport from the Swiss National Science Foundation.
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# 1 Introduction
## 1 Introduction
Quantum chromodynamics (QCD) on the lattice predicts a phase transition of confined hadronic matter into deconfined quark and gluon matter (called the Quark Gluon Plasma state – QGP) at a critical temperature $`T_c`$ $``$ 150 MeV, respectively at energy density $`ϵ_c`$ $``$ 1 GeV/fm<sup>3</sup> . The order of the transition is parameter dependent . An investigation of relevant observables in heavy ion collisions as a function of energy and/or the impact parameter<sup>2</sup><sup>2</sup>2Provided that there is a unique assignement between impact parameter and QGP phase transition. of the collision could reveal this transition, through a discontinuous behaviour of many QGP signatures at the transition point.
An anomalous suppression of the $`J/\mathrm{\Psi }`$ meson predicted to be a signature of Quark Gluon Plasma formation has been measured to occur in the ratio of the $`J/\mathrm{\Psi }`$ over the Drell Yan (DY) process in Pb+Pb collisions at $`\sqrt{s}`$=17 GeV investigated as a function of transverse energy ($`E_T`$) . In S+U collisions at $`\sqrt{s}`$=19 GeV and in the most peripheral Pb+Pb collisions, the $`(J/\mathrm{\Psi })/DY`$ ratio agrees with expectations . The ratio $`(J/\mathrm{\Psi })/DY`$ is relevant for the investigation of the $`J/\mathrm{\Psi }`$ suppression, under the assumption that $`J/\mathrm{\Psi }`$ production in these collisions is a hard process.
Recent measurements of the dimuon invariant mass m($`\mu ^+\mu ^{}`$) spectrum between the $`\varphi `$ and the $`J/\mathrm{\Psi }`$ mass (Intermediate Mass Region=IMR) revealed a dimuon enhancement above expectation, which is increasing with the number of nucleons participating in the collision (N) . This enhancement can be understood as due to an excess of $`D\overline{D}`$ production<sup>3</sup><sup>3</sup>3With ’$`D\overline{D}`$’ we denote the number of $`D`$ and $`\overline{D}`$ hadrons which were simultaneously found within the acceptance of the NA50 experiment ( $`D\overline{D}=(D+\overline{D})_{acc}/2`$ ). as suggested by several features of the data, e.g. the shape of the mass, rapidity, angular and transverse momentum distributions of the dimuons . The interpretation of the IMR enhancement as due to open charm is not unique though, because the open charm was extracted through a fit to the dimuon continuum and the $`D\overline{D}`$ meson signal was not directly identified. Recent work suggesting that the seen enhancement could be due to rescattering of D mesons in nuclear matter is not supported by the data . An other possible interpretation is that the IMR excess could be due to thermal dimuons . Non-perturbative effects are known to play a role in heavy flavour hadron production in elementary reactions showing up in deviations of data from perturbative QCD calculations . Based on this fact, one could expect that theoretical investigation of non-perturbative effects and different reaction dynamics as in the plasma phase may result in an enhancement of open charm production in nuclear reactions over perturbative QCD expectations.
If the total charm produced in heavy ion collisions indeed deviates from the perturbative QCD expectations for a hard process as suggested by the NA50 data, it follows that the $`(J/\mathrm{\Psi })/DY`$ ratio is not the proper quantity for the search for the $`J/\mathrm{\Psi }`$ suppression as signature of Quark Gluon Plasma formation in nuclear reactions. It is only the ratio $`(J/\mathrm{\Psi })`$/(total $`c\overline{c}`$) that matters. We therefore investigate here first the dependence of the $`J/\mathrm{\Psi }`$ and the $`D\overline{D}`$ yields per collision on N. We further investigate the dependence of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio on N and on the length of the nuclear matter traversed by the $`J/\mathrm{\Psi }`$, as well as the dependence of both charm and strangeness production on the initial energy density reached in the collision.
## 2 N dependence of open and closed charm yields in Pb+Pb collisions at $`\sqrt{s}`$=17 GeV
### 2.1 N dependence of the Drell Yan yield
Calculation details
The dependence of the Drell Yan (DY) yield per nucleus-nucleus collision in arbitrary units produced in Pb+Pb collisions at $`\sqrt{s}`$=17 GeV per N+N collision, on the transverse energy of the collision has been measured by the NA50 collaboration (figure 7 in ). In this figure it is shown that the theoretically expected DY yield, assuming DY production is a hard process, does partly reproduce the measured one; the deviations at low transverse energy are understood to result from properties of the lead (Pb) nucleus, in particular from the different radii of the proton and neutron distributions in Pb . It is therefore justified to use the theoretically calculated $`E_T`$ dependence of DY yield per collision, to estimate the dependence of $`J/\mathrm{\Psi }`$ and $`D\overline{D}`$ yields on the number of participating nucleons in the collision.
However, since the deviations of the very low statistics measurement of DY yield at high $`E_T`$, from the theoretical estimated DY yield seen in figure 7 of cannot be understood, it would be important to measure the DY yield per collision with high statistics in the high $`E_T`$ region. In this way the last drop in the $`(J/\mathrm{\Psi })/DY`$ ratio above N $``$ 360 can be experimentally verified.
### 2.2 N dependence of the $`D\overline{D}`$ yield
Calculation details
The NA50 collaboration observed an excess (E) of the measured over the expected $`D\overline{D}/DY`$ ratio in S+U and Pb+Pb collisions at $`\sqrt{s}`$=19, 17 GeV, which increases with the number of participants N (figure 12 and table 4 in ). If we fit the S+U and Pb+Pb E points of the above figure to a function $`f=cN^\alpha `$, we find that the excess is increasing with N as $`N^{(\alpha =0.45\pm 0.11)}`$ ($`\chi ^2`$/Degrees Of Freedom=1.7, DOF=7). The $`N`$ dependence of the excess E of the $`D\overline{D}/DY`$ production in S+U collisions at $`\sqrt{s}`$=19 GeV and Pb+Pb collisions at $`\sqrt{s}`$=17 GeV over expectations, reflects the $`N`$ dependence of the $`D\overline{D}/DY`$ ratio. This results from the fact, that all other quantities involved in the definition of E , do not depend on $`N`$. Therefore the $`N`$ dependence of the $`D\overline{D}`$ production yield is given by the $`N`$ dependence of the quantity
$$n_{D\overline{D}}=En_{DY}(D\overline{D}/DY)n_{DY}$$
(1)
where $`n_{D\overline{D}}`$, $`n_{DY}`$ denote the yields of $`D\overline{D}`$ and DY per collision in arbitrary units. The arbitrary units are due to the fact that NA50 did not published absolute yields per collision of the $`J/\mathrm{\Psi }`$, $`DY`$ and $`D\overline{D}`$ separately, corrected for losses due to e.g. acceptance, as a function of N, $`E_T`$. We suggest that it would be important to do so.
The DY yield used for the above calculation has been extracted from the theoretical curve shown in figure 7 in , at the transverse energy ($`E_T`$) points in which the $`D\overline{D}`$ excess factor E has been measured. The $`E_T`$ points corresponding to the excess factor E were extracted, by interpolating between the $`E_T`$ values given in table 1 of as a function of the mean impact parameter, at the values of mean impact parameter for which the factor E has been measured (listed in table 1 of ). For the most central and the most peripheral points, for which no mean b are given in table 1 of , we estimated the values of b, from the values of $`N`$ as a function of b for Pb+Pb collisions calculated by . These calculations agree with the values (N,b) estimated by NA50, when compared in their common range.
Though the $`D\overline{D}`$ measured by NA50 represents the joint probability that a $`D`$ and a $`\overline{D}`$ are both found in the experimental acceptance, the N dependence of it, is expected to be the same or similar to the N dependence of the sum $`D+\overline{D}`$, respectively of the $`D`$, and of the $`\overline{D}`$ hadrons. It is however possible that the above N-dependences deviate from eachother, in case of a higher than 1 charm pair multiplicity per event with charm (see for a numerical estimation of the charm multiplicity per event in central Pb+Pb collisions at 158 A GeV). Therefore, for clarity, the N dependence of the total extrapolated $`D+\overline{D}`$ yield in nuclear collisions should be estimated including acceptance corrections by NA50.
Results and discussion
The resulting $`D\overline{D}`$ yield in arbitrary units (figure 1) increases as $`N^{(\alpha =1.70\pm 0.12)}`$ ($`\chi ^2/DOF`$ = 2.5, DOF=7). This N behaviour indicates that $`D\overline{D}`$ production in Pb+Pb collisions at $`\sqrt{s}`$=17 GeV, did not establish yet equilibrium, in which case a proportionality with $`N`$ –assuming N to be proportional to the volume of the source <sup>4</sup><sup>4</sup>4 The assumption N $``$ V, is based on the observation that the freeze out volume V of the particle source is found to be proportional to N .– is expected ($`\alpha `$= 1). This appears justified because the temperature in the collision zone –assuming local thermalization of light particles– drops with time and the mean temperature expected to be reached in these collisions of the order $``$ $`10^2`$ MeV, is much lower than the mass of charm quarks and/or charmed hadrons.
### 2.3 N dependence of the $`J/\mathrm{\Psi }`$ yield
Calculation details
In the following we estimate the $`J/\mathrm{\Psi }`$ yield per collision as a function of $`N`$, at the same $`N`$ values where the $`D\overline{D}`$ was measured. The $`N`$ dependence of the $`J/\mathrm{\Psi }`$ yield per collision is given by the N dependence of the quantity :
$$n_{J/\mathrm{\Psi }}=((J/\mathrm{\Psi })/DY)n_{DY}$$
(2)
where $`n_{J/\mathrm{\Psi }}`$, $`n_{DY}`$ denote the yields of $`J/\mathrm{\Psi }`$ and DY in arbitrary common units. The $`(J/\mathrm{\Psi })/DY`$ values have been extracted from figure 4 of at the $`E_T`$ values where the E factor has been measured, interpolating between the different points.
Results and discussion
The resulting $`J/\mathrm{\Psi }`$ yield per collision produced in Pb+Pb collisions in arbitrary units (figure 2) increases like $`N^{(\alpha =0.70\pm 0.04)}`$ ($`\chi ^2/DOF`$=1.43, DOF=7). This N dependence indicates an increasing $`J/\mathrm{\Psi }`$ dissociation with higher centrality. The strength of this dissociation as measured by the $`\alpha `$ parameter, is higher for the $`J/\mathrm{\Psi }`$ as compared to any other hadron<sup>5</sup><sup>5</sup>5Deuterons have an even smaller $`\alpha `$ parameter, but they are not elementary hadrons and are weekly bound (see discussion in ). produced in these collisions for example as compared to antiprotons. For the latter, a large annihilation cross section with baryons is expected and there is indeed experimental evidence that they are absorbed with increasing centrality in Pb+Pb collisions ($`\alpha `$($`\overline{p}`$)=0.80 $`\pm `$ 0.04, ($`\chi ^2/DOF`$=1.0, $`DOF`$=3) at y=3.7, $`p_T`$=0 <sup>6</sup><sup>6</sup>6We extracted here the $`\alpha `$ parameter for $`\overline{p}`$, after quadratically adding the statistical and a 5% systematic error.).
The $`J/\mathrm{\Psi }`$ multiplicity as a function of N extracted with an other method , agrees with the here presented results within the errors.
## 3 The $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in nuclear collisions
### 3.1 The N dependence of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in nuclear collisions
Calculation details
Assuming that the IMR excess is due partly or solely to open charm, allows us to search for an anomalous suppression of $`J/\mathrm{\Psi }`$ as compared to the open charm production. The $`N`$ dependence of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in Pb+Pb and S+U collisions at $`\sqrt{s}`$ of 17 and 19 GeV, estimated as:
$$(J/\mathrm{\Psi })/D\overline{D}((J/\mathrm{\Psi })/DY)/(D\overline{D}/DY)((J/\mathrm{\Psi })/DY)/E$$
(3)
in arbitrary units due to the E factor in equation (3), is a decreasing function of $`N`$ (figure 3). The $`(J/\mathrm{\Psi })/DY`$ in S+U collisions was taken from . Note that possible deviations of the DY yield from its theoretical calculation (as seen in figure 7 in ), do not drop out in the $`((J/\mathrm{\Psi })/DY)/(D\overline{D}/DY)`$ ratio shown here, because the $`D\overline{D}/DY`$ –unlike the $`(J/\mathrm{\Psi })/DY`$– was calculated by NA50 not using the minimum bias theoretical DY yield values but the measured ones.
In order to show the influence of the very last drop of $`(J/\mathrm{\Psi })/DY`$, on the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio, the $`(J/\mathrm{\Psi })/DY`$ ratio divided by $`N^{0.45\pm 0.11}`$ is plotted as a function of N in figure 4. This quantity resembles the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio
$$(J/\mathrm{\Psi })/D\overline{D}((J/\mathrm{\Psi })/DY)/N^{0.45\pm 0.11}$$
(4)
in arbitrary units, because $`N^{0.45\pm 0.11}`$ is the found N dependence of the $`D\overline{D}/DY`$ ratio. The open points of figure 4 are extracted from the ’minimum bias analysis’ results of figure 4 in <sup>7</sup><sup>7</sup>7’Minimum bias analysis’ in NA50, means that the DY for the $`(J/\mathrm{\Psi })/DY`$ ratio, was determined using the theoretically estimated DY yield per collision as a function of $`E_T`$ and the measured $`dN/dE_T`$ vs $`E_T`$ spectrum of minimum bias trigger events (see ).. The closed points show the $`((J/\mathrm{\Psi })/DY)/N^{0.45}`$ calculated here, at the N values where the $`D\overline{D}`$ excess factor was measured.
Results and discussion
The $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in Pb+Pb collisions found as shown in equation (3), decreases with $`N`$ as $`N^{(\alpha =0.79\pm 0.14)}`$ ($`\chi ^2/DOF`$ = 3.3, $`DOF`$=7), respectively like $`N^{(\alpha =1.17\pm 0.14)}`$ ($`\chi ^2/DOF`$=1.54, DOF=6) when the first point is not fitted<sup>8</sup><sup>8</sup>8The first point of the $`D\overline{D}/DY`$ enhancement factor E lyes significantly above the $`N^\alpha `$ function fit to the E distribution (figure 12 in ).. The $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in S+U collisions (figure 3) decreases with $`N`$ as $`N^{(\alpha =0.62\pm 0.22)}`$ ($`\chi ^2/DOF`$=0.69, $`DOF`$=3). The $`(J/\mathrm{\Psi })/D\overline{D}`$ found as shown in equation (4) when fitted to the function $`N^\alpha `$ until N=380, gives $`N^{(\alpha =1.07\pm 0.07)}`$, ($`\chi /DOF`$=0.93, DOF=7).
If the $`J/\mathrm{\Psi }`$ is completely dissociated in a quark gluon plasma and is formed later mainly through $`c`$ and $`\overline{c}`$ quark coalescence, we expect that the N dependence of the ratio $`(J/\mathrm{\Psi })/D\overline{D}`$ –rather than the $`(J/\mathrm{\Psi })/(D\overline{D})^2`$– reflects the N dependence of the volume of the charm environment . This is due to the expectation that because of the very low cross section of charm production at these energies, there is most often just one $`c\overline{c}`$ pair per event containing charm, whatever N. Then the probability to form a $`J/\mathrm{\Psi }`$ from coalescence is proportional to $`(J/\mathrm{\Psi })/D\overline{D}`$ and inversely proportional to the volume of the particle source –made up by $`u\overline{u}d\overline{d}s\overline{s}`$ quarks and gluons– within which the $`c`$ and $`\overline{c}`$ quarks scatter. Assuming this volume is proportional to N (see footnote 3), one would expect that $`(J/\mathrm{\Psi })/D\overline{D}`$ decreases as $`N^1`$, as actually observed.
In this case, one can use the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio to extract the absolute value of the volume of its environment with a coalescence model. The ’charm’ coalescence volume would reflect partly the QGP hot spot volume and partly the hadronic source volume from which hadrons with charm and anticharm can also form a $`J/\mathrm{\Psi }`$. If the absolute yields per collision of $`J/\mathrm{\Psi }`$ and $`D\overline{D}`$ as a function of N, needed for this calculation would be published by NA50, the charm coalescence volume could be calculated.
Figure 4 suggests that the coalescence picture could hold for the full N range of Pb+Pb collisions up to N=380. Obviously, it would be better to use the $`D\overline{D}`$ data themselves instead of the N parametrization, if high enough statistics would be available.
On the other hand, if the multiplicity of charm quarks is high enough that often more than 1 charm quark pair per event with charm is produced, then it is the ratio $`(J/\mathrm{\Psi })/(D\overline{D})^2`$ which is expected to be inversely proportional to the volume of the charm source (this is exactly the case if $`d`$ coalescence out of $`p`$ and $`n`$ is investigated in a baryon rich source). The N dependence of the $`(J/\mathrm{\Psi })/(D\overline{D})^2`$ ratio, which would be relevant in the above discussed case is N<sup>(α=-2.26±0.48)</sup>, $`(\chi ^2/DOF`$=2.5, DOF=7) respectively N<sup>(α=-3.1±0.24)</sup>, $`(\chi ^2/DOF`$=1.2, DOF=6) if the first point is not fitted. The question on the absolute multiplicity of charm in nuclear reactions, should be answered by experiment.
### 3.2 The L dependence of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in nuclear collisions
Calculation details
The two distributions of $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio for S+U and Pb+Pb collisions in figure 3, are measured at different energies, and therefore they cannot be compared in terms of their absolute yields but only with respect to their shapes. In order to compare their absolute yields, the data from figure 5 of will be used. There the $`(J/\mathrm{\Psi })/DY`$ ratio in p+A, S+U and Pb+Pb collisions is shown as a function of L, all normalised to the same energy ($`\sqrt{s}`$=19 GeV) and corrected for the isospin dependence of DY production. The parameter L, is the length that the $`J/\mathrm{\Psi }`$ traverses through nuclear matter. In order to convert figure 5 of to the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio as a function of the L parameter, the $`(J/\mathrm{\Psi })/DY`$ ratio data points have been divided by the E factor as described in equation (3). The correlation of the L parameter with N and b for Pb+Pb collisions, has been estimated using the theoretical calculation of .
Results and discussion
The L dependence of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in arbitrary units in p+A, S+U and Pb+Pb collisions calculated here, is shown in figure 5, together with the $`(J/\mathrm{\Psi })/DY`$ ratio published in . The closed points show the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in S+U and Pb+Pb collisions extracted as indicated in equation 3. The open squares and circles show the $`(J/\mathrm{\Psi })/DY`$ ratio in S+U and Pb+Pb collisions from . The open stars, show both the L dependence of the $`(J/\mathrm{\Psi })/DY`$ as well as the L dependence of the $`(J/\mathrm{\Psi })/D\overline{D}`$ in p+A collisions which are the same, since the factor E has the value 1 for the latter.
The $`J/\mathrm{\Psi }`$ over the $`D\overline{D}`$ production investigated as a function of the volume through which the $`J/\mathrm{\Psi }`$ traverses (V $``$ L<sup>3</sup>) is suppressed as compared to the shape of the exponential fit going through the $`(J/\mathrm{\Psi })/D\overline{D}`$ p+A data, in both the S+U and Pb+Pb collisions at all $`L`$ points (figure 5).
The energy density of the lowest S+U point has been estimated to be $``$ 1.1 GeV/fm<sup>3</sup> , which is comparable to the predicted critical energy density for the QGP phase transition of $``$ 1 GeV/fm<sup>3</sup>. A similar energy density of 1.2 GeV/fm<sup>3</sup> has been estimated to be reached in the most peripheral Pb+Pb collisions measured by NA50.
In the following we investigate the initial energy density, rather than only the volume of the particle source (V $``$ $`L^3`$), as a critical parameter for the appearance of the QGP phase transition.
## 4 The $`ϵ`$ dependence of charm and strangeness in nuclear collisions
### 4.1 Charm
Calculation details
We estimate here the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio as a function of the energy density $`ϵ`$. For this purpose we use part of the data shown in figure 7 in . There the ratio of $`((J/\mathrm{\Psi })/DY)_{measured}`$ over the $`((J/\mathrm{\Psi })/DY)_{expected}`$ is shown. The ’$`((J/\mathrm{\Psi })/DY)_{expected}`$’, is taken to be the exponential fit seen in figure 5, which represents the ’normal’ $`J/\mathrm{\Psi }`$ dissociation (i.e. understood without invoking QGP formation). Dividing these data points by $`N^{0.45\pm 0.11}`$, and normalising the distribution of S+U and Pb+Pb points to the p+p and p+A data as in figure 5, we estimate the $`((J/\mathrm{\Psi })/D\overline{D})`$ ratio over the expectation expressed by the above mentioned exponential curve, which fits the $`((J/\mathrm{\Psi })/D\overline{D})`$ data points for p+p, p+d and p+A collisions.
Results and discussion
The result of this calculation is shown in figure 6 in logarithmic scale and in figure 7 (a) in linear scale. It demonstrates a deviation of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio both in S+U and Pb+Pb collisions, from the p+p and p+A expectation curve, occuring above $`ϵ`$ $``$ 1 GeV/fm<sup>3</sup>. The logarithmic scale is shown to reveal small changes in the slope of the $`((J/\mathrm{\Psi })/D\overline{D})`$ distribution as a function of $`ϵ`$, appearing at $`ϵ`$ $``$ 2.2 and 3.2 GeV/fm<sup>3</sup>.
### 4.2 Strangeness
Figure 7 compares the two QGP signatures of $`J/\mathrm{\Psi }`$ suppression and of strangeness enhancement. For this purpose we represent all data points as a function of the estimated energy density. Note that the energy density as critical scale variable, has the advantage that unlike the temperature, it is defined irrespective of whether equilibrium is reached in the collisions studied.
Figure 7 (b) shows the multiplicity of kaons per event ($`K^+`$, but also some $`K_s^0`$ data scaled to $`K^+`$ are shown) divided by the effective volume of the particle source at thermal freeze out in the center of mass frame, as a function of the initial energy density. The effective volume represents the part of the real source volume, within which pions are correlated with each other (the so called ’homogenity’ volume in the literature ). The effective volume is smaller than but proportional to the real source volume. For a more precise calculation of the freeze out source volume a detailled model is needed. Here we estimate the effective volume at thermal freeze out $`V_{thermal}`$ based on measurements. The (smaller) effective volume at chemical freeze out $`V_{chemical}`$, is not experimentally measured, we give however an estimate of the ratio $`V_{chemical}/V_{thermal}`$. Note that we compare the kaon data without rescaling for the different energy between AGS and SPS<sup>9</sup><sup>9</sup>9 The total $`K^+`$ multiplicity in p+p collisions, increases by a factor of $``$ 5, from 11.1 to 158 GeV per nucleon ..
Calculation details
The effective volume V of the particle source has been estimated in the center of mass frame, assuming cylindrical shape of the source:
$`V=(\pi Radius_{cylinder}^2)Length_{cylinder}`$ $``$
$`V=(\pi 4R_{side}^2)(\sqrt{12}R_{long})`$
where $`R_{side}`$ is a measure of the transverse radius and $`R_{long}`$ is a measure of the longitudinal radius of the particle source, and the factors $`4`$ and $`\sqrt{12}`$ arise from the definition of $`R_{side}`$, $`R_{long}`$ . The $`R_{side}`$ and $`R_{long}`$ values for central Au+Au collisions at 10.8 A GeV and for central Pb+Pb collisions at 158 A GeV have been taken from . We dont use the more elaborated estimation of the homogenity volume given in , because the $`R_{ol}`$ component is not given in . Based on the data of we estimated the effective volume of the source at thermal freeze out in central Au+Au collisions at 10.8 A GeV ($`V1949fm^3`$) and central Pb+Pb collisions at 158 A GeV ($`V6532fm^3`$). The effective volume increases by a factor of 3.35 from AGS to SPS energy.
Based on the temperature at thermal and chemical freeze out which has been estimated from measurements using thermal models , and the above estimated volumes at thermal freeze out, we can further estimate the volume at chemical freeze out. For this we assume that the relation $`VT^3`$, which holds in the universe for massless particles in thermal equilibrium and for adiabatic expansion , holds approximately for heavy ion collisions at AGS and SPS energy. Then from the temperature values at thermal and chemical freeze out given in averaged over all models, we find that the ratio $`V_{chemical}/V_{thermal}`$(AGS Si+Au 14.6 A GeV) =0.45 and $`V_{chemical}/V_{thermal}`$(SPS Pb+Pb 158 A GeV) = 0.28. Using the volume at chemical freeze out as estimated above, would stretch the $`K/V`$ ratio in figure 7 (b) between SPS and AGS, by a factor $``$ 1.6 apart. We dont use these values in figure 7, because the above calculation is model dependent, e.g. the assumption of massless particles is not met, while the assumption of thermal equilibrium may not be true.
The ratio K/N is expected to be proportional to the number density of kaons $``$ K/V, (V=volume), assuming V $``$ N (for justification of this assumption see footnote 3, and ). Based on this expectation, we estimated here the K/V ratios from the K/N ratios, by normalising the K/N ratios to the K/V value of the most central Au+Au events of E866 respectively Pb+Pb events of NA49, for which the value of the volume has been estimated above.
The kaon data from Au+Au collisions at 11.1 A GeV (E866 and E802 experiments) and from Pb+Pb collisions at 158 A GeV (NA49 experiment) are kaon multiplicities extrapolated to full acceptance. Therefore NA49 and E866 data are absolutely normalised. We estimated K/N from the NA49 experiment using the kaon multiplicities from and the number of wounded nucleons from as available<sup>10</sup><sup>10</sup>10 We take N equal to the number of wounded nucleons for NA49, because it is used in all other experiments presented here, and allows for a straightforward geometrical interpretation of N., otherwise we used the N estimated from the experimental baryon distribution .
The data from NA52 and WA97 have been measured in a small phase space acceptance and have been scaled here arbitrarily, in order to match the NA49 data in figure 7. This scaling is justified since all NA52, NA49 and WA97 measurements are kaons produced in Pb+Pb collisions at 158 A GeV, and ’extrapolates’ the NA52 and WA97 data to the NA49 full acceptance multiplicities, allowing for comparison of the shapes of the distributions. It is assumed, that the N and $`ϵ`$ dependence of kaons does not change significantly with the phase space acceptance.
In order to calculate the energy density we have performed the following steps. The energy density for all colliding systems has been estimated using the Bjorken formula and data given in . The transverse radius of the overlapping region of the colliding nuclei is found as: $`R_{trans}=1.13(N/2)^{1/3}`$, where N is the total number of participant nucleons. The formation time was taken 1 fm/c .
For the E866 experiment, lacking $`E_T`$ values, (but with measured E<sub>forward</sub>), we used instead of $`(dE_T/d\eta )_{ycm}`$, ($`ycm`$=midrapidity) the total energy of the nucleons participating in the collision ($`E_{tot,part}=N_{projectile,participants}E_{beam}`$), assuming the proportionality $`(dE_T/d\eta )_{ycm}E_{tot,part}`$, and we further normalised the results in such a way that the maximum energy density of our estimate matches the absolute value of the maximal achieved energy density in the most central Au+Au events at this energy of 1.3 GeV/fm<sup>3</sup>, given in .
NA52 measures $`E_T`$ near midrapidity (y $``$ 3.3). These values were used to estimate the energy density and the results have been normalised to the maximum energy density reached in Pb+Pb collisions at the same centrality of $`ϵ_{max}`$=3.2 GeV/fm<sup>3</sup>, extracted by NA49 in . Parametrizing the dependence of the energy density on the number of participants found from the NA52 data as described above, we estimated the energy density corresponding to the N values of the WA97 and the NA49 kaon measurements, given in . Data from S+S collisions taken from and are also shown.
To estimate the systematic error on the energy density found with the above methods, we calculated the energy density in Pb+Pb collisions at 158 A GeV, using the VENUS 4.12 event generator. We estimated with VENUS the $`(dE_T/d\eta )_{ycm}`$ at $`ycm`$=midrapidity and the number of participant nucleons and used them to find the energy density from the Bjorken formula . The deviation of the energy density calculated with VENUS $`(dE_T/d\eta )_{ycm}`$ from the energy density found using the NA52 transverse energy measurements is $``$ 30% of the latter. The deviation of the energy density calculated with VENUS $`dE_T/d\eta `$ from the energy density found using the total energy of the participant nucleons and of the newly produced particles, (which is similar to the method used to estimate the energy density for the AGS data), over the latter energy density, is at the same level.
In this context, it appears important for a more precise comparison of data as a function of $`ϵ`$, that experiments publish together with the number of participants also the $`dE_T/d\eta `$ at midrapidity for each centrality region, for both nucleus+nucleus and for p+p collisions, estimated by models or measured if available (e.g. in NA49).
Results and discussion
Figure 7 (b) suggests that kaons below $`ϵ`$ $``$ 1 GeV/fm<sup>3</sup> did not reach equilibrium, while this seems to be the case above. Indeed kaons produced in Au+Au collisions at 11.1 A GeV and in very peripheral Pb+Pb collisions at 158 A GeV , increase faster than linear with N, indicating non thermal kaon production, while they increase nearly proportional to N above $`ϵ`$ $``$ 1 GeV/fm<sup>3</sup> . The connection of strangeness equilibrium and the QGP phase transition has been discussed e.g. in . There it is shown, that strangeness in heavy ion collisions is expected to reach equilibrium values if the system runs through a QGP phase, while this is less probable in a purely hadronic system.
Figure 7 demonstrates that both the $`J/\mathrm{\Psi }`$ and kaon production exhibit a dramatic change above the energy density of $``$ 1 GeV/fm<sup>3</sup>. While the equilibration of strange particles as suggested by their $`N^1`$ dependence above 1 GeV/fm<sup>3</sup>, could in principle also be due to equilibrium reached in a hadronic environment, the combined appearance of this effect and of the $`(J/\mathrm{\Psi })/D\overline{D}`$ suppression at the same energy density value is a striking result, indicating a change of phase above $`ϵ_c`$=1 GeV/fm<sup>3</sup>.
The expectation for the shape of the $`J/\mathrm{\Psi }`$ suppression as a function of energy density are three succesive drops of the $`J/\mathrm{\Psi }`$ ; a drop by $``$ 8$`\%`$ due to $`\psi ^{^{}}`$ dissociation, a drop by $``$ 32$`\%`$ due to the $`\chi _c`$ dissociation and a drop by $``$ 100$`\%`$ due to the $`J/\mathrm{\Psi }`$ dissociation. All these without taking into account regeneration of $`J/\mathrm{\Psi }`$ through other processes. These can be e.g. coalescence of charm quarks or $`J/\mathrm{\Psi }`$ not travelling through the plasma. The $`\psi ^{^{}}`$ feeds only 8$`\%`$ of the total $`J/\mathrm{\Psi }`$’s and can therefore hardly be observed as a break in the $`J/\mathrm{\Psi }`$ production.
The absolute value of the energy density $`ϵ`$ and therefore of the N values at which these changes could be observed is not exactly given by the models. The critical energy densities for the dissociation of the states $`\mathrm{\Psi }^{^{}}`$, $`\chi _c`$ and $`J/\mathrm{\Psi }`$ could even be so near to each other that no clear multistep behaviour is seen in $`(J/\mathrm{\Psi })/D\overline{D}`$.
Figure 7 suggests that the breaks in the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio at $`ϵ`$ $``$ 2.2 and 3.2 GeV/fm<sup>3</sup>, are less dramatic than the change above $`ϵ`$ $``$ 1 GeV/fm<sup>3</sup>. Therefore, all bound $`c\overline{c}`$ states could be dissociated at similar energy densities, which lye near 1 GeV/fm<sup>3</sup>.
Alternatively, the $`\psi ^{^{}}`$ and the $`\chi _c`$ could dissociate above $`ϵ`$ $``$ 1 GeV/fm<sup>3</sup> and the dissociation of the $`J/\mathrm{\Psi }`$ could start at $`ϵ`$=2.2 GeV/fm<sup>3</sup>, if we interpret the change in the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio, below and above $`ϵ`$=2.2 GeV/fm<sup>3</sup>, as a step behaviour. In this context, the steep drop of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in the bin(s) of largest N (figures 4, 6, 7), cannot be interpreted in a natural way. The steps of $`(J/\mathrm{\Psi })/D\overline{D}`$ remain to be established through a direct measurement of $`J/\mathrm{\Psi }`$ and $`D\overline{D}`$ absolute yields as a function of ($`E_T`$, N, $`ϵ`$).
In the above discussed picture, three QGP signatures appear in nuclear collisions at energy density larger than $``$ 1 GeV/fm<sup>3</sup>:
a) $`J/\mathrm{\Psi }`$ suppression (figure 7 (a)) –which could be due to bound $`c\overline{c}`$ states dissociation–
b) strangeness enhancement (figure 7 (b)), possibly due to equilibration of $`s\overline{s}`$ in QGP as opposed to hadrons,
c) the invariant mass m($`e^+e^{}`$) excess at m below the $`\rho `$ mass , possibly due to a $`\rho `$ change and/or to increased production of the lowest mass glueball state in QGP .
This coincidence of QGP signatures, suggests a change of phase at $`ϵ`$ $``$ 1 GeV/fm<sup>3</sup> as expected .
From the above discussion it follows, that a direct measurement of open charm production in nuclear collisions appears essential for the physics of the Quark Gluon Plasma phase transition. Furthermore, if enhanced over expectations, open charm in nuclear collisions defies theoretical understanding.
## 5 Possibilities for future measurements
A measurement of open and closed charm production in Pb+Pb collisions as a function of energy below the SPS top energy of $`\sqrt{s}=17`$ GeV searching for the disappearance of the seen $`J/\mathrm{\Psi }`$ suppression in central Pb+Pb collisions at a certain $`\sqrt{s}`$, could prove clearly the QGP phase transition. Using the same nuclei at different $`\sqrt{s}`$ and looking only at central collisions, differences due to different nuclear profiles drop out. No currently existing or planned experiment at SPS is however able to perform this measurement without major upgrades, though one proposal (NA6i) could significantly improve the identification of open charm production through a better determination of the decay vertex . An upgrade of NA50/NA6i, or alternatively a completely new experiment, could possibly achieve this goal. The study could also be performed at the Relativistic Heavy Ion Collider (RHIC) using lower energy and/or large and small nuclei, and in fixed target experiments at RHIC favoured because of higher luminosity as compared to the collider mode, important for a low energy scan.
It would be also important (and easier than the above) to measure the $`J/\mathrm{\Psi }`$, $`D\overline{D}`$ and $`DY`$ absolute yields per collision, below $`ϵ`$=1 GeV/fm<sup>3</sup>, by using the most peripheral (not yet investigated) Pb+Pb collisions or collisions of lighter nuclei at the highest beam energy at SPS ($`\sqrt{s}`$ = 17,19 GeV).
An other piece of information important for the understanding of charm production in nuclear collisions would be the direct comparison of the $`(J/\mathrm{\Psi })/DY`$ and the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratios in nuclear collisions at $`\sqrt{s}<`$ 19 GeV and in $`p+\overline{p}`$ collisions at the Tevatron. Tevatron reaches an energy density similar to or larger than the one estimated in very central S+S collisions at 200 A GeV . Therefore it would supply a comparison for these points and a continuation of the absorption line fitted through the p+p and p+A data measured by NA50 (figure 5), or otherwise. Differences due to the change of dominant production mechanisms of charm in $`p\overline{p}`$ collisions as compared to A+B, p+p, can be accounted for theoretically. A high $`E_T`$ cut could additionally help in sorting out ’central’ $`p+\overline{p}`$ collisions. This comparison should be done possibly in the very same dimuon mass region for all processes (also DY) e.g. using Monte Carlo’s tuned to $`p+\overline{p}`$ Tevatron data.
This comparison would answer the question, if the energy density is indeed the only critical variable for the appearance of a thermalised QGP state with 3 effective flavours u,d,s, or whether there is also a critical volume (e.g. as measured by the L variable: V $``$ L<sup>3</sup>). Furthermore, at present the comparison of nuclear collisions to p+p and p+A data is done at the same energy and not at the same energy density. This issue is important, since if for example the energy density is the only critical scale variable, the QGP should be formed also in elementary collisions like $`p\overline{p}`$ at a higher beam energy and the same energy density.
Further it is important to search for thresholds in the production of many particles e.g. $`\mathrm{\Omega }`$, which was found to be enhanced by a factor 15 above p+A data in Pb+Pb collisions at 158 A GeV in the energy density region corresponding to the green stars in figure 7 (b). Similarly interesting would be a measurement of the invariant mass of $`e^+e^{}`$ in low energy densities.
## 6 Conclusions
In this letter, consequences resulting from the viable possibility that the dimuon invariant mass (m($`\mu ^+\mu ^{}`$)) enhancement, measured by the NA50 experiment in the intermediate mass region (IMR): $`m(\varphi )<m(\mu ^+\mu ^{})<m(J/\mathrm{\Psi })`$, in S+U and Pb+Pb collisions at $`\sqrt{s}`$ 19, 17 GeV, reflects a $`D\overline{D}`$ enhancement over expectations, are worked out.
The dependence of the $`J/\mathrm{\Psi }`$ and the $`D\overline{D}`$ yields per collision in Pb+Pb collisions on the mean number of participants has been estimated. This dependence reveals the non thermal features of charm production at this energy. The $``$ $`N^{0.7}`$ dependence of the $`J/\mathrm{\Psi }`$ yield (figure 2) suggests strong dissociation of $`J/\mathrm{\Psi }`$ with higher centrality. The dissociation is stronger than the absorption seen in any other hadron, e.g. $`\overline{p}`$ in Pb+Pb collisions. The N dependence of the $`D\overline{D}`$ yield of $`N^{1.7}`$ (figure 1) indicates also non-thermal open charm production at this energy, showing up in an excess rather than reduction as compared to the thermal expectation.
If the dimuon excess observed by NA50 is partly or solely due to open charm, it is appropriate to search for an anomalous suppression of $`J/\mathrm{\Psi }`$ as compared to the total open charm production, rather than to the DY process. We therefore investigated here the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio in Pb+Pb collisions and we find it to decrease approximately as $``$ $`N^1`$ (figures 3, 4). This is the N dependence expected for the $`J/\mathrm{\Psi }`$ if it were completely dissociated in quark gluon matter and were later dominantly formed through $`c\overline{c}`$ quark coalescence, assuming N $``$ volume of the $`c\overline{c}`$ environment and charm quark multiplicity of one in events with charm. In that case, based on coalescence arguments, the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio could be used to estimate the volume of the charm environment, which may reflect partly the size of the quark gluon plasma. This is probable under the assumption that the final measured $`J/\mathrm{\Psi }`$ is dominated by the $`J/\mathrm{\Psi }`$ originating from $`c\overline{c}`$ pairs which travel through the plasma volume, an assumption which may hold only for large plasma volumes, i.e. for the most central collisions.
A further consequence of a possible open charm enhancement is that the $`J/\mathrm{\Psi }`$ over the $`D\overline{D}`$ ratio appears to be suppressed already in S+U collisions as compared to p+A collisions, unlike the $`(J/\mathrm{\Psi })/DY`$ ratio (figure 5). The $`\psi ^{^{}}/D\overline{D}`$ ratio would also be additionally suppressed as compared to the $`\psi ^{^{}}/DY`$ in both S+U and Pb+Pb collisions. These phenomena could be interpreted as onset of dissociation of bound charm states above energy density $`ϵ`$ $``$ 1 GeV/fm<sup>3</sup>.
We estimated and compared the dependence of the $`(J/\mathrm{\Psi })/D\overline{D}`$ ratio and of the kaon multiplicity per volume (assuming K/N $``$ K/V = kaon number density) in several collisions and $`\sqrt{s}`$ as a function of the initial energy density. We find that both the kaon number density and the ratio $`(J/\mathrm{\Psi })/D\overline{D}`$ exhibit dramatic changes at the energy density of 1 GeV/fm<sup>3</sup>, as demonstrated in figure 7. This is the main result of this paper.
It follows that three major QGP signatures ($`s\overline{s}`$ enhancement, $`\rho `$ changes and $`J/\mathrm{\Psi }`$ suppression) all appear above the energy density of $``$ 1 GeV/$`fm^3`$, which is the critical energy density for the QGP phase transition according to lattice QCD.
This discussion underlines the importance of a direct measurement of open charm production in nuclear collisions, and of other experimental investigations proposed in section 5, for the understanding of ultrarelativistic nuclear reactions and the dynamics of the Quark Gluon Plasma phase transition.
Acknowledgments
I would like to thank Prof. P. Minkowski, Prof. K. Pretzl, Prof. U. Heinz and Prof. J. Rafelski for stimulating discussions, and Dr. C. Cicalo, Dr. O. Drapier, Dr. C. Gerschel, Dr. E. Scomparin, Dr. P. Seyboth, Dr. F. Sikler, Dr. U. Wiedemann, and especially Dr. J.Y. Ollitrault for clarifying discussions on their data and/or for communicating their results to me.
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# Finite-temperature properties of Pb(Zr1-xTix)O3 alloys from first principles
\[
## Abstract
A first-principles-derived approach is developed to study finite-temperature properties of Pb(Zr<sub>1-x</sub>Ti<sub>x</sub>)O<sub>3</sub> (PZT) solid solutions near the morphotropic phase boundary (MPB). Structural and piezoelectric predictions are in excellent agreement with experimental data and direct first-principles results. A low-temperature monoclinic phase is confirmed to exist, and is demonstrated to act as a bridge between the well-known tetragonal and rhombohedral phases delimiting the MPB. A successful explanation for the large piezoelectricity found in PZT ceramics is also provided.
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Ferroelectric perovskite A(BB<sup>′′</sup>)O<sub>3</sub> alloys are of growing importance for a variety of device applications , and are also of great current fundamental interest since little is known about the effects responsible for their anomalous properties. A good example of an A(B,B<sup>′′</sup>)O<sub>3</sub> solid solution that is of both fundamental and technological importance is the Pb(Zr<sub>1-x</sub>Ti<sub>x</sub>)O<sub>3</sub> system. Usually denoted as PZT, this mixed-cation alloy is currently in widespread use in piezoelectric transducers and actuators . Its phase diagram exhibits a morphotropic phase boundary (MPB) separating a region with a tetragonal ground state ($`x`$ $`>`$ 0.52) from a region with rhombohedral symmetry ($`x`$ $`<`$ 0.45) .
High piezoelectric response is experimentally found in ceramics of PZT around the MPB. The origins of this large piezoelectric response are unclear. On the one hand, semi-empirical simulations predict that the large experimental value of the $`d_{33}`$ piezoelectric coefficient results mainly from the large value of $`d_{33}`$ that a single-crystal PZT would exhibit . On the other hand, recent first-principles calculations have found that the $`d_{33}`$ coefficient of a tetragonal single crystal of Pb(Zr<sub>0.5</sub>Ti<sub>0.5</sub>)O<sub>3</sub> are estimated to be three times smaller than the experimental value obtained for ceramics at low temperature.
Furthermore, recent synchrotron x-ray powder diffraction studies have revealed the existence of an unexpected low-temperature monoclinic phase of PZT at $`x`$=0.48 , which implies that the phase diagram of PZT is more complex than previously thought. This monoclinic phase may act as a second-order transitional bridge between the tetragonal phase, for which the electrical polarization $`𝐏`$ lies along the pseudo-cubic direction, and the rhombohedral phase, for which $`𝐏`$ is along the pseudo-cubic direction. If this is indeed the case, the polarization of the monoclinic phase continuously rotates as the composition $`x`$ decreases in the MPB region . Such a continuous rotation has yet to be observed.
Obviously, accurate simulations are needed to understand the properties of perovskite alloys in general, and of PZT in particular. Since the beginning of the present decade, first-principles methods have emerged as a powerful tool for investigating properties of ferroelectric systems theoretically (see and references therein). However, these methods are essentially restricted to the study of the zero-temperature properties of small cells, while accurate and interesting predictions of alloy properties would require calculations on much larger cells at finite temperature. Ideally one desires a computational scheme with the capability of predicting the properties of “real” perovskite alloy systems at finite temperature, with the accuracy of the first-principles methods.
The purpose of this letter is to demonstrate that it is possible to develop such a scheme, and to apply it to study the finite-temperature behavior of PZT in the vicinity of the MPB. Remarkably, we find that the existence of an intermediate monoclinic phase emerges naturally from this approach. Moreover, the theory provides a novel and successful explanation for the large piezoelectric response of PZT near the MPB, thereby explaining and resolving the previous theoretical difficulties in obtaining agreement with the known experimental values of the piezoelectric coefficients.
Our scheme is based on the construction of an effective Hamiltonian from first-principles calculations. A ferroelectric effective Hamiltonian must include the ferroelectric local soft mode and the strain variables, since ferroelectric transitions are accompanied by a softening of a polar phonon mode and by the appearance of a strain. An alloy effective Hamiltonian must also include the compositional degrees of freedom. We propose to incorporate all such degrees of freedom by writing the total energy $`E`$ as a sum of two energies,
$`E(\{𝐮_i\},\{𝐯_i\},\eta _H`$ , $`\{\sigma _j\})=E_{\mathrm{ave}}(\{𝐮_i\},\{𝐯_i\},\eta _H)`$ (2)
$`+E_{\mathrm{loc}}(\{𝐮_i\},\{𝐯_i\},\{\sigma _j\}),`$
where $`𝐮_i`$ is the local soft mode in unit cell $`i`$; $`\{𝐯_i\}`$ are the dimensionless local displacements which are related to the inhomogeneous strain variables inside each cell ; $`\eta _H`$ is the homogeneous strain tensor; and the $`\{\sigma _j\}`$ characterize the atomic configuration of the alloy. That is, $`\sigma _j`$=+1 or $`1`$ corresponds to the presence of a B or B<sup>′′</sup> atom, respectively, at lattice site $`j`$ of the A(B<sub>1-x</sub>B<sup>′′</sup><sub>x</sub>)O<sub>3</sub> alloy. The energy $`E_{\mathrm{ave}}`$ depends only on the soft mode and strain variables. The $`\{\sigma _j\}`$ parameters are thus incorporated into the second energy term $`E_{\mathrm{loc}}`$.
For $`E_{\mathrm{ave}}`$, we generalize the analytical expression successfully used in Ref. for simple ABO<sub>3</sub> systems to the case of an alloy. This generalization simply consists in using the virtual crystal approximation (VCA) , i.e., in replacing A(B<sub>1-x</sub>,B<sup>′′</sup><sub>x</sub>)O<sub>3</sub> by a uniform but composition-dependent “virtual” ABO<sub>3</sub> system. $`E_{\mathrm{ave}}`$ thus consists of five parts: a local-mode self-energy, a long-range dipole-dipole interaction, a short range interaction between soft modes, an elastic energy, and an interaction between the local modes and local strain .
For $`E_{\mathrm{loc}}`$, we propose an expression that includes (i) the onsite effect of alloying on the self-energy up to the fourth order in the local-mode vector $`𝐮_i`$; and (ii) intersite contributions which are linear in $`𝐮_i`$ and in $`𝐯_i`$:
$`E_{\mathrm{loc}}`$ $`(\{𝐮_i\},\{𝐯_i\},\{\sigma _j\})=`$ (5)
$`{\displaystyle \underset{i}{}}[\mathrm{\Delta }\alpha (\sigma _i)u_i^4+\mathrm{\Delta }\gamma (\sigma _i)(u_{\mathrm{𝑖𝑥}}^2u_{\mathrm{𝑖𝑦}}^2+u_{\mathrm{𝑖𝑦}}^2u_{\mathrm{𝑖𝑧}}^2+u_{\mathrm{𝑖𝑧}}^2u_{\mathrm{𝑖𝑥}}^2)]`$
$`+{\displaystyle \underset{ij}{}}[Q_{j,i}(\sigma _j)𝐞_{\mathrm{𝑗𝑖}}𝐮_i+R_{j,i}(\sigma _j)𝐟_{\mathrm{𝑗𝑖}}𝐯_i)],`$
where the sums over $`i`$ and $`j`$ run over unit cells and mixed sublattice sites, respectively. $`\mathrm{\Delta }\alpha (\sigma _i)`$ and $`\mathrm{\Delta }\gamma (\sigma _i)`$ characterize the onsite contribution of alloying, while $`Q_{j,i}(\sigma _j)`$ and $`R_{j,i}(\sigma _j)`$ are related to alloying-induced intersite interactions. Here $`𝐞_{\mathrm{𝑗𝑖}}`$ is a unit vector joining the site $`j`$ to the center of the soft-mode vector $`𝐮_i`$, and $`𝐟_{\mathrm{𝑗𝑖}}`$ is a unit vector joining the site $`j`$ to the origin of the displacement $`𝐯_i`$. In principle, terms involving higher powers of $`𝐮_i`$ and $`𝐯_i`$ might be included to improve the quality of the expansion, but as will be shown below, we find this level of expansion to give a very good account of experimental findings. We also find that $`Q_{j,i}(\sigma _j)`$ and $`R_{j,i}(\sigma _j)`$ decrease rapidly as the distance between $`i`$ and $`j`$ increases. As a result, we included contributions up to the third neighbors for $`Q_{j,i}(\sigma _j)`$, and over the first-neighbor shell for $`R_{j,i}(\sigma _j)`$.
All the parameters of Eqs. (1) and (2) are derived from first principles. The 18 parameters of $`E_{\mathrm{ave}}`$ (see Table II of Ref. ) are determined by fitting the results of first-principles VCA calculations. On the other hand, $`\mathrm{\Delta }\alpha (\sigma _i)`$, $`\mathrm{\Delta }\gamma (\sigma _i)`$, $`Q_{j,i}(\sigma _j)`$ and $`R_{j,i}(\sigma _j)`$ are derived by performing first-principles calculations in which a true atom \[e.g., Ti or Zr in Pb(Zr,Ti)O<sub>3</sub>\] is surrounded by VCA atoms. The first-principles method used in the present study is the plane-wave ultrasoft-pseudopotential method within the local-density approximation (LDA) . The VCA approach adopted averages the B and B<sup>′′</sup> pseudopotentials, and is the one of Ref. .
Once our effective Hamiltonian is fully specified, the total energy of Eq. (1) is used in Monte-Carlo simulations to compute finite-temperature properties of ferroelectric alloys. We typically use a $`12\times 12\times 12`$ supercell (8640 atoms), since this choice yields well-converged results. The $`\{\sigma _j\}`$ variables are chosen randomly in order to mimic maximal compositional disorder – consistent with experimental reality – and are kept fixed during the Monte-Carlo simulations. We find that averaging our results over a couple of different realizations of the disorder leads to well-converged statistical properties. The outputs of the Monte-Carlo procedure are the local mode vectors $`𝐮`$ (directly related to the electrical polarization), and the homogeneous strain tensor $`\eta _H`$. We use the correlation-function approach of Ref. to derive the piezoelectric response from these Monte-carlo simulations. Up to 10<sup>6</sup> Monte-Carlo sweeps are first performed to equilibrate the system, and then 2$`\times `$10<sup>4</sup> sweeps are used to get the various statistical averages. The temperature is decreased in small steps.
Figure 1 shows the largest, middle and smallest cartesian coordinates ($`u_1`$, $`u_2`$ and $`u_3`$) of the supercell average of the local mode vectors in Pb(Zr<sub>0.5</sub>Ti<sub>0.5</sub>)O<sub>3</sub> as a function of the temperature, as predicted by our approach described by Eqs. (1) and (2). Each coordinate is close to zero at high temperature, characterizing a paraelectric cubic phase. As the system is cooled down, $`u_1`$ drastically increases while $`u_2`$ and $`u_3`$ remain nearly null. This indicates a transition to a ferroelectric tetragonal phase, consistent with measurements . We predict that the spontaneous polarization reaches 0.79 C/m<sup>2</sup> at very low temperature, which compares well with the first-principles results of 0.70 and 0.74 C/m<sup>2</sup> . The tetragonal axial ratio $`c/a`$ ranges from 1 close to the transition region to 1.02 for lower temperature. This is in good agreement both with the experimental value of 1.02 – 1.025 obtained for disordered samples, and with the first-principles result of 1.03 obtained for an ordered alloy . Figure 1 also shows the predictions of the VCA alloy theory, corresponding to the neglect of $`E_{\mathrm{loc}}`$ in Eq. (1). Interestingly, one sees that $`E_{\mathrm{loc}}`$ has no effect on the phase transition sequence, which is consistent with recent findings that the VCA approach can reproduce some structural properties of PZT . Whether or not $`E_{\mathrm{loc}}`$ is included in the total energy, we find a Curie temperature $`T_c`$ that is higher than the experimental value of 640 K . This difficulty of reproducing $`T_c`$ is a general feature of the effective-Hamiltonian approach , and may be due to higher perturbative terms neglected in the analytical expression for the total energy. In order to compare our results with experimental data, we will henceforth rescale our temperature as in Ref. so that the theoretical $`T_c`$ is forced to match the experimental one.
Figure 2 shows the piezoelectric coefficients predicted for a tetragonal single crystal of Pb(Zr<sub>0.5</sub>Ti<sub>0.5</sub>)O<sub>3</sub> as a function of the rescaled temperature, when neglecting or incorporating $`E_{\mathrm{loc}}`$ in Eq. (1). The independent coefficients for the 4mm point group are $`d_{33}`$, $`d_{31}`$ and $`d_{15}`$. One can notice that inclusion of $`E_{\mathrm{loc}}`$ has only a small effect on $`d_{31}`$ and $`d_{33}`$: $`d_{31}`$ is rather small for any temperature except near the transition, and $`d_{33}`$ is around 50 – 55 pC/N at room temperature in both simulations. Using the experimental values of the elastic compliances to compute $`e_{33}`$ from our calculated $`d_{31}`$ and $`d_{33}`$, we find that our alloy effective hamiltonian leads to an $`e_{33}`$ of 4.3 C/m<sup>2</sup>, while neglecting $`E_{\mathrm{loc}}`$ in Eq. (1) yields a similar $`e_{33}`$ of 3.8 C/m<sup>2</sup> at low temperature. Both predictions agree well with the first-principles results ranging between 3.4 and 4.8 C/m<sup>2</sup> , confirming that the VCA can reproduce the $`e_{33}`$ coefficient of PZT .
Figure 2 also demonstrates that incorporating $`E_{\mathrm{loc}}`$ in the total energy leads to a large enhancement of the $`d_{15}`$ coefficient, which is consistent with recent measurements revealing that the piezoelectric elongation of the tetragonal unit cell of PZT does not occur along the polar direction . This enhancement is highly relevant for the piezoelectric response $`d_{33}`$ in ceramic samples, denoted $`d_{33,c}`$, which involves an average of the form
$$d_{33,c}=_0^{\pi /2}[(d_{31}+d_{15})\mathrm{sin}^2\theta +d_{33}\mathrm{cos}^2\theta ]\mathrm{sin}\theta \mathrm{cos}\theta d\theta $$
(6)
over the single-crystal coefficients. The true alloy approach of Eqs. (1)-(2) leads to a $`d_{33,c}`$ of 163 pC/N at room temperature, in excellent agreement with the experimental value of 170 pC/N . On the other hand, neglecting $`E_{\mathrm{loc}}`$ leads to a smaller $`d_{33,c}`$ of 90 pC/N. This difference clearly demonstrates the necessity of incorporating the local alloying effect into the total energy for understanding the large piezoelectric response of PZT ceramics near the MPB.
We now use our alloy effective Hamiltonian to investigate the low-temperature phases of Pb(Zr<sub>1-x</sub>Ti<sub>x</sub>)O<sub>3</sub> near the MPB. We choose a constant temperature of 50 K in the Monte-Carlo simulations, and vary the overall composition of the solid solution. This compositional variation affects two quantities: (i) the populations of $`\sigma _j`$ equal to $`+1`$ or $`1`$, and (ii) the alloy-related parameters. For the latter, only the parameters entering the local-mode self-energy of $`E_{\mathrm{ave}}`$ , and the $`\mathrm{\Delta }\alpha (\sigma _i)`$ and $`\mathrm{\Delta }\gamma (\sigma _i)`$ in Eq. (2), are allowed to be composition-dependent. This composition-dependence is assumed to be linear, and is determined by performing first-principles simulations on two different compositional cells. Such a linear composition-dependence approach is only realistic when exploring a narrow range of compositions, as done in the present study.
Figure 3 shows that the local mode, and hence the polarization, is parallel to the pseudo-cubic direction for Ti compositions larger than 50%, which is consistent with a tetragonal phase. For compositions lower than 47%, the polarization becomes parallel to the pseudo-cubic direction, indicating the ”high temperature” rhombohedral phase . The most interesting feature of Fig. 3 is the behavior of the local mode for the compositional range between 47.5% and 49.5%: as $`x`$ decreases, $`u_1`$ decreases, while $`u_2`$ and $`u_3`$ increase and remain nearly equal to each other. This behavior is characteristic of an intermediate phase that is neither tetragonal nor rhombohedral. The strain tensor given by our simulations indicates that this intermediate phase is the monoclinic phase experimentally found by Noheda et al. We further predict that the monoclinic phase for $`x48\%`$ can be characterized by an angle of 90.7 and by lattice vectors a<sub>m</sub> = $`a_0`$($``$1.005,$``$1.005,$``$0.009), b<sub>m</sub> = $`a_0`$(1.002,$``$1.002,0.000), and c<sub>m</sub> = $`a_0`$(0.004,0.004,1.018), where a<sub>0</sub> is a cubic lattice constant. All these predictions are in excellent quantitative agreement with the experimental results of Ref. . Figure 3 clearly demonstrates that the monoclinic phase acts as a bridge between the rhombohedral and tetragonal phases, as indicated by the continuous rotation of the polarization as a function of composition. Our computational scheme is also able to reproduce the compositional range narrowing of the monoclinic phase observed when increasing the temperature . It should be noted that a VCA-only calculation (i.e., neglecting $`E_{\mathrm{loc}}`$) does not reveal a monoclinic phase. This finding demonstrates once again the need for incorporating the local effect of alloying into the total energy to study subtle effects.
In summary, we have developed a first-principles derived computational scheme to study finite-temperature properties of Pb(Zr<sub>1-x</sub>Ti<sub>x</sub>)O<sub>3</sub> solid solutions near the MPB as a function of composition and temperature. We find that there is a low-temperature monoclinic phase acting as a bridge between the rhombohedral phase, existing for $`x<0.47`$, and the tetragonal phase, occurring for $`x`$ larger than 0.50. The predicted structural data are in very good agreement with measurements, as well as with direct first-principles calculations. The use of this approach also provides an explanation for the large experimental value of $`d_{33}`$ in tetragonal ceramics of PZT near the MPB . This large piezoelectricity is simply due to the very large value of the $`d_{15}`$ coefficient predicted to occur in the single crystal.
L.B. thanks the financial assistance provided by the Arkansas Science and Technology Authority (grant N99-B-21), and the National Science Foundation (grant DMR-9983678). A.G. acknowledges support from the Spanish Ministry of Education (grant PB97-0598). D.V. acknowledges the financial support of Office of Naval Research grant N00014-97-1-0048. We wish to thank B. Noheda, H. Chen, M. Cohen, E. Cross, T. Egami and Q. Zhang for very useful discussions.
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# SPONTANEOUS CP VIOLATION IN THE LEFT–RIGHT–SYMMETRIC MODEL
## 1 Introduction
In this talk I report on a recent investigation $`^{\mathrm{?},\mathrm{?}}`$ of a non-standard mechanism of CP violation in an attractive extension of the Standard Model (SM), the spontaneously broken left–right–symmetric model.$`^\mathrm{?}`$ Available data on the mass differences and the measured amount of CP violation in the K and B system seriously constrain the model’s parameter space. Even if current theoretical uncertainties persist in the K system, the expected experimental progress in B physics $`^\mathrm{?}`$ will soon bring conclusive tests of the model in the form of precise values of the CP violating asymmetries in the decays $`B_dJ/\psi K_S`$ and $`B_sJ/\psi \varphi `$, which we predict to sizeably deviate from their SM expectations.
It is well known that CP is a natural symmetry of pure gauge theories with massless fermions. Consequently, CP violation actually probes the least known sector of unified theories, namely the scalar and Yukawa couplings. The current development of dedicated accelerators to probe CP violation in the B system prompts studies of possible departures from the SM. Left–right (LR) symmetric models based on the extended gauge group SU(2)$`{}_{\mathrm{L}}{}^{}\times `$SU(2)$`{}_{\mathrm{R}}{}^{}\times `$U(1) offer the advantage of a well-defined, and actually quite constraining context, largely testable experimentally, while presenting a structure significantly different from the SM. Whereas in general LR models CP violation can come from several sources, we consider here a restricted model which exhibits the aesthetically attractive feature of spontaneous CP violation. This means that the Lagrangian exhibits manifest CP symmetry, which at low energies is broken by ”misaligned” phases of the symmetry-breaking vacuum expectation value (VEV). In this model, spontaneous breakdown is the only source of CP violation. Under this hypothesis, Ecker and Grimus have shown $`^\mathrm{?}`$ that, except for an exceptional case (which will not be considered here), the Yukawa couplings can be parametrized in terms of two real symmetrical matrices. As a consequence, all CP-violating phases of the model can be be calculated exactly, as they are related to a unique phase (denoted $`\alpha `$ below) which affects the SU(2)$`{}_{\mathrm{L}}{}^{}\times `$SU(2)<sub>R</sub> breaking VEV. This point is important, as it links baryogenesis, and in particular the sign of the matter–antimatter asymmetry, to low-energy CP violation.$`^\mathrm{?}`$
In practice, the above-defined ”Spontaneously Broken Left–Right model” (SB–LR) adds a very economical four parameters to the SM: 2 boson masses and 2 parameters describing the VEV that breaks SU(2)$`{}_{\mathrm{L}}{}^{}\times `$SU(2)<sub>R</sub>. Despite being an extension of the SM, the SB–LR is in some sense more restrictive than the SM itself. Indeed, while the SM is a subset of general LR obtained by sending the R-sector masses to infinity, a similar procedure applied to the SB–LR yields additional constraints since the CKM phase $`\delta `$ is no longer independent, but predicted within the model. To be specific, we find that in the SB–LR $`\delta `$ is too small, $`|\delta |<0.25`$ or $`|\delta \pi |<0.25`$, whereas a recent global fit $`^\mathrm{?}`$ yields $`\delta =1.0\pm 0.2`$. Hence the SM limit of the SB–LR is inconsistent by 3.5$`\sigma `$ with current experiments. This has the important consequence that the SB–LR is actually testable, and distinct from the SM: experimental bounds cannot be indefinitely evaded by simply sending the R-sector to infinite masses: scalars and vectors in the range (2–20) TeV are definitely needed.
Experimental constraints on the SB–LR, mainly from the K system, have been thoroughly investigated in the late 80’s.$`^\mathrm{?}`$ Since then, many SM parameters, in particular the CKM angles and the top quark mass, have been measured much more accurately, and also theory has progressed, as exact relations for the CP-violating phases in the quark mixing-matrices have been derived,$`^\mathrm{?}`$ which supersede the previously used small-phase approximation $`^\mathrm{?}`$ that breaks down for $`b`$ decays due to the large top-quark mass. The perspective of finding non-standard CP violation in the B system at the B factories, the Tevatron and the LHC has prompted us to undertake a new comprehensive analysis of restrictions on the SB–LR from available experimental data. The main results can be summarized as follows:
* the role of the Higgs bosons, neglected in most analyses, is crucial;
* the decoupling limit of the model, where the extra boson masses $`M_2`$ and $`M_H`$ are sent to infinity, is experimentally excluded, which implies upper bounds on $`M_2`$ and $`M_H`$;
* neglecting uncertainties of quark masses and CKM angles, the SB–LR favours opposite signs of the CP-violating observables $`\mathrm{Re}ϵ`$ and $`a_{\mathrm{CP}}(B_dJ/\psi K_S)`$, which are both expected to be positive in the SM; hence, the model cannot accommodate both the experimentally measured $`ϵ`$ and the SM expectation $`a_{\mathrm{CP}}^{\mathrm{SM}}(B_dJ/\psi K_S)0.75`$ and is excluded if $`a_{\mathrm{CP}}`$ will be measured close to its SM expectation;
* the CP asymmetry in $`B_sJ/\psi \varphi `$, which is negligible in the SM, can be as large as 35% in SB–LR.
## 2 The Left-Right-Symmetric Model with Spontaneous CP Violation
Before discussing its predictions for CP-violating phenomena, let us explain very shortly the essential features of the SB–LR. As already mentioned, it is based on the gauge group SU(2)$`{}_{\mathrm{L}}{}^{}\times `$SU(2)$`{}_{\mathrm{R}}{}^{}\times `$U(1), which cascades down to the unbroken electromagnetic subgroup U(1)<sub>em</sub> through the following simple symmetry-breaking pattern:
$$\underset{\underset{\text{U(1)}_{\mathrm{em}}}{\underset{}{\text{SU(2)}_\mathrm{L}\times \text{U(1)}}}}{\underset{}{\text{SU(2)}_\mathrm{R}\times \text{SU(2)}_\mathrm{L}}}\times \text{U(1)}$$
The scalar sector is highly model-dependent; for the generation of quark masses, there has to be at least one scalar bidoublet $`\mathrm{\Phi }`$, i.e. a doublet under both SU(2), which, by spontaneous breakdown of SU(2)$`{}_{\mathrm{R}}{}^{}\times `$SU(2)<sub>L</sub>, acquires the VEV
$$\mathrm{\Phi }=\frac{1}{\sqrt{2}}\left(\begin{array}{cc}v& 0\\ 0& we^{i\alpha }\end{array}\right).$$
(1)
Here, $`v`$ and $`w`$ are real and the phase $`\alpha `$ is the (only) source of CP violation in the model. The particle content of $`\mathrm{\Phi }`$ corresponds to four particles, one analogue of the SM Higgs, two flavour-changing neutral Higgs bosons, and one flavour-changing charged Higgs. The masses of these new Higgs particles can be assumed to be degenerate to good accuracy; they will be denoted by $`M_H`$ below.
LR symmetry implies that the left-handed quark sector of the SM gets complemented by a right-handed one, with quark mixing matrices $`V_\mathrm{L}`$ and $`V_\mathrm{R}`$, respectively, and $`|V_\mathrm{L}|=|V_\mathrm{R}|`$ (but $`V_\mathrm{L}V_\mathrm{R}`$ due to different complex phases!). In the standard Maiani convention, $`V_\mathrm{L}`$ contains one, $`V_\mathrm{R}`$ five complex phases, which depend on the three generalized Cabibbo-type angles (“CKM angles”), the quark masses, and the VEV (1). The presence of such a large number of weak phases, calculable in terms of only one non-SM variable,<sup>a</sup><sup>a</sup>aI.e. the variable $`\beta `$ introduced in the next section; note also that there is a 64-fold discrete ambiguity of the phases due to the signs of the quark masses, which are physical in LR models. is a feature that makes the investigation of CP-violating phenomena in the SB–LR very interesting. The left- and right-handed charged gauge-bosons $`W_\mathrm{L}`$ and $`W_\mathrm{R}`$ mix with each other; the mass eigenstates are denoted by $`W_1`$ and $`W_2`$. The mixing angle $`\zeta `$, obtained as
$$\zeta =\frac{2|vw|}{|v|^2+|w|^2}\left(\frac{M_1}{M_2}\right)^2,$$
(2)
is rather small: as the ratio $`|v|/|w|`$ is smaller than 1,<sup>b</sup><sup>b</sup>bWhich can always be achieved by a redefinition of the Higgs bidoublet $`\mathrm{\Phi }\sigma _2\mathrm{\Phi }^{}\sigma _2`$. one has $`\zeta <(M_1/M_2)^210^3`$ (assuming $`M_2`$ in the TeV range as indicated by the experimental absence of right-handed weak currents). There are, however, arguments according to which a small ratio $`|v|/|w|𝒪(m_b/m_t)`$ would naturally explain the observed smallness of the CKM angles; $`^\mathrm{?}`$ in this case, $`\zeta 10^5`$. An experimental bound on $`\zeta `$ can in principle be obtained from the upper bound on the electromagnetic dipole moment of the neutron, which is induced by L–R mixing; existing theoretical calculations are, however, very sensitive to the precise values of only poorly known nucleon matrix elements; $`^\mathrm{?}`$<sup>c</sup><sup>c</sup>cIn addition, the Higgs contributions to the dipole moment are usually not included. the present status of an experimental bound on $`\zeta `$ is thus not quite clear, although large values of $`\zeta 10^4`$ appear to be disfavoured.
The fact that the SB–LR has no perceptible impact on SM tree-level amplitudes (no experimental indication of right-handed weak interactions or large flavour-changing neutral currents!) implies that the new boson masses must be in the TeV range. The SB–LR effects thus manifest themselves mostly in
* $`W_\mathrm{L}`$$`W_\mathrm{R}`$ mixing in top-dominated penguin diagrams, enhanced by large quark-mass terms from spin-flips, $`\zeta \zeta m_t/m_b`$ (similar for penguins with charged-Higgs particles);
* SM amplitudes that are forbidden or heavily suppressed (e.g. electromagnetic dipole moment of the neutron);
* mixing of neutral K and B mesons, where the suppression factor $`(M_1/M_2)^2`$ is partially compensated by large Wilson-coefficients or hadronic matrix-elements (chiral enhancement in K mixing), and to which the flavour-changing Higgs bosons contribute at tree level.
The SB–LR does, however, not significantly modify the decay amplitudes of the “gold-plated” decay mode $`B_dJ/\psi K_S`$ or the decay $`B_sJ/\psi \varphi `$: <sup>d</sup><sup>d</sup>dRecall that, in the SM and using the standard “generalized Wolfenstein parametrization” of the CKM matrix, the amplitudes of these decays carry only small or zero weak phases and the CP-violating asymmetry is essentially given by the B–$`\overline{\text{B}}`$ mixing-phase. these decays are dominated by one single CKM amplitude ($`bccs`$) with contributions from colour-suppressed tree and penguin-topologies with internal $`c`$ or $`t`$ quarks. The “gold-plated” mode $`B_dJ/\psi K_S`$ is the standard example for a special type of decays of a neutral B meson into a CP-eigenstate, whose time-dependent CP-asymmetry takes a particularly simple form and allows the direct extraction of a weak CKM phase without hadronic uncertainties:
$$𝒜_{\mathrm{CP}}(B_dJ/\psi K_S)=\frac{\mathrm{\Gamma }(t)\overline{\mathrm{\Gamma }}(t)}{\mathrm{\Gamma }(t)+\overline{\mathrm{\Gamma }}(t)}=\mathrm{sin}\varphi _{\mathrm{weak}}\mathrm{sin}\mathrm{\Delta }M_dt,$$
(3)
where $`\overline{\mathrm{\Gamma }}(t)`$ denotes the decay-rate of $`\overline{B}_d^0(t)J/\psi K_S`$ and $`\mathrm{\Delta }M_d`$ is the mass-difference in the $`B_d^0`$$`\overline{B}_d^0`$ system. The LR contribution to the tree-topology is given by $`W_L`$$`W_R`$ mixing, $`W_R`$ or neutral-Higgs exchange and suppressed by $`(M_1/M_2)^210^3`$ or more with respect to the SM contribution. As for the penguins, internal $`W_R`$ or charged-Higgs exchange are suppressed by the same order of magnitude as for the tree-topology, and the only potentially relevant contribution comes from $`W_L`$$`W_R`$ mixing: the corresponding top-penguin topology is enhanced by a spin-flip factor $`\zeta m_t/m_b`$, which is at most 5% in the most unfavourable case $`|v|/|w|=1`$ and in the range of permille for the preferred value $`|v|/|w|𝒪(m_b/m_t)`$. Consequently, the SB–LR contributions to the amplitudes of the “gold-plated” B decays are small and do neither scratch the lustre of the golden plates nor yield sizeable direct CP violation: <sup>e</sup><sup>e</sup>eWhich would show up as terms in $`\mathrm{cos}\mathrm{\Delta }M_dt`$ in (3). the main impact of the model on these decays is from B–$`\overline{\text{B}}`$ mixing and modifies the extracted value of the weak phase $`\varphi _{\mathrm{weak}}`$. Sizeable LR-effects on amplitudes are, however, to be expected in theoretically less “clean” channels like $`B_d\pi \pi `$ and $`bs\gamma `$.
## 3 Phenomenological Analysis
As mentioned above, the main impact of the SB–LR on channels with available experiemental data is to modify the SM pattern of neutral K and B meson mixing. The relevant quantity to be calculated is the matrix-element
$$M^0_{\mathrm{eff}}^{|\mathrm{\Delta }F|=2}\overline{M}^0=2m_MM_{12}$$
with the effective weak Hamiltonian $`_{\mathrm{eff}}^{|\mathrm{\Delta }F|=2}`$. Experimental observables which restrict SB–LR contributions to $`M_{12}^{B,K}`$ are
$$\mathrm{\Delta }M_{B,K}=2\left|M_{12}^{B,K}\right|,ϵ_K\frac{1}{2\sqrt{2}}e^{i\pi /4}\mathrm{sin}(\mathrm{arg}M_{12}^K),a_{\mathrm{CP}}(B_dJ/\psi K_S)=\mathrm{sin}(\mathrm{arg}M_{12}^{B_d}).$$
(4)
In addition, we also consider constraints posed by the smallness of direct CP violation in the K system encoded in the value of $`ϵ^{}/ϵ`$, but in view of the theoretical uncertainties associated with this observable $`^\mathrm{?}`$ we only require the sign of $`ϵ^{}/ϵ`$ to be correctly reproduced. $`|M_{12}^K|`$ is affected by long-distance QCD uncertainties which are also present in the SM, so that in our analysis, instead of attempting a full calculation of that quantity, we impose the reasonable requirement that SB–LR effects be smaller than the experimental mass difference: $`2|M_{12}^{K,SBLR}|<\mathrm{\Delta }M_K`$. The quantities $`\mathrm{arg}M_{12}^K`$ and $`M_{12}^B`$, on the other hand, are short-distance dominated and can be reliably calculated in the SB–LR. For the technicalities I refer to Ref. $`^\mathrm{?}`$; here I only specify the set of SM input parameters,<sup>f</sup><sup>f</sup>fNote that we do not take into account uncertainties associated with these parameters.
$$\begin{array}{ccccccccc}\overline{m}_t(\overline{m}_t)\hfill & =& 170\mathrm{GeV},\hfill & \overline{m}_b(\overline{m}_b)\hfill & =& 4.25\mathrm{GeV},\hfill & & & \\ \overline{m}_c(\overline{m}_c)\hfill & =& 1.33\mathrm{GeV},\hfill & \overline{m}_s(2\mathrm{GeV})\hfill & =& 110\mathrm{MeV},\hfill & & & \\ m_s/m_d\hfill & =& 20.1,\hfill & m_u/m_d\hfill & =& 0.56,\hfill & & & \\ \multicolumn{6}{c}{|V_{us}|=0.2219,|V_{ub}|=0.004,|V_{cb}|=0.04,}& & & \end{array}$$
(5)
and remind that $`M_{12}`$ depends in addition on the SB–LR parameters $`M_2`$, $`M_H`$ and, as shown in Ref. $`^\mathrm{?}`$, the variable $`\beta `$ defined as
$$\beta =\mathrm{arctan}\frac{2|wv|\mathrm{sin}\alpha }{|v|^2|w|^2}.$$
(6)
The combined analysis of the experimental data on $`\mathrm{\Delta }M_{K,B_d,B_s}`$, $`ϵ_K`$, the sign of $`ϵ^{}/ϵ`$ and $`a_{\mathrm{CP}}^{\mathrm{exp}}(B_dJ/\psi K_S)=0.79_{0.41}^{+0.44}`$ (see the talk $`^\mathrm{?}`$) yields the allowed region for $`M_2`$ and $`M_H`$ shown in Fig. 1 and the correlations between values of $`a_{\mathrm{CP}}`$ and $`ϵ_K`$ shown in Fig. 2.
The evident preference of the SB–LR for opposite signs of $`a_{\mathrm{CP}}`$ and $`ϵ_K`$, which is in contradiction to experiment at 98% CL, $`^\mathrm{?}`$ actually helps to resolve the 64-fold discrete ambiguity of the CKM phases mentioned in the previous section: only one of these 64 solutions can reproduce a positive $`ϵ_K`$ compatible with the experimental result and $`a_{\mathrm{CP}}>0`$. The resulting value of $`a_{\mathrm{CP}}`$ is, however, rather smallish, $`a_{\mathrm{CP}}<0.1`$, and at variance with the SM expectation $`a_{\mathrm{CP}}^{\mathrm{SM}}0.75`$, but in agreement with the present experimental result within 2$`\sigma `$. I thus quote as a first specific and testable prediction of the SB–LR:
$$\overline{)a_{\mathrm{CP}}^{\mathrm{SB}\mathrm{LR}}(B_dJ/\psi K_S)<0.1a_{\mathrm{CP}}^{\mathrm{SM}}(B_dJ/\psi K_S)0.75}$$
With the parameters $`M_2`$, $`M_H`$ and $`\beta `$ being constrained, we can now predict the allowed range for mixing-induced CP violation in the $`B_s`$ system. In the simple case of $`bccs`$ dominated $`B_sf`$ transitions into a final state $`f`$ that is a CP eigenstate (e.g. $`f=D_s^+D_s^{}`$, $`J/\psi \eta (^{})`$), the CP asymmetry is completely analogous to the $`B_dJ/\psi K_S`$ case:
$$a_{\mathrm{CP}}(B_sD_s^+D_s^{},J/\psi \eta (^{}))=\mathrm{sin}(\mathrm{arg}M_{12}^{B_s});$$
the corresponding correlation with $`a_{\mathrm{CP}}(B_dJ/\psi K_S)`$ is plotted in Fig. 3(a). The situation is a bit more complicated for the decay $`B_sJ/\psi \varphi `$, as here the final state is a superposition of CP-even and odd states. The standard time-dependent CP asymmetry as defined in (3) does now contain poorly known hadronic matrix-elements. These can, in principle, be extracted by an analysis of the angular correlations between the observed decay products $`(J/\psi )\mathrm{}^+\mathrm{}^{}`$ and $`(\varphi )K^+K^{}`$ $`^\mathrm{?}`$. Such an analysis requires, however, large statistics and is probably possible only at the LHC.$`^\mathrm{?}`$
In Fig. 3(b), we show the correlation between the $`B_s^0`$$`\overline{B_s^0}`$ mass and width difference $`\mathrm{\Delta }\mathrm{\Gamma }_s`$ in the SB–LR. The reduction of $`\mathrm{\Delta }\mathrm{\Gamma }_s`$ through new-physics effects is not very effective in this case, whereas the mass difference $`\mathrm{\Delta }M_s`$ may be reduced significantly. Although, at first glance, values of $`\mathrm{\Delta }M_s`$ as small as $`0.55\mathrm{\Delta }M_s^{\mathrm{SM}}`$ may seem to be at variance with the experimental bound $`\mathrm{\Delta }M_s>14.3\mathrm{ps}^1`$ at 95% CL $`^\mathrm{?}`$, this is actually not the case: with the hadronic parameters from $`^\mathrm{?}`$ and $`|V_{ts}|=0.04`$ with the generalized Cabibbo-angles fixed from (5), one has the theoretical prediction (see $`^\mathrm{?}`$, e.g., for the full formula)
$$\mathrm{\Delta }M_s^{\mathrm{SM}}=(14.5\pm 6.3)\mathrm{ps}^1.$$
Combining this with the experimental bound, one has
$$\frac{\mathrm{\Delta }M_s^{\mathrm{LR}}}{\mathrm{\Delta }M_s^{\mathrm{SM}}}>\frac{14.3}{14.5+2\times 6.3}=0.53.$$
A pattern of $`B_s`$ mass and decay width differences like that emerging in the SB–LR would be in favour of experimental studies of the $`B_s`$ decays at hadron machines, where small values of $`\mathrm{\Delta }M_s`$ and large values of $`\mathrm{\Delta }\mathrm{\Gamma }_s`$ would be desirable.
Let us finally illustrate the CP-violating asymmetry of the decay $`B_sJ/\psi \varphi `$:
$$𝒜_{\mathrm{CP}}(B_s(t)J/\psi \varphi )\frac{\mathrm{\Gamma }(t)\overline{\mathrm{\Gamma }}(t)}{\mathrm{\Gamma }(t)+\overline{\mathrm{\Gamma }}(t)}=\left[\frac{1D}{F_+(t)+DF_{}(t)}\right]\mathrm{sin}(\mathrm{\Delta }M_st)\mathrm{sin}(\mathrm{arg}M_{12}^{B_s}),$$
(7)
where $`\mathrm{\Gamma }(t)`$ and $`\overline{\mathrm{\Gamma }}(t)`$ denote the time-dependent rates for decays of initially, i.e. at $`t=0`$, present $`B_s^0`$ and $`\overline{B_s^0}`$ mesons into $`J/\psi \varphi `$ final states, respectively. The remaining quantities are defined as
$$D\frac{|A_{}(0)|^2}{|A_0(0)|^2+|A_{}(0)|^2},$$
(8)
and
$$F_\pm (t)\frac{1}{2}\left[\left(1\pm \mathrm{cos}(\mathrm{arg}M_{12}^{B_s})\right)e^{+\mathrm{\Delta }\mathrm{\Gamma }_st/2}+\left(1\mathrm{cos}(\mathrm{arg}M_{12}^{B_s})\right)e^{\mathrm{\Delta }\mathrm{\Gamma }_st/2}\right].$$
(9)
Here $`A_0(t)`$, $`A_{}(t)`$ and $`A_{}(t)`$ are linear polarization amplitudes that describe the CP-even and odd final-state configurations $`^\mathrm{?}`$. Note that we have $`F_+(t)=F_{}(t)=1`$ for a negligible width difference $`\mathrm{\Delta }\mathrm{\Gamma }_s`$. Obviously, the advantage of the “integrated” observable (7) is that it can be measured without performing an angular analysis and is thus accessible at HERA-B and Tevatron Run II. The disadvantage is of course that it also depends on the hadronic quantity $`D`$, which precludes a theoretically clean extraction of $`\mathrm{arg}M_{12}^{B_s}`$ from (7). However, this feature does not limit the power of this CP asymmetry to search for indications of new physics, which would be provided by a sizeable measured value of (7). Model calculations of $`D`$, making use of the factorization hypothesis, typically give $`D=0.1\mathrm{}0.5`$ $`^\mathrm{?}`$, which is also in agreement with a recent analysis of the $`B_sJ/\psi \varphi `$ polarization amplitudes performed by the CDF collaboration $`^\mathrm{?}`$. A recent calculation of the relevant hadronic form factors from QCD sum rules on the light-cone $`^\mathrm{?}`$ yields $`D=0.33`$ in the factorization approximation. Consequently, the CP-odd contributions proportional to $`|A_{}(0)|^2`$ may have a significant impact on (7). In Fig. 3(c), we plot this CP asymmetry as a function of $`t`$, for fixed values of $`D=0.3`$, $`\mathrm{sin}\mathrm{arg}M_{12}^{B_s}=0.38`$, $`\mathrm{\Delta }\mathrm{\Gamma }_s/\mathrm{\Gamma }_s=0.14`$ and $`\mathrm{\Delta }M_s=14.5\mathrm{ps}^1`$. Although the $`B_s^0`$$`\overline{B_s^0}`$ oscillations are very rapid, as can be seen in this figure, it should be possible to resolve them experimentally, for example at the LHC. The first extremal value of (7), corresponding to $`\mathrm{\Delta }M_st=\pi /2`$, is given to a very good approximation by
$$A_{\mathrm{CP}}(B_sJ/\psi \varphi )=\left(\frac{1D}{1+D}\right)\mathrm{sin}(\mathrm{arg}M_{12}^{B_s}),$$
(10)
which would also fix the magnitude of the $`B_sJ/\psi \varphi `$ CP asymmetry (7) in the case of a negligible width difference $`\mathrm{\Delta }\mathrm{\Gamma }_s`$. In Fig. 3(d), we show the prediction of the SB–LR for (10) as a function of the hadronic parameter $`D`$. For a value of $`D=0.3`$, the CP asymmetry may be as large as $`25\%`$. The dilution through the hadronic parameter $`D`$ is not effective in the case of the CP-violating observables of the $`B_sJ/\psi [l^+l^{}]\varphi [K^+K^{}]`$ angular distribution, which allow one to probe $`\mathrm{sin}(\mathrm{arg}M_{12}^{B_s})`$ directly. I thus quote as a second specific and testable prediction of the SB–LR:
$$\overline{)a_{\mathrm{CP}}^{\mathrm{SB}\mathrm{LR}}(B_sJ/\psi \varphi )(1040)\%a_{\mathrm{CP}}^{\mathrm{SM}}(B_sJ/\psi \varphi )10^2}$$
## 4 Summary
In this talk I have presented a detailed investigation of the present status of the left–right symmetrical model with spontaneous CP violation, based on the gauge group SU(2)$`{}_{\mathrm{L}}{}^{}\times `$SU(2)$`{}_{\mathrm{R}}{}^{}\times `$U(1). The parameter space of this model includes the masses of the predominantly right-handed charged gauge boson, $`M_2`$, and of FC neutral and charged Higgs bosons, which we have assumed to be degenerate with a common mass $`M_H`$. Also included are the parameter $`\beta `$, which measures the size of the VEV of the Higgs bidoublet $`\mathrm{\Phi }`$ that characterizes the spontaneous breakdown of CP symmetry, and the 64-fold discrete ambiguity of the CKM phases due to different quark mass signatures. In contrast to previous publications, in which the constraints on the model from K and B physics were treated separately, our paper $`^\mathrm{?}`$ is the first one to consider them in a coherent way and to use the exact results for the CKM phases instead of the small phase approximation. We have concentrated on experimental constraints imposed by the mass differences $`\mathrm{\Delta }M_{K,B}`$ and observables describing CP violation, i.e. $`ϵ_K`$, $`ϵ^{}/ϵ`$ and $`a_{\mathrm{CP}}(B_dJ/\psi K_S)`$. In view of the large theoretical uncertainties, we only use the sign, but not the absolute value of Re$`(ϵ^{}/ϵ)`$ as a constraint, and we do not use the electric dipole moment of the neutron. Our main finding is that, although the K and B constraints can be met separately by a large range of input parameters, it is their combination that severely restricts the model. We find in particular that the CP violating observables $`ϵ_K`$ and $`a_{\mathrm{CP}}(B_dJ/\psi K_S)`$ are crucial: the sets of input parameters that pass the constraints imposed by the meson mass differences $`\mathrm{\Delta }M_{K,B}`$ yield to a large majority opposite signs of $`ϵ_K`$ and $`a_{\mathrm{CP}}(B_dJ/\psi K_S)`$.
We have also performed an analysis of mixing-induced CP-violating effects in $`B_sD_s^+D_s^{}`$, $`J/\psi \eta ^{(^{})}`$, $`J/\psi \varphi `$ decays and have demonstrated that the corresponding CP asymmetries may be as large as $`𝒪(40\%)`$, whereas the SM predicts vanishingly small values. Since the decay amplitudes of these modes are not significantly affected in the SB–LR, direct CP violation remains negligible, as in the SM. From an experimental point of view, $`B_sJ/\psi \varphi `$ is a particularly promising mode, which is very accessible at B physics experiments at hadron machines. We have proposed a simple strategy to search for indications of new physics in this transition, which does not require an angular analysis of the $`J/\psi [l^+l^{}]`$ and $`\varphi [K^+K^{}]`$ decay products. In contrast to the large mixing-induced CP asymmetries in the $`B_s`$ channels, the SB–LR predicts a small value for $`a_{\mathrm{CP}}(B_dJ/\psi K_\mathrm{S})`$ below 10%. Since the $`B_s`$ decays cannot be explored at the asymmetric $`e^+`$$`e^{}`$ B factories operating at the $`\mathrm{{\rm Y}}(4S)`$ resonance, such a pattern would be in favour of hadronic B experiments. We look forward to experimental data to check whether this scenario is actually realized by Nature.
We would like to stress that our study does not claim to be exhaustive as we did not allow the most crucial SM input parameters, the CKM angles and quark masses, to float within their presently allowed ranges. Taking into account these uncertainties would certainly affect the phases of the CKM matrices and thus mainly show up in the CP violating observables, which, as we have shown, are crucial. It is therefore not to be excluded that an analysis of the input parameter uncertainties would result in increasing the viable LR parameter ranges, but we doubt that it will change the anticorrelation between the signs of $`ϵ_K`$ and $`a_{\mathrm{CP}}(B_dJ/\psi K_S)`$ which implies a small maximum value of $`a_{\mathrm{CP}}^{\mathrm{SB}\mathrm{LR}}(B_dJ/\psi K_S)`$ attainable in the model.
## Acknowledgements
It is a pleasure to thank R. Fleischer, J.-M. Frère and J. Matias for enjoyable collaboration on the work presented here. I also gratefully acknowledge financial support from the organizers of the Moriond meeting. This work has been supported by DFG through a Heisenberg fellowship.
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# 1 Introduction
## 1 Introduction
Quite some time ago ’t Hooft pointed out some very special features which arise while formulating gauge theories on a torus (a review can be found in reference ). They are related to the freedom in choosing boundary conditions for the gauge potential: since only local, gauge invariant quantities are required to be periodic, periodicity of the gauge potential has to be satisfied only up to a gauge transformation. Under displacement of the gauge potential by a torus period $`l_\nu `$ in direction $`\widehat{\nu }`$
$`𝐀_\mu (x+l_\nu \widehat{\nu })=𝛀_\nu (x)𝐀_\mu (x)𝛀_\nu ^{}(x)ı𝛀_\nu (x)_\mu 𝛀_\nu ^{}(x),\mu =1,\mathrm{},4,`$ (1)
with $`𝛀_\nu (x_{\mu \nu })`$ SU(N) matrices, also known as the twist matrices. The choice of such matrices is arbitrary up to a consistency condition derived from the requirement of single-valuedness of the gauge potential. Built in two ways from $`𝐀_\mu (x_\nu ,x_\rho )`$ (through $`𝐀_\mu (x_\nu +l_\nu ,x_\rho )`$ and $`𝐀_\mu (x_\nu ,x_\rho +l_\rho )`$) single-valuedness for $`𝐀_\mu (x_\nu +l_\nu ,x_\rho +l_\rho )`$ implies
$`𝛀_\mu (x_\nu +l_\nu )𝛀_\nu (x_\mu )=𝛀_\nu (x_\mu +l_\mu )𝛀_\mu (x_\nu )\mathrm{Exp}\left({\displaystyle \frac{ı2\pi n_{\mu \nu }}{\mathrm{N}}}\right)`$ (2)
with $`n_{\mu \nu }`$ a gauge invariant antisymmetric tensor of integers, defined modulo N and independent of $`x`$ (this twist factor is allowed due to the invariance of $`𝐀_\mu `$ under a gauge transformation with an element of the center $`ZZ_\mathrm{N}`$ of SU(N)). Indeed, the actual choice of the twist matrices is irrelevant and only the consistency conditions given by $`n_{\mu \nu }`$ matter. Whenever $`n_{\mu \nu }0`$ (mod N), for some $`\mu `$, $`\nu `$, we say the boundary conditions are twisted. The twist is reflected in a gauge invariant way through the non-trivial periodicity of the Polyakov loops, defined on the torus as
$$_\mu (x)=\frac{1}{\mathrm{N}}\mathrm{Tr}(𝐋_\mu )=\frac{1}{\mathrm{N}}\mathrm{Tr}\left(\mathrm{Texp}\left\{ı_{\gamma _{\mu (x,x^{})}}𝐀_\nu 𝑑x^\nu \right\}𝛀_\mu (x^{})\mathrm{Texp}\left\{ı_{\gamma _{\mu (x^{},x)}}𝐀_\nu 𝑑x^\nu \right\}\right)$$
(3)
with $`\gamma _{\mu (a,b)}`$ a straight line in the positive $`\mu `$ direction starting at $`a`$ and ending at $`b`$ and $`x^{}`$ the border of the torus patch. Periodicity holds only up to the twist factors, i.e.
$$_\mu (x+l_\nu \widehat{\nu })=\mathrm{Exp}\left(\frac{ı2\pi n_{\mu \nu }}{\mathrm{N}}\right)_\mu (x).$$
(4)
With twisted boundary conditions the topological charge Q is no longer necessarily an integer:
$$Q=\frac{1}{16\pi ^2}\mathrm{Tr}\left(𝐅_{\mu \nu }\stackrel{~}{𝐅}_{\mu \nu }\right)d_4x=\nu \frac{\kappa }{\mathrm{N}},\mathrm{with}\nu ,\kappa ZZ,$$
(5)
$`\kappa `$ is associated to the matrix of integers $`(n_{\mu \nu })`$ through
$$\kappa =\frac{1}{4}n_{\mu \nu }\stackrel{~}{n}_{\mu \nu }=\stackrel{}{k}\stackrel{}{m},$$
(6)
with $`k_i=n_{0i}`$, $`n_{ij}=ϵ_{ijk}m_k`$. The mathematical proof of this relation for the topological charge can be found in . Then $`Q`$ is fractional and proportional to $`1/\mathrm{N}`$ whenever $`\stackrel{}{k}\stackrel{}{m}0\mathrm{modulo}\mathrm{N}`$ (non-orthogonal twist). This implies, through Schwarz-inequality, that the action of any configuration is bounded from below by
$$S=\frac{1}{2}\mathrm{Tr}\left(𝐅_{\mu \nu }𝐅_{\mu \nu }\right)d_4x8\pi ^2|Q|=8\pi ^2\left|\nu \frac{\kappa }{\mathrm{N}}\right|$$
(7)
with the bound saturated for self or anti-self dual configurations ($`𝐅_{\mu \nu }=\pm \stackrel{~}{𝐅}_{\mu \nu }`$). It is clear that whenever $`\kappa 0(\mathrm{modulo}\mathrm{N})`$ there is an obstruction for zero-action configurations. Minimal action is attained in such cases if $`|\nu \kappa /\mathrm{N}|=1/\mathrm{N}`$ with $`S=8\pi ^2/`$N. These are in fact the kind of solutions we will describe in this paper.
Some of these fractional charge solutions have already been found either analytically or numerically. ’t Hooft has explicitly constructed non-abelian solutions with constant field strength which turn out to be (anti-)self-dual whenever the sides of the torus satisfy certain relations (see for details or the appendix at the end of this paper). There are also a few numerical studies of solutions with non constant field strength. The first one is presented in reference and it is obtained there the fractional charge solution with $`|Q|=1/2`$ and $`S=4\pi ^2`$ for the $`SU(2)`$ group, on a $`L^3\times T`$ torus with $`TL`$ and satisfying twisted boundary conditions given by the twist vectors $`\stackrel{}{m}=(1,1,1)`$ and $`\stackrel{}{k}=(1,1,1)`$. A full parametrization of the field strength $`𝐅_{\mu \nu }`$, and of the gauge field $`𝐀_\mu `$ for this solution is presented in reference . Another SU(2) solution is presented in , in this case the fractional charge solution with $`|Q|=1/2`$ and $`S=4\pi ^2`$, on a torus $`L^2\times T^2`$ with $`TL`$ and satisfying the appropriate boundary conditions to have the properties of a vortex. The same kind of solution for the SU(3) group, with $`|Q|=1/3`$ and $`S=8\pi ^2/3`$, is presented in reference , and the generalization to SU(N) group with $`N>3`$ can be found in . In this article we present a numerical study of SU(N) solutions, with charge $`|Q|=1/N`$ and action $`S=8\pi ^2/N`$, and living on a $`L^3\times T`$ torus with $`TL`$. Some preliminary results have been presented in .
These configurations are interesting by itself from a mathematical point of view, and physically interesting for their possible relevance in low energy phenomena like the confinement property or the breaking of the chiral symmetry. As has been pointed out in these solutions may play a role in the properties of the theory in the limit of large number of colors, N. One of the arguments to question the contribution of instantons to long-distance phenomena as confinement is based on the large N expansion . Any instanton mediated interaction is suppressed by the semi-classical factor exp$`(8\pi ^2/g^2)`$, since the large N limit is achieved while keeping $`g^2`$N fixed, integer charge instantons are (at least in the dilute gas picture) naively suppressed by exp(-N). The argument no longer holds for twisted instantons with action $`8\pi ^2/`$N. Another interesting point is the possible relation between the center vortex picture of confinement, proposed in and now being investigated , and fractional charge solutions. As have been pointed out in , it is possible to built vortex configurations in $`R^4`$ from solutions of the Yang Mills equations of motion in $`T^4`$. We also want to mention the model of confinement based in fractional charge objects presented in reference , and some favourable results shown in .
The paper is structured as follows. In section two the numerical method to obtain the solutions is briefly described. We will be interested in solutions living on a $`L^3\times T`$ torus with $`TL`$ which, in the limit $`T\mathrm{}`$, represent vacuum to vacuum tunneling. The analysis will be restricted to spatial twist $`\stackrel{}{m}=(1,1,1)`$. Section three presents a detailed analysis of these solutions. Our conclusions are presented in section four. Finally, we include an appendix with the analytic solutions obtained through ’t Hooft construction. Their relation with the cases we have studied is discussed through the text.
## 2 Numerical minimization of the action
To generate numerically the minimal action configurations we follow the method which has allowed to successfully extract these kind of solutions for other sizes of the torus and values of the number of colors in references . We use the standard discretization of Yang-Mills theories on the lattice . We work on $`N_s^3\times N_t`$ lattices, $`N_tN_s`$, with variables defined on each link of the lattice taking values on N $`\times `$ N unitary matrices $`\widehat{𝐔}_\mu (n)`$. The lattice action used is the Wilson action,
$$S_W=\underset{n,\mu ,\nu }{}\mathrm{Tr}\left(1\widehat{𝐔}_\mu (n)\widehat{𝐔}_\nu (n+\widehat{\mu })\widehat{𝐔}_\mu ^{}(n+\widehat{\nu })\widehat{𝐔}_\nu ^{}(n)\right),$$
(8)
where $`\mu `$ and $`\nu `$ specify directions (from 1 to 4) and $`\widehat{\mu }`$ , $`\widehat{\nu }`$ are unit vectors along the corresponding direction, $`n_\mu =1,\mathrm{},N_\mu `$.
The link variables $`\widehat{𝐔}_\mu (n)`$ satisfy the (twisted) boundary conditions,
$$\widehat{𝐔}_\mu (n+N_\nu \widehat{\nu })=𝛀_\nu (n)\widehat{𝐔}_\mu (n)𝛀_\nu ^{}(n+\widehat{\mu })),$$
(9)
where $`N_4=N_t`$, $`N_i=N_s,i=1,2,3`$ and $`𝛀_\mu `$ are the twist matrices with consistency condition,
$$𝛀_\mu (n+N_\nu \widehat{\nu })𝛀_\nu (n)=𝛀_\nu (n+N_\mu \widehat{\mu })𝛀_\mu (n)\mathrm{Exp}(2\pi ın_{\mu \nu }/\mathrm{N}).$$
(10)
It is possible to make a change of variables
$$𝐔_\mu (N_\mu ,n_\nu )=\widehat{𝐔}_\mu (N_\mu ,n_\nu )𝛀_\mu (n_\nu )𝐔_\mu (n_\mu N_\mu ,n_\nu )=\widehat{𝐔}_\mu (n_\mu N_\mu ,n_\nu )$$
(11)
such that the new link variables are strictly periodic. In terms of the new links
$$S_W=\underset{n,\mu ,\nu }{}\mathrm{Tr}\left(1Z_{\mu \nu }^{}(n)𝐔_\mu (n)𝐔_\nu (n+\widehat{\mu })𝐔_\mu ^{}(n+\widehat{\nu })𝐔_\nu ^{}(n)\right),$$
(12)
where $`Z_{\mu \nu }(n)ZZ_\mathrm{N}`$ take the values: $`Z_{\mu \nu }(n)=1`$ for all plaquettes except the one at the top-right corner in the $`(\mu ,\nu )`$ plane which is equal to $`\mathrm{Exp}(2\pi in_{\mu \nu }/\mathrm{N})`$.
The strategy to obtain the solution is minimize the lattice action with respect to the variable $`𝐔_\mu (n)`$ (this minimization procedure is usually known as cooling). We use the Cabibbo-Marinari-Okawa algorithm in which each link variable is updated in the way: $`𝐔_\mu (n)\mathrm{𝐀𝐔}_\mu (n)`$, where $`𝐀`$ is a SU(N) matrix built from a SU(2) matrix $`𝐚`$ which is embedded into one of the $`N(N1)/2`$ subgroups of SU(N). Once we obtain the matrix $`𝐚`$ minimizing the new action $`S_W(\mathrm{𝐀𝐔}_\mu (n))`$, we update the link variable $`𝐔_\mu (n)`$, and repeat the procedure for all the $`N(N1)/2`$ subgroups of SU(N) and for all lattice sites. This constitutes one cooling sweep. We iterate this procedure up to we obtain that the Wilson action is stable with a given precision (in this work, the eight relevant digit) and close to the value of the expected continuum action: $`S=8\pi ^2/N`$.
## 3 The solutions
As mentioned in the introduction, we are interested in solutions with minimal non-trivial action,
$$S=8\pi ^2|Q|=\frac{8\pi ^2}{N},$$
(13)
on a volume $`[L/2,L/2]^3\times [T/2,T/2]`$, with $`TL`$. When $`T\mathrm{}`$ these solutions represent vacuum to vacuum tunneling.
We have restricted our analysis to the following non-orthogonal twist tensors:
1. Spatial twist, always $`\stackrel{}{m}=(1,1,1)`$.
2. Temporal twist, two cases:
* $`\stackrel{}{k}=(1,0,0)`$ for $`N=3,4,5,8,10`$ the solution is in this case anti-self-dual, $`Q=1/N`$.
* $`\stackrel{}{k}=(n,n,n)`$ for $`N=3n+1=4,7,10,13,19,25`$ the solution is here self-dual, $`Q=1/N`$.
A list of all the lattices analyzed is presented in Table 1.
From the lattice configurations we can easily derive information concerning continuum quantities. Part of it can be extracted in a gauge invariant way, such is the case for instance of the eigenvalues of the field strength or the Polyakov loops. However to derive information concerning the gauge potential gauge fixing is needed.
The continuum field strength tensor is extracted, up to $`𝒪(a^2)`$, from the clover average of the plaquette:
$$𝐐_{\mu \nu }(n)=\frac{1}{4}\text{ }\text{}$$
(14)
through
$$𝐅_{\mu \nu }(na)=\frac{1}{a^2}\frac{1}{2i}\left[𝐐_{\mu \nu }(n)𝐐_{\mu \nu }^{}(n)\frac{11}{\mathrm{N}}\mathrm{Tr}\left(𝐐_{\mu \nu }(n)𝐐_{\mu \nu }^{}(n)\right)\right]$$
(15)
In terms of the gauge fixed links the gauge potential is,
$$𝐀_\mu \left[(n+\frac{1}{2})a\right]=\frac{1}{a}\frac{1}{2i}\left[𝐔_\mu ^{gf}(n)𝐔_\mu ^{gf}(n)\frac{11}{\mathrm{N}}\mathrm{Tr}\left(𝐔_\mu ^{gf}(n)𝐔_\mu ^{gf}(n)\right)\right].$$
(16)
The Polyakov loops $`𝐋_\mu (x)`$ are simply given on the lattice by the ordered product of the $`\mu `$-links corresponding to the path $`\gamma _\mu (x)`$ in Eq. (3).
Non-gauge invariant information about the configurations will be presented in the temporal gauge: $`𝐀_4=0`$. In addition we fix $`𝐀_i(t=\mathrm{})=0`$ which is allowed because in the $`T\mathrm{}`$ limit fractional instantons describe vacuum to vacuum tunneling. In this gauge $`𝐀_i(t=\mathrm{})=ı𝛀_4_i𝛀_4^{}`$, with $`𝛀_4`$ the temporal twist matrix, and the spatial twist matrices are constant. This is not yet a complete gauge fixing, we still have the freedom to make a global gauge transformation and also to multiply the twist matrices by an element of the center of the group. We have made use of the global gauge transformation to bring the spatial twist matrices to a particular form. An explicit construction of constant spatial twist matrices compatible with the twist $`\stackrel{}{m}=(1,1,1)`$ can be easily found, following ’t Hooft :
$$𝛀_3=𝐐𝛀_2=𝐏^{N1}𝛀_1=e^{\frac{i2\pi p}{N}}\mathrm{𝐏𝐐}^{N1}$$
(17)
with p an integer number taking values $`p=1,2,\mathrm{}\mathrm{N}`$, and $`𝐏`$, $`𝐐`$ the matrices,
$$𝐏=\left(\begin{array}{cccc}0& 1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& 1\\ (1)^{N+1}& 0& \mathrm{}& 0\end{array}\right)𝐐=e^{i\pi (1N)/N}\left(\begin{array}{cccc}\varphi _0& 0& \mathrm{}& 0\\ 0& \varphi _1& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& \mathrm{}& \varphi _{N1}\end{array}\right)$$
(18)
where $`\varphi _n=\mathrm{exp}(i2\pi n/N)`$ with $`n=0,1,\mathrm{},\mathrm{N}1`$.
The invariance under multiplication by an element of the center of the group is fixed by imposing that the Polyakov loops take the value $`Ae^{i\pi }`$ at the position where the energy density of the solution is maximal.
On the lattice the $`𝐀_4=0`$, $`𝐀_i(t=T/2)=0`$ gauge is implemented by transforming the corresponding link variables to the identity. The gauge transformation, $`\omega (n)`$, which implements the change, is constructed in the following way: choose a point in the time slice $`t=T/2`$, i.e. $`n^0=(n_t=1,\stackrel{}{n}^0)`$; $`\omega (n)`$ is the product of the link variables along a certain path connecting $`n^0`$ with $`n=(n_t,n_x,n_y,n_z)`$. In particular we choose $`n^0=(1,1,1,1)`$ and the path such that it reaches the point $`n`$ first in the $`x`$ direction up to $`n_x`$, then in the $`y`$ direction up to $`n_y`$, in the $`z`$ direction up to $`n_z`$ and finally in the $`t`$ direction up to $`n_t`$. In this gauge the information about the twist matrices is encoded in the links $`𝐔_0(n_t=N_t)`$, $`𝐔_i(n_t=1,n_i=N_s)`$, the latter are rotated to the form indicated in Eq. (17) .
### 3.1 Gauge-invariant quantities
1. Global quantities. In Table 1 we give the values obtained for the action $`S`$, electric and magnetic parts of the action, $`S_e`$ and $`S_b`$ respectively, and topological charge $`Q`$:
$`S={\displaystyle \frac{1}{2}}{\displaystyle \mathrm{Tr}\left(𝐅_{\mu \nu }𝐅_{\mu \nu }\right)d_4x}={\displaystyle \mathrm{Tr}\left(𝐄_i^2+𝐁_i^2\right)d_4x}`$
$`S_e={\displaystyle \mathrm{Tr}\left(𝐄_i^2\right)d_4x}S_b={\displaystyle \mathrm{Tr}\left(𝐁_i^2\right)d_4x}`$
$`Q={\displaystyle \frac{1}{16\pi ^2}}{\displaystyle \mathrm{Tr}\left(𝐅_{\mu \nu }\stackrel{~}{𝐅}_{\mu \nu }\right)d_4x}={\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \mathrm{Tr}\left(\stackrel{}{𝐄}\stackrel{}{𝐁}\right)d_4x}`$ (19)
where $`𝐄_i=𝐅_{4i}`$ and $`\frac{1}{2}ϵ_{ijk}𝐅_{ij}=𝐁_k`$. We can see from the data that the configurations obtained are (anti-)self-dual to a very good degree, being therefore solutions of the Euclidean equations of motion. Those values are very near to the continuum expected values $`S\mathrm{N}/8\pi ^2=1`$, $`S_e\mathrm{N}/8\pi ^2=0.5`$, $`S_b\mathrm{N}/8\pi ^2=0.5`$ and $`Q\mathrm{N}=\pm 1`$.
The first thing we should check is the scaling of the solutions. Without loss of generality we set the spatial length to $`l_s=1`$, being then the lattice spacing $`a=1/N_s`$. To see how the continuum limit $`a0`$ is approached we vary the lattice spacing while keeping all other parameters fixed (among them the ratio $`N_t/N_s`$). We fit the $`N`$ and $`a^2`$ dependence of the action to the expression $`S\mathrm{N}/8\pi ^2=1\mathrm{\Delta }a^2/(\mathrm{N}\sqrt{\mathrm{N}})`$ and obtain that for the value $`\mathrm{\Delta }=8.893`$ the data in table 1 are well described with errors smaller than the $`0.2\%`$ . From this fit we understand how the continuum limit is approached for any value of $`\mathrm{N}`$, and also that the $`\mathrm{N}`$ dependence is such that the lattice corrections decrease with increasing $`\mathrm{N}`$. This property will be discussed further later on, it implies that for large $`\mathrm{N}`$ we can obtain good continuum results already on rather coarse lattices, this is a rather general property which, as we will see, affects other quantities apart from the integrated action and charge densities.
2. Energy profile. The energy profile, defined as
$$ϵ(t)=\mathrm{Tr}\left(E_i^2(\stackrel{}{x},t)+B_i^2(\stackrel{}{x},t)\right)d_3\stackrel{}{x},$$
(20)
is located on a region of size $`N/3`$ and has only one maximum for all values of N up to $`N=13`$ (instanton profile) and a double peak structure for the values $`N=19,25`$. In Figures 1a and 1b we show the scaling with the lattice spacing for the solutions with $`N=3`$ and $`N=19`$. We can see that points coming from lattices with different sizes describe very similar curves, scaling towards the same continuum function.
For values of N up to $`N=13`$ , $`ϵ(t)`$ is well fitted by
$$\varphi (t)=\frac{1}{Acosh(wt)+Bt^2+C}.$$
(21)
The values obtained for the parameters $`A,B,C,w`$ are given in table 2. For the values $`N=19,25`$ we fit to the expression $`(\varphi (tt_0)+\varphi (t+t_0))/2`$ and also in table 2 we give the results of these fits.
The N dependence is such that
$$ϵ(t)\frac{\varphi (t/\mathrm{N})}{\mathrm{N}^2}$$
(22)
as illustrated in figures 1c and 1d where $`\mathrm{N}^2ϵ(t)`$ is plotted as a function of $`t/\mathrm{N}`$ for $`\stackrel{}{k}=(1,0,0)`$ and $`\mathrm{N}=4,5,8,10`$ in figure 1c and for $`\stackrel{}{k}=(n,n,n)`$ and $`\mathrm{N}=4,7,10,13,19,25`$ in figure 1d. We also plot in figure 1d the energy profile for the abelian solution described in the appendix (equations 50,51 and 52). We plot $`N^2ϵ(t)=24\pi ^2`$ for the values of $`t/N`$ between $`t/N=1/6`$ and $`t/N=1/6`$. This is the profile for the selfdual abelian solution in the $`\mathrm{N}\mathrm{}`$ limit. We can see that the energy profiles of the solutions with $`\stackrel{}{k}=(n,n,n)`$ are approaching the one of the abelian solution in the $`\mathrm{N}\mathrm{}`$ limit.
3. Action density. Defined as,
$$S(\stackrel{}{x},t)=\mathrm{Tr}\left(𝐄_i^2(\stackrel{}{x},t)+𝐁_i^2(\stackrel{}{x},t)\right).$$
(23)
For values of N up to $`N=13`$, the action density has only one maximum what we will call the center of the instanton. We fit the center and their first nearest neighbours to the expression,
$$S(\stackrel{}{x},t)=S_0\left(1\underset{i}{}\frac{(x_ix_i^0)^2}{a_i^2}\frac{(tt^0)^2}{a_t^2}\right)$$
(24)
where $`S_0`$ is the height, $`x_i^0`$,$`t^0`$ the position and $`a_i`$,$`a_t`$ the width of the maximum. The values obtained are shown in table 3. For the values $`N=19,25`$ we observe two maximum in the action density and we make the same fit for each one. The results are also shown in table 3. In both cases, the errors are obtained from the difference between the same quantities calculated for the electric and magnetic part of the action. We can see that with increasing N all maximum become spatially flat. In fact, this is a general and important property of the solutions; when N is large some quantities, among them the action density, are spatially independent. In particular this implies that all the coordinate dependence of the action density comes through the time dependence of the energy profile $`ϵ(t)`$ defined above. This fact allows to easily understand the decrease of the lattice artifacts with increasing N since generally constant fields give rise to a much smoother, continuum-like, behaviour.
4. Eigenvalues of $`𝐅_{\mu \nu }`$.
Since the solution is (anti) self-dual we only give the results for $`𝐁_i`$. The main properties for the eigenvalues of $`𝐅_{\mu \nu }`$ are illustrated in figures 2a, 2b, 2c and 2d. We only show the results for the solutions with twist $`\stackrel{}{k}=(n,n,n)`$ because the same properties are obtained for the solutions with twist $`\stackrel{}{k}=(1,0,0)`$.
In figure 2a we show the eigenvalues of $`𝐁_1`$ for the solutions with $`N=7`$ and $`\stackrel{}{k}=(2,2,2)`$. Very similar results are obtained if we plot $`𝐁_2`$ or $`𝐁_3`$ instead of $`𝐁_1`$. That we plot is the spatial average of each eigenvalue as a function of time. The error bars mean spatial dispersion of the eigenvalues (difference between the maximum and minimum value of the eigenvalue at each temporal point). The first property we observe is spatial independence of the eigenvalues. This property also holds for values of $`N7`$ as shown in figures 2b, 2c and 2d. From figure 2a we observe that we obtain good results from very coarse lattices. In this figure we plot points coming from lattices with the following sizes: $`3^3\times 21`$, $`5^3\times 35`$ and $`8^3\times 56`$, being the results almost independent of the lattice size. This property also holds for bigger values of $`N`$.
The behavior with N is shown in Figure 2b in which we plot the first, second and $`\mathrm{N}^{th}`$ eigenvalues of $`\mathrm{N}\times 𝐁_1`$ as a function of $`t/\mathrm{N}`$ for $`\mathrm{N}=7,13,25`$ and $`\stackrel{}{k}=(n,n,n)`$. The meaning of the points and the error bars is the same as in figure 2a. The first eigenvalue is approximately independent of N and the other N-1 become degenerate for increasing N. This structure is very similar to the one of the selfdual abelian solution described in the appendix (equations 50, 51 and 52) , the first eigenvalue of $`N\times B_1`$ takes the value $`N\times \frac{2\pi }{N}`$ and the other $`N1`$ are equal to the value $`N\times \frac{2\pi }{(N1)N}`$. To compare we plot on figure 2b the first eigenvalue for the abelian solution in the $`N\mathrm{}`$ limit, the value $`2\pi `$ for the interval $`1/6t/N1/6`$. We see that the first eigenvalue of $`B_i`$ for the solutions with $`\stackrel{}{k}=(n,n,n)`$ are approaching the one of the abelian solution in the $`\mathrm{N}\mathrm{}`$ limit.
The symmetry properties of the spatial twist vector $`\stackrel{}{m}`$ suggest us to consider the following combinations of $`B_i`$ fields,
$`𝐁_0={\displaystyle \frac{1}{\sqrt{3}}}\left(𝐁_1+𝐁_2+𝐁_3\right)`$
$`𝐁_L={\displaystyle \frac{1}{\sqrt{6}}}\left(2𝐁_1+𝐁_2+𝐁_3\right)`$
$`𝐁_T={\displaystyle \frac{1}{\sqrt{2}}}\left(𝐁_2+𝐁_3\right)`$ (25)
one parallel to the $`\stackrel{}{m}`$ vector and the other two perpendicular. Also interesting is that, if there is a common component in color space for $`B_i`$ fields, we will see this component appearing in $`B_0`$ and not in $`B_L`$ and $`B_T`$.
In figure 2c we show the first, second and $`\mathrm{N}^{th}`$ eigenvalues of $`N\times 𝐁_0`$ as a function of $`t/N`$ for the same solutions appearing in figure 2b (with the same meaning for points and error bars). We observe the same structure seen in $`𝐁_i`$ and the expected property if there is one common component in color space, the dominant eigenvalue is bigger than the one for $`𝐁_i`$. We can also observe that this eigenvalue of $`N\times 𝐁_0`$ is approaching to the shape of the one for the abelian selfdual solution in the $`N\mathrm{}`$ limit, in this case $`2\pi \sqrt{3}`$.
In figure 2d we show the first, second, $`(N1)^{th}`$ and $`N^{th}`$ eigenvalues of $`N\times 𝐁_L`$ as a function of $`t/N`$ for the same solutions of figure 2b (also with the same meaning for points and error bars). The same results are obtained if we show $`𝐁_T`$ instead of $`𝐁_L`$. The eigenvalue structure is completely different to the ones shown before. In this case the eigenvalues are distributed in pairs, each pair with two opposite values. Another interesting property is that at $`t/N=0`$ the eigenvalues go to zero very fast for large $`N`$, being $`𝐁_0`$ the only one non trivial in this limit.
5. Colour orientation of the field strength. This is studied by calculating the commutators of $`𝐁_i`$’s fields. We show in Figure 3a the eigenvalues of $`N^{3/2}[𝐁_1,𝐁_2]`$ as a function of $`t/\mathrm{N}`$ for the solutions of Figure 2b,c,d (again the error bars represent spatial dispersion). The same results are obtained if we plot the eigenvalues of $`N^{3/2}[𝐁_2,𝐁_3]`$ or $`N^{3/2}[𝐁_1,𝐁_3]`$ instead of $`N^{3/2}[𝐁_1,𝐁_2]`$. Only two eigenvalues are relevant being the other N-2 very close to zero. We can see a nice scaling with $`N^{3/2}`$ of the two relevant eigenvalues at points with $`|t|/N>0.1`$, but not at the center of the solution in which the approach to zero is faster. For these quantities the spatial independence for large N also holds.
To study the abelian content of $`𝐅_{\mu \nu }`$ we calculate the quantity:
$$cos^2\left(\alpha _{ij}\right)=\frac{tr(𝐁_i𝐁_j)^2}{tr(𝐁_i)^2tr(𝐁_j)^2}$$
(26)
Whenever $`cos^2\left(\alpha _{ij}\right)=1`$ the solution is abelian. In Figure 3b we plot $`cos^2\left(\alpha _{12}\right)`$ as a function of $`t/\mathrm{N}`$ for the solutions with $`\mathrm{N}=4,7,10,13,19,25`$ and $`\stackrel{}{k}=(n,n,n)`$. As in other figures, the points mean the spatial average of the quantity and errors the spatial dispersion. For large N the solutions become abelian at the instanton center (t=0). $`cos^2\left(\alpha _{ij}\right)`$ shows again $`\stackrel{}{x}`$ independence for large N.
For the solutions with $`\stackrel{}{k}=(1,0,0)`$ we obtain the same behaviour near the instanton center, but for $`T\pm \mathrm{}`$, $`cos^2\left(\alpha _{ij}\right)`$ goes towards a $`\mathrm{N}`$ depending constant.
6. Polyakov Loops and structure of vacuum. In the gauge we have chosen the relationship between the Polyakov loops and the twist matrices is specially clear. We have $`𝐀_4=0`$ and hence the temporal Polyakov loop directly provides the twist matrix $`𝛀_4(\stackrel{}{x})`$. At $`t=\mathrm{}`$ $`𝐀_i`$ is also fixed to zero and in consequence
$$𝐋_i(\stackrel{}{x},t=\mathrm{})=𝛀_i$$
(27)
Since the spatial twist matrices are constant, compatibility with the spatial boundary conditions for $`\stackrel{}{m}=(1,1,1)`$ (see Eq. (4)), implies that $`\mathrm{N}_\mu (t=\mathrm{})=\mathrm{Tr}(𝐋_\mu (t=\mathrm{}))=0`$. From the boundary conditions in the time direction:
$$_i(\stackrel{}{x},t=T/2)=\mathrm{exp}\left(ik_i\frac{2\pi }{\mathrm{N}}\right)_i(\stackrel{}{x},t=T/2),$$
it is clear that also at $`t=\mathrm{}`$ the spatial Polyakov loops are zero. To characterize the vacuum states between which the configurations interpolate we need thus an additional quantity provided by $`\mathrm{N}_{zyx}=\mathrm{Tr}(𝐋_z𝐋_y𝐋_x)`$. Using the twist matrices given in Eq. (17) the values for $`_{zyx}`$ in a vacuum are,
$$_{zyx}=\mathrm{exp}(ı2\pi p/\mathrm{N})$$
(28)
where $`p`$ takes the values $`p=1,\mathrm{},\mathrm{N}`$. There are therefore N different vacua labeled by the value of $`_{zyx}`$. Our solutions interpolate between two of them as can be seen from the boundary condition in the time direction for $`_{zyx}`$:
$$_{zyx}(x_j,t=T/2)=\mathrm{exp}\left(ı\underset{i}{}k_i\frac{2\pi }{\mathrm{N}}\right)_{zyx}(x_j,t=T/2).$$
(29)
We can parametrize the data obtained for $`_i`$, $`_0`$, $`_{zyx}`$ as,
$`_i(x_{ji},t)=f_i(t)e^{i\alpha _i(t)}\mathrm{exp}\left\{{\displaystyle \frac{i2\pi }{\mathrm{N}}}(\stackrel{}{m}\times \stackrel{}{r})_i\right\}`$ (30)
$`_0(\stackrel{}{x})=Ae^{i\alpha _0}\mathrm{exp}\left\{{\displaystyle \frac{i2\pi }{\mathrm{N}}}\stackrel{}{k}\stackrel{}{r}\right\}`$ (31)
$`_{zyx}(\stackrel{}{x},t)=f_{zyx}(t)e^{i\alpha _{zyx}(t)}`$ (32)
note that with the gauge fixing condition for the Polyakov loops those functions take the values: $`\alpha _0=\pi `$, $`\alpha _i(t=0)=\pi `$ ($`\stackrel{}{r}=\stackrel{}{0}`$ is the maximum of the solution).
For $`_i`$ we make a fit at each temporal point to the spatial dependence of equation 30. The values of $`\sqrt{\chi ^2/N_{s}^{}{}_{}{}^{3}}`$ obtained for the solutions with temporal twist vector $`\stackrel{}{k}=(n,n,n)`$ are always smaller than $`1.24^o`$, $`0.27^o`$, $`0.07^o`$ and $`0.004^o`$ for the values $`N=4,7,10`$ and $`13`$ respectively. For the solutions with temporal twist vector $`\stackrel{}{k}=(1,0,0)`$ these values are $`2.7^o`$, $`1.28^o`$, $`0.81^o`$, $`0.13^o`$ and $`0.05^o`$ for $`N=3,4,5,8`$ and $`10`$ respectively. We extract the values of $`_i`$ at the spatial maximum of the solution and make a fit to the expression:
$`f_i(t)={\displaystyle \frac{B}{N}}cosh(wt/N)`$ (33)
the values of $`B`$ and $`w`$ are given in table 4.
The functions $`f_{zyx}`$, $`\alpha _{zyx}^{(1,0,0)}`$ and $`\alpha _{zyx}^{(n,n,n)}`$ parametrizing $`_{zyx}`$ in equation 32 are well fitted by:
$`f_{zyx}(t)=1{\displaystyle \frac{A_{zyx}}{Ncosh(w_{zyx}t/N)}}`$
$`\alpha _{zyx}^{(1,0,0)}(t)={\displaystyle \frac{\pi }{N}}\left(1+tagh(v_{zyx}{\displaystyle \frac{t}{N}})\right)`$
$`\alpha _{zyx}^{(n,n,n)}(t)={\displaystyle \frac{\pi }{N}}\left(1tagh(v_{zyx}{\displaystyle \frac{t}{N}})\right)`$ (34)
the values obtained for $`A_{zyx}`$, $`w_{zyx}`$, $`v_{zyx}`$ are given in table 4. And finally $`_0`$ only needs the constant $`A`$ to be fitted, their values are given in table 4.
We compare our results with the Polyakov loops for the abelian selfdual solution described in the appendix,
$`_i(x_{ji},t)={\displaystyle \frac{1}{N}}e^{i\pi }Exp\left(i{\displaystyle \frac{2\pi }{N}}{\displaystyle \frac{N1}{3}}{\displaystyle \frac{t}{T}}\right)Exp\left(i{\displaystyle \frac{2\pi }{N}}(\stackrel{}{m}\times \stackrel{}{x})_i\right)`$
$`_0(\stackrel{}{x})={\displaystyle \frac{1}{N}}e^{i\pi }Exp\left\{{\displaystyle \frac{i2\pi }{\mathrm{N}}}\stackrel{}{k}\stackrel{}{x}\right\}`$
where $`T/2tT/2`$ and $`0.5x_i0.5`$. We can see that the Polyakov loops in the $`N\mathrm{}`$ limit for the solutions with $`\stackrel{}{k}=(n,n,n)`$ are the same as the ones for the abelian selfdual solution.
### 3.2 Gauge-dependent quantities
1. Eigenvalues of $`𝐀_i`$. After gauge fixing to the gauge described previously, we calculate the eigenvalues of $`𝐀_i`$. The main properties of these quantities are shown in figures 4a, 4b, 4c and 4d. We only show the results for the solutions with twist $`\stackrel{}{k}=(n,n,n)`$, because the same figures are obtained for the solutions with twist $`\stackrel{}{k}=(1,0,0)`$ changing the sign of the gauge field, $`𝐀_i𝐀_i`$.
In figure 4a we show how the eigenvalues of $`𝐀_1`$ scale towards the continuum limit. We plot these quantities for the solutions with $`N=7`$, temporal twist vector $`\stackrel{}{k}=(2,2,2)`$ and lattice sizes $`3^3\times 21`$, $`5^3\times 35`$ and $`8^3\times 56`$. As in previous figures, points values are the spatial average of the eigenvalues and error bars mean spatial dispersion. We can see that the discretization errors still are important for the solution with size $`N_s=3`$ but are very small for bigger sizes. For these quantities we consider in the following lattice sizes with $`N_s4`$.
The behavior with N is shown in Figure 4b in which we plot the first, second and $`\mathrm{N}^{th}`$ eigenvalues of $`𝐀_1`$ as a function of $`t/\mathrm{N}`$ for $`\mathrm{N}=7,10,13`$ and $`\stackrel{}{k}=(n,n,n)`$ (points and error bars have the same meaning as before). Very similar results are obtained if we plot $`𝐀_2`$ or $`𝐀_3`$ instead of $`𝐀_1`$. One of the eigenvalues is approximately independent of N and the other N-1 become degenerate and approach zero for increasing $`N`$. We also show the biggest eigenvalue for the abelian selfdual solution described in section 3 of the appendix in the $`N\mathrm{}`$ limit. In this limit the function describing this eigenvalue is $`2\pi (\frac{t}{N}+\frac{1}{6})`$ for the values $`\frac{1}{6}\frac{t}{N}\frac{1}{6}`$. We can see that the first eigenvalue of $`A_1`$ is approaching the one of the abelian selfdual solution in the $`N\mathrm{}`$ limit.
As for $`𝐁_i`$ fields, we consider the following combinations of $`𝐀_i`$ fields,
$`𝐀_0={\displaystyle \frac{1}{\sqrt{3}}}\left(𝐀_1+𝐀_2+𝐀_3\right)`$
$`𝐀_L={\displaystyle \frac{1}{\sqrt{6}}}\left(2𝐀_1+𝐀_2+𝐀_3\right)`$
$`𝐀_T={\displaystyle \frac{1}{\sqrt{2}}}\left(𝐀_2+𝐀_3\right)`$ . (35)
In figure 4c we show the first, second and $`N^{th}`$ eigenvalues of $`𝐀_0`$ for the same solutions appearing in figure 4b (points and error bars have the same meaning as before). We observe that the eigenvalue structure is the same one seen for $`𝐀_1`$ and the expected result if there is a common component in color space for $`𝐀_1`$, $`𝐀_2`$ and $`𝐀_3`$, the first eigenvalue of $`𝐀_0`$ is bigger than the first one for $`𝐀_i`$. We also show the first eigenvalue of $`𝐀_0`$ for the abelian selfdual solution in the $`\mathrm{N}\mathrm{}`$ limit, in this case the function $`2\pi \sqrt{3}(\frac{t}{N}+\frac{1}{6})`$ for points $`\frac{1}{6}\frac{t}{N}\frac{1}{6}`$. The first eigenvalue of $`A_0`$ is approaching the one of the abelian selfdual solution in the $`N\mathrm{}`$ limit.
In figure 4d we show the first, second, $`(N1)^{th}`$ and $`N^{th}`$ eigenvalues of $`\sqrt{N}\times 𝐀_L`$ for the same solutions appearing in figure 4b (points and error bars have the same meaning as before). Very similar results are obtained if we plot the eigenvalues of $`𝐀_T`$ instead of the ones for $`𝐀_L`$. The eigenvalue structure is completely different to the one shown for $`𝐀_0`$, $`𝐀_i`$. The eigenvalues are distributed in pairs, each pair with two opposite values. As can be seen from the figure these eigenvalues goes to zero as $`1/\sqrt{N}`$. This means that $`𝐀_L`$ and $`𝐀_T`$ go to zero for large N while $`𝐀_0`$ is independent of $`\mathrm{N}`$.
2. Colour orientation of the potential. As for the field strength, this is studied by calculating the commutators of $`𝐀_i`$’s fields. We show in Figure 5 the eigenvalues of $`N^{1/2}[𝐀_1,𝐀_2]`$ as a function of $`t/\mathrm{N}`$ for the solutions of Figures 4b, 4c and 4d (again the error bars represent spatial dispersion). Only two eigenvalues are relevant being the other N-2 very close to zero and it is also clear the scaling with $`N^{1/2}`$ of the two relevant eigenvalues. From this figure we conclude that the gauge field $`𝐀_i`$ become abelian in the large N limit, compatible with the properties presented before, $`𝐀_0`$ is the remaining component in this limit while $`𝐀_L`$, $`𝐀_T`$ goes to zero with $`N^{1/2}`$.
3. Twist matrix $`𝛀_\mathrm{𝟒}`$. The temporal twist matrix behaves very differently depending of the temporal twist used. Note that their trace is the temporal Polyakov loop given before. We study their eigenvalues and obtain the following. For temporal twist $`\stackrel{}{k}=(n,n,n)`$ we fit the eigenvalues to the expression:
$`\lambda _j=Exp\left({\displaystyle \frac{ı2\pi \stackrel{}{k}\stackrel{}{r}}{N(N1)}}+ı{\displaystyle \frac{\pi }{3}}+ıN\pi \right)j=1,\mathrm{},(N1)/3`$
$`\lambda _j=Exp({\displaystyle \frac{ı2\pi \stackrel{}{k}\stackrel{}{r}}{N(N1)}}+ı\pi )+ıN\pi )j=(N1)/3+1,\mathrm{},2(N1)/3`$
$`\lambda _j=Exp\left({\displaystyle \frac{ı2\pi \stackrel{}{k}\stackrel{}{r}}{N(N1)}}+ı{\displaystyle \frac{5\pi }{3}}+ıN\pi \right)j=2(N1)/3+1,\mathrm{},N1`$
$`\lambda _N=Exp\left({\displaystyle \frac{i2\pi \stackrel{}{k}\stackrel{}{r}}{N}}+i\pi \right)`$ (36)
This expression was obtained from the numerical data for the eigenvalues of the twist matrix $`𝛀_\mathrm{𝟒}`$. This is a good parametrization for all values of $`N`$ studied ($`N=4,7,10,13,19,25`$) and the fits to this expression are better for bigger values of N. The interesting point is that this expression for the eigenvalues is the one for the eigenvalues of the twist matrix $`𝛀_\mathrm{𝟒}`$ given in 52 for the abelian solution described in the appendix.
For temporal twist $`\stackrel{}{k}=(1,0,0)`$ we obtain for large $`N`$ that the eigenvalues only depend on the $`x_1`$ coordinate, but we have not found a good parametrization in this case.
## 4 Conclusions
In this paper we have presented a set of solutions of the SU(N) Yang Mills equations of motion. These solutions are selfdual or antiselfdual, have fractional topological charge $`Q=1/N`$ and live on the four dimensional torus, $`T^4`$. We have studied the case when the lengths of the torus are $`L^3\times T`$ with $`TL`$ and with twist vectors $`\stackrel{}{m}=(1,1,1)`$ and $`\stackrel{}{k}=\frac{N1}{3}(1,1,1)`$ , $`\stackrel{}{k}=(1,0,0)`$. Now we summarize the main results we have obtained.
The obtained results show a clear tendency to describe continuum functions, being the different lattice sizes used enough to observe independence of the number of lattice points. This property indicates that the obtained configurations describe continuum Yang-Mills fields.
For each value of $`N`$ and twist we always obtain the same solution up to a gauge transformation and a spatial translation. This means that we can repeat the procedure to obtain another configuration and the differences observed will be a gauge transformation and a spatial translation.
The main characteristic of the solutions are the following
* The obtained solutions are selfdual or antiselfdual in all the studied cases. We observe numericaly that this property is satisfied with a very high precision. This guarantees that these configurations are solutions of the equations of motion.
* The size of the solutions is approximately $`N/3`$. By size we understand the length of the region in the temporal direction in which the core of the solution is included. We can see this property in all quantities shown in section 3. For example, we can look at the energy profile and check that the most relevant part is located in a region of size $`N/3`$. Also in the same region the eigenvalues of the field strength take their maximum values and go to zero out of this zone. We can say similar assertions for all quantities calculated in this paper.
* The action density has only one maximum for values of $`N13`$ and a double peak structure for bigger values. The spatial dependence of the action density disappears with increasing $`N`$ being only dependent of the temporal coordinate.
* The orientation of the field strength $`𝐅_{\mu \nu }`$ in color space is very dependent on the value of $`N`$. For smaller values of $`N`$ the different components of $`𝐅_{\mu \nu }`$ are built from different components in color space while for bigger values of $`N`$ the same component in color space gives the main contribution to the field strength. This property also holds for the gauge field $`𝐀_\mu `$. This means that, in the $`N\mathrm{}`$ limit, these solutions are abelians.
* In the $`N\mathrm{}`$ limit one eigenvalue gives the most important contribution to some quantities calculated in this article. In this limit, each component of the field strength $`𝐅_{\mu \nu }`$ has one eigenvalue which is approximately $`N`$ times bigger than the other $`N1`$ eigenvalues. This property also holds for each component of the gauge field $`𝐀_\mu `$.
* Independence of the temporal twist vector $`\stackrel{}{k}`$ for some of the calculated quantities. This is an $`N`$ independent property which holds when the length in the temporal direction goes to $`\mathrm{}`$. This property can be seen, for example, in the field strength $`𝐅_{\mu \nu }`$; if we change the temporal twist vectors used, $`\stackrel{}{k}_1=(1,0,0)`$ and $`\stackrel{}{k}_2=(n,n,n)`$ we obtain that $`𝐁_𝐢^{\stackrel{}{𝐤}_\mathrm{𝟏}}=𝐁_𝐢^{\stackrel{}{𝐤}_\mathrm{𝟐}}`$ and $`𝐄_𝐢^{\stackrel{}{𝐤}_\mathrm{𝟏}}=𝐄_𝐢^{\stackrel{}{𝐤}_\mathrm{𝟐}}`$. A similar property is held for the gauge field, under the change of the twist vectors we obtain the relation $`𝐀_𝐢^{\stackrel{}{𝐤}_\mathrm{𝟏}}=𝐀_𝐢^{\stackrel{}{𝐤}_\mathrm{𝟐}}`$.
We have not succeeded in finding an analytic expression describing the properties of the studied solutions. Obviously, the first requirement for an ansatz prepared to find the analytical expression is that this ansatz satisfies the previously shown properties. The most promising approach seems to be an ansatz based on the similarity of the solution in the large $`\mathrm{N}`$ limit with the abelian solution presented in the appendix. Nevertheless, if the solutions in this limit coincide with the abelian solution, something singular must happen at points $`|t/N|=1/6`$, being therefore not ease to use this similarity to find the analytic expression. To conclude, we hope that all the numerical data presented will be helpful for other attempts to find the analytical expression of the solutions presented in this paper.
## A. Analytic solutions.
’t Hooft construction. The $`𝐀_\mu `$ and $`𝐅_{\mu \nu }`$ fields are built from a diagonal matrix $`𝐓`$, in the following way
$$𝐀_\mu (x)=\frac{\pi }{N}\underset{\nu }{}\frac{\alpha _{\mu \nu }x_\nu }{l_\mu l_\nu }𝐓,𝐅_{\mu \nu }(x)=\frac{2\pi }{N}\frac{\alpha _{\mu \nu }}{l_\mu l_\nu }𝐓.$$
(37)
where $`\alpha _{\mu \nu }`$ is an antisymmetric tensor and $`l_\mu `$ the length of the torus in the $`\mu `$ direction. The matrix $`𝐓`$ has the form,
$$𝐓=\left(\begin{array}{cc}l11_{k\times k}& \mathrm{𝟎}_{k\times l}\\ \mathrm{𝟎}_{l\times k}& k11_{l\times l}\end{array}\right)$$
(38)
being $`k`$ and $`l`$ integer numbers ($`k+l=N`$). To build the twist matrices we use the $`𝐏`$ and $`𝐐`$ matrices defined in equation 18. From these matrices we construct another set of matrices,
$$𝐏_1,𝐐_1=\left(\begin{array}{cc}(𝐏,𝐐)_{k\times k}& \mathrm{𝟎}_{k\times l}\\ \mathrm{𝟎}_{l\times k}& 11_{l\times l}\end{array}\right);𝐏_2,𝐐_2=\left(\begin{array}{cc}11_{k\times k}& \mathrm{𝟎}_{k\times l}\\ \mathrm{𝟎}_{l\times k}& (𝐏,𝐐)_{l\times l}\end{array}\right)$$
(39)
satisfying the properties,
$`𝐏_1𝐐_1=𝐐_1𝐏_1Exp\left\{{\displaystyle \frac{i2\pi }{N}}\left(11{\displaystyle \frac{𝐓}{k}}\right)\right\};𝐏_2𝐐_2=𝐐_2𝐏_2Exp\left\{{\displaystyle \frac{i2\pi }{N}}\left(11+{\displaystyle \frac{𝐓}{l}}\right)\right\}.`$ (40)
And the ansatz for the twist matrices is,
$$𝛀_\mu (x)=𝐏_1^{s_\mu }𝐐_1^{t_\mu }𝐏_2^{u_\mu }𝐐_2^{v_\mu }Exp\left\{\frac{i\pi }{N}\underset{\nu }{}\frac{\alpha _{\mu \nu }x_\nu }{l_\nu }𝐓\right\}$$
(41)
where $`s_\mu `$,$`t_\mu `$,$`u_\mu `$ and $`v_\mu `$ are arbitrary integer numbers. These matrices must satisfy the consistency condition,
$`𝛀_\mu (x_\nu +l_\nu )𝛀_\nu (x_\mu )=𝛀_\nu (x_\mu +l_\mu )𝛀_\mu (x_\nu )\mathrm{Exp}\left({\displaystyle \frac{i2\pi n_{\mu \nu }}{\mathrm{N}}}\right).`$ (42)
This condition imposes the following equations for $`s_\mu `$,$`t_\mu `$,$`u_\mu `$ and $`v_\mu `$
$`{\displaystyle \frac{1}{k}}(t_\mu s_\nu t_\nu s_\mu )=l{\displaystyle \frac{\alpha _{\mu \nu }}{N}}+{\displaystyle \frac{n_{\mu \nu }}{N}}+A_{\mu \nu };{\displaystyle \frac{1}{l}}(v_\mu u_\nu v_\nu u_\mu )=k{\displaystyle \frac{\alpha _{\mu \nu }}{N}}+{\displaystyle \frac{n_{\mu \nu }}{N}}+B_{\mu \nu }`$ (43)
where $`A_{\mu \nu }`$ and $`B_{\mu \nu }`$ are integer numbers. To solve these two equations we give to $`\alpha _{\mu \nu }`$ and $`n_{\mu \nu }`$ the form,
$`n_{\mu \nu }=n_{\mu \nu }^{(1)}+n_{\mu \nu }^{(2)};\alpha _{\mu \nu }={\displaystyle \frac{n_{\mu \nu }^{(1)}}{k}}{\displaystyle \frac{n_{\mu \nu }^{(2)}}{l}}`$ (44)
and then equations in formula 43 are transformed to,
$`n_{\mu \nu }^{(1)}=t_\mu s_\nu t_\nu s_\mu +kA_{\mu \nu };n_{\mu \nu }^{(2)}=v_\mu u_\nu v_\nu u_\mu +lB_{\mu \nu }`$ (45)
which can be solved if the following condition is satisfied,
$`n_{\mu \nu }^{(1)}\stackrel{~}{n}_{\mu \nu }^{(1)}=0modk;n_{\mu \nu }^{(2)}\stackrel{~}{n}_{\mu \nu }^{(2)}=0modl`$ (46)
The $`n_{\mu \nu }^{(1)}`$ and $`n_{\mu \nu }^{(2)}`$ tensors are orthogonal twist tensors for a SU(k) and SU(l) group respectively.
The topological charge for these solutions is,
$$Q=\frac{1}{N}\frac{\alpha _{\mu \nu }\stackrel{~}{\alpha }_{\mu \nu }}{4}kl=\frac{1}{N}\left(\stackrel{}{k}^{(1)}\stackrel{}{m}^{(1)}\frac{l}{k}+\stackrel{}{k}^{(2)}\stackrel{}{m}^{(2)}\frac{k}{l}\stackrel{}{k}^{(1)}\stackrel{}{m}^{(2)}\stackrel{}{k}^{(2)}\stackrel{}{m}^{(1)}\right)$$
(47)
where we have defined the vectors $`k_{i}^{}{}_{}{}^{(n)}=n_{0i}^{}{}_{}{}^{(n)}`$, $`m_{i}^{}{}_{}{}^{(n)}=ϵ_{ijk}n_{jk}^{}{}_{}{}^{(n)}/2`$, with $`n=1,2`$.
Some examples. Now we give some examples of solutions built using ’t Hooft construction, and for the torus lengths used in this article: $`l_x=l_y=l_z=1`$ and $`l_t\mathrm{}`$. In fact, our examples will be given for any value of $`l_t`$. The topological charge in our examples is given by equation 47 and the action by the following equation,
$`S={\displaystyle \frac{8\pi ^2}{N}}{\displaystyle \frac{1}{2}}\{{\displaystyle \frac{1}{l_t}}(\stackrel{}{k}^{(1)}\stackrel{}{k}^{(1)}{\displaystyle \frac{l}{k}}+\stackrel{}{k}^{(2)}\stackrel{}{k}^{(2)}{\displaystyle \frac{k}{l}}2\stackrel{}{k}^{(1)}\stackrel{}{k}^{(2)})+`$
$`l_t(\stackrel{}{m}^{(1)}\stackrel{}{m}^{(1)}{\displaystyle \frac{l}{k}}+\stackrel{}{m}^{(2)}\stackrel{}{m}^{(2)}{\displaystyle \frac{k}{l}}2\stackrel{}{m}^{(1)}\stackrel{}{m}^{(2)})\}.`$ (48)
The minimum value for the action, $`S=8\pi ^2|Q|`$, is obtained when the solution is selfdual or antiselfdual. This condition imposes a value for the temporal length $`l_t`$, obtained solving the equation,
$$\frac{\stackrel{}{k}^{(1)}}{l}\frac{\stackrel{}{k}^{(2)}}{k}=\pm l_t\left(\frac{\stackrel{}{m}^{(1)}}{l}\frac{\stackrel{}{m}^{(2)}}{k}\right)$$
(49)
the positive sign for selfdual solutions and the negative sign for antiselfdual solutions. Now we give some examples for the twist vectors used in this article:
1. Solutions for twist vectors $`\stackrel{}{m}=(1,1,1)`$ and $`\stackrel{}{k}=(n,n,n)`$, with $`N=3n+1`$ for $`n=1,2,3,\mathrm{}`$. We choose $`k=N1`$ and $`l=1`$ and the twist vectors in subspaces SU(k) and SU(l) as,
$`\stackrel{}{m}^{(1)}=(1,1,1)\stackrel{}{k}^{(1)}=(n,n,n);\stackrel{}{m}^{(2)}=(0,0,0)\stackrel{}{k}^{(2)}=(0,0,0)`$
The topological charge for this solution is $`Q=1/N`$ and it is selfdual when $`l_t=n`$.
2. Solutions for twist vectors $`\stackrel{}{m}=(1,1,1)`$ and $`\stackrel{}{k}=(1,0,0)`$. We choose the twist vectors in $`SU(k)`$ and $`SU(l)`$ subspaces as,
$`\stackrel{}{m}^{(1)}=(1,1,1)\stackrel{}{k}^{(1)}=(0,0,0);\stackrel{}{m}^{(2)}=(0,0,0)\stackrel{}{k}^{(2)}=(1,0,0)`$
This choice works for any values of $`k,l`$. The topological charge for this solution is $`Q=1/N`$ and in this case it is not possible to solve the selfduality equation.
Changing the gauge. Now we change the gauge for one of the solutions described before to the gauge used in section 3. We choose the solution given in example 1. In this case $`k=N1`$, $`l=1`$, $`n_{\mu \nu }^{(1)}=n_{\mu \nu }`$ and $`n_{\mu \nu }^{(2)}=0`$. The fields are,
$`𝐀_\mu (x)={\displaystyle \frac{\pi }{N(N1)}}{\displaystyle \underset{\nu }{}}{\displaystyle \frac{n_{\mu \nu }x_\nu }{l_\mu l_\nu }}𝐓,𝐅_{\mu \nu }(x)={\displaystyle \frac{2\pi }{N(N1)}}{\displaystyle \frac{n_{\mu \nu }}{l_\mu l_\nu }}𝐓.`$ (50)
we remember that the torus lengths were $`l_x=l_y=l_z=1`$ and $`l_t`$ can take any value, and the twist vectors were $`\stackrel{}{m}=(1,1,1)`$ and $`\stackrel{}{k}=(n,n,n)`$, with $`N=3n+1`$. The action and the topological charge take the values,
$`S={\displaystyle \frac{8\pi ^2}{N}}{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{l_t}{n}}+{\displaystyle \frac{n}{l_t}}\right)Q={\displaystyle \frac{1}{N}}`$
the twist matrices are,
$$𝛀_\mu (x)=𝐏_1^{s_\mu }𝐐_1^{t_\mu }Exp\left\{\frac{i\pi }{N(N1)}\underset{\nu }{}\frac{n_{\mu \nu }x_\nu }{l_\nu }𝐓\right\}$$
and the values for $`s_\mu `$, $`t_\mu `$,
$`s_1=0s_2=1s_3=N2s_4=2(N1)/3`$
$`t_1=1t_2=0t_3=N2t_4=(N1)/3.`$
The gauge used in section 3 is: $`𝐀_4=0`$ and $`𝐀_i(t=\mathrm{})=0`$. For our example we can not use the same gauge because $`t`$ is finite and $`𝐅_{\mu \nu }0`$ for any value of $`t`$ (this is the condition needed to put $`A_i(t)=0`$ at some point $`t`$). The most similar gauge condition is the following one,
$`𝐀_4=0,𝐀_3(t=0)=0,𝐀_2(t=0,z=0)=0,𝐀_1(t=0,z=0,y=0)=0`$
because all links associated with these fields were rotated to the identity. The field $`𝐀_\mu `$ and the twist matrices in this gauge are,
$`𝐀_4=0`$
$`𝐀_3={\displaystyle \frac{2\pi }{N(N1)}}𝐓{\displaystyle \frac{n}{l_t}}t`$
$`𝐀_2={\displaystyle \frac{2\pi }{N(N1)}}𝐓\left({\displaystyle \frac{n}{l_t}}tz\right)`$
$`𝐀_1={\displaystyle \frac{2\pi }{N(N1)}}𝐓\left({\displaystyle \frac{n}{l_t}}t+zy\right)`$ (51)
$`𝛀_4=𝐏_{1}^{}{}_{}{}^{s_4}𝐐_{1}^{}{}_{}{}^{t_4}Exp\left\{i{\displaystyle \frac{2\pi }{N(N1)}}𝐓\stackrel{}{k}\stackrel{}{r}\right\}`$
$`𝛀_3=𝐏_{1}^{}{}_{}{}^{s_3}𝐐_{1}^{}{}_{}{}^{t_3}Exp\left\{i{\displaystyle \frac{2\pi }{N(N1)}}𝐓(yx)\right\}`$
$`𝛀_2=𝐏_{1}^{}{}_{}{}^{s_2}𝐐_{1}^{}{}_{}{}^{t_2}Exp\left\{i{\displaystyle \frac{2\pi }{N(N1)}}𝐓x\right\}`$
$`𝛀_1=𝐏_{1}^{}{}_{}{}^{s_1}𝐐_{1}^{}{}_{}{}^{t_1}`$ (52)
We compare along section 3 the results obtained for large values of N with this solution (with $`l_t=n`$) because some properties are very similar for both solutions in this limit.
Finally, we give the value of the Polyakov loops for this example,
$`_\mu ={\displaystyle \frac{1}{N}}Exp\left(i{\displaystyle \frac{2\pi }{N}}{\displaystyle \underset{\nu }{}}{\displaystyle \frac{n_{\mu \nu }x_\nu }{l_\nu }}\right)`$
note that this quantity is gauge invariant and could be calculated with the two different $`𝐀_\mu `$ given before, obtaining the same result.
## Acknowledgements
The results presented in this article are part of the Ph. D. Thesis of the author . Most of the ideas developed in this work come from suggestions and fruitful discussions with Antonio González-Arroyo. I also acknowledge useful discussions with Margarita García Pérez and Carlos Pena. This work has been supported by the Spanish Ministerio de Educación y Cultura under a postdoctoral Fellowship and by CICYT under grant AEN97-1678.
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# Quantum Bit Escrow
## 1 Introduction
We start with an informal definition of a (very) weak variant of bit commitment. In this variant there is first a commitment stage in which Alice commits a bit $`b`$ to Bob. Later on there is a reveal stage in which Alice reveals the bit and Bob proves he played honestly. The protocol should be binding in the sense that if Alice changes her mind at revealing time then Bob has a good probability of catching her cheating, and sealing in the sense that if Bob learns information about the committed bit then Alice has a good probability of catching him cheating. Thus, the fundamental (and only) difference between this primitive and bit commitment is that in bit commitment Bob can not learn from the encoding any information about $`b`$, while in the weak primitive Bob can learn a lot of information about the encoded bit, but if he does so Alice catches him cheating with a good probability.
###### Definition 1
(Weak bit commitment) A weak bit commitment protocol is a quantum communication protocol between Alice and Bob which consists of two stages, the depositing stage and the revealing stage, and a final classical declaration stage at which both Alice and Bob each declare “accept” or “reject”. The following requirements should hold.
* If both Alice and Bob are honest, then at depositing stage Alice decides on a bit, $`b`$. She then communicates with Bob, where Alice’s protocol depends on $`b`$. At revealing stage Alice and Bob communicate, and during this stage Alice reveals to Bob the deposited bit $`b`$. Both Alice and Bob accept.
* (Binding) If Alice tries to change her mind about the value of $`b`$, then there is non zero probability that an honest Bob would reject.
* (Sealing) If Bob attempts to learn information about the deposited bit $`b`$, then there is non zero probability that an honest Alice would reject.
Later on, we will give more formal definitions of “Alice changing her mind” and “Bob learning information”, and we will quantify the degree to which a protocol is binding or sealing.
Now, consider the following protocol:
###### Protocol 1
(Bit Escrow) For an angle $`\alpha [\pi ,\pi ]`$ define $`\varphi _\alpha =\mathrm{cos}(\alpha )|0+\mathrm{sin}(\alpha )|1`$. Let,
$`\varphi _{b,x}`$ $`=`$ $`\{\begin{array}{cc}\varphi _\theta \hfill & b=0,x=0\hfill \\ \varphi _\theta \hfill & b=0,x=1\hfill \\ \varphi _{\frac{\pi }{2}\theta }\hfill & b=1,x=0\hfill \\ \varphi _{\frac{\pi }{2}+\theta }\hfill & b=1,x=1\hfill \end{array}`$
for some fixed angle $`\theta \frac{\pi }{8}`$, say, $`\theta =\frac{\pi }{8}`$. See Figure 1.
To deposit bit $`b`$, Alice picks a random $`x\{0,1\}`$, and sends $`\varphi _{b,x}`$ to Bob. Later on, one of the following two challenges is issued:
* Either Alice is asked to reveal the deposited bit, and then Alice sends the classical bits $`b`$ and $`x`$ to Bob <sup>1</sup><sup>1</sup>1 This means that when Bob gets the qubit $`q_b`$ that is supposed to carry a classical value for $`b`$, Bob measures $`q_b`$ first in the $`\{|0,|1\}`$ basis. We carry this convention throughout the paper.. Bob measures $`\varphi `$ according to the basis $`\{\varphi _{0,x},\varphi _{1,x}\}`$ and verifies that the result of the measurement is $`\varphi _{b,x}`$.
* Or Bob is asked to return the deposited qubit, he returns a qubit $`q`$, and Alice measures it in the $`\{\varphi _{0,x},\varphi _{1,x}\}`$ basis and verifies that it is $`\varphi _{b,x}`$.
We rigorously define and prove:
###### Theorem 1
Protocol 1 has the following properties:
* The deposited qubit does not reveal, in an information theoretic sense, all the information about the deposited bit $`b`$.
* (Binding) When Bob asks Alice to reveal the classical bit $`b`$ that she deposited, if Alice influences the value of b with advantage $`ϵ`$ then she is detected cheating with probability $`\mathrm{\Omega }(ϵ^2)`$.
* (Sealing) When Alice challenges Bob to return the deposited qubit, then if Bob can predict $`b`$ with advantage $`ϵ`$ then he is detected cheating with probability $`\mathrm{\Omega }(ϵ^2)`$.
Protocol 1 and Theorem 1 do not achieve the goal set in definition 1 of weak bit commitment. Definition 1 asks for a protocol that is both binding and sealing, i.e., a commitment s.t. if either player cheats he is detected cheating with a good probability. Protocol 1 and Theorem 1 only give a commitment that is either binding (if Alice has to reveal) or sealing (if Bob has to return the qubit), but not simultaneously both. We therefore call this protocol a bit escrow protocol. The question of achieving simultaneous binding and sealing i.e. a weak bit commitment protocol, is left open. This question was addressed in , who independently defined the binding and sealing properties, and we discuss it in section 1.2.
We describe soon how to use the first two properties in Theorem 1 to get a biased coin flipping protocol with a constant bias.
### 1.1 Quantum Coin flipping
Alice and Bob are going through a divorce. They want to decide by a coin flip over the phone who is going to keep the car. The problem is that they do not trust each other any more.
###### Definition 2
(Classical coin flipping) A coin flipping protocol with $`\delta `$ bias is one where Alice and Bob communicate and finally decide on a value $`c\{0,1\}`$ s.t. if at least one of the players is honest then for any strategy of the dishonest player $`Prob(c=0)[\frac{1}{2}\delta ,\frac{1}{2}+\delta ]`$.
Classical coin flipping can be implemented either by a trusted party or by assuming players with limited computational power and some cryptographic assumptions. However, if the players have unlimited computational power then no coin flipping protocol is possible in a classical world. This is because any protocol represents a two player game, and therefore game theory tells us that there is a player with an always winning strategy.
By contrast, in the quantum setting coin flipping (without computational assumptions) is not a priori ruled out. This is because any attempt by a player to measure extra information by deviating from the protocol can disturb the quantum state, and therefore be detected by the other player. This leads Lo and Chau and later Mayers et. al. to consider quantum coin flipping. There are several ways to define quantum coin flipping when cheaters can be detected. We define:
###### Definition 3
(quantum coin flipping) A quantum coin flipping protocol with bias $`\delta `$ is one where Alice and Bob communicate and finally each decides on a value $`c\{0,1,err\}`$. Let $`c_A`$ ($`c_B`$) denote Alice’s (Bob’s) result. We require:
* If both players are honest then $`c_A`$ always equals $`c_B`$, $`Prob(c_A=err)=0`$, and $`0`$ and $`1`$ have equal probability: $`Prob(c_A=0)=Prob(c_A=1)=\frac{1}{2}`$.
* If one of the players is honest and the other is not, then for any strategy of the dishonest player, the honest player’s result $`c`$ satisfies for any $`b\{0,1\}`$:
$$Prob(c=b)\frac{1}{2}+\delta $$
Lo and Chau showed that there is no quantum coin flipping protocol with $`0`$ bias, under a certain restriction (“ideal coin flipping”.) Mayers et al generalized their proof to the general $`0`$ bias case. Lo and Chau leave open the question whether non-exact protocols exist. Mayers et al suggest a quantum coin flipping protocol that is based on a biased-coin protocol that is repeated many times. Mayers et al prove that it works well against some strong, natural attacks. However, no general proof is given or claimed for the coin-flipping protocol or the biased-coin sub-protocol.
We give a simple protocol for quantum biased coin flipping, with constant bias. It is a modification of protocol 1:
###### Protocol 2
(A biased coin flipping protocol)
* Alice picks $`b,x_R\{0,1\}`$ and sends Bob $`\varphi _{b,x}`$. We set $`\theta =\frac{\pi }{8}`$.
* Bob chooses $`b^{}_R\{0,1\}`$ and sends it to Alice.
* Alice sends Bob $`b`$ and $`x`$. Bob checks against the qubit she sent in the first step. The result of the game is $`r=err`$ if Alice is caught cheating and $`r=bb^{}`$ otherwise.
Based on the properties of protocol 1 we can prove that no player can fully control the game:
###### Theorem 2
Protocol 2 has $`\delta 0.42`$ bias.
i.e., no player can force his result with probability greater than $`0.92`$. We note that while our protocol is resilient against all powerful malicious quantum players, it requires only simple single qubit operations from the honest player. An intriguing question is whether quantum coin flipping protocols are possible for arbitrarily low biases.
### 1.2 Weak Bit Commitment?
Hardy and Kent (see Section 1.3) noticed that Protocol 1 can be used to give a weak bit commitment protocol if Alice and Bob can access a random independent coin flip. This is done as follows: at revealing time Alice first reveals the bit $`b`$, and then they receive a random independent coin flip. If the coin is $`0`$, Bob is challenged to convince Alice that he hasn’t been cheating, and if the coin flip turns out to be $`1`$, then Alice is challenged. This is still correct if the coin flip is biased, as long as both probabilities for $`0`$ and for $`1`$ are constant.
Since we already have a biased coin flipping protocol, we might consider using this biased coin flipping protocol combined with the bit escrow protocol to give a weak bit commitment protocol. Consider the following protocol (see Figure 2):
###### Protocol 3
To deposit bit $`b`$, Alice picks a random $`x\{0,1\}`$, and sends $`\varphi =\varphi _{b,x}`$ to Bob. To reveal the bit, Alice sends $`b`$ to Bob. Then a biased-coin flipping protocol (Protocol 2) is played.
* If Alice loses she is asked to reveal $`x`$ and Bob measures $`\varphi `$ according to the basis $`\{\varphi _{0,x},\varphi _{1,x}\}`$ and verifies that the result of the measurement is $`\varphi _{b,x}`$.
* If Bob loses he is asked to return the deposited qubit $`q`$, and Alice measures it in the $`\{\varphi _{0,x},\varphi _{1,x}\}`$ basis and verifies that it is $`\varphi _{b,x}`$.
It is left as an open question whether this protocol, or perhaps a protocol which uses a different coin flipping procedure, is actually a weak bit commitment protocol. The main difficulty in proving or disproving such a result is the issue of independence between the coin flipping protocol and the bit escrow protocol. In other words, one has to prove that the cheater cannot use entanglement to correlate the events of being detected cheating in the bit-escrow protocol and winning the biased coin flipping protocol, in such a way that the cheater is never challenged when he (or she) has positive probability of being detected.
It is our hope that our techniques could be extended to give weak bit commitment with $`\mathrm{\Omega }(ϵ^c)`$ binding and sealing for some constant $`c`$. Our results also show that Protocol 3 cannot be more than $`\mathrm{\Omega }(ϵ^2)`$ sealing or binding. It might be interesting to find a protocol that does better, or prove that such a protocol does not exist. It seems that a weak bit commitment protocol with better than quadratic security parameters can be used repeatedly to give a secure coin flipping protocol with unbounded bias.
### 1.3 Related Work
Some of the work presented here was independently done by Hardy and Kent . They independently defined the binding and sealing properties and the weak bit commitment primitive (giving it different names). The protocol they analyze is similar in structure to protocol 3. Hardy and Kent’s result asserts that a protocol similar to Protocol 3 is simultaneously sealing and binding. I.e., if Alice (Bob) uses a strategy that gives her (him) $`ϵ`$ advantage, then Alice (Bob) is detected cheating with some probability which is strictly greater than $`0`$ (they do not analyze the dependence of the detection probability on $`ϵ`$). However, no proof is given regarding the security against a cheater who tries to correlate the two parts of the protocol to his (or her) advantage.
## 2 Preliminaries
The model. Let $`\{e_1,\mathrm{},e_{2^n}\}`$ be an orthonormal basis for $`\text{ }\mathrm{C}^n`$, and let $`|i=|i_1,\mathrm{},i_n`$ be the vector $`e_i`$. A pure state over $`n`$ qubits is a vector $`v\text{ }\mathrm{C}^{2^n}`$ of norm $`1`$. Any pure state $`|v`$ can be expressed as $`|v=\mathrm{\Sigma }_ia_i|i`$, with $`\mathrm{\Sigma }_i|a_i|^2=1`$. A mixed state is a classical distribution over pure states, $`\{p_i,\varphi _i\}`$, where $`0p_i1`$, $`\mathrm{\Sigma }_ip_i=1`$ and $`\varphi _i`$ is a pure state, and the interpretation we give it is that the system is with probability $`p_i`$ in the pure state $`\varphi _i`$. A quantum system is, in general, in a mixed state. The system Alice builds in the first stage of Protocol 2 is in a mixed state that is with probability $`\frac{1}{4}`$ in some pure state $`\varphi _{b,x}`$.
A quantum system can undergo two basic operations: unitary evolution and measurement.
: If a unitary transformation $`U:\text{ }\mathrm{C}^{2^n}\text{ }\mathrm{C}^{2^n}`$ is applied to a pure state $`\varphi `$, then the new state of the system is the pure state $`U\varphi `$. If $`U`$ is applied to the mixture $`\{p_i,\varphi _i\}`$ then the new state of the system is the mixture $`\{p_i,U\varphi _i\}`$. The interpretation we give it is that with probability $`p_i`$ the system was in the pure state $`\varphi _i`$ hence it is now in the pure state $`U\varphi _i`$.
: An orthogonal measurement is a decomposition of the system into orthogonal subspaces. More formally, suppose the system is in a super position $`\varphi \text{ }\mathrm{C}^{2^n}`$. Suppose $`_1,\mathrm{},_k`$ are orthogonal subspaces, and $`\text{ }\mathrm{C}^{2^n}=_1\mathrm{}_k`$. A measurement of $`\varphi `$ according to the decomposition $`_1,\mathrm{},_k`$, will get result $`i`$ (or $`_i`$) with probability $`q_i=|\mathrm{\Pi }__i|\varphi |^2`$ where $`\mathrm{\Pi }__i`$ is the projection on subspace $`_i`$, and then the state will collapse to $`\frac{1}{\sqrt{q_i}}\mathrm{\Pi }__i|\varphi `$. In other words, $`\varphi `$ falls into the subspace $`_i`$ with probability which is the length of the projection squared, and the new vector is the normalized projected vector. An orthogonal measurement can be represented using an Hermitian matrix $`M`$ whose eigenspaces are the subspaces $`_i`$. A measurement of a mixture is the mixture of the measurements of the pure states.
Given a system $`\rho `$ on $`\text{ }\mathrm{C}^{2^n}`$, one can use an ancilla, say $`|0,\mathrm{},0\text{ }\mathrm{C}^{2^m}`$, apply a unitary transformation $`U:\text{ }\mathrm{C}^{2^n}\text{ }\mathrm{C}^{2^m}\text{ }\mathrm{C}^{2^n}\text{ }\mathrm{C}^{2^m}`$, and then an orthogonal measurement on $`\text{ }\mathrm{C}^{2^n}\text{ }\mathrm{C}^{2^m}`$. It turns out that this is the most general measurement possible. There are several equivalent ways to formulate this so called ’generalized measurement’, and we refer the interested reader to .
The Density Matrix. The density matrix of a pure state $`|\varphi `$ is the matrix $`|\varphi \varphi |`$, where $`\varphi |=((\varphi )^t)^{}`$ is the conjugate transpose of $`\varphi `$. For example, the density matrix of $`\varphi _{0,0}`$ is
$`|\varphi _\theta \varphi _\theta |`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}^2(\theta )& \mathrm{cos}(\theta )\mathrm{sin}(\theta )\\ \mathrm{cos}(\theta )\mathrm{sin}(\theta )& \mathrm{sin}^2(\theta )\end{array}\right)`$
The density matrix of a mixed state $`\{p_i,\varphi _i\}`$ is $`\mathrm{\Sigma }_ip_i|\varphi _i\varphi _i|`$. All density matrices are Hermitian, positive semi-definite and have trace $`1`$. If a unitary matrix $`U`$ operates on the system, it transforms the density matrix $`\rho `$ to $`U\rho U^{}`$. A measurement $`M`$ operating on a system whose density matrix is $`\rho `$ results in an expected outcome $`Trace(M\rho )`$.
Distinguishing Between Density Matrices. Given a quantum system $`\rho `$ and a generalized measurement $`𝒪`$ on it, let $`\rho ^𝒪`$ denote the classical distribution on the possible results that we get by measuring $`\rho `$ according to $`𝒪`$. i.e., it is some classical distribution $`p_1,\mathrm{},p_k`$ where we get result $`i`$ with probability $`p_i`$. Given two different mixed states, we can ask how well one can distinguish between the two mixtures. We need a measure for the distance between two classical distributions and we choose the $`l_1`$ norm:
###### Definition 4
Let $`p_1,\mathrm{},p_k`$ and $`q_1,\mathrm{},q_k`$ be two probability distributions over $`\{1,\mathrm{},k\}`$. Then $`|pq|_1=\mathrm{\Sigma }_i|p_iq_i|`$.
A fundamental theorem about distinguishing density matrices tells us:
###### Theorem 3
Let $`\rho _1,\rho _2`$ be two density matrices on the same space $``$. Then for any generalized measurement $`𝒪`$
$$|\rho _1^𝒪\rho _2^𝒪|_1Trace(\sqrt{A^{}A})$$
where $`A=\rho _1\rho _2`$. Furthermore, the bound is tight, and the orthogonal measurement $`𝒪`$ that projects a state on the eigenvectors of $`\rho _1\rho _2`$ achieves this bound.
Theorem 3 shows that the density matrix captures all the accessible information that a quantum state contains. If two different mixtures have the same density matrix (which is quite possible) then physically they are two different systems, but practically (and from a computational point of view) they are indistinguishable.
The quantity $`Trace(\sqrt{A^{}A})`$ is of independent interest. If we define $`A_t=Trace(\sqrt{A^{}A})`$ then $`||||_t`$ defines a norm, and has some additional properties such as $`AB_t=A_tB_t`$, $`A_t=1`$ for any density matrix $`A`$ and $`AB_t,BA_tA_tB_t`$. If $`\varphi _1,\varphi _2`$ are two pure states, and $`\rho _i`$ is the reduced density matrix of $`\varphi _i`$, then $`\rho _0\rho _1_t=2\sqrt{1|\varphi _1|\varphi _2|^2}`$. See for more details.
Locality. We now turn to the local view of a subsystem. Suppose we are in a mixed state $`\rho `$ over $`k+m`$ qubits, where Alice holds the first $`k`$ qubits $`A`$ and Bob holds the last $`m`$ qubits $`B`$. Assume that Alice applies a generalized measurement $`𝒪`$ on her qubits $`A`$. This induces a new density matrix $`\rho _B^𝒪`$ on $`B`$. E.g., if Alice and Bob were in the super position $`\varphi =\frac{1}{\sqrt{2}}(|00+|11)`$ over two qubits and Alice measured the second qubit according to the basis $`\{|0,|1\}`$, then Bob is with probability $`\frac{1}{2}`$ in the super position $`|0`$ and with probability $`\frac{1}{2}`$ in $`|1`$, hence $`\rho _B^𝒪=\left(\begin{array}{cc}\frac{1}{2}& 0\\ 0& \frac{1}{2}\end{array}\right)`$. A fundamental fact from physics, which can also be proven rigorously, tells us that in fact $`\rho _B^𝒪`$ does not depend on $`𝒪`$, but only on the original matrix $`\rho `$. We thus denote it by $`\rho |_B`$, and call it the density matrix $`\rho `$ reduced onto the subsystem $`B`$. Alternatively, we say that the rest of the system is traced out. The physical interpretation of the above result is that a player is guaranteed locality, i.e., a player Bob who holds a subsystem $`B`$ knows that the results he gets from measurements he applies on $`B`$ do not depend on the way the system outside $`B`$ evolves. It is also some kind of commitment. If Alice sends Bob $`k`$ qubits that have reduced density matrix $`\rho _B`$, then whatever Alice later does can not change this reduced density matrix.
Purification. A density matrix on a Hilbert space $`A`$ can always be viewed as a reduced density matrix of a pure state on a larger Hilbert space, a process which is called “purification”. A pure state $`|\varphi _{A,B}`$ is a purification of the density matrix $`\rho _A`$ if the reduced density matrix of $`|\varphi \varphi |_{A,B}`$ to the Hilbert space $`A`$ is $`\rho `$. The most straight forward way to purify a density matrix $`\rho =_iw_i|\varphi _i\varphi _i|`$ is by the state $`|\varphi =_i\sqrt{w_i}|i|\varphi _i`$.
Fidelity.
The fidelity is a way to measure distances between density matrices, which is an alternative to the trace metric. Given two density matrices $`\rho _0,\rho _1`$ on the same Hilbert space $`A`$ the fidelity is defined to be:
$$f(\rho _0,\rho _1)=sup|\varphi _0|\varphi _1|^2$$
(2)
where the supremum is taken over all purifications $`|\varphi _0`$ of $`\rho _0`$ and $`|\varphi _1`$ of $`\rho _1`$ to the same dimensional Hilbert space. We note here a few important properties which can easily be proven:
1. $`0f(\rho _0,\rho _1)1`$
2. $`f(\rho _0,\rho _1)=1\rho _0=\rho _1`$
3. For $`\rho _0`$ which is a pure state, i.e. $`\rho _0=|\varphi _0\varphi _0|`$, we have
$$f(\rho _0,\rho _1)=\varphi _0|\rho _1|\varphi _0.$$
Note that the fidelity increases as the distance between two density matrices decreases. It is also not too difficult to see that the supremum is always achieved, i.e. we can replace the supremum by a maximum; See for more details.
Entanglement. Suppose Alice holds a register $`A`$, Bob holds $`B`$, and the system is in a pure state $`\psi _{AB}`$. If we look at Bob’s system alone then we might see a mixed state, and as we said before, Alice can not change the reduced density matrix of Bob by local operations on her side. On the other hand Alice might gain different aspects of knowledge on the actual result that Bob gets.
###### Example 1
$`\psi _{AB}=\frac{1}{\sqrt{2}}(|00+|11)`$. If Alice measures in the $`\{|0,|1\}`$ basis, then Bob’s system is with probability half in the state $`|0`$, and with probability half in the state $`|1`$, and the register $`A`$ reflects the result Bob gets, i.e., Alice knows whether Bob gets a zero or a one. Now, $`\psi _{AB}`$ can also be represented as $`\frac{1}{\sqrt{2}}(|+,++|,)`$ where $`|+=\frac{1}{\sqrt{2}}(|0+|1)`$ and $`|=\frac{1}{\sqrt{2}}(|0|1)`$. Alice can measure the register $`A`$ in the $`\{|+,|\}`$ basis. Now Bob’s system is with probability $`\frac{1}{2}`$ in the state $`|+`$, and with probability half in the state $`|`$, and the register $`A`$ reflects the result Bob gets, i.e., Alice knows whether Bob gets $`|+`$ or $`|`$. Notice that Bob’s reduced density matrix is the same in both cases.
An important Theorem by Mayers and independently Lo and Chau states:
###### Theorem 4
Suppose the reduced density matrix of $`B`$ is the same in $`\varphi _{AB}`$ and $`\psi _{AB}`$. Then Alice can move from $`\varphi _{AB}`$ to $`\psi _{AB}`$ by applying a local transformation on her side.
I.e., even though Alice can not change Bob’s reduced density matrix, she can determine how to “open” the mixture, and do so in a way that gives her full knowledge of Bob’s result.
## 3 The Binding Property
In Protocol 1 Alice sends a qubit to Bob (we call it a “deposit” step) and later on she tells Bob how to “open” the qubit (the “reveal” step) which also determines the value that is supposed to be in the qubit. Such a protocol is worthless unless the deposit step is “binding” Alice to a pre-determined value. We first define the binding property in a general way. We then analyze how binding Protocol 1 is. Suppose we have a two step protocol:
: Alice prepares a super-position $`\psi _{AB}`$ with two quantum registers $`A`$ and $`B`$. Alice sends the second register $`B`$ to Bob.
: Alice and Bob communicate. Bob follows the protocol and Alice is arbitrary. If Alice wants to create a bias towards $`0`$ she uses one strategy, and if she wants a bias towards $`1`$ she uses a different strategy. Bob decides on a result $`r_B\{0,1,err\}`$.
Let us denote by $`p_0`$ the probability that Alice claims the result is $`0`$ in the zero strategy, by $`p_1`$ the probability that Alice claims the result is $`1`$ in the zero strategy, and by $`p_{err}`$ the probability that Bob decides the answer is $`r_B=err`$ when Alice uses the zero strategy. We similarly define $`q_0,q_1,q_{err}`$ for the one strategy.
###### Definition 5
($`(ϵ,\gamma )`$ binding) A protocol is $`(ϵ,\gamma )`$ binding, if whenever Bob is honest, for any strategy Alice uses, if $`p_{err},q_{err}ϵ`$ then $`|p_0q_0|,|p_1q_1|\gamma `$.
### 3.1 Protocol 1 is quadratically binding
###### Theorem 5
Protocol 1 is $`(ϵ,\gamma =\frac{2\sqrt{ϵ}}{\mathrm{cos}(2\theta )})`$ binding.
###### Proof 3.6.
(of Theorem 5). At deposit time Alice sends Bob one qubit $`B`$, which might be entangled with the qubits $`A`$ that Alice holds. Let us denote the reduced density matrix of $`B`$ by $`\rho `$. At revealing time, Alice may choose whether she wants to bias the result towards $`0`$, in which case she applies the generalized measurement $`M_0`$, or towards $`1`$ in which case she applies $`M_1`$. The measurements $`M_0`$ and $`M_1`$ do not change the reduced density matrix $`\rho `$ of Bob, but rather give different ways to realize $`\rho `$ as a mixture of pure-states, and give Alice information about the value that Bob actually gets to see in this mixture.
Now, we even go further and give Alice complete freedom to choose the way she realizes the reduced density matrix $`\rho `$ of Bob as a mixture, and we give her the knowledge of Bob’s value for free. Let us say that when Alice applies $`M_0`$, the reduced density matrix $`\rho `$ is realized as the mixture $`\{p_i,\varphi _i\}`$, and when Alice applies $`M_1`$ the reduced density matrix $`\rho `$ is realized as the mixture $`\{p_i^{},\varphi _i^{}\}`$.
Now, let us focus on the zero strategy. Say Alice realizes $`\rho `$ as $`\{p_i,\varphi _i\}`$. When the $`i`$’th event happens, Alice’s strategy tells her to send some two qubits $`q_b,q_x`$ to Bob, that are supposed to hold classical $`0,1`$ values for $`b`$ and $`x`$. Bob then measures $`q_b`$ and $`q_x`$ in the $`\{|0,|1\}`$ basis. Now, if one of $`q_b,q_x`$ is not a classical bit, then Alice can measure it herself in the $`\{|0,|1\}`$ basis, and get a mixture over classical bits. Furthermore, we can push all the probabilistic decisions into the mixture $`\{p_i,\varphi _i\}`$. Thus, w.l.o.g, we can assume Alice’s answers $`q_b`$ and $`q_x`$ are classical bits that are determined by the event $`i`$. Let us denote by $`u_i`$ the vector $`\varphi _{b_i,x_i}`$ where $`b_i,x_i`$ are Alice’s answers when event $`i`$ occurs. W.l.o.g we may assume $`u_i\{\varphi _{b,x}\}`$, otherwise we know Bob immediately rejects.
The probability Bob discovers that Alice is cheating is then $`1|\varphi _i|u_i|^2`$ and the overall probability Bob detects Alice is cheating is
$`p_{err}`$ $`=`$ $`\mathrm{\Sigma }_ip_i(1|\varphi _i|u_i|^2)`$
Let us define the density matrix $`\rho _0=\mathrm{\Sigma }_ip_i|u_iu_i|`$.
###### Claim 1.
$`\rho \rho _0_t2\sqrt{p_{err}}`$.
###### Proof 3.7.
$`|\varphi _i\varphi _i||u_iu_i|_t=2\sqrt{1|\varphi _i|u_i|^2}`$. Therefore
$`\rho \rho _0_t`$ $`=`$ $`\mathrm{\Sigma }_ip_i|\varphi _i\varphi _i|\mathrm{\Sigma }_ip_i|u_iu_i|_t`$
$``$ $`\mathrm{\Sigma }_ip_i|\varphi _i\varphi _i||u_iu_i|_t`$
$`=`$ $`2\mathrm{\Sigma }_ip_i\sqrt{1|\varphi _i|u_i|^2}`$
Now, by Cauchy-Schwartz inequality,
$`\mathrm{\Sigma }_ip_i\sqrt{1|\varphi _i|u_i|^2}`$ $`=`$ $`\mathrm{\Sigma }_i\sqrt{p_i}\sqrt{p_i(1|\varphi _i|u_i|^2)}`$
$``$ $`\sqrt{\mathrm{\Sigma }_ip_i}\sqrt{\mathrm{\Sigma }_ip_i(1|\varphi _i|u_i|^2)}`$
$`=`$ $`\sqrt{p_{err}}`$
and the claim follows.
Similarly, if Alice tries to bias the result towards $`1`$, $`B`$ ends up in the mixture $`\{p_i^{},\varphi _i^{}\}`$, and when $`\varphi _i^{}`$ occurs Alice sends $`b^{},x^{}`$ to Bob that correspond to a vector $`u_i^{}\{\varphi _{b,x}\}`$. We define $`\rho _1`$ to be the reduced density matrix $`\rho _1=\mathrm{\Sigma }_ip_i^{}|u_i^{}u_i^{}|`$. As before, $`\rho \rho _1_t2\sqrt{q_{err}}`$. Hence, $`\rho _0\rho _1_t2(\sqrt{p_{err}}+\sqrt{q_{err}})`$.
To conclude the proof, we establish the following claim:
###### Claim 2.
Let $`\rho _0`$ and $`\rho _1`$ be density matrices corresponding to mixtures over $`\{\varphi _{b,x}\}`$. Let $`p_0`$ be the probability of $`\varphi _{0,0}`$ or $`\varphi _{0,1}`$ in the first mixture, and $`p_1=1p_0`$ be the probability of $`\varphi _{1,0}`$ or $`\varphi _{1,1}`$. Similarly let $`q_0`$ and $`q_1`$ be the corresponding quantities for the second mixture. Then $`\rho _0\rho _1_t2|p_0q_0|\mathrm{cos}2\theta `$.
###### Proof 3.8.
We show that we can distinguish the mixtures with probability at least $`|p_0q_0|\mathrm{cos}2\theta `$ when we measure them according to the basis $`\{|0,|1\}`$. If we do the measurement on a qubit whose state is the reduced density matrix $`\rho _0`$ we get the $`|0`$ answer with probability $`p_0\mathrm{cos}^2(\theta )+p_1\mathrm{sin}^2(\theta )`$, while if we do the measurement on a qubit whose state is the reduced density matrix $`\rho _1`$ we get the $`|0`$ answer with probability $`q_0\mathrm{cos}^2(\theta )+q_1\mathrm{sin}^2(\theta )`$. The difference is $`|p_0\mathrm{cos}^2(\theta )+p_1\mathrm{sin}^2(\theta )(q_0\mathrm{cos}^2(\theta )+q_1\mathrm{sin}^2(\theta ))|=|p_0q_0|(\mathrm{cos}^2(\theta )\mathrm{sin}^2(\theta ))`$, where we used $`p_1q_1=(1p_0)(1q_0)=q_0p_0`$. Altogether we get $`\rho _0\rho _1_t2|p_0q_0|(\mathrm{cos}^2(\theta )\mathrm{sin}^2(\theta ))`$ as desired.
Putting it together:
$$2\mathrm{cos}(2\theta )|p_0q_0|\rho _1\rho _1_t2(\sqrt{p_{err}}+\sqrt{q_{err}})4\sqrt{ϵ}$$
I.e., $`|p_0q_0|\frac{2\sqrt{ϵ}}{\mathrm{cos}(2\theta )}`$.
### 3.2 A Quadratic Strategy for Alice
We now show that Alice has a quadratic strategy for Protocol 1, and thus Theorem 5 is essentially tight. In fact, we show the quadratic bound for a more general family of protocols. Let $`\rho _0,\rho _1`$ be two density matrices of the same dimension, $`\rho _0`$ can be realized as the mixture $`\{p_i^0,|\alpha _i^0\}`$, and $`\rho _1`$ as $`\{p_i^1,|\alpha _i^1\}`$. To encode $`b`$, honest Alice picks $`|\alpha _i^b`$ with probability $`p_i`$ and sends it to Bob. At revealing time Alice sends $`b`$ and $`i`$ to Bob, and Bob tests whether Alice is cheating by projecting his state on $`|\alpha _i^b`$.
###### Theorem 3.9.
Let $`f`$ be the fidelity $`f(\rho _0,\rho _1)`$. For any $`0\alpha \pi /4`$ there exists a strategy for Alice with advantage $`\sqrt{f}sin(2\alpha )/2`$ and probability of detection at most $`\frac{(1f)sin^2(\alpha )}{2}`$.
On first reading of the next proof the reader might want to check the proof in the simpler case where $`\rho _0`$ and $`\rho _1`$ represent pure states, i.e., $`\rho _b=|\psi _b\psi _b|`$.
###### Proof 3.10.
We first represent the strategy of a honest Alice in quantum language. Consider two maximally parallel purifications $`|\psi _0`$ and $`|\psi _1`$ of $`\rho _0`$ and $`\rho _1`$, where $`\rho _0`$ and $`\rho _1`$ are density matrices of the register $`B`$, and the purifications are states on a larger Hilbert space $`AB`$. By , $`|\psi _0|\psi _1|^2=f(\rho _0,\rho _1)`$. At preparation time, Alice prepares the state
$`|\beta `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(|0,\psi _0+|1,\psi _1)`$
on $`AB`$ and one extra qubit $`C`$. Alice then sends the register $`B`$ to Bob. At revealing time, Alice measures the qubit $`C`$ in the $`|0,|1`$ basis, to get a bit $`b`$. The state of registers $`A,B`$ is now $`|\psi _b`$. Alice then applies a unitary transformation $`U_b`$ on register $`A`$, which rotates her state $`|\psi _b`$ to the state
$`|\psi _b^{}`$ $`=`$ $`{\displaystyle \underset{j}{}}\sqrt{p_j^b}|j_A|\alpha _j^b_B`$
This is possible by Theorem 4. After applying $`U_b`$, Alice measures register $`A`$ in the computational basis and sends Bob the bit $`b`$ and the outcome of the second measurement, $`j`$. This strategy is similar to the honest strategy, except for that Alice does not know what bit and state is sent until revealing time.
We can also assume w.l.o.g. that the maximally parallel purifications satisfy that $`\psi _0|\psi _1`$ is real and positive. This can be assumed since otherwise we could multiply $`|\psi _0`$ by an overall phase without changing the reduced density matrix and the absolute value of the inner product.
To cheat, Alice creates the encoding $`|\beta _{CAB}`$ and sends register $`B`$ to Bob. Alice’s one strategy is also as described above. The zero strategy, on the other hand, is a slight modification of the honest strategy. At revealing time, Alice measures the control qubit $`C`$ in the $`\{|\varphi _\alpha ,|\varphi _\alpha ^{}\}`$ basis, where
$`|\varphi _\alpha `$ $`=`$ $`c|0+s|1,`$ (3)
$`|\varphi _\alpha ^{}`$ $`=`$ $`s|0+c|1,`$
and $`s=\mathrm{sin}(\alpha )`$, $`c=\mathrm{cos}(\alpha )`$. If the outcome is a projection on $`|\varphi _\alpha `$ Alice sends $`b=0`$ and proceeds according to the $`b=0`$ honest protocol, i.e. applies $`U_0`$ to register $`A`$, measures in the computational basis and sends the result to Bob. If the outcome is a projection on $`|\varphi _\alpha ^{}`$, Alice proceeds according to the $`b=1`$ honest protocol. Let us now compute Alice’s advantage and Alice’s probability of getting caught cheating.
We can express $`|\beta `$ as:
$`|\beta `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(c|\varphi _\alpha ,\psi _0s|\varphi _\alpha ^{},\psi _0)+`$
$`{\displaystyle \frac{1}{\sqrt{2}}}(s|\varphi _\alpha ,\psi _1+c|\varphi _\alpha ^{},\psi _1).`$
Hence, the probability Alice sends $`b=0`$ in the zero strategy is $`\frac{1}{2}|c\psi _0+s\psi _1|^2=\frac{1}{2}(c^2+s^2+2cs\psi _0|\psi _1)=\frac{1}{2}(1+2cs\sqrt{f})`$. We conclude:
###### Claim 3.
Alice’s advantage is $`\frac{\sqrt{f}\mathrm{sin}(2\alpha )}{2}`$.
We now prove that the detection probability is at most $`(1f)s^2`$. The state of $`AB`$ conditioned that the first measurement yields $`|\varphi _\alpha `$ can be written as $`\frac{1}{\sqrt{Pr(b=0)}}\frac{1}{\sqrt{2}}(c|\psi _0+s|\psi _1)`$ where $`Pr(b=0)`$ is the probability Alice sends $`b=0`$ in the zero strategy. The above state can be written as
$$\frac{1}{\sqrt{Pr(b=0)}}\frac{1}{\sqrt{2}}(c+\sqrt{f}s)|\psi _0+\sqrt{1f}s|\psi _0^{}$$
The rest of the protocol involves Alice’s rotation of the state by $`U_0`$, then Alice’s measurement of the register $`A`$ and Bob’s measurement of the register $`B`$. The entire process can be treated as a generalized measurement on this state, where this measurement is a projection onto one of two subspaces, the “cheating Alice” and the “Honest Alice” subspaces. We know that $`|\psi _0`$ lies entirely in the honest Alice subspace, and thus the probability that Alice is caught, conditioned that $`C`$ was projected on $`\varphi _\alpha `$, is at most $`\frac{1}{Pr(b=0)}\frac{1}{2}(1f)s^2`$.
In the same way, when we condition on a projection on $`\varphi _\alpha ^{}`$, Alice’s state can be written as $`\frac{1}{\sqrt{Pr(b=1)}}\frac{1}{\sqrt{2}}((c\sqrt{f}s)|\psi _1\sqrt{1f}s|\psi _1^{})`$. which gives a probability of detection which is at most $`\frac{1}{Pr(b=1)}\frac{1}{2}(1f)s^2`$. Adding the conditional probabilities together we get that the detection probability is at most $`\frac{(1f)s^2}{2}`$.
## 4 The Sealing Property
###### Definition 4.11.
($`(ϵ,p)`$ sealing) A bit escrow protocol is $`(ϵ,p)`$ sealing, if whenever Alice is honest and deposits a bit $`b`$ s.t. , for any strategy Bob uses and a value $`c`$ Bob learns, it holds that either
* $`\mathrm{Pr}_{b_R\{0,1\},protocol}(c=b)\frac{1}{2}+ϵ`$, or
* $`\mathrm{Pr}_{b_R\{0,1\},protocol}(r_A=err)p`$
The probability is taken over $`b`$ taken uniformly from $`\{0,1\}`$ and the protocol.
We show here that protocol 1 is quadratically sealing. This means that whatever Bob does, he will always be detected cheating with probability which is at least the square of his advantage. Later, we show that this is tight.
### 4.1 Protocol 1 is Quadratically Sealing
###### Theorem 4.12.
Protocol 1 is $`(ϵ=O(\frac{\sqrt{p}}{\mathrm{sin}(2\theta )}),p)`$ sealing.
###### Proof 4.13.
We first describe a general scenario. Alice is honest and sends $`|\varphi _{b,x}_A`$ to Bob. Bob has an ancilla $`|0_C`$. Bob applies some unitary transformation $`U`$ acting on the registers $`A`$ and $`C`$. Let us denote
$`|\alpha _{b,x}`$ $`=`$ $`U(|\varphi _{b,x},0_{AC})`$
Bob then sends register $`A`$ to Alice, and keeps register $`C`$ to himself. We want to show that if $`C`$ contains much information about $`b`$ then Alice detects Bob cheating with a good probability.
We can express $`\alpha _{b,x}`$ as a superposition,
$`|\alpha _{b,x}`$ $`=`$ $`|\varphi _{b,x},w_{b,x}+|\varphi _{\neg b,x},w_{b,x}^{}`$ (4)
where we have used the basis $`|\varphi _{b,x}`$, $`|\varphi _{\neg b,x}`$, for $`A`$. In this representation, the probability $`p`$ Bob is caught cheating is:
$`p`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{b,x}{}}w_{b,x}^{}^2`$ (5)
which in particular implies that $`w_{b,x}^{}2\sqrt{p}`$.
We now want to express Bob’s advantage. Let $`\rho _0`$ ($`\rho _1`$) be the reduced density matrix of the register $`B`$ conditioned on the event that $`b=0`$ ($`b=1`$). Then,
$$\rho _b=\underset{x}{}Pr(x)(|w_{b,x}w_{b,x}|+|w_{b,x}^{}w_{b,x}^{}|)$$
(6)
Bob’s advantage is at most the trace distance between $`\rho _0`$ and $`\rho _1`$, and we want to bound it from above. Triangle inequality gives: $`\rho _0\rho _1_t\frac{1}{2}(|w_{0,0}w_{0,0}||w_{1,1}w_{1,1}|_t+|w_{0,1}w_{0,1}||w_{1,0}w_{1,0}|_t+_{b,x}||w_{b,x}^{}w_{b,x}^{}|||_t)`$.
As the trace norm of two pure states $`a`$ and $`b`$ is $`2\sqrt{1|a|b|^2}`$, and using Equation 5, we get:
$`\rho _0\rho _1_t`$ $``$ $`\sqrt{1|w_{0,0}|w_{1,1}|^2}+`$
$`\sqrt{1|w_{0,1}|w_{1,0}|^2}+2p`$
We now claim;
###### Lemma 4.14.
$`|w_{0,0}|w_{1,1}|,|w_{0,1}|w_{1,0}|1O(ctg^2(2\theta )+4)p`$.
Thus, altogether, $`\rho _0\rho _1_tO(ctg(2\theta )\sqrt{p})`$ which completes the proof.
We now turn to the proof of Lemma 4.14.
###### Proof 4.15.
(of Lemma).
We will prove that all the unprimed $`w`$ vectors lie in one bunch of small width, using the unitarity of $`U`$. The unitarity of $`U`$ implies that $`\varphi _{b,x}|\varphi _{b^{},x^{}}=\alpha _{b,x}|\alpha _{b^{},x^{}}`$. We can express $`\alpha _{b,x}`$ as in Equation 4. We get:
$`\varphi _{b,x}|\varphi _{b^{},x^{}}`$ $`=`$ $`\varphi _{b,x}|\varphi _{b^{},x^{}}w_{b,x}|w_{b^{},x^{}}+`$
$`\varphi _{b,x}|\varphi _{\neg b^{},x^{}}w_{b,x}|w_{b^{},x^{}}^{}+`$
$`\varphi _{\neg b,x}|\varphi _{b^{},x^{}}w_{b,x}^{}|w_{b^{},x^{}}+`$
$`\varphi _{\neg b,x}|\varphi _{\neg b^{},x^{}}w_{b,x}^{}|w_{b^{},x^{}}^{}`$
Substituting the values $`b,x,b^{},x^{}`$ for actual values, and noticing that $`|w_{b,x}^{}|w_{b^{},x^{}}^{}|4p`$, we in particular get the following equations:
$`w_{b,x}|w_{b,x}`$ $`=_{4p}`$ $`1`$ (7)
$`w_{0,0}|w_{1,0}^{}+w_{0,0}^{}|w_{1,0}`$ $`=`$ $`0`$ (8)
$`w_{0,1}|w_{1,1}^{}+w_{0,1}^{}|w_{1,1}`$ $`=`$ $`0`$ (9)
$`w_{1,0}|w_{1,1}^{}w_{1,0}^{}|w_{1,1}=_{4cp/s}{\displaystyle \frac{c}{s}}(1w_{1,0}|w_{1,1})`$ (10)
$`w_{0,1}|w_{1,0}^{}+w_{0,1}^{}|w_{1,0}=_{4sp/c}{\displaystyle \frac{s}{c}}(1w_{0,1}|w_{1,0})`$ (11)
$`w_{0,0}|w_{1,1}^{}w_{0,0}^{}|w_{1,1}=_{4sp/c}{\displaystyle \frac{s}{c}}(1w_{0,0}|w_{1,1})`$ (12)
$`w_{0,0}^{}|w_{0,1}w_{0,0}|w_{0,1}^{}=_{4cp/s}{\displaystyle \frac{c}{s}}(1w_{0,0}|w_{0,1})`$ (13)
where $`c=\mathrm{cos}(2\theta )`$, $`s=\mathrm{sin}(2\theta )`$ and we write $`x=_qy`$ if $`|xy|q`$. A partial result can already be derived from what we have so far. By equation 5, we note that the length of the primed $`w`$ vectors is at most $`2\sqrt{p}`$. Inserting this to equations 11 and 12, we get that $`|w_{0,0}|w_{1,1}|`$ and similarly $`|w_{0,1}|w_{1,0}|`$ are close to $`1`$ up to terms of order $`\sqrt{p}`$. This is a weaker than the result which we want to achieve in lemma 4.14, which is closeness to $`1`$ up to order $`p`$ terms. If we stop here, the closeness of the unprimed $`w`$ vectors up to order $`\sqrt{p}`$ implies that Bob’s information is at most of the order of $`{}_{}{}^{4}\sqrt{p}`$. Note, however, that so far all we have used is unitarity, and we have not used the particular properties of the set of vectors we use in the protocol. In the rest of the proof, we will use the symmetry in protocol 1 to improve on this partial result, and to show that Bob’s information is at most of the order of $`\sqrt{p}`$. Basically, the symmetry which we will use is the fact that the vectors in the protocol can be paired into orthogonal vectors.
We proceed as follows. The idea is to express equations 10-13 as inequalities involving only the distances between two $`w`$ vectors, $`w_{b,x}w_{b^{},x^{}}`$ and then to solve the set of the four inequalities to give an upper bound on the pairwise distances. This will imply a bound on the inner products, $`w_{b,x}|w_{b^{},x^{}}`$, by the following connection:
###### Claim 5.
$`1Re(w_{b,x}|w_{b^{},x^{}})\frac{w_{b,x}w_{b^{},x^{}}^2}{2}`$.
where $`Re(z)`$ denotes the real part of the complex number $`z`$.
###### Proof 4.16.
$`w_{b,x}w_{b^{},x^{}}^2=w_{b,x}w_{b^{},x^{}}|w_{b,x}w_{b^{},x^{}}22Re(w_{b,x}|w_{b^{},x^{}})`$.
We denote:
$`a`$ $`=`$ $`w_{0,0}w_{0,1}`$
$`b`$ $`=`$ $`w_{0,0}w_{1,1}`$
$`c`$ $`=`$ $`w_{0,1}w_{1,0}`$
$`d`$ $`=`$ $`w_{1,0}w_{1,1}`$
Let $`LHS`$ ($`RHS`$) be the sum of the left (right) hand side of the last four equations.
###### Claim 6.
$`Re(RHS)\frac{c}{2s}(a^2+d^2)+\frac{s}{2c}(b^2+c^2)`$.
###### Proof 4.17.
$`Re(RHS)`$ $`=`$ $`{\displaystyle \frac{c}{s}}(2Re(w_{0,0}|w_{0,1})Re(w_{1,0}|w_{1,1}))`$
$`+`$ $`{\displaystyle \frac{s}{c}}(2Re(w_{0,0}|w_{1,1})Re(w_{0,1}|w_{1,0}))`$
and now we can apply claim 5.
Expressing the left hand side of the equations in terms of $`a,b,c`$ and $`d`$ might look a bit more complicated, and this is where we invoke the symmetric properties of the protocol, namely equations 8 and 9.
###### Claim 7.
$`Re(LHS)4\sqrt{p}(a+b+c+d)`$
###### Proof 4.18.
We first look at the LHS of Equation 12 \+ Equation 13. By adding $`w_{0,1}|w_{1,1}^{}+w_{0,1}^{}|w_{1,1}=0`$ (due to Equation 9) and by using the fact that $`Re(\alpha |\beta )=Re(\beta |\alpha )`$ we get that the LHS of these two equations contributes $`Re(w_{0,0}^{}|w_{0,1}w_{1,1})+Re(w_{0,1}^{}|w_{1,1}w_{0,0})+Re(w_{1,1}^{}|w_{0,1}w_{0,0})2\sqrt{p}(c+d)+2\sqrt{p}b+2\sqrt{p}a`$.
Similarly, the LHS of Equation 10 +Equation 11 is $`Re(w_{1,0}^{}|w_{0,1}w_{1,1})+Re(w_{0,1}^{}|w_{1,0}w_{1,1})+Re(w_{1,1}^{}|w_{1,0}w_{0,1})2\sqrt{p}(a+b+d+c)`$.
Altogether, $`Re(LHS)4\sqrt{p}(a+b+c+d)`$.
Combining Claims 7 and 6 with our knowledge that $`Re(RHS)Re(LHS)+\frac{8cp}{s}+\frac{8sp}{c}`$ we get:
$`{\displaystyle \frac{c}{2s}}(a^2+d^2)+{\displaystyle \frac{s}{2c}}(b^2+c^2)`$
$`4\sqrt{p}(a+b+c+d)+{\displaystyle \frac{8cp}{s}}+{\displaystyle \frac{8sp}{c}}`$
We want to show that $`a,b,c,d`$ are all of the order of $`\sqrt{p}`$. Define $`\mathrm{\Delta }=a+b+c+d`$. For $`0\theta \frac{\pi }{8}`$, $`ctg(2\theta )tg(2\theta )`$. Since all terms in the left hand side are positive, we have for each of $`a,b,c,d`$ an upper bound in terms of $`\mathrm{\Delta }`$:
$`a^2,b^2,c^2,d^2`$ $``$ $`{\displaystyle \frac{8\mathrm{\Delta }\sqrt{p}}{s/c}}+16p(1+({\displaystyle \frac{c}{s}})^2)`$
Thus, $`\mathrm{\Delta }=a+b+c+d4\sqrt{\frac{8\mathrm{\Delta }\sqrt{p}c}{s}+\frac{16p}{s^2}}`$.
Solving the quadratic equation
$`\mathrm{\Delta }^2{\displaystyle \frac{2^7\sqrt{p}c}{s}}\mathrm{\Delta }{\displaystyle \frac{2^8p}{s^2}}`$ $``$ $`0`$
for $`\mathrm{\Delta }`$ we get
$`\mathrm{\Delta }`$ $``$ $`132\sqrt{p}ctg(2\theta )`$
Finally,
$`|w_{0,0}|w_{1,1}|`$ $``$ $`|Re(w_{0,0}|w_{1,1})|`$
$`=`$ $`{\displaystyle \frac{w_{0,0}^2+w_{1,1}^2b^2}{2}}`$
$``$ $`{\displaystyle \frac{2b^28p}{2}}`$
$``$ $`1(2^{15}ctg^2(2\theta )+4)p`$
where the third inequality is true due to equation 7. Similarly, we have the same lower bound for $`|w_{0,1}|w_{1,0}|`$, which implies lemma 4.14.
Thus, our bit escrow protocol gives quadratic sealing.
###### Remark 4.19.
Protocol 1 is sealing even if we modify it a little bit, as follows: at revealing time Alice first reveals $`b`$ and then Bob returns the qubit $`q`$. In other words, if Bob has learned $`ϵ`$ information about $`b`$ after the deposit stage, then even if later on he gets to know $`b`$, he cannot avoid being detected with probability $`\mathrm{\Omega }(ϵ^2)`$. To see this, we use linearity. If Bob has a strategy which gives him detection probability $`p`$ in the modified protocol, then w.l.o.g. his strategy is to apply the identity if $`b=0`$ and some unitary operation $`U`$ if $`b=1`$. However, since the $`b=1,x=0`$ and $`b=1,x=1`$ cases are linear combinations of the $`b=0,x=0`$ and $`b=0,x=1`$ cases, one can show that if Bob’s probability for detection is $`p`$ in the $`b=0`$ case, then it is also $`O(p)`$ in the $`b=1`$ case, and therefore Bob does not have to apply $`U`$ in the first place. This means that if he has a cheating strategy for the modified protocol, then he also has a cheating strategy with about the same parameters for protocol 1, and so by Theorem 4.12 the modified protocol is also quadratically secure.
###### Remark 4.20.
One might suspect that this quadratic gap will always be the case for any reasonable set of vectors for Alice. This is not correct. If Alice only uses $`\varphi _{0,1}`$ and $`\varphi _{1,0}`$, then Bob has a strategy which gives him $`{}_{}{}^{4}\sqrt{p}`$ advantage. We will not elaborate on this in this paper.
### 4.2 A Quadratic Strategy for Bob
###### Theorem 4.21.
Let $`\rho _0,\rho _1`$ be two density matrices of the same dimension, such that $`\rho _0\rho _1_t=t`$. Consider the following protocol. Alice tosses a random bit $`b`$. She chooses a pure state from the mixture $`\rho _b`$, and sends it to Bob. Then Bob returns to Alice the state, and Alice projects it on the original state to test whether Bob has manipulated it. We claim that for any $`1p0`$, there is a strategy for Bob such that he learns $`b`$ with advantage $`t\sqrt{p}`$, and his probability of detection is at most $`\frac{1}{2}(1\sqrt{1p})`$, which is $`\mathrm{\Theta }(p)`$ for small $`p`$.
proof: Alice prepares an encoding $`\psi _b`$ of $`b\{0,1\}`$ in register $`B`$, and sends register $`B`$ to Bob. Let $`\rho _b`$ be the reduced density matrix of $`\psi _b`$ to register $`B`$. We denote $`t=\rho _0\rho _1_t`$. By Theorem 3 we know that if Bob is interested in learning information about $`b`$, and is not concerned with being detected cheating, the best he can do is a measurement according to the eigenvalue basis of $`\rho _0\rho _1`$. Given, any $`0p1`$ we modify this strategy to a strategy where the detection probability is at most $`p`$, and yet, Bob gets much information.
Let us consider more precisely Bob’s best strategy for learning $`b`$ if he is not concerned with being caught. Let $`\{e_1,\mathrm{},e_K\}`$ be the eigenvector basis of $`\rho _0\rho _1`$. Let $`V^+`$ ($`V^{}`$) be the set of eigenvectors $`e`$ with non-negative (negative) eigenvalues. The measurement $`M`$ is defined by the Hermitian matrix for which $`V^+`$ is an eigenspace of eigenvalue $`0`$ and $`V^{}`$ is an eigenspace of eigenvalue $`1`$. By Theorem 3
$`|Trace(\rho _0M)Trace(\rho _1M)|`$ $`=`$ $`{\displaystyle \frac{t}{2}}`$ (14)
To apply a weak form of the measurement $`M`$, Bob takes a one qubit ancilla $`C`$. He applies a unitary transformation $`U`$ on the received message and the ancilla, as follows:
$`U|e,0`$ $`=`$ $`\{\begin{array}{cc}|e,0\hfill & \text{If }eV^+\hfill \\ |e|v\hfill & \text{If }eV^{}\hfill \end{array}`$
where $`|v=\sqrt{1p}|0+\sqrt{p}|1`$ and $`U`$ is completed to a unitary transformation. After applying $`U`$ Bob returns register $`B`$ to Alice, and keeps the ancilla $`C`$ for himself. Notice that the special case $`p=1`$ is equivalent to the measurement $`M`$.
###### Lemma 4.22.
$`||U\rho _0|_CU\rho _1|_C||_t=t\sqrt{p}`$.
###### Proof 4.23.
We will show
###### Claim 8.
$`U\rho _0|_C`$ $`=`$ $`Trace(\rho _0M)|00|+(1Trace(\rho _0M))|vv|)`$
$`U\rho _1|_C`$ $`=`$ $`Trace(\rho _1M)|00|+(1Trace(\rho _1M))|vv|)`$
Thus, $`U\rho _0|_CU\rho _1|_C=(Trace(\rho _0M)Trace(\rho _1M))(|00||vv|)=\pm \frac{t}{2}(|00||vv|)`$, where the last equality is due to Equation 14. Since, $`|00||vv|_t=2\sqrt{10|v|^2}=2\sqrt{p}`$ we get $`||U\rho _0|_CU\rho _1|_C||_t=t\sqrt{p}`$ as desired.
We now prove Claim 8.
###### Proof 4.24.
(of Claim 8). We express $`\rho _0=_jw_j|\alpha _j\alpha _j|`$, where $`\alpha _j`$ is a pure state.We further express each $`\alpha _j`$ in the eigenbasis $`\{e_i\}`$:
$`|\alpha _j`$ $`=`$ $`{\displaystyle \underset{i+}{}}a_{ij}^+|e_i^++{\displaystyle \underset{i}{}}a_{ij}^{}|e_i^{}`$
Applying $`U`$, this state is taken to:
$`U|\alpha _j,0`$ $`=`$ $`{\displaystyle \underset{i+}{}}a_{ij}^+|e_i^+|0+{\displaystyle \underset{i}{}}a_{ij}^{}|e_i^{}|v`$
The reduced density matrix to the register $`C`$, in case of event $`|\alpha _j`$ is:
$`{\displaystyle \underset{i+}{}}|a_{ij}^+|^2|00|+{\displaystyle \underset{i}{}}|a_{ij}^{}|^2|vv|`$
and altogether, $`U\rho _0|_C=_jw_j(_{i+}|a_{ij}^+|^2)|00|+`$
$`_jw_j(_i|a_{ij}^{}|^2)|vv|`$. To complete the proof we just notice that $`_jw_j(_{i+}|a_{ij}^+|^2)=Trace(\rho _0M)`$. The proof for $`U\rho _0|_C`$ is similar.
We now analyze the error detection probability.
###### Lemma 4.25.
$`Prob(err)\frac{1}{2}(1\sqrt{1p})`$
###### Proof 4.26.
Say Alice sent Bob the state $`|w`$. We can express it as $`|w=a|w^++b|w^{}`$ where $`|w^+Span(V^+)`$ and $`|w^{}Span(V^{})`$. Bob applies $`U`$ on $`w`$ and gets
$`U|w`$ $`=`$ $`a|w^+,0+b|w^{},v`$
$`=`$ $`a|w^+,0+\sqrt{1p}b|w^{},0+\sqrt{p}b|w^{},1`$
Therefore, if we measure the last qubit, then with probability $`pb^2`$ we end up in $`|w^{}`$ and with probability $`1pb^2`$ we end up in $`a|w^++\sqrt{1p}b|w^{}`$ normalized. Thus the density matrix of $`U|w`$ after tracing out the last qubit is:
$`\rho `$ $`=`$ $`\left(\begin{array}{cc}|a|^2& b\overline{a}\sqrt{1p}\\ \overline{b}a\sqrt{1p}& |b|^2\end{array}\right)`$
To find out the probability for Alice not to detect Bob cheating, we calculate $`w|\rho |w`$. We get:
$`Pr(\neg Err)`$ $`=`$ $`|a|^4+2|ab|^2\sqrt{1p}+|b|^4`$
$`=`$ $`12|ab|^2(1\sqrt{1p})`$
The probability of Alice detecting an error is thus $`2|ab|^2(1\sqrt{1p})\frac{1}{2}(1\sqrt{1p})`$.
###### Remark 4.27.
The average of $`|ab|`$ can tend to $`0.5`$, even when $`t`$ tends to $`0`$. This can be seen by taking $`\rho _0`$ to be composed of two states which are the basis states $`|0`$ and $`|1`$ rotated by $`\theta `$ towards each other, whereas $`\rho _1`$ is the mixture of the basis states rotated by $`\theta `$ outwards. As $`\theta `$ tends to $`0`$, $`t`$ tends to $`0`$, but $`|ab|`$ tend to $`0.5`$.
## 5 Proof of Theorem 2
We show that no cheater can control the game.
:
Suppose Alice is honest and Bob is arbitrary. Let us look at the mixture that Alice generates at the first step of Protocol 2. Let $`\rho _{b=0}`$ be the density matrix in the case $`b=0`$, and $`\rho _{b=1}`$ in the case $`b=1`$. Then $`\rho _{b=0}\rho _{b=1}_t=2\mathrm{cos}(2\theta )`$. It follows from Theorem 3 that whatever Bob does, the probability that $`b^{}=b`$ and Bob wins is at most $`\mathrm{Pr}(b^{}=b)\frac{1}{2}+\frac{\mathrm{cos}(2\theta )}{2}=\mathrm{cos}^2(\theta )`$ which for $`\theta =\frac{\pi }{8}`$ is at most $`0.86`$.
:
Now, suppose Bob is honest and Alice is arbitrary. , which is at most $`\frac{p_0+q_1}{2}`$, whereas the probability that Alice loses is at least $`\frac{p_1+q_0}{2}`$. The difference $`|x(1x)|`$ is at most $`\frac{p_0q_0+q_1p_1}{2}\frac{|p_0q_0|+|p_1q_1|}{2}=|p_0q_0|`$, i.e., $`x\frac{1+|p_0q_0|}{2}`$.
Also, $`\frac{p_{err}+q_{err}}{2}1x`$, as whenever Alice is caught cheating she loses. This implies that $`\sqrt{p_{err}}+\sqrt{q_{err}}2\sqrt{1x}`$ as the maximum is obtained when $`p_{err}=q_{err}=1x`$.
Finally, from the proof of Theorem 5 we have $`|p_0q_0|\frac{\sqrt{p_{err}}+\sqrt{q_{err}}}{cos(2\theta )}`$. Putting it all together we get:
$`x`$ $``$ $`{\displaystyle \frac{1+|p_0q_0|}{2}}`$
$``$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{\sqrt{p_{err}}+\sqrt{q_{err}}}{2cos(2\theta )}}`$
$``$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{\sqrt{1x}}{\mathrm{cos}(2\theta )}}`$
For $`\theta =\frac{\pi }{8}`$ we get the quadratic equation $`4x^2+4x70`$. Solving it we get $`x\frac{\sqrt{8}1}{2}0.9143`$.
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# ISO Spectroscopy of the Young Bipolar Nebulae S106 IR and Cep A East Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) and with the participation of ISAS and NASA.
## 1 Introduction
The infrared emission-line spectrum of a Young Stellar Object (YSO) is dominated by the interaction of the central star with the remnants of the clouds from which it formed. The intense UV radiation generated by accretion as well as by the central star itself causes dissociation of molecular material close to the YSO and ionizes much of the atomic material, giving rise to typical nebular and recombination lines. Furthermore, the interaction of the UV field with the remnants of the star’s natal cloud will produce a photodissociation region (PDR; see Hollenbach & Tielens 1999 for a comprehensive review). In a PDR, heating of the gas occurs by collisions with electrons, photoelectrically ejected from grain surfaces. The gas in the surface regions of PDRs reaches temperatures of typically 500 K (e.g. Tielens & Hollenbach 1985), producing a distinctive infrared spectrum due to collisionally excited low-lying levels of many molecular and neutral species.
Another phenomenon associated with YSOs are strong stellar winds, often collimated into a bipolar outflow. These will cause a shock as they drive a supersonic wave into the surrounding molecular cloud, heating it in the process. Shocks are usually divided into two distinct categories: J- or Jump-shocks, and C- or Continuous-shocks. In J-shocks viscous heating of the neutrals occurs in a thin shock front in which radiative cooling is insignificant, while the post-shock gas is heated to several times 10<sup>4</sup> degrees, dissociating all molecular material (e.g. Hollenbach & McKee 1989). Cooling of the post-shock gas occurs through atomic fine-structure lines as well as through re-formation of molecules. In contrast, C-shocks are magnetized, non-dissociative shocks in which the physical conditions change more gradually from their pre- to post-shock values and cooling is mainly through radiation from molecular material (e.g. Kaufman & Neufeld 1996). If the temperatures in a C-shock become sufficiently high to start to dissociate molecules, the cooling through the molecular lines diminishes, and the shock temperature increases until it turns into a J-shock. Shocks with a shock velocity larger than 40 km s<sup>-1</sup> are usually J-shocks, whereas slower shocks are usually of C-type (Chernoff et al. 1982). In this paper we will illustrate the power of infrared spectroscopy to study the above mentioned phenomena using two young bipolar nebulae, S106 and Cep A.
The bipolar nebula S106 is one of the most studied H ii regions in our galaxy. The exciting source of the region is a very young massive stellar object, known as either IRS4 in the terminology of Gehrz et al. (1982) or IRS3 following Pipher et al. (1976). Around this central source, Hodapp & Rayner (1991) found a cluster of about 160 stars, embedded in the molecular cloud surrounding S106. Observed radio emission and H i recombination lines from the region have been suggested to arise in a strong ($`\dot{M}`$ $``$ 2 $`\times `$ 10<sup>-5</sup> M yr<sup>-1</sup>), fast ($``$ 200 km s<sup>-1</sup>), stellar wind, driving a shock into the surrounding extended molecular cloud (Hippelein & Münch 1981). A dark lane, largely devoid of any optical or infrared emission, separates the two lobes of the H ii region. This dark lane has been quoted many times as a prime example of a large (30″) circumstellar disk, consisting of cool gas and dust, seen nearly edge-on (Bieging 1984; Mezger et al. 1987). However, Loushin et al. (1990) showed this structure to be an expanding ring of molecular material rather than a protoplanetary disk. Near-infrared ro-vibrational emission of H<sub>2</sub> is seen to arise just outside the H ii region, suggesting an origin in a PDR (Hayashi et al. 1990).
Cep A is another well-known site of recent star formation. It contains a luminous (2.5 $`\times `$ 10<sup>4</sup> L; Evans et al. 1981) far-infrared source, as well as several fainter infrared sources. The core of the Cep A region remains obscured at optical and infrared wavelengths due to massive amounts of extinction, and is the source of an energetic, complex molecular outflow (Bally & Lane 1982). Radio-observations show that it contains two strings of about 13 ultracompact H ii regions, arranged in a Y-shape (Hughes 1988). The source IRS 6a (Lenzen et al. 1984), located in the infrared nebula to the east of the core, is the dominant source of the eastern lobe at 20 $`\mu `$m (Ellis et al. 1990), but polarization measurements show that the nebula is illuminated by the compact radio source HW-2 (Hughes & Wouterloot 1984), located $``$ 5″ south of IRS 6a (Casement & McLean 1996). The fact that it is not visible at 20 $`\mu `$m implies an extremely large extinction towards HW-2. H<sub>2</sub> emission occurs in both the eastern and western parts of the nebula (Doyon & Nadeau 1988). Hartigan et al. (1996) showed that the molecular emission to the east appears as a regular jet, whereas that to the west concentrates in shells, which they proposed to arise in wakes from bow shocks surrounding Herbig-Haro objects in the region. ISO observations of the western part of the region were discussed by Wright et al. (1996), who modelled the observed infrared fine structure lines as arising in a planar J-shock with shock velocity 70–80 km s<sup>-1</sup>. The same authors modelled the observed molecular hydrogen emission as arising in a combination of several C-shocks. This multitude of shocks in the region agrees with the explanation of Narayanan & Walker (1996), who reported the presence of multiple episodes of outflow activity from the region. Recently, Goetz et al. (1998) presented strong evidence for this scenario through near-infrared H<sub>2</sub> and \[Fe ii\] images of Cep A East showing several distinct regions containing shocks.
In this paper we will present new Infrared Space Observatory (ISO; Kessler et al. 1996) data on infrared fine-structure and molecular emission lines from the central region of S106, centered on IRS4, and the eastern part of Cep A, centered on IRS 6a. These data give for the first time access to a large spectral range that was unaccessible from ground-based observations and contain lines which are decisive probes of the excitation conditions. We will show that the observed emission-line spectrum can be well explained as arising in the combination of a PDR and a H ii region in the case of S106, whereas this emission is shock-excited in the case of Cep A.
## 2 Observations and data
ISO Short Wavelength (2.4–45 $`\mu `$m) Spectrometer (SWS; de Graauw et al. 1996) full grating scans (“AOT S01”) of S106 IR and Cep A East were obtained in ISO revolutions 335 (JD 2450373.740) and 843 (JD 2450880.050), respectively. Each of these observations took 2124 seconds of observing time. In addition to this, deeper SWS grating scans on selected molecular and fine transition lines in the two objects (“AOT S02”) were obtained in revolutions 134 (at JD 2450172.705; S106), 220 (JD 2450258.463; Cep A) and 566 (JD 2450603.706; Cep A). In revolutions 558 (JD 2450596.103) and 580 (JD 2450617.838), scans at the full SWS grating resolution (“AOT S06”) covering the wavelength ranges of 3.0–3.5, 4.0–6.8, 7.0–7.6 and 12.1–16.5 $`\mu `$m were obtained for S106 as well. ISO Long Wavelength (43–197 $`\mu `$m) Spectrometer (LWS; Clegg et al. 1996) grating scans (“AOT L01”) of S106 IR and Cep A East were obtained in revolutions 134 (JD 2450172.726) and 566 (JD 2450603.681), respectively. Data were reduced in a standard fashion using calibration files corresponding to OLP version 7.0, after which they were corrected for remaining fringing and glitches. To increase the S/N in the final spectra, the detectors were aligned and statistical outliers were removed, after which the spectra were rebinned to a lower spectral resolution.
The SWS full grating scans consist of twelve different grating scans, each covering a small wavelength region, which were joined to form one single spectrum. Because of the variation of the diffraction limit of the telescope with wavelength, different SWS bands use apertures of different sizes. This is illustrated in Figs. 1 and 2, where we show their relative position and size, overlaid on K-band (2.2 $`\mu `$m) images of S106 and Cep A East by Hodapp & Rayner (1991) and Hodapp (1994). For a source that is not point-like, one may see a discontinuity in the spectra at the wavelengths where a change in aperture occurs, which can indeed be seen in both the S106 and Cep A spectra. The relative discontinuities in S106 are close to the maximum possible values, indicating that the source is extended across the entire SWS aperture at the longer wavelengths (33″$`\times `$ 20″). In Cep A the discontinuities are smaller, pointing to a smaller size of the far-infrared source.
Since both spectrometers on board ISO use entrance apertures that are smaller than the beam size, a wavelength dependent correction has to be applied to the standard flux calibration when observing extended sources. We determined these correction factors by convolving the K-band images of S106 and Cep A East, with the beam profile and applied these to the data before further analysis. The maximum correction to the flux was 8% in the region around 40 $`\mu `$m. If the sources are more extended at the longer wavelengths or in specific lines than the K-band images we employed in estimating the diffraction losses, we will have underestimated the flux. The maximum error this could introduce in the flux calibration is $``$15%.
Figures 3 and 4 show the resulting SWS and LWS full grating scans for S106 IR and Cep A East. Scanned lines and measured line fluxes or upper limits (total flux for line with peak flux 3$`\sigma `$) are listed in Table 1. The errors listed in Table 1 include the errors in the absolute flux calibration; errors in line ratios used in the remainder of this paper may be lower. Plots of all detected lines, rebinned to a resolution $`\lambda /\mathrm{\Delta }\lambda `$ of 1500 with an oversampling factor of four (SWS), or to a resolution of 500 (LWS), are given in Figs. 5 and 6. In the same figures we also show the line profiles expected for an unresolved point-source and an unresolved extended source filling the entire aperture. All detected lines in both S106 and Cep A are compatible with these profiles. For lines measured with the full SWS grating resolution (AOT S02 or S06), this indicates that they have a FWHM of less than $``$ 150 km s<sup>-1</sup>.
## 3 Solid-state features
The SWS spectrum of S106 (Fig. 3) consists of a relatively smooth continuum, with numerous strong emission lines superimposed. The continuum appears to consist of two components, with the break between the two occurring around 10 $`\mu `$m. This is similar to the situation seen in other regions of recent star formation, where these two components are commonly ascribed to emission due to dust close to the star (2–10 $`\mu `$m) and small dust grains in the wide circumstellar environment (10–200 $`\mu `$m; e.g. Cesarsky et al. 1996).
The familiar UIR bands at 3.3, 3.4, 6.2, 7.6, 7.8, 8.6, 11.3 and 12.7 $`\mu `$m, usually attributed to polycyclic aromatic hydrocarbons (PAHs), are present in emission and strong. As first demonstrated by ground-based spatially resolved CVF spectroscopy of the 3.3 $`\mu `$m band by Felli et al. (1984), these features do not come from the central point-source IRS4, but originate in the wide circumstellar environment. The band-strengths of the UIR features in our ISO spectra appear consistent with the ground-based measurements present in literature (Hefele & Hölzle 1980; Felli et al. 1984; Geballe et al. 1985).
An absorption band around 3.0 $`\mu `$m due to the O–H stretch mode of H<sub>2</sub>O ice, previously undetected in S106, is present in the SWS spectrum. Although partly filled in by emission from the 7.6, 7.8, 8.6 and 11.3 $`\mu `$m UIR bands, a 9.7 $`\mu `$m absorption feature due to the Si–O stretching mode in amorphous silicates also appears present in Fig. 3. The combination of the contamination by the UIR bands and the limited baseline can explain why this feature has not been detected from the ground. A weak dip around 18–19 $`\mu `$m in our ISO data may be identified with the Si–O bending mode of amorphous silicate, confirming the presence of silicate absorption in S106.
Since extinction in the continuum surrounding the 9.7 $`\mu `$m feature is small compared to the extinction within this feature, the extinction $`A_\lambda `$ at wavelength $`\lambda `$ across a non-saturated 9.7 $`\mu `$m feature can simply be obtained from the relation $`A_\lambda `$ = $`2.5\mathrm{log}(I/I_0)`$. Using an average interstellar extinction law which includes the silicate feature (Fluks et al. 1994), we can then convert these values of $`A_\lambda `$ to a visual extinction, resulting in a value of $`A_V`$ = 13$`\stackrel{m}{.}`$7 toward S106 IR. This value is in excellent agreement with that of 13$`\stackrel{m}{.}`$4 $`\pm `$ 2.7, derived towards stars in the S106 embedded cluster (Hodapp & Rayner 1991), or the value of 15<sup>m</sup>$`\pm `$ 3<sup>m</sup>predicted from infrared H i lines by Alonso-Kosta & Kwan (1989). It is also in agreement with the extinction of $`A_V`$ = 12<sup>m</sup> towards the northern lobe of the S106 nebula (Felli et al. 1984). It is somewhat larger than that towards the southern lobe ($`A_V`$ = 8<sup>m</sup>; Felli et al. 1984) and much smaller than that towards the central source ($`A_V`$ = 21<sup>m</sup>; Eiroa et al. 1979).
For Cep A East the SWS spectrum (Fig. 4) is dominated by absorption bands from a variety of ices. The O–H bend and stretch modes of H<sub>2</sub>O at 3.0 and 6.0 $`\mu `$m are strong, as are the 4.27 $`\mu `$m <sup>12</sup>C=O stretch and the 15.3 $`\mu `$m O=C=O bending mode of CO<sub>2</sub>. The <sup>12</sup>C=O stretch of CO at 4.67 $`\mu `$m is clearly detected. The absorption band around 7 $`\mu `$m, attributed to solid CH<sub>4</sub> (Boogert et al. 1996, 1997) is present as well. The unidentified absorption feature around 6.8 $`\mu `$m, also observed toward NGC 7538 IRS9 and RAFGL 7009S (Schutte et al. 1996; d’Hendecourt et al. 1996) is also present in Cep A. The 9.7 $`\mu `$m amorphous silicate feature consists of a very deep, saturated, absorption. From the non-saturated wings of this feature we can again derive a visual extinction, resulting in a value of $`A_V`$ = 270<sup>m</sup> for Cep A East. This value is within the range of $`A_V`$ = 75–1000<sup>m</sup> extinction for the central source and nebula reported by Lenzen et al. (1984). PAHs appear absent in Cep A East.
From the integrated optical depth $`\tau (\nu )𝑑\nu `$ of a non-saturated absorption feature we can compute a column density $`N`$ using an intrinsic band strength $`A_m`$. For H<sub>2</sub>O, CO, CO<sub>2</sub> and CH<sub>4</sub> ices, values of $`A_m`$ were measured by Gerakines et al. (1995) and Boogert et al. (1997). For silicates, $`A_m`$ is taken from Tielens & Allamandola (1987), based on the lab measurements by Day (1979). Integrated optical depth and column density values are listed in Table 2. The derived abundances of 100:6:15 for H<sub>2</sub>O:CO:CO<sub>2</sub> in Cep A East are within the range of values observed in other lines of sight (Whittet et al. 1996; d’Hendecourt et al. 1996).
## 4 Hydrogen recombination lines
A rich spectrum of H i recombination lines is present in the SWS data of S106 (Table 1). In contrast, H i lines are absent from the infrared spectra of Cep A East. It is not clear whether the H i lines observed with the rather large (27″$`\times `$ 14″) SWS beam in S106 are dominated by the strong ($`\dot{M}`$ $``$ 10<sup>-5</sup> M yr<sup>-1</sup>) stellar wind (Hippelein & Münch 1981; Felli et al. 1984), or could be due to the extended H ii region seen in H$`\alpha `$ and Br$`\gamma `$ (Bally et al. 1998; Hayashi et al. 1990; Greene & Rayner 1994; Maillard et al. 1999).
A clue to what is the case could come from the observed line profiles. As can be seen from Fig. 5, the H i lines observed in S106 are wider than the instrumental line profile for a point source. Unfortunately, the effect of line broadening of a compact source with FWHM of a few hundred km s<sup>-1</sup>, as expected for the stellar wind (Felli et al. 1985; Drew et al. 1993), would be about the same as that of observing narrow recombination lines from an extended H ii region. However, from the fact that our Br$`\alpha `$ line flux of 1.7 $`\times `$ 10<sup>-14</sup> W m<sup>-2</sup> is identical to that observed by other authors in smaller apertures (Garden & Geballe 1986; Persson et al. 1988; Drew et al. 1993), we conclude that at least for Br$`\alpha `$ the dominant component must be the compact wind. In view of the limited spectral resolution of our data, a full analysis of the wind structure is beyond the scope of the present paper.
## 5 Molecular hydrogen emission
Both S106 and Cep A East are well-known sources of extended molecular hydrogen emission (Longmore et al. 1986; Hayashi et al. 1990; Bally & Lane 1982; Doyon & Nadeau 1988; Goetz et al. 1998). All ground-based H<sub>2</sub> measurements in literature refer to the ro-vibrational transitions, with upper energy levels $`E((\mathrm{J})/\mathrm{k}`$ above 5000 K. These do not probe the bulk of the H<sub>2</sub>, which is expected to be at much lower temperatures. However, in both S106 and Cep A East we have detected many pure rotational lines of H<sub>2</sub> (Table 1), which have much lower upper energy levels and are therefore able to directly probe the physical conditions in the dominant chemical species. And since their transition probabilities are quite small, these lines are optically thin and the excitation temperature will be close to the kinetic temperature of the gas.
From the H<sub>2</sub> line fluxes $`I(J)`$ listed in Table 1 it is possible to calculate the apparent column densities of molecular hydrogen in the upper J levels, averaged over the SWS beam, $`N(J)`$, using $`N(J)=\frac{4\pi I(J)}{A}\frac{\lambda }{hc}`$, with $`\lambda `$ the wavelength, $`h`$ Planck’s constant and $`c`$ the speed of light. The transition probabilities $`A`$ were taken from Turner et al. (1977). Line fluxes were corrected for extinction using the average interstellar extinction law by Fluks et al. (1994). For S106 IR we adopted the value of $`A_V`$ = 13$`\stackrel{m}{.}`$7 derived in Section 3. The fact that the 0–0 S(3) line of H<sub>2</sub>, with a wavelength near the center of the 9.7 $`\mu `$m amorphous silicate feature, was detected in Cep A indicates that it cannot suffer from the same extinction as the continuum; the H<sub>2</sub> emission must originate in a spatially separate region from the continuum. Therefore we adopt a much smaller value of $`A_V`$ = 2<sup>m</sup>, expected to be a reasonable value for an origin in either a shock or PDR, for the emission lines observed toward Cep A East. Since the extinction correction at these mid-infrared wavelengths is small, a mis-estimate of $`A_V`$ by a few magnitudes will not affect our results.
A useful representation of the H<sub>2</sub> data is to plot the log of $`N(\mathrm{J})/\mathrm{g}`$, the apparent column density for a given J upper level divided by the statistical weight, versus the energy of the upper level, taken from Dabrowski (1984). These plots are shown in Fig. 7. Also plotted in these figures are several measurements of H<sub>2</sub> ro-vibrational lines in S106 and Cep A from literature (Longmore et al. 1986; Tanaka et al. 1989; Bally & Lane 1982). Although these measurements were taken at the same positions on the sky as ours, the beam sizes used were different and hence one may expect to see systematic differences in the derived specific intensities if the source of H<sub>2</sub> radiation is more extended than the beam size. In fact only the S106 measurements by Longmore et al. (1986), taken with the smallest beam diameter, 12″, differ systematically from the other measurements, indicating that the extent of the H<sub>2</sub> emitting region in S106 is probably somewhat larger than their beam size. This agrees well with H<sub>2</sub> 1–0 S(1) images of the region (Hayashi et al. 1990; Maillard et al. 1999). In Fig. 7 the Longmore et al. measurements of S106 were scaled to the Tanaka et al. data.
The statistical weight $`g`$ used in Fig. 7 is a combination of the rotational and nuclear spin components. We have assumed the high temperature equilibrium relative abundances of 3:1 for the ortho and para forms of H<sub>2</sub> (Burton et al. 1992). For a Boltzmann distribution, the points in the plot shown in Fig. 7 should form a nearly straight line whose slope is inversely proportional to the excitation temperature, while its intercept is a measure of the total column density of warm gas. Using the formula by Parmar et al. (1991) and the rotational constants by Dabrowski (1984) to fit our data points in the low-lying pure rotational levels to a Boltzmann distribution, we arrive at values of 500 K and 9 $`\times `$ 10<sup>19</sup> cm<sup>-2</sup> and 730 K and 5 $`\times `$ 10<sup>19</sup> cm<sup>-2</sup> for S106 IR and Cep A East, respectively. Using the distances to S106 and Cep A of 1200 and 690 pc (Rayner 1994; Mel’nikov et al. 1995), this corresponds to total molecular hydrogen masses of 0.008 and 0.003 M within the SWS beam. Using the alternative distance estimate towards S106 of 600 pc (Staude et al. 1982) would decrease the S106 H<sub>2</sub> mass to 0.004 M. The fitted Boltzmann distributions are shown as the leftmost dashed lines in Fig. 7. The fact that the points for the ortho and para form of H<sub>2</sub> lie on the same line proves that our assumption on their relative abundances is correct.
As can be seen from Fig. 7, both the pure-rotational and ro-vibrational lines of H<sub>2</sub> with upper level energies higher than 5000 K deviate significantly from the leftmost dashed line. In both cases a Boltzmann distribution was fitted to these lines as well, shown as the rightmost dashed line in both figures. In the case of S106, the relative location of the 1–0 and 2–1 S(1)-lines in Fig. 7 are indicative of fluorescent excitation through UV pumping (Draine & Bertoldi 1996; Black & van Dishoeck 1997). The fitted Boltzmann distribution to the higher energy levels has thus no physical meaning, but may still be useful to provide a simple parametrization of the relative population of the energy levels. In the case of Cep A, the lines may indicate the presence of a smaller column (3 $`\times `$ 10<sup>17</sup> cm<sup>-2</sup>) of hot (a few 1000 K) molecular hydrogen, in addition to the large column of warm gas.
Employing predictions of H<sub>2</sub> emission from PDR, J-shock and C-shock models by Burton et al. (1992), Hollenbach & McKee (1989) and Kaufman & Neufeld (1996), we determined the excitation temperature $`T_{\mathrm{rot}}`$ from the low-lying pure rotational levels from these models as a function of density $`n`$ and either incident FUV flux $`G`$ (in units of the average interstellar FUV field G<sub>0</sub> = 1.2 $`\times `$ 10<sup>-4</sup> erg cm<sup>-2</sup> s<sup>-1</sup> sr<sup>-1</sup>; Habing 1968) or shock velocity $`v_s`$ in an identical way as was done for the observations. The resulting relations between $`T_{\mathrm{rot}}`$ and $`n`$ or $`v_s`$ are shown in Fig. 8. As can be seen from these plots, the PDR and J-shock models predict a fairly small (200–540 K) range of resulting excitation temperatures, whereas in the C-shocks this range is much larger (100–1500 K). Furthermore, we see that in the model predictions for shocks the resulting $`T_{\mathrm{rot}}`$ does not depend much on density, whereas for PDRs it does not depend much on $`G`$, suggesting that once the mechanism of the H<sub>2</sub> emission is established, it can be used to constrain $`v_s`$ or $`n`$ in a straightforward way.
Comparing the excitation temperatures of 500 and 730 K for S106 and Cep A with those plotted in Fig. 8, we note that for S106 this falls well within the range of PDR- and C-shock model predictions, but are too high compared to the ones expected from J-shocks. The observed bright PAH emission features (Sect. 3) and the atomic fine-structure line spectrum (Sect. 6) point towards a PDR origin for the H<sub>2</sub> emission in S106. The higher temperature for Cep A can only be reproduced by the C-shock models. Therefore we tentatively conclude that a dense ($``$ 10<sup>6</sup> cm<sup>-3</sup>) PDR seems to be the best candidate to explain the observed H<sub>2</sub> emission in S106 IR and a slow ($``$ 20 km s<sup>-1</sup>) non-dissociative shock can explain the observed warm column of H<sub>2</sub> in Cep A East. Since the regions we are looking at probably only fill part of the SWS beam, the absolute intensity of the H<sub>2</sub> emission listed in Table 1 can also be reproduced by these same models by varying the beam filling factor.
## 6 Carbon-monoxide emission lines
In both S106 IR and Cep A East several ro-vibrational emission lines due to gas-phase CO were detected in the long-wavelength part of the LWS spectra (Figs. 5 and 6). Similar to what was done for the H<sub>2</sub> emission in the previous section, we constructed a CO excitation diagrams, using molecular data from Kirby-Docken & Liu (1978). They are shown in Fig. 9. The temperature and column of CO resulting from the Boltzmann fit to this excitation diagram are 420 K and 5.6 $`\times `$ 10<sup>18</sup> cm<sup>-2</sup> and 340 K and 1.4 $`\times `$ 10<sup>18</sup> cm<sup>-2</sup> for S106 and Cep A East, respectively. Corresponding CO masses are 1 $`\times `$ 10<sup>-2</sup> M (S106 at 1200 pc; for the 600 pc distance estimate this mass estimate would decrease to 3 $`\times `$ 10<sup>-3</sup> M) and 1 $`\times `$ 10<sup>-3</sup> M (Cep A East).
The CO excitation temperature of 420 K for S106 is somewhat lower than that found from the H<sub>2</sub> lines, in agreement with what is expected from a PDR, in which the CO emission arises in deeper embedded regions than the H<sub>2</sub>. We conclude that a PDR is the most likely candidate for the source of the gas-phase CO emission in S106. For Cep A East the CO excitation temperature of 340 K is much lower than that found from the H<sub>2</sub> lines. This behaviour is hard to reconcile with an origin in the C-shock invoked to explain the H<sub>2</sub> observations. However, a J-shock, necessary to explain the ionic lines which we will discuss in the next section, may produce the observed CO emission while only emitting little warm H<sub>2</sub>. Since the LWS CO observations were made with a much larger beam than the SWS H<sub>2</sub> observations, an alternative explanation may be that we are observing emission from two spatially distinct regions. The relatively large mass of warm CO as compared to that seen in H<sub>2</sub> may also easily be explained in this scenario.
The observed CO lines have critical densities of around 10<sup>6</sup> cm<sup>-3</sup>. Therefore the detection of these lines also implies densities of this order of magnitude or higher in the originating region. For the shock seen in Cep A East this may not be unreasonable. However, assuming that these densities would exist in the entire S106 PDR would be implausible. If these CO lines do indeed arise in the large-scale environment of S106, the PDR must therefore have a clumpy structure (e.g. Burton et al. 1990). An alternative explanation of the presence of these lines, would be to identify the originating region with the surface of the extended disk-like structure surrounding S106 which could act like a PDR. With the present data-set we cannot make a distinction between these two possibilities.
## 7 Atomic fine structure lines
Important constraints on the physical conditions in the line emitting region come from the observed fine structure lines. As a first step towards identifying the mechanism responsible for the observed emission we can look at the mere presence of certain lines. The observed lines with high ionization potentials in S106, such as \[O iii\], \[Ar iii\], \[Ne iii\], \[Ni iii\], \[Fe iii\] and \[S iv\] can only originate in the H ii region surrounding S106. The PAH emission as well as the molecular lines observed towards S106 are indicative of the presence of a PDR as well. This PDR might contribute to the observed \[Fe ii\], \[Ni ii\], \[Si ii\], \[O i\] and \[C ii\] emission. We thus need to model the fine-structure lines in S106 IR as arising in the combination of an H ii model and a PDR.
We have used the photo-ionization code cloudy (version 90.04; Ferland 1996) to generate model predictions for line strengths for the H ii region surrounding S106 IR, assuming a spherical geometry and constant hydrogen density throughout the region. We generated a grid of models in which the input spectrum, taken from Kurucz (1991) models for a stellar photosphere, and the electron density in the H ii region were varied. The total luminosity of the model was fixed to the total bolometric luminosity of S106 of 4.2 $`\times `$ 10<sup>4</sup> L (computed from the spectral energy distribution using data from literature and assuming a distance of 1.2 kpc). The line ratios of the observed \[O iii\] and \[S iii\] fine structure lines are expected to depend mainly on density, and hardly on the temperature. Their behaviour in our H ii region model together with the observed line ratios in S106 are shown in Fig. 10. From the fact that both line ratios agree on the thus obtained value for the electron density in the H ii region, 1.3–2.5 $`\times `$ 10<sup>3</sup> cm<sup>-3</sup>, we conclude that our assumption of a constant density is reasonable.
Several other lines ratios (\[Ar ii\]/\[Ar iii\], \[Ne ii\]/\[Ne iii\] and \[S iii\]/\[S iv\]) depend mainly on the effective temperature of the star and do not depend on density much. Their behaviour in the cloudy models for S106 is shown in Fig. 10 as well. From these plots, it can be seen that the source of ionizing photons should have a temperature between 37,000 and 40,000 K, corresponding to a spectral type of O6–O8, in good agreement with earlier determinations (Gehrz et al. 1982; Staude et al. 1982). However, not all observed line ratios yield the same temperature for the central star. Most likely this is due to the fact that the Kurucz models used for the input spectrum do not include the opacity shortward of the He ii ionization limit due to the stellar wind, which is expected to be strong in the case of S106. The true temperature for the central source is therefore expected to be around the lower range of temperatures deduced from the Kurucz synthetic photospheres, i.e. around 37,000 K. For illustrative purposes, we also show the ratio of \[O i\] 63.2 $`\mu `$m/\[O iii\] 51.8 $`\mu `$m for our model H ii region, showing that virtually all the atomic oxygen will be in the form of \[O iii\] and that we can thus safely attribute all of the observed \[O i\] emission to the PDR.
To be able to explain the observed intensity of 1.3 $`\times `$ 10<sup>-2</sup> erg s<sup>-1</sup> cm<sup>-1</sup> sr<sup>-1</sup> for the \[Si ii\] line at 34.8 $`\mu `$m in S106, the Tielens & Hollenbach (1985) PDR models require a density higher than $``$ 10<sup>5</sup> cm<sup>-3</sup> and $`G`$ $``$ 10<sup>5</sup> G<sub>0</sub>. This regime can also reproduce the observed ratios of \[Si ii\] and the \[Fe ii\] and \[O i\] lines, although to reproduce both the exact strength of the \[Si ii\] emission and the relative strength to that to the \[Fe ii\] lines requires a \[Fe ii\] depletion of $``$ 30% higher than the one assumed in the Tielens & Hollenbach models. An alternative explanation could be that the PDR doesn’t fill the SWS aperture at 35 micrometer (33″$`\times `$ 20″) completely, increasing the surface brightness of all lines. The predicted intensities of \[S i\] and \[Fe i\] are sufficiently low to be undetectable, in agreement with the observations.
Towards Cep A East, only fine structure lines that can be produced in shocks were observed. In contrast to both C- and J-shocks, PDRs do not produce significant quantities of \[S i\] emission (Tielens & Hollenbach 1985). Therefore this line must be completely due to one or more shocks. C-shocks only contain trace fractions of ions and hence cannot explain the observed \[Fe ii\], \[Si ii\] and \[C ii\] emission. Hence the simplest hypothesis would be to try to explain the observed fine-structure lines in Cep A in terms of a J-shock model. To explain the \[S i\] surface brightness of 1.0 $`\times `$ 10<sup>-3</sup> erg s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup> requires a moderately dense (10<sup>4</sup>–10<sup>5</sup> cm<sup>-3</sup>) pre-shock gas. If the shock does not fill the ISO beam completely, as is expected, this number will go up. The absence of \[Fe i\] emission does indicate that we are dealing with higher (10<sup>5</sup>–10<sup>6</sup> cm<sup>-3</sup>) densities, so line ratios seem to provide more reliable constraints in this case. However, the aperture sizes for lines measured at different wavelengths are also different in some cases, making this method only reliably applicable to line ratios measured in identical SWS or LWS apertures. Therefore the ratio of \[Fe ii\] 26.0 $`\mu `$m to \[S i\] 25.2 $`\mu `$m might provide the most reliable indicator of physical conditions in the shock. To reproduce the observed value of 0.26 with the Hollenbach & McKee (1989) J-shock models requires either moderately dense (10<sup>5</sup> cm<sup>-3</sup>) gas with a low (30 km s<sup>-1</sup>) shock velocity or a high (10<sup>6</sup> cm<sup>-3</sup>) density with a faster ($``$ 60 km s<sup>-1</sup>) shock. To have the predicted \[Fe i\] emission sufficiently low to explain our non-detections of this species, the moderately dense, slow J-shock model is required. This regime also reproduces the relative \[Si ii\] strength.
As was discussed in the previous section, one or multiple J-shocks cannot reproduce the observed H<sub>2</sub> emission (although a contribution to this can be expected). The success of the J-shock model in explaining the observed fine-structure lines and the absence of PAH emission in Cep A East, excluding the possibility of a significant contribution from a PDR, leads us to pose that a combination of one or more J- and C-shocks must be responsible for the observed emission in Cep A. The presence of more than one type of shock in the region could be linked to the reported multiple episodes of outflow activity from the embedded source (Narayanan & Walker 1996). In the presence of both a J- and a C-shock, the \[Si ii\] and \[Fe ii\] emission would originate completely in the J-shock, whereas both the C- and the J-shock would contribute to the observed \[S i\] and H<sub>2</sub> spectra. With the J-shock parameters derived from the fine-structure lines, an additional C-shock component with a shock-velocity of about 20 km s<sup>-1</sup> is required to reproduce the observed pure rotational H<sub>2</sub> emission. The observed hot column of H<sub>2</sub> in Cep A may then be due to formation pumping (the effect that H<sub>2</sub> gets re-formed with non-zero energy in the post-shock gas after being dissociated in the shock front) in the J-shock. To have the ro-vibrational lines of comparable intensity as the rotational lines (through collisional excitation) requires pre-shock densities of the order of 10<sup>6</sup> cm<sup>-3</sup>. For both S106 IR and Cep A, predicted lines strengths of the best fit models are also listed in Table 1.
## 8 Discussion and conclusions
In the previous sections we saw that the infrared emission-line spectrum of S106 could be well explained by the presence of an H ii region and a clumpy PDR. The average density in the H ii region should be around 2 $`\times `$ 10<sup>3</sup> cm<sup>-3</sup>. The density in the PDR, which was estimated to be 10<sup>5</sup>–10<sup>6</sup> cm<sup>-3</sup>, is several orders of magnitude higher, in agreement with one would expect in a scenario in which a young massive star has photo-ionized its natal cloud in all but the densest clumps.
In Section 7 we concluded that the incident FUV flux on the S106 PDR should be rather high ($``$ 10<sup>5</sup> G<sub>0</sub>). We also showed that the central source in S106 is of spectral type O8, with a total luminosity of 4.2 $`\times `$ 10<sup>4</sup> L. From a Kurucz (1991) model for a stellar photosphere with $`T_{\mathrm{eff}}`$ = 37,000 K and $`\mathrm{log}g`$ = 4.0, we compute that such a star emits a total FUV (6–13.6 eV) flux of 3.6 $`\times `$ 10<sup>36</sup> erg s<sup>-1</sup> sr<sup>-1</sup>. To dilute this stellar FUV field to a value of 10<sup>5</sup>–10<sup>6</sup> G<sub>0</sub>, the PDR must be at a location 3–10 $`\times `$ 10<sup>3</sup> AU away from the central star, corresponding to a projected distance of 5–17″. This projected distance is independent of the assumed distance towards S106. It is within the range allowed by the SWS entrance apertures and is compatible with an origin in the $``$ 30″ diameter region of H<sub>2</sub> emission in the H<sub>2</sub> 1–0 S(1) image of S106 by Hayashi et al. (1990). We conclude that a central O8 star is sufficient to produce the FUV radiation field reaching the S106 PDR.
These results agree well with other determinations. Schneider et al. (2000) have spatially resolved the 157.7 $`\mu `$m \[C ii\] and 63.2 $`\mu `$m \[O i\] emission in S106 and concluded that there is a PDR region with density $``$ 10<sup>6</sup> cm<sup>-3</sup> and strong (up to 8 $`\times `$ 10<sup>5</sup> G<sub>0</sub>) UV field in the immediate environment of S106 and a more extended dense (10<sup>5</sup>–10<sup>6</sup> G<sub>0</sub>) PDR region exposed to a lower intensity (300–500 G<sub>0</sub>) UV field. We are clearly dominated by their first component. This same component may be identified with the high-density (10<sup>5</sup> cm<sup>-3</sup>) fluorescent H<sub>2</sub> regions identified by Hayashi et al. (1990).
Both the presence of \[S i\] 25.2 $`\mu `$m emission and the H<sub>2</sub> spectrum point unambiguously to shocked gas as the origin of the infrared emission lines in Cep A East. A combination of a slow (20 km s<sup>-1</sup>) non-dissociative shock and a faster (30–60 km s<sup>-1</sup>) dissociative shock are required to explain our ISO spectra. The density of the pre-shock gas in Cep A East is similar to that in the S106 PDR: 10<sup>5</sup>–10<sup>6</sup> cm<sup>-3</sup>. Our J-shock component may be identified with the \[Fe ii\] 1.64 $`\mu `$m clumps observed by Goetz et al. (1998). The C-shocks could be due to their H<sub>2</sub> knots which do not show ionic emission.
It is interesting to compare the results derived here for the emission from the eastern lobe of Cep A with the results by Wright et al. (1996), obtained using the same instruments on board ISO, for the western part of the nebula. They derived a $`T_{\mathrm{rot}}`$ of 700 $`\pm `$ 30 K for the H<sub>2</sub> lines with upper level energies up to 7000 K, and a temperature range up to 11,000 K for the higher upper level energies. They explained this H<sub>2</sub> emission as arising from a combination of at least two C-shocks with different pre-shock density, shock velocity and covering factor. In addition to this, they also reported emission from \[Ne ii\], \[S i\] and \[Si ii\], which they explained as arising in a planar J-shock with pre-shock density of 10<sup>3</sup>–10<sup>4</sup> cm<sup>-3</sup> and shock velocity 70–80 km s<sup>-1</sup>. Qualitatively, the detected lines in the eastern and western part of the nebula are well in agreement. The main difference between the observed lines seems to be that they are much more intense in the western part of Cep A, indicating that the densities there are a factor of 100 higher than those obtained for the eastern lobe, but the shock velocities are comparable in both parts of the nebula. The similarity between these two parts of Cep A are in agreement with a scenario in which the driving source of the molecular outflow, HW-2, went through multiple episodes of outflow activity, as suggested by Narayanan & Walker (1996). In this scenario, the J-shock component could be due to the most recent period of enhanced mass loss, whereas one or more C-shocks could be due to older generations of outflows. Alternatively, the stellar wind material could produce a (fast) J-shock while the surrounding molecular cloud might be swept up by a slower C-shock.
As was already remarked by Staude & Elsässer (1993), the differences between the environments of these two massive embedded YSOs, S106 and Cep A, are remarkable. The mid- and far-infrared observations presented in this paper show that the difference between these two sources cannot be due to a different orientation of their nebulae; it can only be a reflection of their different evolutionary status, with Cep A being the younger of the two. In S106 the stellar wind and UV radiation of the exciting source have cleared and excited a sufficiently large region to create strong PDR emission, whereas in the case of Cep A, the central source is still heavily embedded and we only observe the interaction of the outflow with its surroundings. In due time, Cep A will clear its surroundings, ionize the hydrogen, and evolve into a bipolar nebula quite similar to S106.
###### Acknowledgements.
The authors would like to thank John Rayner and Klaus-Werner Hodapp for providing us with the K-band images of S106 and Cep A shown in Figs. 1 and 2. Peter van Hoof is kindly acknowledged for providing us with a copy of the cloudy computer code. Frank Molster is thanked for providing the authors with valuable input on a draft of the paper. MvdA acknowledges financial support from NWO grant 614.41.003 and through a NWO Pionier grant to L.B.F.M. Waters. This research has made use of the Simbad data base, operated at CDS, Strasbourg, France.
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# References
1. Gauge fixing in Yang–Mills theories is well understood in perturbation theory. But outside of perturbation theory the situation is more complicated, because of the existence of Gribov copies . These complications arise as soon as the theory is nonlinear, which is the case for nonabelian gauge theories, but also, for instance, with a nonlinear gauge-fixing condition in an abelian theory.
This does not mean that Yang–Mills theories do not exist outside perturbation theory. It is well known that they can be defined on the lattice, using the compactness of the group to dispense with gauge fixing altogether. However, nonperturbative gauge fixing is interesting in a variety of contexts, ranging from the need to make contact between lattice and continuum calculations (see for example Ref. ) to the definition of chiral gauge theories on the lattice (see for example Ref. ). For a discussion of the numerical implementation of the gauge-fixing method discussed in this paper, see Ref. .
One can try to extend the standard BRST construction to the lattice. When one does this, one finds, quite generally and rigorously, that the BRST gauge-fixed partition function vanishes . This result can be heuristically explained as the result of pairwise cancellations of lattice Gribov copies, which occur with opposite signs of the Faddeev–Popov determinant. (For some recent work on overcoming this problem, see Refs. .)
A different nonperturbative method for gauge fixing has been proposed some time ago by Jona-Lasinio and Parrinello and by Zwanziger (JLPZ) (see also Ref. ). Starting from the euclidean ungauged partition function
$$Z=𝒟_\mathrm{H}A\mathrm{exp}\left(S_{\mathrm{inv}}(A)\right),$$
(1)
one inserts one in the form
$$1=\frac{𝒟_\mathrm{H}h\mathrm{exp}\left(S_{\mathrm{ni}}(A^h)\right)}{𝒟_\mathrm{H}h\mathrm{exp}\left(S_{\mathrm{ni}}(A^h)\right)}$$
(2)
into the functional integral. Here $`𝒟_\mathrm{H}A`$ and $`𝒟_\mathrm{H}h`$ denote the invariant measures over the gauge field $`A`$ and the group-valued scalar field $`h`$, respectively; $`A^h`$ is the gauge transform of $`A`$ under a gauge transformation $`h`$. $`S_{\mathrm{ni}}(A)`$ is any gauge noninvariant local functional of $`A`$. Because of the invariance of $`𝒟_\mathrm{H}A`$ and $`S_{\mathrm{inv}}(A)`$, we can perform a gauge transformation in the numerator of Eq. (1) with Eq. (2) inserted, and, dropping a trivial factor $`𝒟_\mathrm{H}h`$, we obtain
$$Z=𝒟_\mathrm{H}A\frac{\mathrm{exp}\left(S_{\mathrm{inv}}(A)S_{\mathrm{ni}}(A)\right)}{𝒟_\mathrm{H}h\mathrm{exp}\left(S_{\mathrm{ni}}(A^h)\right)}.$$
(3)
This procedure is completely rigorous on the lattice, and sources coupled to gauge-invariant operators can be added without changing the argument. We also note that the Boltzmann weights for both integrations contained in Eq. (3) are positive.
While this method is conceptually very simple, much less is known about how this nonperturbative method works in perturbation theory, as compared with the standard gauge fixing based on BRST invariance. Note, however, that there are some similarities between the two procedures. A noninvariant term $`S_{\mathrm{ni}}(A)`$ is added to the action, and (for a suitable choice) this will make the quadratic part of the action invertible. The contributions of different orbits are weighted properly because of the group integral in the denominator of Eq. (3), which plays a role similar to that of the Faddeev–Popov (FP) determinant in the usual case. In the FP case, however, the determinant can be expressed as an integral over ghosts, and the gauge-fixed action including the ghost terms is local, whereas here this is not the case. This is an important difference, because locality is a key ingredient in power-counting arguments, and thus at the heart of the usual perturbative analysis of renormalization.
It is therefore of interest to find out whether perturbation theory can be systematically developed for a JLPZ gauge-fixed Yang–Mills theory, and its relation to the usual FP results. This question has been addressed previously by Fachin , who analyzed the vacuum polarization for the choice $`S_{\mathrm{ni}}(A)=\mathrm{tr}(M^2A_\mu ^2)`$ at one loop. He concluded that for $`M\mathrm{}`$ the transverse part of the vacuum polarization agrees with that obtained using the FP method; the longitudinal part vanishes for $`M\mathrm{}`$, as in Landau gauge. This equivalence for $`M\mathrm{}`$ at fixed cutoff was already derived formally in Ref. , and another formal discussion appeared recently in Ref. . Here, we are interested in considering the situation where $`M`$ is chosen to be of the same order of magnitude as the cutoff. The transverse part of the one-loop vacuum polarization with JLPZ gauge fixing, for example, was found to contain terms proportional to $`p^2\delta _{\mu \nu }p_\mu p_\nu `$ times $`p^2/M^2`$ or $`p^4/M^4`$ , so, if we choose $`M\mathrm{\Lambda }`$ (with $`\mathrm{\Lambda }`$ the cutoff), such terms are of order $`1/\mathrm{\Lambda }^2`$, and vanish when we take $`\mathrm{\Lambda }\mathrm{}`$. We present below a general argument that, if we choose $`M`$ to be fixed in units of the cutoff, perturbation theory for a JLPZ gauge-fixed theory is equivalent to that of the same theory gauge-fixed using the standard FP method. We discuss in some detail what “equivalent” means in this context. We also show that our arguments for the equivalence of JLPZ and FP gauge fixing can be used to construct a set of local Feynman rules for the JLPZ gauge-fixed gauge theory.
2. We begin with the JLPZ gauge-fixed partition function in the presence of a source for the gauge field $`A_\mu `$ (working in euclidean space-time):
$$Z(J)=𝒟_\mathrm{H}A\frac{\mathrm{exp}\left(S_{\mathrm{inv}}(A)S_{\mathrm{ni}}(A)+J_\mu A_\mu \right)}{𝒟_\mathrm{H}h\mathrm{exp}\left(S_{\mathrm{ni}}(A^h)\right)},$$
(4)
where $`S_{\mathrm{inv}}(A)`$ is the gauge-invariant classical action, and $`S_{\mathrm{ni}}(A)`$ is not invariant. As noted before, $`A^h`$ is the gauge transform of $`A`$ under a finite gauge transformation,
$$A_\mu ^h=h(A_\mu +\frac{i}{g}_\mu )h^{}.$$
(5)
We will take
$`S_{\mathrm{inv}}(A)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{tr}(F_{\mu \nu }^2),`$ (6)
$`S_{\mathrm{ni}}(A)`$ $`=`$ $`\mathrm{tr}(M^2A_\mu ^2),`$
where we do not indicate the integration $`d^4x`$ explicitly, and $`M`$ a parameter with the dimension of a mass. We will assume that $`M`$ is proportional to the cutoff for the case of regulators such as the lattice, in which case $`M`$ is proportional to the inverse lattice spacing. However, it is sufficient for our arguments that $`M`$ is chosen large compared to all physical scales. This assumption applies in the case of dimensional regularization. Also, $`A_\mu =A_\mu ^aT^a`$, with $`T^a`$ the generators of the gauge group, with
$`\mathrm{tr}(T^aT^b)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta ^{ab},`$ (7)
$`[T^a,T^b]`$ $`=`$ $`if^{abc}T^c.`$ (8)
If we now consider perturbation theory, and use a regulator without power-like divergences such as dimensional regularization, we may extend the range of integration of the variables $`A_\mu ^a`$ from $`\mathrm{}`$ to $`\mathrm{}`$. The invariant measure is just the flat measure
$$𝒟_\mathrm{H}A=𝒟A=\underset{x,\mu ,a}{}dA_\mu ^a(x),$$
(9)
and similarly for $`𝒟_\mathrm{H}h`$ (after parametrizing $`h`$ as in Eq. (11) below).
We wish to show that this gauge-fixed partition function is equivalent (in a way to be discussed in more detail in the next section), under certain assumptions which are valid in perturbation theory, to the standard FP gauge-fixed form for the covariant gauge $`_\mu A_\mu =0`$. The derivation is based on two ingredients. The first ingredient is that we insert a constant into the partition function, written in the form
$$\text{constant}=\mathrm{det}(\mathrm{})𝒟\eta \delta (_\mu A_\mu \frac{1}{M}\mathrm{}\eta ),$$
(10)
where $`𝒟\eta `$ denotes the (flat) measure for $`\eta `$, a new field which takes values in the Lie algebra of the gauge group. We may choose boundary conditions such as to avoid the trivial zero mode of $`\mathrm{}`$. The second ingredient is to make use of the fact that the physics at scales below $`M`$ is not altered by adding or changing terms which are irrelevant in the sense of Wilson’s renormalization group.
First, let us expand $`S_{\mathrm{ni}}(A^h)`$ in $`g/M`$, writing
$$h=\mathrm{exp}(i\frac{g}{M}\theta ).$$
(11)
We find
$$S_{\mathrm{ni}}(A^h)=\frac{1}{2}M^2(A_\mu ^a)^2M\theta ^a_\mu A_\mu ^a+\frac{1}{2}_\mu \theta ^a(D_\mu (A)\theta )^a+O\left(\frac{g}{M}\right),$$
(12)
where
$$(D_\mu (A)\theta )^a=_\mu \theta ^a+gf^{abc}A_\mu ^b\theta ^c.$$
(13)
Note that $`_\mu D_\mu (A)`$ is just the FP operator for the covariant gauge.
We then shift the vector field
$$A_\mu A_{}^{}{}_{\mu }{}^{}=A_\mu \frac{1}{M}_\mu \eta .$$
(14)
This gives
$`S_{\mathrm{ni}}(A^h)`$ $`=`$ $`{\displaystyle \frac{1}{2}}M^2(A_{}^{}{}_{\mu }{}^{a})^2+MA_{}^{}{}_{\mu }{}^{a}_\mu \eta ^a+{\displaystyle \frac{1}{2}}(_\mu \eta ^a)^2`$
$`M\theta ^a_\mu A_{}^{}{}_{\mu }{}^{a}\theta ^a\mathrm{}\eta ^a+{\displaystyle \frac{1}{2}}_\mu \theta ^a(D_\mu (A^{})\theta )^a+O\left({\displaystyle \frac{g}{M}}\right),`$
and
$$S_{\mathrm{inv}}(A)+S_{\mathrm{ni}}(A)=S_{\mathrm{inv}}(A^{})+\frac{1}{2}M^2(A_{}^{}{}_{\mu }{}^{a})^2+MA_{}^{}{}_{\mu }{}^{a}_\mu \eta ^a+\frac{1}{2}(_\mu \eta ^a)^2+O\left(\frac{g}{M}\right).$$
(16)
The terms $`\frac{1}{2}M^2(A_{}^{}{}_{\mu }{}^{a})^2+MA_{}^{}{}_{\mu }{}^{a}_\mu \eta ^a+\frac{1}{2}(_\mu \eta ^a)^2`$ cancel between the numerator and denominator of the integrand in Eq. (4). The terms $`M\theta ^a_\mu A_{}^{}{}_{\mu }{}^{a}`$ and $`MA_{}^{}{}_{\mu }{}^{a}_\mu \eta ^a`$ (which is equal to $`M\eta ^a_\mu A_{}^{}{}_{\mu }{}^{a}`$ by partial integration) may be dropped because of the $`\delta `$-function, $`\delta (_\mu A_{}^{}{}_{\mu }{}^{})`$, in Eq. (10).
All the $`O(g/M)`$ terms are irrelevant, and may therefore be omitted without changing the (renormalized) theory. Doing this, we can perform the $`h`$-integral in the denominator of the integrand in Eq. (4), obtaining
$`Z(J)`$ $`=`$ $`\mathrm{det}(\mathrm{}){\displaystyle 𝒟\eta 𝒟A^{}\delta (_\mu A_{}^{}{}_{\mu }{}^{})\mathrm{det}^{1/2}(_\mu D_\mu (A^{}))}`$
$`\times \mathrm{exp}\left(S_{\mathrm{inv}}(A^{})+J_\mu A_\mu {\displaystyle \frac{1}{2}}\eta \mathrm{}(_\mu D_\mu (A^{}))^1\mathrm{}\eta \right).`$
Using the fact that physical quantities do not change when we replace the source term $`J_\mu A_\mu `$ by a new source term $`J_\mu ^{}A_{}^{}{}_{\mu }{}^{}`$ coupling to $`A_{}^{}{}_{\mu }{}^{}`$ instead of to $`A_\mu `$ ($`A_\mu `$ and $`A_{}^{}{}_{\mu }{}^{}`$ are equivalent interpolating fields), we can now also perform the $`\eta `$ integral. Dropping the primes, we obtain
$$Z(J)=𝒟A\delta (_\mu A_\mu )\mathrm{det}(_\mu D_\mu (A))\mathrm{exp}\left(S_{\mathrm{inv}}(A)+J_\mu A_\mu \right).$$
(18)
Representing the $`\delta `$-function as
$$\delta (_\mu A_\mu ^a)\underset{\xi 0}{lim}\mathrm{exp}\left(\frac{1}{2\xi }(_\mu A_\mu ^a)^2\right),$$
(19)
and introducing an algebraic field $`B`$ and ghost fields $`c`$ and $`\overline{c}`$, this can be recast as
$`Z(J)=\underset{\xi 0}{lim}{\displaystyle 𝒟A𝒟B𝒟\overline{c}𝒟c}`$ (20)
$`\times \mathrm{exp}\left(S_{\mathrm{inv}}(A){\displaystyle \frac{1}{2}}\xi B^2+iB_\mu A_\mu \overline{c}(_\mu D_\mu (A))c+J_\mu A_\mu \right).`$
This is the standard BRST-invariant form of the FP gauge-fixed partition function in Landau gauge. Of course, correlation functions of gauge-invariant operators are independent of $`\xi `$. We observe that, if we would take the limit $`M\mathrm{}`$ before removing the cutoff, our argument constitutes an alternative derivation of the equivalence of Eq. (4) to Landau gauge given in Ref. . Note that the standard way of gauge fixing employed in lattice QCD computations is formally equivalent to this limit. What is new here, is that we do not take the limit $`M\mathrm{}`$ first, but keep it at the order of the cutoff. Nevertheless, the parameter $`M`$ has disappeared from Eq. (20).
The derivation given above is valid only in perturbation theory. We assumed that the $`\theta `$\- and $`\eta `$-integrals converge, i.e. that the FP operator $`_\mu D_\mu (A)`$ has only positive eigenvalues. This is not in general the case, but it is true in perturbation theory.
If we use a regulator with a hard cutoff, such as the lattice, additional subtractions will be needed in order to remove power-like divergences, which may appear as a consequence of dropping irrelevant terms. Also note that, even though Eq. (14) is linear, in general the invariant measure written in terms of the Lie-algebra valued fields is nonlinear for such a regulator, and it would therefore change under this transformation. However, this nonlinearity is proportional to the coupling constant $`g`$, and therefore the effects of this shift are of order $`g/M`$.
3. In this section, we will discuss in more detail what we mean by “equivalent.” It is clear that, in general, correlation functions of the form $`A_\mu ^a(x)A_\nu ^b(y)\mathrm{}`$ are not the same in the JLPZ and FP versions. A trivial example is the $`O(g^0)`$ two-point function $`A_\mu (p)A_\nu (q)=\delta (p+q)G_{\mu \nu }(p)`$, with $`G_{\mu \nu }(p)`$ equal to
$`G_{\mu \nu }^{\mathrm{JLPZ}}(p)`$ $`=`$ $`{\displaystyle \frac{1}{p^2}}\left(\delta _{\mu \nu }{\displaystyle \frac{p_\mu p_\nu }{p^2}}\right)+{\displaystyle \frac{1}{M^2}}{\displaystyle \frac{p_\mu p_\nu }{p^2}},`$ (21)
$`G_{\mu \nu }^{\mathrm{FP}}(p)`$ $`=`$ $`{\displaystyle \frac{1}{p^2}}\left(\delta _{\mu \nu }{\displaystyle \frac{p_\mu p_\nu }{p^2}}\right)+{\displaystyle \frac{\xi }{p^2}}{\displaystyle \frac{p_\mu p_\nu }{p^2}},\xi 0,`$
in the theories defined by Eqs. (4) and (20), respectively. This is not in contradiction with the general argument presented above, because of the change of interpolating field, $`A_\mu A_{}^{}{}_{\mu }{}^{}`$.
The equality of the transverse part of the two-point function at tree level (Eq. (21)) suggests in what sense the two theories defined by Eqs. (4) and (20) are equivalent. These two theories are not identical, because of the fact that we dropped irrelevant terms in going from the JLPZ to the FP version. The correct statement is that physical quantities (which are necessarily extracted from correlation functions of gauge-invariant operators) in the JLPZ version can be mapped into those of the FP version by a finite renormalization of the bare coupling constant $`g`$. In the FP version of the theory the fact that only a coupling constant renormalization is needed follows from BRST invariance. In the JLPZ version, the same follows from the observation that JLPZ gauge fixing can be “undone” by multiplying Eq. (4) (for $`J_\mu =0`$) by a constant in the form $`𝒟_\mathrm{H}g`$, and transforming $`A_\mu A_\mu ^g`$, removing $`S_{\mathrm{ni}}`$ from the partition function. Wave-function renormalizations are not necessary for physical quantities.
In other words, renormalized perturbation theory for physical quantities is the same in both versions for a suitable definition of the renormalized coupling constant, but the relation between the renormalized and bare coupling constants is, in general, different in the FP and JLPZ versions of the theory.
4. In the case of noninvariant correlation functions, it is well known that, in the FP version of the theory, only a universal coupling-constant renormalization and multiplicative wave-function renormalizations are necessary as a consequence of BRST invariance. The situation is less clear in the JLPZ version of the theory. Some progress can be made however, by observing that a set of local and renormalizable Feynman rules can be constructed for the JLPZ partition function, Eq. (4). Start with diagonalizing the quadratic form in $`\theta `$ and $`\eta `$ in Eq. (S0.Ex2) by a shift $`\theta =\theta ^{}\eta `$. After this shift, Eq. (S0.Ex2) can be written as the sum of two parts,
$$S_{\mathrm{ni}}(A^h)=S_{\mathrm{ni}}^\eta (A^{},\eta )+S_{\mathrm{ni}}^\theta (A^{},\eta ,\theta ^{}),$$
(22)
with $`S_{\mathrm{ni}}^\theta (A^{},\eta ,\theta ^{}=0)=0`$. To order $`g/M`$ we obtain
$`S_{\mathrm{ni}}^\eta (A^{},\eta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}M^2(A_{}^{}{}_{\mu }{}^{a})^2+{\displaystyle \frac{1}{2}}gf^{abc}_\mu \eta ^aA_\mu ^b\eta ^c+O\left({\displaystyle \frac{g}{M}}\right),`$ (23)
$`S_{\mathrm{ni}}^\theta (A^{},\eta ,\theta ^{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \theta _{}^{}{}_{}{}^{a}(D_\mu (A^{})\theta ^{})^agf^{abc}_\mu \theta _{}^{}{}_{}{}^{a}A_{}^{}{}_{\mu }{}^{b}\eta ^c+O\left({\displaystyle \frac{g}{M}}\right)`$
upon using $`_\mu A_{}^{}{}_{\mu }{}^{}=0`$. Next, we define the following action (dropping the primes on $`A_\mu `$ and $`\theta `$):
$`S_{\mathrm{pert}}^{\mathrm{JLPZ}}`$ $`=`$ $`S_{\mathrm{inv}}(A+\eta /M)+S_{\mathrm{ni}}(A+\eta /M)S_{\mathrm{ni}}^\eta (A,\eta )+S_{\mathrm{ni}}^\theta (A,\eta ,\theta )`$
$`=`$ $`{\displaystyle \frac{1}{4}}(F_{\mu \nu }^a)^2+{\displaystyle \frac{1}{2}}(_\mu \eta ^a)^2{\displaystyle \frac{1}{2}}gf^{abc}_\mu \eta ^aA_\mu ^b\eta ^c`$
$`+{\displaystyle \frac{1}{2}}_\mu \theta ^a(D_\mu (A)\theta )^agf^{abc}_\mu \theta ^aA_\mu ^b\eta ^c+O\left({\displaystyle \frac{g}{M}}\right).`$
Here “$`O(g/M)`$” indicates all higher order terms in the expansion in $`g`$, including those coming from the shift Eq. (14), and they should be kept (up to the order in $`g`$ of interest). This action will give rise to a correct set of Feynman rules if we add the rule that a factor $`1`$ be applied for each connected $`\theta `$-subdiagram without external $`\theta `$ lines. This is reminiscent of the minus sign for ghost loops in the FP case: the integral over $`\theta `$ in the denominator of Eq. (4) gives rise to an effective action $`S_{\mathrm{eff}}(A,\eta )`$,
$$\mathrm{exp}(S_{\mathrm{eff}}(A,\eta ))=𝒟\theta \mathrm{exp}\left(S_{\mathrm{ni}}^\theta (A,\eta ,\theta )\right),$$
(25)
the vertices of which correspond precisely to these connected $`\theta `$-subdiagrams. The additional minus sign corresponds to the fact that this integral appears in the denominator of the integrand in Eq. (4). If we work with a regulator in which it is important to keep the nonlinear terms in the measure, this can be taken into account by treating these nonlinear terms as part of the action.
Finally, there is still the $`\delta `$-function of Eq. (10), through which the field $`\eta `$ was introduced. After the shift Eq. (14), this is just $`\delta (_\mu A_\mu )`$, which again can be represented as in Eq. (19). This leads to a set of Feynman rules for a local theory, Eq. (S0.Ex8), which is renormalizable by power counting for any value of $`\xi `$, with the original JLPZ version (Eq. (4)) corresponding to the limit $`\xi 0`$ (at fixed cutoff).
It is instructive to compare our local Feynman rules with the non-local Feynman rules of Ref. . Our Feynman rules are defined by Eq. (S0.Ex8), with the minus-sign rule for connected $`\theta `$-subdiagrams, and the $`\delta `$-function constraint on $`A_\mu `$, enforced by taking $`\xi 0`$ in Eq. (19). Note that the longitudinal part of the gauge field is represented by $`\eta `$, through Eq. (10). The Feynman rules of Ref. were derived by integrating out the field $`\theta `$, or equivalently $`h`$, in Eq. (4), and without introducing the field $`\eta `$ to represent the longitudinal part of the gauge field. This leads to non-local vertices, where the non-locality arises from the integration over $`\theta `$, as well as from the non-local relation between $`\eta `$ and the longitudinal part of $`A_\mu `$.
It is straightforward to obtain these non-local Feynman rules from our Feynman rules by (perturbatively) integrating out $`\theta `$ in Eq. (S0.Ex8), and by converting the Feynman rules for $`\eta `$ to Feynman rules for the longitudinal part $`A_{L\mu }`$ of $`A_\mu `$. This is obvious, because both our Feynman rules and those of Ref. are derived from the same partition function, Eq. (4). In momentum space, the relation between $`\eta `$ and the longitudinal part $`A_{L\mu }`$ of $`A_\mu `$ is
$$\eta (p)=iM\frac{p_\mu }{p^2}A_{L\mu }(p),$$
(26)
and from $`\eta (p)\eta (q)=\delta (p+q)/p^2`$ we thus find
$$\frac{p_\mu p_\nu }{p^2}A_{L\mu }(p)A_{L\nu }(q)=\frac{1}{M^2}\delta (p+q),$$
(27)
in accordance with the first equation in Eq. (21), as well as with Ref. . The simplest non-local vertex of Ref. corresponds to the $`\eta A\eta `$ three-point vertex of Eq. (S0.Ex8), which reads in momentum space, with $`k`$ the momentum of $`A_\mu `$ and $`p,q`$ the $`\eta `$ momenta,
$$\frac{1}{2}gf^{abc}i(p_\mu q_\mu ).$$
(28)
Using Eq. (26), this translates into a non-local three-point vertex for the gauge fields:
$$\frac{1}{6}M^2gf^{abc}\left[i(p_\mu q_\mu )\frac{p_\kappa q_\lambda }{p^2q^2}+i(q_\kappa k_\kappa )\frac{q_\lambda k_\mu }{q^2k^2}+i(k_\lambda p_\lambda )\frac{k_\mu p_\kappa }{k^2p^2}\right],$$
(29)
where we symmetrized in the three gauge-field lines. This is exactly the non-local three-point gauge-field vertex of Ref. . The other non-local vertices of the Feynman rules of Ref. can be obtained in a similar way.
It follows that, for correlation functions involving only gauge fields on the external lines, one obtains the same result from either our Feynman rules or those of Ref. . This is true in particular for the one-loop vacuum polarization, which was calculated explicitly in Ref. . The result found there agrees with the vacuum polarization in Landau gauge, i.e. as calculated from Eq. (20).
We conclude this section by noting that other field redefinitions of the field $`\theta `$ can be used, for instance
$$\mathrm{exp}\left(i\frac{g}{M}\theta \right)=\mathrm{exp}\left(i\frac{g}{M}\theta ^{}\right)\mathrm{exp}\left(i\frac{g}{M}\eta \right),$$
(30)
or
$$\mathrm{exp}\left(i\frac{g}{M}\theta \right)=\mathrm{exp}\left(i\frac{g}{M}\theta ^{}\right)\left(1i\frac{g}{M}\eta \right).$$
(31)
These examples differ only by terms of order $`g/M`$ from the one employed above, and thus lead to a different specific form of the $`O(g/M)`$ terms in Eq. (S0.Ex8).
5. We argued that the gauge-fixing procedure proposed in Refs. , with the choice of Eq. (6) for $`S_{\mathrm{ni}}(A)`$, is perturbatively equivalent to the standard gauge-fixing procedure with Fadeev–Popov ghosts. In our derivation, we chose $`M`$ to be of the order of the cutoff, thus extending earlier arguments in which the (formal) limit $`M\mathrm{}`$ was considered.
“Equivalent” here means that perturbatively calculated relations between physical quantities will be the same in both versions of the theory. Since, in addition, the JLPZ method leads to a weighting of the integration over orbits which takes Gribov copies correctly into account, this method may be the preferred one for nonperturbative calculations in gauge-fixed Yang–Mills theories.
We also derived local Feynman rules for the JLPZ version of the theory, Eq. (4), from which it can be seen that the theory is renormalizable by power counting. The choice of $`M`$ at the order of the cutoff is a key ingredient here. By construction, correlation functions of gauge-invariant operators are the same when calculated perturbatively from either Eq. (4) or Eq. (S0.Ex8), after the field redefinition of Eq. (14) is taken into account. In order to renormalize correlation functions of gauge noninvariant operators, counterterms may have to be added to Eq. (S0.Ex8), in addition to those needed to renormalize gauge-invariant quantities. The locality guarantees that all counterterms necessary for renormalization are local, and the renormalizability guarantees that only a finite number, all with mass dimension less than or equal to four, will be needed. It is also clear that this can be done in such a way that the invariance of gauge-invariant correlation functions is maintained, because of the gauge invariance of the original formulation of Eq. (4). Hence, we believe that no problems will be encountered in carrying out this program order by order in perturbation theory, for the theory of Eq. (S0.Ex8). But it is not clear how this would then “translate back” to a JLPZ-like formulation as in Eq. (4). What is lacking is a tool similar to BRST symmetry in the FP version of the theory, which could be used to further control the form of the counterterms. It would be interesting and useful if such a mechanism could be found.
We end with a comment on our use of the specific form of the action $`S_{\mathrm{ni}}(A)`$ in Eq. (6). This choice is the “most” (and only) relevant local operator in the sense of the renormalization group. Our analysis does not work when $`S_{\mathrm{ni}}(A)`$ is chosen to be a marginal operator. We expect that they will work if we would add a marginal operator to $`S_{\mathrm{ni}}(A)`$ of Eq. (6). For instance, a term of the form $`c_\mu \mathrm{tr}(A_\mu ^4)`$, with $`c`$ a constant of order $`g^2`$ (which is natural on the lattice), can be removed by a field redefinition of the form $`A_\mu A_\mu +(c/2)A_\mu ^3/M^2`$. Since the nonlinear term of this field redefinition is of order $`1/M^2`$, this will just remove the term $`_\mu \mathrm{tr}(A_\mu ^4)`$, without introducing any other marginal terms.
Acknowledgements
MG would like to thank Giancarlo Rossi, Massimo Testa and Arjan van der Sijs for useful discussions, and the Physics Department of the University of Rome II “Tor Vergata” for hospitality. MG and MO are supported in part by the US Department of Energy, and YS is supported in part by the Israel Academy of Science.
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# Status of the solution to the solar neutrino problem based on non-standard neutrino interactions
## I Introduction
In recent years the accuracy with which the solar neutrino flux is being measured has been improved significantly . Better statistics and calibration of the pioneering experiments, as well as the first next-generation experiment SuperKamiokande, measuring the solar neutrino spectrum and the event rate as a function of the zenith angle with unprecedented precision, have provided a lot of new information about the solar neutrino problem . On the theoretical side several substantial improvements have been made in the standard solar model (SSM) which now includes diffusion of helium and heavy elements and updated low energy nuclear cross sections relevant to the solar neutrino production . Furthermore, the SSM has received an important independent confirmation by the excellent agreement between its predicted sound speeds and recent helioseismological observations .
All five solar neutrino experiments observe a solar neutrino flux which is smaller than predicted by the SSMs. In order to understand this discrepancy it has been suggested that neutrinos are endowed with properties which are not present in the standard electroweak theory . These new properties allow the electron neutrinos to be converted along their way from the center of the sun to the detectors on earth into different neutrino flavors, i.e. into muon, tau, or possibly sterile neutrinos. The fact that the terrestrial experiments are less sensitive to these neutrino flavors explains the observed lower counting rates. The most plausible solution is that neutrinos are massive and there is mixing in the lepton sector. Then neutrino oscillations in vacuum or matter (the Mikheyev-Smirnov-Wolfenstein (MSW) effect) can explain the deficit of observed neutrinos with respect to the predictions of the SSM .
In his seminal paper Wolfenstein observed that non-standard neutrino interactions (NSNI) with matter can also generate neutrino oscillations. In particular this mechanism could be relevant to solar neutrinos interacting with the dense solar matter along their path from the core of the sun to its surface . In this case the flavor changing neutrino interactions (FCNI) are responsible for the off-diagonal elements in the neutrino propagation matrix (similar to the $`\mathrm{\Delta }m^2\mathrm{sin}^22\theta `$ term induced by vacuum mixing). For massless neutrinos resonantly enhanced conversions can occur due to an interplay between the standard electroweak neutrino interactions and non-universal flavor diagonal neutrino interactions (FDNI) with matter .
While many extensions of the standard model allow for massive neutrinos, it is important to stress that also many New Physics models predict new neutrino interactions. The minimal supersymmetric standard model without $`R`$-parity has been evoked as an explicit model that could provide the FCNI and FDNI needed for this mechanism. Systematic studies of the data demonstrated that resonantly enhanced oscillations induced by FCNI and FDNI for massless neutrinos , or FCNI in combination with massive neutrinos can solve the solar neutrino problem.
In this paper we investigate the current status of the solution to the solar neutrino problem based on NSNI, which is briefly reviewed in section II. In the first part of our study (section III) we present a comprehensive statistical analysis of this solution. Our analysis comprises both the measured total rates of Homestake , GALLEX , SAGE and SuperKamiokande and, for the first time in the context of NSNI, the full SuperKamiokande data set (corresponding to 825 effective days of operation) including the recoil electron spectrum and the day-night asymmetry. We have not included in our $`\chi ^2`$ analysis the seasonal variation but we will comment on this effect. For the solar input we take the solar neutrino fluxes and their uncertainties as predicted in the standard solar model by Bahcall and Pinsonneault (hereafter BP98 SSM) . The BP98 SSM includes helium and heavy elements diffusion, as well as the new recommended value for the low energy $`S`$-factor, $`S_{17}=19_2^{+4}`$ eV b. We also study the dependence of the allowed parameter space on the high energy <sup>8</sup>B neutrino flux, by varying the flux normalization as a free parameter.
In the second part of our study for the first time a systematic, model-independent investigation of the phenomenological constraints on FCNI and new non-universal FDNI relevant for solar neutrinos is presented (section IV). Our two main goals are: a) to find out whether NSNI can be sufficiently large to provide a viable solution to the solar neutrino problem, and b) to study various kinds of new interactions in order to single out those New Physics models that can provide such interactions. Since the typical energy scales relevant for solar neutrinos are lower than the weak interaction scale and therefore lower than any New Physics scale, it is sufficient to discuss the effective operators induced by heavy boson exchange that allow for non-standard neutrino scattering off quarks or electrons. These operators are related by the $`SU(2)_L`$ symmetry of the standard electroweak theory to operators that induce anomalous contributions to leptonic decays. Since $`SU(2)_L`$ violation cannot be large for New Physics at or above the weak scale, one can use the upper bounds on lepton flavor violating decays or on lepton universality violation to put model-independent bounds on the relevant non-standard neutrino interactions.
We find that non-standard neutrino interactions can provide a good fit to the solar neutrino data if there are rather large non-universal FDNI (of order $`0.5G_F`$) and small FCNI (of order a few times $`10^3G_F`$). Our phenomenological analysis indicates that FCNI could only be large enough to provide $`\nu _e\nu _\tau `$ transitions, while $`\nu _e\nu _\mu `$ transitions are not relevant for the solution of the solar neutrino problem, because of strong experimental constraints. Large FDNI can only be induced by an intermediate doublet of $`SU(2)_L`$ (a scalar or a vector boson) or by a neutral vector singlet. We conclude that the minimal supersymmetric model with broken $`R`$-parity is the favorite model for this scenario.
In section V we discuss how to to confirm or exclude the solution to the solar neutrino problem based on non-standard neutrino interactions by future experiments. We argue that the magnitudes of FCNI parameters necessary for $`\nu _e\nu _\tau `$ conversion in the sun could be tested independently by the upcoming $`B`$-factories. Finally, we discuss briefly the possibility of distinguishing this solution from the others by future solar and long-baseline neutrino oscillation experiments.
## II Neutrino flavor conversion induced by non-standard neutrino interactions
Any model beyond the standard electroweak theory that gives rise to the processes
$`\nu _ef`$ $``$ $`\nu _{\mathrm{}}f,`$ (1)
$`\nu _\alpha f`$ $``$ $`\nu _\alpha f,`$ (2)
where (here and below) $`f=u,d,e`$ and $`\mathrm{}=\mu ,\tau `$ and $`\alpha =e,\mu ,\tau `$, is potentially relevant for neutrino oscillations in the sun, since these processes modify the effective mass of neutrinos propagating in dense matter.
The evolution equations for massless neutrinos that interact with matter via the standard weak interactions and the non-standard interactions in (1) and (2) is given by :
$$i\frac{d}{dr}\left(\begin{array}{c}A_e(r)\\ A_{\mathrm{}}(r)\end{array}\right)=\sqrt{2}G_F\left(\begin{array}{cc}n_e(r)& ϵ_\nu _{\mathrm{}}^fn_f(r)\\ ϵ_\nu _{\mathrm{}}^fn_f(r)& ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}n_f(r)\end{array}\right)\left(\begin{array}{c}A_e(r)\\ A_{\mathrm{}}(r)\end{array}\right),$$
(3)
where $`A_e(r)`$ and $`A_{\mathrm{}}(r)`$ are, respectively, the probability amplitudes to detect a $`\nu _e`$ and $`\nu _{\mathrm{}}`$ at position $`r`$. For neutrinos that have been coherently produced as $`\nu _e`$ in the solar core at position $`r_0`$, the equations in (3) are subject to the boundary conditions $`A_e(r_0)=1`$ and $`A_{\mathrm{}}(r_0)=0`$. While $`W`$-exchange of $`\nu _e`$ with the background electrons gives rise to the well known forward scattering amplitude $`\sqrt{2}G_Fn_e(r)`$, the FCNI in (1) induce a flavor changing forward scattering amplitude $`\sqrt{2}G_Fϵ_\nu _{\mathrm{}}^fn_f(r)`$ and the non-universal FDNI are responsible for the flavor diagonal entry $`\sqrt{2}G_Fϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}n_f(r)`$ in eq. (3). Here
$$n_f(r)=\{\begin{array}{cc}n_n(r)+2n_p(r)& f=u\\ 2n_n(r)+n_p(r)& f=d\end{array}$$
(4)
is the respective fermion number density at position $`r`$ in terms of the proton \[neutron\] number density $`n_p(r)`$ \[$`n_n(r)`$\] and
$$\epsilon =ϵ_\nu _{\mathrm{}}^f\frac{G_{\nu _e\nu _{\mathrm{}}}^f}{G_F}\text{and}\epsilon ^{}=ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}\frac{G_{\nu _{\mathrm{}}\nu _{\mathrm{}}}^fG_{\nu _e\nu _e}^f}{G_F},$$
(5)
describe, respectively, the relative strength of the FCNI in (1), and the new flavor diagonal, but non-universal interactions in (2). $`G_{\nu _\alpha \nu _\beta }^f`$ ($`\alpha ,\beta =e,\mu ,\tau `$) denotes the effective coupling of the four-fermion operator
$$𝒪_\nu ^f(\overline{\nu _\alpha }\nu _\beta )(\overline{f}f).$$
(6)
that gives rise to such interactions. The Lorentz structure of $`𝒪_\nu ^f`$ depends on the New Physics that induces this operator. Operators which involve only left-handed neutrinos (and which conserve total lepton number $`L`$) can be decomposed into a $`(VA)(VA)`$ and a $`(VA)(V+A)`$ component. (Any single New Physics contribution that is induced by chiral interactions yields only one of these two components.) It is, however, important to note that only the vector part of the background fermion current affects the neutrino propagation for an unpolarized medium at rest . Hence only the $`(VA)(V)`$ part of $`𝒪_\nu ^f`$ is relevant for neutrino oscillations in normal matter. One mechanism to induce such operators is due to the exchange of heavy bosons that appear in various extensions of the standard model. An alternative mechanism arises when extending the fermionic sector of the standard model and is due to $`Z`$-induced flavor-changing neutral currents (FCNCs). For a discussion of $`Z`$-induced FCNC effects on solar neutrinos, see Refs. .
A resonance occurs when the diagonal entries of the evolution matrix in eq. (3) coincide at some point $`r_{res}`$ along the trajectory of the neutrino, leading to the resonance condition
$$ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}n_f(r_{res})=n_e(r_{res}).$$
(7)
An immediate consequence is that new FDNI for $`f=e`$ alone cannot induce resonant neutrino flavor conversions.
As we will see in section IV only $`\nu _e\nu _\tau `$ conversions are compatible with the existing phenomenological constraints on $`ϵ_\nu _{\mathrm{}}^f`$ and $`ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}`$. We note that in the minimal supersymmetric standard model with broken $`R`$-parity the relevant parameters are given by
$$ϵ_{\nu _\tau }^d=\frac{\lambda _{331}^{}{}_{}{}^{}\lambda _{131}^{}}{4M_{\stackrel{~}{b}}^2\sqrt{2}G_F},\text{and}ϵ_{}^{}{}_{\nu _\tau }{}^{d}=\frac{|\lambda _{331}^{}|^2|\lambda _{131}^{}|^2}{4M_{\stackrel{~}{b}}^2\sqrt{2}G_F},$$
(8)
in terms of the trilinear couplings $`\lambda _{ijk}^{}`$ and the bottom squark mass $`M_{\stackrel{~}{b}}`$.
The neutrino evolution matrix in eq. (3) vanishes in vacuum and is negligibly small for the matter densities of the earth’s atmosphere. Therefore the probability of finding an electron neutrino arriving at the detector during day time is easily obtained by evolving the equations in (3) from the neutrino production point to the solar surface. Furthermore, typically there are many oscillations between the neutrino production and detection point and a resonance. Therefore the phase information before and after the resonance is usually lost after integration over the production and detection region and one may use classical survival probabilities. Then at day time we have
$$P_{\nu _e\nu _e}^{\mathrm{day}}=|A_e(r_s)|^2\frac{1}{2}+(\frac{1}{2}P_c)\mathrm{cos}2\theta _m^p\mathrm{cos}2\theta _m^s,$$
(9)
where $`r_s`$ is the solar surface position and in the analytic expression in eq. (9) we denote by $`\theta _m^p`$ and $`\theta _m^s`$, respectively, the effective, matter-induced mixing at the neutrino production point and at the solar surface. In terms of the New Physics parameters $`\epsilon `$, $`\epsilon ^{}`$ and the fermion densities the effective mixing is given by
$$\mathrm{tan}2\theta _m=\frac{2ϵ_\nu _{\mathrm{}}^fn_f}{ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}n_fn_e}.$$
(10)
Note that $`\mathrm{tan}2\theta _m=2ϵ_\nu _{\mathrm{}}^e/(ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{e}1)`$ is constant for $`f=e`$. $`P_c`$ is the level crossing probability. The approximate Landau-Zener expression is
$$P_c=\mathrm{exp}\left[\pi \gamma /2\right]\mathrm{with}\gamma =4\sqrt{2}G_F\left|\frac{(ϵ_\nu _{\mathrm{}}^f/ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f})^2}{ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}}\frac{n_e}{\frac{d}{dx}\left(\frac{n_f}{n_e}\right)}\right|_{res}.$$
(11)
When neutrinos arrive at the detector during the night, a modification of the survival probability has to be introduced since the non-standard neutrino interactions with the terrestrial matter may regenerate electron neutrinos that have been transformed in the sun. Assuming that the neutrinos reach the Earth as an incoherent mixture of the effective mass-eigenstates $`\nu _1`$ and $`\nu _2`$ the survival probability during night-time can be written as :
$$P_{\nu _e\nu _e}^{\mathrm{night}}=\frac{P_{\nu _e\nu _e}^{\mathrm{day}}\mathrm{sin}^2\theta _m^s+P_{2e}(12P_{\nu _e\nu _e}^{\mathrm{day}})}{\mathrm{cos}2\theta _m^s}.$$
(12)
Here $`P_{2e}`$ is the probability of a transition from the state $`\nu _2`$ to the flavor eigenstate $`\nu _e`$ along the neutrino path in the Earth.
For our analysis we assume a step function profile for the Earth matter density, which has been shown to be a good approximation in other contexts (see e.g. Ref. for a recent analysis of matter effects for atmospheric neutrinos). Then the earth matter effects on the neutrino propagation correspond to a parametric resonance and can be calculated analytically ,
$$P_{2e}=\mathrm{sin}^2\theta _m^s+W_1^2\mathrm{cos}2\theta _m^s+W_1W_3\mathrm{sin}2\theta _m^s,$$
(13)
where the parameters $`W_1`$ and $`W_3`$ contain all the information of the Earth density and are defined in Ref. . (The only difference is that in our case also the off-diagonal element of the neutrino evolution matrix varies when the neutrino propagates through the earth matter.)
It is this interaction with the terrestrial matter that can produce a day-night variation of the solar neutrino flux and, consequently, a seasonal modulation of the data. (Note that this seasonal variation is of a different nature than the one expected for vacuum oscillations from the change of the baseline due to the eccentricity of the earth’s orbit around the sun.)
## III Analysis of the solar neutrino data
In this section we present our analysis of the solution to the solar neutrino problem based on neutrino flavor conversions induced by NSNI in matter. Our main goal is to determine the values of $`\epsilon `$ and $`\epsilon ^{}`$ that can explain the experimental observations without modifying the standard solar model predictions.
### A Rates
First we consider the data on the total event rate measured by the Chlorine (Cl) experiment , the Gallium (Ga) detectors GALLEX and SAGE and the water Cherenkov experiment SuperKamiokande (SK) . We compute the allowed regions in parameter space according to the BP98 SSM and compare the results with the regions obtained for an arbitrary normalization $`f_B`$ of the high energy neutrino <sup>8</sup>B neutrino fluxes.
We use the minimal $`\chi ^2`$ statistical treatment of the data following the analyses of Refs. . Our $`\chi ^2`$ function is defined as follows:
$$\chi _R^2(\epsilon ,\epsilon ^{}[,f_B])=\underset{i,j=1,\mathrm{},4}{}(R_i^{th}(\epsilon ,\epsilon ^{}[,f_B])R_i^{obs})\left[\sigma _R^2\right]_{ij}^1(R_j^{th}(\epsilon ,\epsilon ^{}[,f_B])R_j^{obs}),$$
(14)
where $`R_i^{th}`$ and $`R_i^{obs}`$ denote, respectively, the predicted and the measured value for the event rates of the four solar experiments ($`i`$ = Cl, GALLEX, SAGE, SK). The error matrix $`\sigma _R`$ contains both the experimental (systematic and statistical) and the theoretical errors.
In Fig. 1 the allowed regions in the parameter space of $`ϵ_\nu ^d`$ and $`ϵ_{}^{}{}_{\nu }{}^{d}`$ for neutrino scattering off $`d`$-quarks are shown at 90, 95 and 99 % confidence level (CL). In Fig. 1a, the <sup>8</sup>B flux is fixed by the BP98 SSM prediction ($`f_B=1`$). The best fit point of this analysis is found at
$$ϵ_\nu ^d=3.2\times 10^3\text{and}ϵ_{}^{}{}_{\nu }{}^{d}=0.61,$$
(15)
with $`\chi _{min}^2=2.44`$ for $`42=2`$ degrees of freedom (DOF). Allowing an arbitrary <sup>8</sup>B flux normalization, a different best fit point is obtained for $`(ϵ_\nu ^d,ϵ_{}^{}{}_{\nu }{}^{d})=(2.2\times 10^2,0.59)`$ and $`f_B=1.36`$ with $`\chi _{min}^2=0.91`$ for $`43=1`$ DOF. The result of this analysis is shown in Fig. 1b. (Effects due to deviations of the hep neutrino flux from the standard solar model prediction are expected to be less significant and we do not consider them in this work.)
In Fig. 2 the allowed regions in the parameter space of $`ϵ_\nu ^u`$ and $`ϵ_{}^{}{}_{\nu }{}^{u}`$ for neutrino scattering off $`u`$-quarks are shown at 90, 95 and 99 % CL. In Fig. 2a, the <sup>8</sup>B flux is fixed by the BP98 SSM prediction. The best fit point of this analysis is found at
$$ϵ_\nu ^u=1.32\times 10^3\text{and}ϵ_{}^{}{}_{\nu }{}^{u}=0.43,$$
(16)
with $`\chi _{min}^2=2.64`$ for two DOF. Allowing an arbitrary <sup>8</sup>B flux normalization $`f_B`$, a different best fit point is obtained for $`(ϵ_\nu ^u,ϵ_{}^{}{}_{\nu }{}^{u})=(5.8\times 10^3,0.425)`$ and $`f_B=1.34`$ with $`\chi _{min}^2=0.96`$ for one DOF. The result of this analysis is shown in Fig. 2b.
It is remarkable that the neutrino flavor conversion mechanism based on NSNI provides quite a good fit to the total rates despite the fact that the conversion probabilities (9) and (12) do not depend on the neutrino energy. This is unlike the case of the vacuum and the MSW conversion mechanisms which provide the appropriate energy dependence to yield a good fit. For NSNI the only way to distinguish between neutrinos of different energies is via the position of the resonance $`r_{res}`$. Note that according to eq. (7), $`r_{res}`$ is a function of $`\epsilon ^{}`$ only. As can be seen in Fig. 3, $`n_e/n_f`$ ($`f=d,u`$) is a smooth and monotonic function of the distance from the solar center $`r`$, allowing to uniquely determine $`r_{res}`$ for a given value of $`\epsilon ^{}`$. From Fig. 3 it follows that a resonance can only occur if $`ϵ_d^{}[0.50,0.77]`$ for NSNI with $`d`$-quarks or $`ϵ_u^{}[0.40,0.46]`$ for NSNI with $`u`$-quarks. For both cases the major part of these intervals corresponds to $`r_{res}<0.2R_{}`$ ($`R_{}`$ being the solar radius). For $`ϵ_{d,u}^{}`$ within the 90 % CL regions (indicated in Fig. 1 and Fig. 2) we find $`r_{res}0.1R_{}`$. Since the nuclear reactions that produce neutrinos with higher energies in general take place closer to the solar center (see chapter 6 of Ref. for the various spatial distributions of the neutrino production reactions), a resonance position close to the solar center implies that predominantly the high energy neutrinos are converted by a resonant transition. For $`r_{res}0.1R_{}`$ practically all <sup>8</sup>B-neutrinos cross the resonance layer, fewer <sup>7</sup>Be-neutrinos pass through the resonance, while most of the $`pp`$-neutrinos are not be affected by the resonance since their production region extends well beyond the resonance layer. Therefore for most of the allowed region in Fig. 1 and Fig. 2 we have
$$P(^8\text{B})<P(^7\text{Be})<P(pp).$$
(17)
We note that the above relation is still valid when taking into account that a significant fraction of the $`pp`$ neutrinos crosses the resonance layer twice, if they are produced just outside resonance. This is – roughly speaking – because a $`\nu _e`$ which undergoes a resonant flavor transition when entering the solar interior at $`r_{res}`$ is reconverted into a $`\nu _e`$ at the second resonance when it emerges again from the solar core. In our numerical calculations we properly take into account the effects of such double resonances.
An immediate consequence of the relation in eq. (17) is that as long as $`f_B=1`$ the NSNI solution predicts that $`R_{SK}<R_{Cl}<R_{Ga}`$, which is inconsistent with the observed hierarchy of the rates, $`R_{Cl}<R_{SK}<R_{Ga}`$, leading to a somewhat worse fit than the standard MSW solutions. However when treating $`f_B`$ as a free parameter, for $`f_B1.31.4`$ the SK rate is sufficiently enhanced to give the correct relation between the rates. In this case also the neutral current contribution from $`\nu _{\mu ,\tau }e^{}`$ scattering is increased due to a larger $`\nu _{\mu ,\tau }`$ flux, which is consistent with the Super-Kamiokande observations. We find that for the best fit points for $`(\epsilon ,\epsilon ^{})`$ in Figs. 1b and 2b and $`f_B1.35`$ the survival probability for <sup>8</sup>B, <sup>7</sup>Be and $`pp`$-neutrinos are $`0.24`$, $`0.4`$ and $`0.7`$, respectively.
In Fig. 1b and Fig. 2b the <sup>8</sup>B neutrino flux normalization $`f_B`$ has been varied as a free parameter in order to study the dependence of the allowed parameter space on the high energy neutrino flux. It is interesting to note that the allowed regions in these figures do not completely contain those in Fig. 1a and Fig. 2a where the boron flux and its uncertainty is determined by the BP98 SSM. In order to explain this apparently inconsistent result we have have plotted $`\chi ^2`$ as a function of $`f_B`$ in Fig. 4 allowing $`f_B`$ to vary within a sufficiently broad interval ($`0<f_B<100`$) for every point in the $`(\epsilon ,\epsilon ^{})`$ parameter space. The horizontal lines indicate the 68, 90 and 99 % CL limits for two DOF. The intersection of these lines with the $`\chi ^2`$ curve determine the relevant ranges of the boron flux allowed by the experimental data. The vertical dotted lines indicate the $`1\sigma `$ and $`3\sigma `$ ranges of the boron neutrino flux in the BP98 SSM.
Note that the $`\chi ^2`$ minima are obtained for a boron flux significantly larger ($`f_B1.35`$) than the one predicted by the SSM ($`f_B=1.0`$), as we already anticipated in the discussion of eq. (17). Moreover, from Fig. 4 it follows that for $`f_B<1`$ the fit to the experimental data imposes stronger constraints on the boron flux than the SSM. For $`f_B>1`$ the situation is exactly the opposite. Therefore the effect of relaxing $`f_B`$ from its SSM value is that regions in the $`(\epsilon ,\epsilon ^{})`$ parameter space where the averaged survival probability for <sup>8</sup>B neutrinos, $`P(^8\text{B})`$ is smaller can be easily compensated by a larger boron flux and obtain a lower value for $`\chi _R^2`$. On the other hand regions where $`P(^8\text{B})`$ is rather large require a small boron flux which is more difficult to achieve when eliminating the SSM constraint on $`f_B`$. This is the main mechanism behind the changes of the allowed regions upon relaxing the SSM constraint on $`f_B`$. It explains why the regions with large $`\epsilon `$ are allowed in Fig. 1a and Fig. 2a and are ruled out in Fig. 1b and Fig. 2b. Here $`P(^8\text{B})`$ is rather large and a small boron flux $`f_B1`$ like in the SSM is preferred to explain the data. The opposite occurs in the area between the two disconnected regions in Fig. 1a and Fig. 2a. Here $`P(^8\text{B})`$ is comparatively small and therefore a larger boron flux increases $`\chi _R^2`$ in this region leading to the merging of the separated contours in Fig. 1b and Fig. 2b when $`f_B`$ is treated as a free parameter.
### B Zenith Angle Data
Next, we consider the zenith angle dependence of the solar neutrino data of the SuperKamiokande experiment. As mentioned above, NSNI with matter may affect the neutrino propagation through the earth resulting in a difference between the event rates during day and night time. The data obtained by the SuperKamiokande collaboration are divided into five bins containing the events observed at night and one bin for the events collected during the day and have been averaged over the period of SuperKamiokande operation: 403.2 effective days for the day events and 421.5 effective days for the night events. The experimental results suggest an asymmetry between the total data collected during the day ($`D`$) and the total data observed during the night ($`N`$:
$$A=2\frac{ND}{N+D}=0.065\pm 0.031(stat.)\pm 0.013(syst.).$$
(18)
In order to take into account the earth matter effect we define the following $`\chi ^2`$-function that characterizes the deviations of the six measured ($`Z_i^{obs}`$) from the predicted ($`Z_i^{th}`$) values of the rate as a function of zenith angle:
$$\chi _Z^2(\epsilon ,\epsilon ^{},\alpha _Z)=\underset{i=1,\mathrm{},6}{}\frac{\left(\alpha _ZZ_i^{th}(\epsilon ,\epsilon ^{})Z_i^{obs}\right)^2}{\sigma _{Z,i}^2}.$$
(19)
Here $`\sigma _{Z,i}`$ refers to the total error associated with each zenith angle bin and we have neglected possible correlations between the systematic errors of these bins. Since we are only interested in the shape of the zenith angle distribution, we have introduced an overall normalization factor, $`\alpha _Z`$, which is treated as a free parameter and determined from the fit. (Using this procedure also prevents over-counting the data on the total event rate when combining all available data in section III E). Note that the experimental value of the day-night asymmetry in eq. (18) is not used in the fit, since the six zenith angle bins already include consistently all the available information about the earth effects.
In Fig. 5 we show the allowed regions in the $`(\epsilon ,\epsilon ^{})`$ parameter space for neutrino scattering off $`d`$\- and $`u`$-quarks, respectively. The contours in Fig. 5 correspond to the allowed regions at 90, 95 and 99 % CL. The best fit (indicated by the open circle) is obtained for $`(ϵ_\nu ^d,ϵ_{}^{}{}_{\nu }{}^{d})=(0.251,0.62`$) and $`\alpha _Z=0.819`$ with $`\chi _{min}^2=1.10`$ for neutrino scattering off $`d`$-quarks and at $`(ϵ_\nu ^u,ϵ_{}^{}{}_{\nu }{}^{u})=(0.229,0.690`$) and $`\alpha _Z=0.685`$ with $`\chi _{min}^2=1.44`$ for neutrino scattering off $`u`$-quarks (having $`63=3`$ DOF in both cases).
Finally, in Fig. 6 we show the expected zenith angle distributions for SuperKamiokande using the values of $`(ϵ,ϵ^{})`$ determined by the best fit. For comparison, we also present in this figure the expected zenith angle distributions for the best fit values of $`(ϵ,ϵ^{})`$ found in the combined analysis (that will be discussed in section III E).
### C Recoil Electron Spectrum
We also consider the measurements of the recoil electron spectrum by SuperKamiokande . The available data, after 825 days of operation, are divided into 18 bins. 17 of these bins have a width of 0.5 MeV and are grouped into two bins for a super low energy analysis with energies between 5.5 MeV and 6.5 MeV and 15 bins with energies ranging from 6.5 MeV (the low energy limit) to 14 MeV. The last bin includes all the events with energies larger than 14 MeV.
Since the electron neutrino survival probability does not depend on the neutrino energy in the NSNI scenario, the spectral distortion of the recoil electrons from <sup>8</sup>B neutrino due to the presence of a $`\nu _{\mu ,\tau }`$ component in the neutrino flux is expected to be very small and therefore, even a relatively small spectral distortion (such as the one expected in small mixing angle MSW solution) could rule out this solution.
The $`\chi ^2`$-function that characterizes the deviations of the measured ($`S_i^{obs}`$) from the predicted ($`S_i^{th}`$) values for the electron recoil spectrum therefore provides an important test of the NSNI solution. It is defined as:
$$\chi _S^2(\epsilon ,\epsilon ^{},\alpha _S)=\underset{i,j=1,\mathrm{},18}{}\left(\alpha _SS_i^{th}(\epsilon ,\epsilon ^{})S_i^{obs}\right)\left[\sigma _S^2\right]_{ij}^1\left(\alpha _SS_j^{th}(\epsilon ,\epsilon ^{})S_j^{obs}\right),$$
(20)
where the error matrix (squared)
$$\left[\sigma _S^2\right]_{ij}=\delta _{ij}[\sigma _i^2(stat.)+\sigma _i^2(uncorr.)]+\sigma _i(corr.)\sigma _j(corr.)+\sigma _i(theor.)\sigma _j(theor.)$$
(21)
includes statistical \[$`\sigma _i(stat.)`$\] and systematic experimental errors (including both the uncorrelated \[$`\sigma _i(uncorr.)`$\] and the correlated \[$`\sigma _i(corr.)`$\] contributions) as well as the theoretical errors \[$`\sigma _i(theor.)`$\] (see Refs. for more details.). Again, as in the analysis for the zenith angle dependence, we introduce an overall normalization factor $`\alpha _S`$, which is taken as a free parameter and determined from the fit, in order to avoid over-counting the data on the total event rate. Fitting the present data to our scenario we obtain $`\chi _{min}^2=20.0`$ for $`181=17`$ DOF, which is still acceptable at the 27 % CL.
### D Seasonal Variations
The earth matter effects on neutrino flavor transitions induce a seasonal variation of the data (beyond the expected variation of the solar neutrino flux due to the eccentricity of the earth’s orbit) due to the variation of the day and night time during the year. Since these variations can be relevant to other neutrino oscillation scenarios , a positive signal could help to distinguish the various solutions and it is worthwhile to analyze the effects of such a variation in the NSNI scenario.
The present SK solar neutrino data do not provide any conclusive evidence in favor of such a variation, but indicate only that the variation seems to be larger for recoil electron energies above 11.5 MeV. In our scenario, however, we do not expect any correlation between the seasonal variation and the recoil electron energies, since the electron neutrino survival probability does not depend on the neutrino energy. Therefore any range of parameters that leads to a considerable seasonal modulation for energies above 11.5 MeV is disfavored by the data for lower energies. However, for the range of parameters ($`\epsilon ,\epsilon ^{}`$) that can solve the solar neutrino problem, earth regeneration effects are never strong enough to induce a significant seasonal variation. Hence taking into account the data on seasonal variations neither changes the shape of the allowed region, nor the best fit points.
### E Combined Analysis
Our final result is the fit derived from the combined analysis of all presently available solar neutrino data. In Fig. 7 and Fig. 8 we show the allowed regions for $`(ϵ_\nu ^d,ϵ_{}^{}{}_{\nu }{}^{d})`$ and $`(ϵ_\nu ^u,ϵ_{}^{}{}_{\nu }{}^{u})`$, respectively, using both the results from the total rates from the Chlorine, GALLEX, SAGE and SuperKamiokande solar neutrino experiments together with the 6 bins from the SuperKamiokande zenith angle data discussed previously. Although adding the spectral information to our analysis does not change the shape of allowed regions nor the best fit points, it is included in order to determine the quality of the global fit. However, we do not take into account the seasonal variation in our combined $`\chi ^2`$ analysis, since the effect is negligible.
For neutrino scattering off $`d`$-quarks the best fit for the combined data is obtained for
$$ϵ_\nu ^d=0.028\text{and}ϵ_{}^{}{}_{\nu }{}^{d}=0.585$$
(22)
with $`\chi _{min}^2=29.05`$ for $`284=24`$ DOF, corresponding to a solution at the 22 % CL (see Fig. 7a). Allowing $`f_B1`$, the best fit is found at $`(ϵ_\nu ^d,ϵ_{}^{}{}_{\nu }{}^{d})=(0.018,0.585)`$ and $`f_B=1.38`$ with $`\chi _{min}^2=26.62`$ for $`285=23`$ DOF, corresponding to a solution at the 27 % CL (see Fig. 7b). For neutrino scattering off $`u`$-quarks the best fit for the combined data is obtained for
$$ϵ_\nu ^u=0.0083\text{and}ϵ_{}^{}{}_{\nu }{}^{u}=0.425$$
(23)
with $`\chi _{min}^2=28.45`$ for $`284=24`$ DOF corresponding to a solution at the 24 % CL (see Fig. 8a). Allowing $`f_B1`$, the best fit is obtained for $`(ϵ_\nu ^u,ϵ_{}^{}{}_{\nu }{}^{u})=(0.0063,0.426)`$ and $`f_B=1.34`$ with $`\chi _{min}^2=26.59`$ for $`285=23`$ DOF, corresponding to a solution at the 27 % CL (see Fig. 8b). These results have to be compared with the fit for standard model neutrinos, that do not oscillate (where the CL is smaller than $`10^7`$), as well as to the standard solutions of the solar neutrino problem in terms of usual neutrino oscillations (36 % CL) .
Finally, in Fig. 6 we show the expected zenith angle distributions for SuperKamiokande using the best fitted values of $`(ϵ,ϵ^{})`$ from the combined analysis.
## IV Phenomenological constraints on $`ϵ`$ and $`ϵ^{}`$
In this section we investigate whether the allowed regions for the parameters $`ϵ_\nu _{\mathrm{}}^f`$ and $`ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}`$ are at all phenomenologically viable. The analysis of non-standard neutrino interactions that could be relevant for the solar neutrino problem is similar to the discussions in Refs. , where the possibility that FCNI explain the LSND results or the atmospheric neutrino anomaly was discussed.
Generically, extensions of the Standard Model include additional fields that can induce new interactions: A heavy boson $``$ that couples weakly to some fermion bilinears $`B_{ij}`$ with the trilinear couplings $`\lambda _{ij}`$, where $`i,j=1,2,3`$ refer to fermion generations, induces the four fermion operator $`B_{ij}^{}B_{kl}`$ at tree-level. The effective coupling is given by
$$G_N^{B^{}B}=\frac{\lambda _{ij}^{}\lambda _{kl}}{4\sqrt{2}M_{}^2},$$
(24)
for energies well below the boson mass $`M_{}`$. Thus, in terms of the trilinear coupling $`\lambda _{\alpha f}`$ that describes the coupling of some heavy boson $``$ to $`\nu _\alpha `$ ($`\alpha =e,\mu ,\tau `$) and a charged fermion $`f=u,d,e`$ the effective parameters in (5) are given by
$$ϵ_\nu _{\mathrm{}}^f=\frac{\lambda _\mathrm{}f^{}\lambda _{ef}}{4\sqrt{2}M_{}^2G_F}\text{and}ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}=\frac{|\lambda _\mathrm{}f|^2|\lambda _{ef}|^2}{4\sqrt{2}M_{}^2G_F}.$$
(25)
Since any viable extension of the Standard Model has to contain the SM gauge symmetry, the effective theory approach presented in Refs. is completely sufficient to describe any New Physics effect for the energy scales typical to present neutrino oscillation experiments. Even though the effective theory obviously does not contain all the information inherent in the full high-energy theory, the parameters of the effective theory are all of what is accessible at low energies, when the “heavy degrees of freedom” are integrated out.
The crucial point for our analysis is the following: Since the SM neutrinos are components of $`SU(2)_L`$ doublets, the same trilinear couplings $`\lambda _{\alpha f}`$ that give rise to non-zero $`ϵ_\nu _{\mathrm{}}^f`$ or $`ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}`$ also induce other four-fermion operators. These operators involve the $`SU(2)_L`$ partners of the neutrinos, i.e. the charged leptons, and can be used to constrain the relevant couplings.
Noting that Lorentz invariance implies that any fermionic bilinear $`B_{ij}`$ can couple to either a scalar ($`𝒮`$) or a vector ($`𝒱`$) boson it is straightforward to write down all gauge invariant trilinear couplings between the bilinears (that contain SM fermions) and arbitrary bosons $`𝒮`$ and $`𝒱`$ that might appear in a generic extension of the Standard Model (see Tabs. 1–3 of Ref. ). From these couplings one then obtains all the effective four fermion operators relevant to the solution to the solar neutrino problem in terms of NSNI as well as the $`SU(2)_L`$-related operators that are used to constrain their effective couplings. (We do not consider here operators that violate total lepton number which can be induced if there is mixing between the intermediate bosons .)
While we refer the reader to Refs. for the details of this model-independent approach, we present here two explicit examples relevant to solar neutrinos to demonstrate how $`SU(2)_L`$ related processes can be used to constrain the parameters $`ϵ_\nu _{\mathrm{}}^f`$ or $`ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}`$. First, consider the bilinear $`\overline{L}f_R`$ (where $`L`$ denotes the lepton doublet and $`f=e,u,d`$) that couples via a scalar doublet to its hermitian conjugate $`\overline{f_R}L`$. In terms of the component fields the effective interaction is
$`{\displaystyle \frac{\lambda _{\alpha f}^{}\lambda _{\beta f}}{M_1^2}}(\overline{\nu _\alpha }f_R)(\overline{f_R}\nu _\beta )+{\displaystyle \frac{\lambda _{\alpha f}^{}\lambda _{\beta f}}{M_2^2}}(\overline{l_\alpha }f_R)(\overline{f_R}l_\beta )=`$ (26)
$``$ $`{\displaystyle \frac{\lambda _{\alpha f}^{}\lambda _{\beta f}}{2M_1^2}}(\overline{\nu _\alpha }\gamma ^\mu \nu _\beta )(\overline{f_R}\gamma _\mu f_R){\displaystyle \frac{\lambda _{\alpha f}^{}\lambda _{\beta f}}{2M_2^2}}(\overline{l_\alpha }\gamma ^\mu l_\beta )(\overline{f_R}\gamma _\mu f_R),`$ (27)
where $`l_\alpha =e_L,\mu _L,\tau _L`$ for $`\alpha =e,\mu ,\tau `$. $`\lambda _{\alpha f}`$ is the trilinear coupling of $`\overline{L_\alpha }f_R`$ to the scalar doublet and $`M_{1,2}`$ denote the masses of its $`SU(2)_L`$ components. The important point is that the scalar doublet exchange not only gives rise to the four-Fermi operator $`𝒪_\nu ^f`$ in (6) (with $`(VA)(V+A)`$ structure), but also produces the $`SU(2)_L`$ related operator
$$𝒪_l^f(\overline{l_\alpha }l_\beta )(\overline{f}f),$$
(28)
which has the same Lorentz structure as $`𝒪_\nu ^f`$, with the neutrinos replaced by their charged lepton partners. Moreover, the effective coupling of $`𝒪_l^f`$, that we denote by $`G_{\alpha \beta }^f`$, is related to $`G_{\nu _\alpha \nu _\beta }^f`$ by
$$G_{\nu _\alpha \nu _\beta }^f=G_{\alpha \beta }^f\frac{M_1^2}{M_2^2}.$$
(29)
Constructing all the relevant four fermion operators that are induced by the couplings between the bilinears listed in Tabs. 1–3 of Ref. , one finds that in general $`𝒪_l^f`$ is generated together with $`𝒪_\nu ^f^{}`$. Here $`f^{}`$ can be different from $`f`$ only for interactions with quarks, that is in some cases $`𝒪_l^u`$ ($`𝒪_l^d`$) is generated together with $`𝒪_\nu ^d`$ ($`𝒪_\nu ^u`$). The leptonic operator $`𝒪_l^e`$ is always generated together with $`𝒪_\nu ^e`$ unless the interaction is mediated by an intermediate scalar $`SU(2)_L`$ singlet that couples to
$$(L_{\mathrm{}}L_e)_s=\frac{1}{\sqrt{2}}(\overline{\nu _{\mathrm{}}^c}e_L\overline{\mathrm{}_L^c}\nu _e).$$
(30)
Note that the singlet only couples between two different flavors, since the coupling has to be antisymmetric in flavor space. Consequently a singlet that couples to the bilinear $`(L_{\mathrm{}}L_e)_s`$ cannot induce a non-zero $`ϵ_\nu _{\mathrm{}}^e`$. The fact that the resulting four fermion operators only mediate FDNI is true because for the solar neutrinos we only care about $`\nu _e\nu _{\mathrm{}}`$ transitions. (For atmospheric neutrinos also $`\nu _\mu \nu _\tau `$ transitions induced by non-standard neutrino interactions with the electrons are of interest. In this case the coupling of $`(L_\mu L_e)_s`$ to $`(L_\tau L_e)_s^{}`$ via singlet exchange inducing FCNI is possible .)
The effective interactions that are mediated by a scalar singlet of mass $`M`$ that couples to $`(L_{\mathrm{}}L_e)_s`$ with the elementary coupling $`\lambda _\mathrm{}e`$ are given by
$`{\displaystyle \frac{|\lambda _\mathrm{}e|^2}{M^2}}\left[(\overline{e_L}\nu _{\mathrm{}}^c)(\overline{\nu _{\mathrm{}}^c}e_L)(\overline{e_L}\nu _{\mathrm{}}^c)(\overline{\mathrm{}_L^c}\nu _e)+(\overline{\nu _e}\mathrm{}_L^c)(\overline{\mathrm{}_L^c}\nu _e)(\overline{\nu _e}\mathrm{}_L^c)(\overline{\nu _{\mathrm{}}^c}e_L)\right]=`$ (31)
$`{\displaystyle \frac{|\lambda _\mathrm{}e|^2}{M^2}}\left[(\overline{e_L}\gamma ^\mu e_L)(\overline{\nu _{\mathrm{}}}\gamma _\mu \nu _{\mathrm{}})(\overline{e_L}\gamma ^\mu \nu _e)(\overline{\nu _{\mathrm{}}}\gamma _\mu \mathrm{}_L)+(\overline{\nu _e}\gamma ^\mu \nu _e)(\overline{\mathrm{}_L}\gamma _\mu \mathrm{}_L)(\overline{\nu _e}\gamma ^\mu e_L)(\overline{\mathrm{}_L}\gamma _\mu \nu _{\mathrm{}})\right],`$ (32)
where we used a Fierz transformation and the identity $`\overline{A^c}\gamma ^\mu B^c=\overline{B}\gamma ^\mu A`$ to obtain (32). One can see that in this case $`𝒪_\nu ^e`$ is generated together with three more operators that have the same effective coupling (up to a sign). However, unlike for the case of intermediate doublets (or triplets), all these operators involve two charged leptons and two neutrinos.
### A Experimental constraints
#### 1 Flavor changing neutrino interactions
There is no experimental evidence for any non-vanishing $`G_e\mathrm{}^f`$ . Therefore, whenever $`𝒪_l^f`$ is generated together with $`𝒪_\nu ^f`$, one can use the upper bounds on $`G_e\mathrm{}^f`$ to derive constraints on $`G_{\nu _e\nu _{\mathrm{}}}^f`$. The most stringent constraints on $`G_e\mathrm{}^e`$ are due to the upper bounds on $`\mu ^{}e^{}e^+e^{}`$ and $`\tau ^{}e^{}e^+e^{}`$ :
$`\text{BR}(\mu ^{}e^{}e^+e^{})`$ $`<`$ $`1.0\times 10^{12},`$ (33)
$`\text{BR}(\tau ^{}e^{}e^+e^{})`$ $`<`$ $`2.9\times 10^6.`$ (34)
Normalizing the above bounds to the measured rates of the related lepton flavor conserving decays, $`\text{BR}(\mu ^{}e^{}\overline{\nu }_e\nu _\mu )100\%`$ and $`\text{BR}(\tau ^{}e^{}\overline{\nu }_e\nu _\tau )=0.18`$ , we obtain
$`ϵ_\mu ^eG_{e\mu }^e/G_F`$ $`<`$ $`1.0\times 10^6,`$ (35)
$`ϵ_\tau ^eG_{e\tau }^e/G_F`$ $`<`$ $`4.2\times 10^3.`$ (36)
Note that the bounds on $`ϵ_{\mathrm{}}^f`$ do only coincide with those for $`ϵ_\nu _{\mathrm{}}^f`$ in the $`SU(2)`$ symmetric limit. We will comment on possible relaxations due to $`SU(2)_L`$ breaking effects later in section IV B.
To constrain $`G_{e\mu }^q`$ we use the upper bounds on $`\mu e`$ conversion from muon scattering off nuclei ,
$`{\displaystyle \frac{\sigma (\mu ^{}\text{Ti}e^{}\text{Ti})}{\sigma (\mu ^{}\text{Ti}\text{capture})}}`$ $`<`$ $`4.3\times 10^{12},`$ (37)
$`{\displaystyle \frac{\sigma (\mu ^{}\text{Pb}e^{}\text{Pb})}{\sigma (\mu ^{}\text{Pb}\text{capture})}}`$ $`<`$ $`4.6\times 10^{11},`$ (38)
$`{\displaystyle \frac{\sigma (\mu ^{}{}_{}{}^{32}\text{S}e^{}{}_{}{}^{32}\text{S})}{\sigma (\mu ^{}{}_{}{}^{32}\text{S}\nu _\mu ^{32}\text{P}^{})}}`$ $`<`$ $`7\times 10^{11},`$ (39)
concluding that
$$ϵ_\mu ^qG_{e\mu }^q/G_F<10^5.$$
(40)
is a conservative upper bound irrespective of the inherent hadronic uncertainties for such an estimate.
To constrain $`G_{e\tau }^q`$ we may use the upper bounds on various semi-hadronic tau decays that violate lepton flavor :
$`\text{BR}(\tau ^{}e^{}\pi ^0)`$ $`<`$ $`3.7\times 10^6,`$ (41)
$`\text{BR}(\tau ^{}e^{}\rho ^0)`$ $`<`$ $`2.0\times 10^6,`$ (42)
$`\text{BR}(\tau ^{}e^{}\eta )`$ $`<`$ $`8.2\times 10^6,`$ (43)
$`\text{BR}(\tau ^{}e^{}\pi ^+\pi ^{})`$ $`<`$ $`1.9\times 10^6.`$ (44)
Let us first consider the tau decays into $`\pi ^0`$ and $`\rho ^0`$. Since these mesons belong to an isospin triplet we can use the isospin symmetry to normalize the above bounds (41) and (42) by the measured rates of related lepton flavor conserving decays. Using $`\text{BR}(\tau ^{}\nu _\tau \pi ^{})=0.11`$ and $`\text{BR}(\tau ^{}\nu _\tau \rho ^{})=0.22`$ we obtain
$$G_{e\tau }^q(\pi )<8.2\times 10^3G_F,\text{and}G_{e\tau }^q(\rho )<4.2\times 10^3G_F.$$
(45)
Since the $`\pi `$ ($`\rho `$) is a pseudoscalar (vector) meson its decay probes the axial-vector (vector) part of the quark current.
In general, any semi-hadronic operator $`𝒪_l^q`$ can be decomposed into an $`I=0`$ and an $`I=1`$ isospin component. Only the effective coupling of the latter can be constrained by the upper bounds on the decays into final states with isovector mesons, like the $`\pi `$ and the $`\rho `$. If the resulting operator is dominated by the $`I=0`$ component, the bounds in (45) do not hold. But in this case we can use the upper bound on $`\text{BR}(\tau ^{}e^{}\eta )`$ in (43). Since the $`\eta `$ is an isosinglet, isospin symmetry is of no use for the normalization. However, we can estimate the proper normalization using the relation between the $`\eta `$ and $`\pi `$ hadronic matrix elements, which is just the ratio of the respective decay constants, $`f_\eta /f_\pi 1.3`$ . Taking into account the phase space effects, we obtain from (43) that
$$G_{e\tau }^q(\eta )<1.1\times 10^2G_F.$$
(46)
Since the $`\eta `$ is a pseudoscalar meson its decay probes the axial-vector part of the $`I=0`$ component of the quark current, while the neutrino propagation is only affected by the vector part. As we have already mentioned, for any single chiral New Physics contribution the vector and axial-vector parts have the same magnitude and we can use (46) to constrain the isosinglet component of $`𝒪_l^q`$. In case there are several contributions, whose axial-vector parts cancel each other , the $`I=0`$ component could still be constrained by the upper bound on $`\text{BR}(\tau ^{}e^{}\pi ^+\pi ^{})`$ in (44). While the calculation of the rate is uncertain due to our ignorance of the spectra and the decay constants of the isosinglet scalar resonances, we expect that the normalization will be similar to that of the $`\pi `$, $`\rho `$ and $`\eta `$ discussed before. Finally we note that the decay $`\tau ^{}e^{}\omega `$ would be ideal to constrain the $`I=0`$ vector part, but at present no upper bound on its rate is available.
While one can always fine-tune some parameters in order to avoid our bounds, our basic assumption is that this is not the case. Thus from (45) and (46) we conclude that
$$ϵ_\tau ^qG_{e\tau }^q/G_F<10^2.$$
(47)
#### 2 Flavor diagonal neutrino interactions
So far we have only discussed the upper bounds on FCNI. However, if the neutrinos are massless then in addition to the FCNI that induce an off-diagonal term in the effective neutrino mass matrix, also non-universal flavor diagonal interactions are needed to generate the required splitting between the diagonal terms.
In general any operator that induces such FDNI is related to other lepton flavor conserving operators, that give additional contributions to SM allowed processes, and therefore violate the lepton universality of the SM. Then the upper bounds on lepton universality violation can be used to constrain these operators. Using the relation to the operators that induce the FDNI one may also constrain the latter.
As we mentioned already for massless neutrinos only a non-zero $`ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{q}`$ ($`q=u,d`$) can lead to a resonance effect, while FDNI that allow for scattering off electrons alone are insufficient to solve the solar neutrino problem. Therefore we only need to discuss the effective flavor diagonal operator
$$𝒪_\nu ^q(\overline{\nu _\alpha }\gamma _\mu \nu _\alpha )(\overline{q}\gamma ^\mu q),$$
(48)
where $`\alpha =e,\mu ,\tau `$ and $`q=u_{L,R},d_{L,R}`$.
It is easy to check that for FDNI induced by heavy boson exchange $`𝒪_\nu ^q`$ is always induced together with
$$𝒪_l^q(\overline{l_\alpha }\gamma _\mu l_\alpha )(\overline{q^{}}\gamma ^\mu q^{}),$$
(49)
where $`q^{}=u_{L,R},d_{L,R}`$ can be different from $`q`$. Moreover, intermediate scalar singlets and triplets (that couple to $`QL`$) as well as charged vector singlets and triplets (that couple to $`\overline{Q}L`$) also give rise to
$$𝒪_{l\nu }^q(\overline{l_\alpha }\gamma _\mu \nu _\alpha )(\overline{q_L}\gamma ^\mu q_L^{}),$$
(50)
where $`q^{}=u,d`$ for $`q=d,u`$. (See (32) as an example for FDNI mediated by an intermediate scalar singlet.) Since $`𝒪_{l\nu }^q`$ induces an additional contribution to the SM weak decay $`\tau _L\pi \nu _\tau `$ for $`\alpha =\tau `$ and to $`\pi l_\alpha \overline{\nu }_\alpha `$ for $`\alpha =e,\mu `$, the relevant effective coupling $`G_{l\nu }^q`$ can be constrained by the upper bounds on lepton universality violation in semi-hadronic decays. The latter leads to a deviation of the parameters
$`R_{e/\mu }^\pi `$ $``$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle \frac{\mathrm{\Gamma }(\pi ^{}e^{}\overline{\nu }_e)}{\mathrm{\Gamma }(\pi ^{}\mu ^{}\overline{\nu }_\mu )}}}1+{\displaystyle \frac{G_{e\nu _e}^qG_{\mu \nu _\mu }^q}{G_F}}`$ (51)
$`R_{\tau /\mu }^\pi `$ $``$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle \frac{\mathrm{\Gamma }(\tau ^{}\nu _\tau \pi ^{})}{\mathrm{\Gamma }(\pi ^{}\mu ^{}\overline{\nu }_\mu )}}}1+{\displaystyle \frac{G_{\tau \nu _\tau }^qG_{\mu \nu _\mu }^q}{G_F}}`$ (52)
$`R_{\tau /e}^W`$ $``$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle \frac{\mathrm{\Gamma }(W^{}\tau \overline{\nu }_\tau )}{\mathrm{\Gamma }(W^{}e\overline{\nu }_e)}}}1+{\displaystyle \frac{G_{\tau \nu _\tau }^qG_{e\nu _e}^q}{G_F}}`$ (53)
from unity. Here $`N`$ denotes a normalization factor, which is just the ratio of the above two rates in the SM such that $`R_{\alpha /\beta }=1`$ if $`G_{l_\alpha \nu _\alpha }^q=G_{l_\beta \nu _\beta }^q`$. In the approximation we assume that $`G_{l\nu }^qG_F`$. From the most recent experimental data it follows that
$`R_{e/\mu }^\pi `$ $`=`$ $`1.0017\pm 0.0015,`$ (54)
$`R_{\tau /\mu }^\pi `$ $`=`$ $`1.005\pm 0.005,`$ (55)
$`R_{\tau /e}^W`$ $`=`$ $`0.987\pm 0.023,`$ (56)
implying that
$$ϵ_{}^{}{}_{\mathrm{}}{}^{q}\frac{G_\mathrm{}\nu _{\mathrm{}}^fG_{e\nu _e}^q}{G_F}<10^2$$
(57)
is a conservative upper bound. If $`𝒪_{l\nu }^q`$ is induced together with $`𝒪_\nu ^q`$, then in the $`SU(2)_L`$ symmetric limit $`ϵ_{}^{}{}_{\mathrm{}}{}^{q}=ϵ_{}^{}{}_{\nu }{}^{q}`$, but a modest relaxation due to $`SU(2)_L`$ breaking effects is possible (see section IV B).
It is essential to realize that not all New Physics operators that induce the FDNI relevant to solar neutrinos are related to $`𝒪_{l\nu }^q`$. For an intermediate $`SU(2)_L`$ scalar doublet (see eq. (27) for $`f=q`$) or a vector doublet (that couples to $`\overline{q^c}L`$) or a neutral vector singlet, only $`𝒪_l^q`$ is induced together with $`𝒪_\nu ^q`$. In this case one may only use the upper bound on $`G_{ll}^q`$ that are due to the constraints on compositeness. The present data from $`p\overline{p}e^+e^{},\mu ^+\mu ^{}+X`$ imply an upper limit on the scale of compositeness $`\mathrm{\Lambda }(qql_\alpha l_\alpha )>1.6`$ TeV, which translates into
$$G_{l_\alpha l_\alpha }^q<10^1G_F$$
(58)
as a conservative estimate for $`\alpha =e,\mu `$. (One-loop contributions to the $`Z`$ width due to $`𝒪_l^q`$ lead to a similar constraint on $`G_{ee}^q`$ .) However, no upper bound on $`\mathrm{\Lambda }(qq\tau \tau )`$ is available.
For a neutral vector singlet $`G_{l_\alpha l_\alpha }^q=G_{\nu _\alpha \nu _\alpha }^q`$ and from (5) and (58) it follows that
$$ϵ_{}^{}{}_{\nu _\mu }{}^{q}<10^1$$
(59)
while there is no model-independent bound on $`ϵ_{}^{}{}_{\nu _\tau }{}^{q}`$. For intermediate $`SU(2)_L`$ doublets the bound in (59) could be relaxed somewhat, since the effective couplings of the relevant operators may differ due to $`SU(2)_L`$ breaking effects, which we discuss next.
### B Constraining $`SU(2)_L`$ breaking effects
The excellent agreement between the SM predictions and the electroweak precision data implies that $`SU(2)_L`$ breaking effects cannot be large. To show that the upper bounds on $`G_{ll}^f`$ (or $`G_{l\nu }^f`$) translate into similar bounds for $`G_{\nu \nu }^f`$ if their related operators stem from the same $`SU(2)_L`$ invariant coupling, we recall from eq. (29) that in general the ratio of the couplings, $`G_{\nu _\alpha \nu _\beta }^f/G_{\alpha \beta }^f`$ (or $`G_{\nu _\alpha \nu _\beta }^f/G_{\alpha \nu _\beta }^f`$), is given by ratio $`M_1^2/M_2^2`$. Here $`M_1`$ and $`M_2`$ are the masses of the particles belonging to the $`SU(2)_L`$ multiplet that mediate the processes described by $`G_{\alpha \beta }^f`$ ($`G_{\alpha \nu _\beta }^f`$) and $`G_{\nu _\alpha \nu _\beta }^f`$, respectively. If $`M_1M_2`$ this multiplet will contribute to the oblique parameters $`S,U`$ and, most importantly, $`T`$ . A fit to the most recent precision data performed in Ref. determined the maximally allowed ratio $`(M_1/M_2)_{\mathrm{max}}^2`$ to be at most 6.8 (at 90 % CL) for intermediate scalars. (Vector bosons in general are expected to have even stronger bounds for the mass ratio). Consequently the upper limits on the effective couplings $`G_{\nu \nu }`$ agree with those we derived for the corresponding $`G_{ll}^f`$ (or $`G_{l\nu }^f`$) within an order of magnitude even for maximal $`SU(2)_L`$ breaking. Thus, barring fine-tuned cancellations.
$$ϵ_\nu _{\mathrm{}}^f<6.8ϵ_{\mathrm{}}^f\text{and}ϵ_{}^{}{}_{\nu _{\mathrm{}}}{}^{f}<6.8ϵ_{}^{}{}_{\mathrm{}}{}^{f},$$
(60)
at 90 % CL.
## V Implications for Future Experiments
In this section we discuss how to test the solution to the solar neutrino problem based on non-standard neutrino interactions in future neutrino experiments.
Let us consider first the possibility of obtaining stronger constraints on New Physics from future laboratory experiments. Our phenomenological analysis shows that FCNI could only be large enough to provide $`\nu _e\nu _\tau `$ transitions while, model-independently, $`\nu _e\nu _\mu `$ transitions are irrelevant for solar neutrinos. Even for $`\nu _e\nu _\tau `$ transitions the required effective coupling has to be close to its current upper bound, which we derived from limits on anomalous tau decays. Therefore the solution to the solar neutrino problem studied in this paper could be tested by the upcoming $`B`$-factories that are expected to improve the present experimental bounds on several rare $`\tau `$ decays. For example, assuming an integrated luminosity of 30 fb<sup>-1</sup> (corresponding to $`3\times 10^7\tau `$ pairs) for the BaBar experiment, the upper limits on the branching ratios in (41)–(44) could be reduced by one order of magnitude. This would decrease the bound on $`ϵ_\tau ^q`$ in (47) to a value close to the smallest possible best fit values for $`\epsilon `$ (c.f. Fig. 7 and Fig. 8) ruling out a large region of the parameter space and making the NSNI solution increasingly fine-tuned.
Next we consider the implications for future solar neutrino experiments . In Tab. 1 we present the expected ranges for the event rates (normalized to the SSM expectation in the absence of neutrino flavor transitions) of those experiments, if the solar neutrino problem is explained by NSNI. In Fig. 9 the predicted rates are presented graphically. The ranges correspond to the 95 % CL regions for $`(\epsilon ,\epsilon ^{})`$ in Figs. 7b and 8b. As before we use the BP98 SSM predictions for the initial neutrino fluxes and the survival probability in eq. (9) to compute the expected rates for each of the five detectors. Specifically, there are three types of detectors: (a) The Sudbury Neutrino Observatory (SNO) , which is measuring the <sup>8</sup>B neutrino charged current (CC) rate, (b) the BOREXINO and KamLAND experiments that are designed to observe the <sup>7</sup>Be neutrino signal and (c) the HELLAZ and HERON experiments dedicated to a precise measurement of the low-energy $`pp`$ neutrino flux.
Tab. 1: Future solar neutrino experiments and their rates predicted by the NSNI solution
| Experiment | Start of operation | Main neutrino source | Rate predicted by NSNI |
| --- | --- | --- | --- |
| SNO | 1999 | <sup>8</sup>B | $`0.220.43`$ |
| BOREXINO | 2001 | <sup>7</sup>Be | $`0.300.52`$ |
| KamLAND | 2001 | <sup>7</sup>Be | $`0.300.52`$ |
| HELLAZ | $`>2002`$ | $`pp`$ | $`0.520.83`$ |
| HERON | $`>2000`$ | $`pp`$ | $`0.520.82`$ |
The predictions for the rates in Tab. 1 reflect the relation between the predominant neutrino fluxes that we presented in eq. (17), i.e. neutrinos with higher energies are in general produced closer to the solar center and therefore more likely to pass through a resonance and undergo flavor conversion.
As can be seen from Fig. 9 the suppression pattern of the NSNI solution is clearly different from the one predicted by the small angle MSW solution (c.f. Fig. 1 of Ref. ). But there is a striking similarity between the NSNI solution and the LMA solution, including the preference for large $`f_B`$ (c.f. Fig. 7 of Ref. ), the absence of a <sup>8</sup>B spectral distortion and the modest day night effect. Consequently, using solar neutrino data, it will be difficult to distinguish the NSNI scenario from the LMA MSW solution. We note, however, that the KamLAND experiment will provide an independent test of the oscillation parameters of the LMA MSW solution by observing anti electron neutrinos from several nuclear reactors around the Kamioka mine in Japan. Thus, if KamLAND would indeed confirm the LMA MSW solution, then the NSNI solution discussed in this paper will be irrelevant.
Since SNO already started taking data and is expected to have some results soon, let us consider some implication for this experiment. As we have pointed out one of the important features of the NSNI conversion mechanism is the absence of any distortion in the solar neutrino spectrum even though the averaged survival probabilities of neutrinos from different nuclear reactions in the sun are not equal. Due to this feature, the following simple relation between the SuperKamiokande solar neutrino event rate $`R_{SK}`$ and the SNO CC event rate $`R_{SNO}^{CC}`$ (both normalized by the SSM predictions) holds:
$$R_{SK}=R_{SNO}^{CC}(1r)+rf_B,$$
(61)
where $`R_{SK}`$ and $`R_{SNO}^{CC}`$ are defined exactly as in eqs. (4) and (6) of Ref. and $`r`$ is given by
$$r\frac{{\displaystyle 𝑑E_eR(E_e)𝑑E_\nu \varphi ^{{}_{}{}^{8}\text{B}}(E_\nu )\sigma _{\nu _{\mu ,\tau }e}(E_\nu ,E_e)}}{{\displaystyle 𝑑E_eR(E_e)𝑑E_\nu \varphi ^{{}_{}{}^{8}\text{B}}(E_\nu )\sigma _{\nu _ee}(E_\nu ,E_e)}}\frac{1}{7}.$$
(62)
Here, $`E_e`$ and $`E_\nu `$ are the electron and neutrino energy, respectively, $`R(E_e)`$ is the SuperKamiokande resolution and efficiency function, $`\varphi ^{{}_{}{}^{8}\text{B}}`$ is the <sup>8</sup>B neutrino flux, and $`\sigma _{\nu _ee}`$ and $`\sigma _{\nu _{\mu ,\tau }e}`$ denote the elastic scattering cross sections for $`\nu _ee^{}\nu _ee^{}`$ and $`\nu _{\mu ,\tau }e^{}\nu _{\mu ,\tau }e^{}`$, respectively.
We note that $`R_{SK}`$ and $`R_{SNO}^{CC}`$ are defined such that $`R_{SK}=R_{SNO}^{CC}`$ in the absence of neutrino flavor transitions including the case where $`f_B1`$. (Strictly speaking, a slight violation of the equality in eq. (61) could be induced by the earth matter effect on these two experiments, since they are located at somewhat different latitudes.) Using the relation (61), the true flux of the <sup>8</sup>B neutrino flux $`(\varphi ^{{}_{}{}^{8}\text{B}})_{true}=f_B(\varphi ^{{}_{}{}^{8}\text{B}})_{SSM}`$ could be precisely determined by combining SuperKamiokande and SNO solar neutrino measurements, if the solar neutrino problem is indeed due to NSNI.
Finally let us discuss shortly the possibility of testing the solution studied in this paper by future long-baseline neutrino oscillation experiments. Since only $`\nu _e\nu _\tau `$ transitions are viable, an independent test would require a $`\nu _\tau `$ ($`\overline{\nu }_\tau `$) appearance experiment using an intense beam of $`\nu _e`$ ($`\overline{\nu }_e`$), which could be created at future neutrino factories (see, e.g., Ref. ).
Assuming a constant density and using the approximation that $`n_dn_u3n_e`$ in the earth, the conversion probability for a neutrino which travels a distance $`L`$ in the earth is given by:
$$P(\nu _e\nu _\tau ;L)\frac{36\epsilon ^2}{36\epsilon ^2+(13\epsilon ^{})^2}\mathrm{sin}^2\left[\frac{1}{2}\sqrt{36\epsilon ^2+(13\epsilon ^{})^2}\sqrt{2}G_Fn_eL\right].$$
(63)
Numerically, the oscillation length in the earth matter can be estimated to be
$$L_{osc}8.1\times 10^3\left[\frac{2\text{mol/cc}}{n_e}\right]\left[\frac{1}{\sqrt{36\epsilon ^2+(13\epsilon ^{})^2}}\right]\text{km}.$$
(64)
Using eqs. (63) and (64) and the approximation $`n_e2`$ mol/cc (which is valid close to the earth surface), we find that, for the case of non-standard neutrino scattering off $`d`$-quark, $`P`$ few $`\times 10^4`$ for K2K ($`L=250`$ km) and $`P`$ few $`\times 10^3`$ for MINOS ($`L=732`$ km) for our best fit parameters. Similarly for $`u`$-quark, $`P`$ few $`\times 10^5`$ for K2K and $`P`$ few $`\times 10^4`$ for MINOS for the best fit parameters. These estimates imply that it would be hard but not impossible, at least for the case of scattering off $`d`$-quarks, to obtain some signal of $`\nu _e\nu _\tau `$ conversions due to NSNI interactions by using an intense $`\nu _e`$ beam which can be created by a muon storage ring .
## VI Conclusions
According to our $`\chi ^2`$ analysis non-standard neutrino interactions (NSNI) can provide a good fit to the solar neutrino data provided that there are rather large non-universal FDNI (of order $`0.5G_F`$) and small FCNI (of order $`10^210^3G_F`$). The fit to the observed total rate, day-night asymmetry, seasonal variation and spectrum distortion of the recoil electron spectrum is comparable in quality to the one for standard neutrino oscillations.
From the model-independent analysis we learn that NSNI induced by the exchange of heavy bosons cannot provide large enough $`\nu _e\nu _\mu `$ transitions, while $`\nu _e\nu _\tau `$ FCNI in principle could be sufficiently strong. However, the current bounds will be improved by the up-coming $`B`$-factories, providing an independent test of the NSNI solution. The required large non-universal FDNI (for $`\nu _e`$ transitions into both $`\nu _\mu `$ and $`\nu _\tau `$) can be ruled out by the upper bounds on lepton universality, unless they are induced by an intermediate doublet of $`SU(2)_L`$ (a scalar or a vector boson) or by a neutral vector singlet. For $`\nu _e\nu _\mu `$ there exists a bound due to the limit on compositeness in this case, but for $`\nu _e\nu _\tau `$ there is no significant constraint at present.
Generically only very few models can fulfill the requirements needed for the solution discussed in this paper: massless neutrinos, small FCNI and relatively large non-universal FDNI. As for the vector bosons the most attractive scenario is to evoke an additional $`U(1)_{B3L_\tau }`$ gauge symmetry (where $`B`$ is the baryon number and $`L_\tau `$ denotes the tau lepton number), which would introduce an additional vector singlet that only couples to the third generation leptons and quarks . Among the attractive theories beyond the standard model where neutrinos are naturally massless as a result of a protecting symmetry, are supersymmetric $`SU(5)`$ models that conserve $`BL`$, and theories with an extended gauge structure such as $`SU(3)_CSU(3)_LU(1)_N`$ models , where a chiral symmetry prevents the neutrino from getting a mass. These particular models, however, do not contribute significantly to the specific interactions we are interested in this paper. $`SU(5)`$ models have negligible NSNI since they are mediated by vector bosons which have masses at the GUT scale. $`SU(3)_CSU(3)_LU(1)_N`$ models can provide large $`ϵ_e`$ and $`ϵ_e^{}`$, but these models do not induce NSNI with quarks. From eq. (7) it follows that no resonant conversion can occur in this case.
Therefore we conclude that the best candidate for the scenario we studied are supersymmetric models with broken $`R`$-parity, where the relevant NSNI are mediated by a scalar doublet, namely the “left-handed” bottom squark. Although in this model neutrino masses are not naturally protected from acquiring a mass, one may either evoke an additional symmetry or assume that non-zero neutrino masses are not in a range that would spoil the solution in terms of the non-standard neutrino oscillations we have studied in this paper.
Even though we consider the conventional oscillation mechanisms as the most plausible solutions to the solar neutrino problem, it is important to realize that in general New Physics in the neutrino sector include neutrino masses and mixing, as well as new neutrino interactions. While it is difficult to explain the atmospheric neutrino problem and the LSND anomalies by NSNI , we have shown in this paper that a solution of the solar neutrino problem in terms of NSNI is still viable. The ultimate goal is of course a direct experimental test of this solution. The upcoming solar neutrino experiments will provide a lot of new information which hopefully will reveal the true nature of the solar neutrino problem.
###### Acknowledgements.
This work was partially supported by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP), Conselho Nacional de Desenvolvimento Científico e Tecnologógico (CNPq) and NSF grant PHY-9605140. We thank Y. Grossman, Y. Nir and C. Peña-Garay for helpful discussions. One of us (MMG) would like to thank the Physics Department, University of Wisconsin at Madison, where part of this work was realized, for the hospitality.
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# 1 Introduction
## 1 Introduction
The ability to embed a lower-dimensional theory in a higher-dimensional one has proved to be an extremely useful one in string theory. One can, for example, re-interpret lower-dimensional $`p`$-brane solitons as solutions of the ten-dimensional string, or eleven-dimensional M-theory. A crucial aspect of this picture is that the Kaluza-Klein reduction must be a consistent one, in the sense that all solutions of the lower-dimensional theory must also be solutions of the original higher-dimensional theory. This consistency is guaranteed in a standard toroidal reduction, but it is far less clear-cut when a reduction on a manifold such as a sphere is considered.
Kaluza-Klein reductions on spheres are of great interest in the framework of string theory, because they can give rise to lower-dimensional gauged supergravities that are relevant for discussing the AdS/CFT correspondence . The generic structure of these gauged supergravities comprises gravity coupled to a set of Yang-Mills gauge fields, and a set of scalar fields with a non-trivial potential, together, possibly, with additional antisymmetric tensor fields. A particular class of solution that can be studied is extremal domain walls, which can be viewed as charged black holes or black $`p`$-branes in the gauged theory, in the extremal limit for which the electric or magnetic charges actually vanish. Thus these solutions are supported entirely by the metric and certain scalar fields within the gauged supergravity.
It therefore becomes of interest to study the circumstances under which a higher-dimensional theory can admit a consistent $`n`$-sphere reduction in which just gravity and appropriate scalar fields are retained. In some cases this may be viewed as a subset of a larger consistent reduction of a gauged supergravity, in which the starting point is supergravity in ten or eleven dimensions. However, the question can also be posed in a more general framework, where the starting point need not necessarily even be a supersymmetric theory.
Before discussing the possible new cases, let us review what is known at present. It is natural, when considering an $`n`$-sphere reduction, to try to retain all the $`SO(n+1)`$ Yang-Mills fields as part of the consistent reduction. Usually, however, this is not possible. It was recently shown in that the cases where this can be done, starting from a $`D`$-dimensional theory of gravity, $`n`$-form field strength and dilaton, are as follows. One can start with $`(D,n)=(11,4)`$, and reduce on $`S^4`$ or $`S^7`$ to seven or four dimensions respectively; another possibility is to start from $`(D,n)=(10,5)`$, and reduce on $`S^5`$ to five dimensions. In these cases, the system has no dilaton. Including a dilaton, with a specific coupling, one can also start with $`n=3`$ and $`D`$ arbitrary, reducing on $`S^3`$ or $`S^{D3}`$; or finally one can start with $`n=2`$ and $`D`$ arbitrary, and reduce on $`S^2`$. In all cases one must also include scalar fields $`T_{ij}`$ in the reduction Ansatz, corresponding to the coset $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$. Additionally, for the $`S^4`$ reduction of $`D=11`$ one must include five 3-form field strengths in the Ansatz, while for the $`S^7`$ reduction one must also include 35 more pseudoscalars.<sup>1</sup><sup>1</sup>1In fact the consistent reductions from $`D=11`$ require, in addition, the inclusion of the $`FFA`$ in the Lagrangian that arises in $`D=11`$ supergravity. In the $`S^5`$ reduction from $`D=10`$, it is necessary to impose the requirement of self-duality on the 5-form field strength. Finally, for the $`S^3`$ reduction in the $`n=3`$ case, one must include a 3-form field strength in the reduction Ansatz.
Our statement of the possible consistent reductions first specified that the $`SO(n+1)`$ Yang-Mills gauge fields were to be included, and then we listed the additional fields that would be needed for consistency. Another way of phrasing the question is to specify which scalar fields will be included in the reduction Ansatz. In fact if we want to include all the scalars $`T_{ij}`$ of the $`SL(n+1,\text{I}\mathrm{R})/SO(n+1)`$ coset, the list of cases where consistent reductions are possible will be the same as the above. The reason for this is that once all the scalars $`T_{ij}`$ are present, they will act as sources for the Yang-Mills gauge fields, and so it would be inconsistent to omit the Yang-Mills fields. However, if we settle for a reduction in which fewer scalars are retained, it becomes possible to omit the Yang-Mills fields and this opens up some further possibilities for consistent reductions, which we shall explore in this paper. These reductions with scalars but no gauge fields will be sufficient for the purpose of constructing the extremal domain-wall solutions in the lower dimension, and then lifting them back to the higher dimension.
As mentioned above, if one includes the full set of scalars $`T_{ij}`$ in a truncation then they will give rise to source terms that require the Yang-Mills fields to be non-zero. Specifically, the source currents are of the form $`T_{k[i}^1_\mu T_{j]k}`$, in the adjoint of $`SO(n+1)`$. If we make a truncation where only the diagonal scalar fields are retained,
$$T_{ij}=\mathrm{diag}(X_1,X_2,\mathrm{}X_{n+1}),$$
(1)
then the currents $`T_{k[i}^1_\mu T_{j]k}`$ will be zero, and thus there is no longer any necessity to include the gauge fields in a consistent truncation. This actually allows a somewhat extended set of $`(D,n)`$ values for which consistent reductions can be achieved, which includes cases that would not allow consistent reductions with $`SO(n+1)`$ gauge fields. The allowed cases are detailed below.
In section 2 we construct an Ansatz for the $`n`$-sphere reduction of the a $`D`$-dimensional theory of gravity, an $`n`$-form field strength, and a dilaton, in which the lower-dimensional fields comprise just gravity and the diagonal scalar fields given by (1). We obtain a complete proof of the consistency of this Kaluza-Klein reduction, showing that it works in all cases where the strength of the coupling of the dilaton to the $`n`$-form in $`D`$ dimensions is appropriate. This requirement on the coupling is a rather stringent one, and the allowable cases turn out to be $`\{n=5,D10\}`$; $`\{n=4,D11\}`$; and $`n3`$ with $`D`$ arbitrary, for $`nD/2`$.
In section 3 we construct $`(n+1)`$-parameter extremal domain-wall solutions in the lower-dimensional theories of gravity plus scalar fields, and then make use of the reduction Ansatz derived in section 2 in order to lift these solutions back to the original $`D`$-dimensional theory. We show that in the higher dimension the lifted solutions admit an interpretation as continuous distributions of $`(Dn2)`$-branes. We discuss and obtain the distribution functions. We obtain the metric of the distributed branes in the dual frame, and show that the structure of these metrics depends only on the dimension $`n`$ of the internal sphere, but is independent of $`D`$. In particular, the metric in the dual frame becomes asymptotically AdS$`\times S^n`$ for $`n3`$, and Minkowski$`\times S^3`$ for $`n=3`$.
In section 4 we analyse the spectrum of excitations of a minimally-coupled scalar in the background of the $`(Dn)`$-dimensional domain-wall solution, showing that it has a universal structure that is characterised by the dimension $`n`$ of the internal sphere used in the dimensional reduction. In the case of the vacuum solutions, where the $`(n+1)`$ parameters in the general solutions are all set to zero, the scalar wave equation can be solved explicitly, allowing a study of the singularity structure. We also analyse the Schrödinger potentials for generic cases, allowing us to determine the structures of the spectra in the various cases.
In an appendix, we show that a single-charge rotating $`p`$-brane in a generic dimension can be dimensionally reduced on the internal (distorted) $`n`$-sphere to give rise to domain-wall black holes with $`[(n+1)/2]`$ electric $`U(1)`$ charges. In the extremal limit, the gauge fields vanish and the balck hole becomes a domain wall that is contained within the set of solutions obtained in this paper.
## 2 Kaluza-Klein sphere reduction
Single-charge $`p`$-branes in supergravity theories in $`D`$ dimensions can be classified as solutions of the theory described by the Einstein-Hilbert action coupled to a dilaton and an $`n`$-form field strength,
$$_D=\widehat{R}\widehat{}\text{1}\mathrm{l}\frac{1}{2}\widehat{}d\widehat{\varphi }d\widehat{\varphi }\frac{1}{2}e^{a\widehat{\varphi }}\widehat{}\widehat{F}_{\left(n\right)}\widehat{F}_{\left(n\right)},$$
(2)
where the constant $`a`$ is given by
$$a^2=4\frac{2(n1)(Dn1)}{D2}.$$
(3)
The requirement that $`a`$ be real puts a strong condition on the possible values for $`n`$, bearing in mind that we must have $`nD`$, and in fact we can always choose a dualisation for $`F_{\left(n\right)}`$ for which $`nD/2`$. From (3), it then follows that the maximum value for $`nD/2`$ is 5. Whe $`n=5`$, the maximal dimension is $`D=10`$, corresponding to the self-dual 5-form in the type IIB theory. For $`n=4`$, the maximal dimension is $`D=11`$, corresponding to 11-dimensional supergravity. In both cases, the constant $`a`$ vanishes. For $`n=0,1,2,3`$, the dimension $`D`$ can be arbitrary. Note that for a given $`n`$ satisfying (3), $`n^{}=Dn`$ satisfies it too. To summarise, the allowed possibilities are
$`n=5,D5:`$ $`D10`$
$`n=4,D4:`$ $`D11`$
$`n=0,1,2,3,D,D1,D2,D3:`$ $`D\text{arbitrary}.`$ (4)
Note that these results come from the requirement only that $`a`$ must be real. If in addition we require that the Lagrangian must be associated with a supersymmetric theory, we get the further restriction that the dimension $`D`$ must be less than or equal to eleven or ten.
The $`p`$-branes for which the first term on the right-hand-side of (3) is 4 can be viewed as the basic building blocks for $`p`$-brane solitons. The $`p`$-branes with values other than 4, (usually $`4/N`$ with $`N`$ an integer) can be viewed as bound states or intersections of these building blocks. For example, for $`D=11`$ and $`D=10`$, our discussion applies to M-branes, the NS-NS string and 5-brane and all the D-branes.
We shall now consider the Kaluza-Klein dimensional reduction of the Lagrangian (2) on $`S^n`$. (The discussion of the reduction instead on $`S^{Dn}`$ can be handled by dualising the $`n`$-form field strength to a $`(Dn)`$-form.) In general, such a reduction is inconsistent if we keep all the massless fields. It was shown , however, that for $`n=2`$ and $`n=3`$ the reduction is always consistent, provided that (3) is satisfied. For $`n=5`$ and $`n=4`$, the reduction is consistent only if additional conditions are satisfied, namely self-duality of the 5-form in $`D=10`$, and the addition of an $`FFA`$ term in $`D=11`$ for the $`n=4`$ case.
In this paper, we shall truncate further to a subset of the massless fields, corresponding to “diagonal inhomogeneous distortions” of the internal $`S^n`$ metric. By this, we mean that we canonically embed the sphere $`S^n`$ in $`n+1`$ dimensional Euclidean space. The round $`S^n`$ metric is given by $`d\mu _id\mu _i`$, where $`\mu _i`$ are Euclidean coordinates satisfying the unit-length constraint $`\mu _i\mu _i=1`$. The diagonal inhomogeneous distortion of the sphere is then achieved by introducing $`(n+1)`$ scalars $`X_i`$, and scaling the coordinate differentials as follows:
$$ds_n^2=\underset{i}{}X_i^1(d\mu _i)^2.$$
(5)
We shall show that for this subset of fields, the Kaluza-Klein reduction is consistent for any of the $`D`$ and $`n`$ values listed in (4), provided that (3) is satisfied.
We find that the Kaluza-Klein reduction Ansatz is given by
$`d\widehat{s}_D^2=Y^{\frac{1}{D2}}\left(\mathrm{\Delta }^{\frac{n1}{D2}}ds_{Dn}^2+g^2\mathrm{\Delta }^{\frac{(Dn1)}{(D2)}}{\displaystyle \underset{i=1}{\overset{n+1}{}}}X_i^1(d\mu _i)^2\right),`$
$`e^{\frac{2}{a}\widehat{\varphi }}=\mathrm{\Delta }^1Y^{\frac{2(Dn1)}{a^2(D2)}},`$ (6)
$`\widehat{F}_{\left(n\right)}=g^{n+1}\mathrm{\Delta }^2UW+g^{n+1}_\nu \left({\displaystyle \frac{X_i\mu ^i}{\mathrm{\Delta }}}\right)dx^\nu Z_i.`$
where
$$\mu _i\mu _i=1,Y=X_i,\mathrm{\Delta }=X_i\mu _i^2,U=2\underset{i}{}X_i^2\mu _i^2\mathrm{\Delta }\underset{i}{}X_i.$$
(7)
The quantities $`W`$ and $`Z_i`$ are respectively the volume-form on the $`n`$-sphere, and a certain $`(n1)`$-form on the $`n`$-sphere:
$`W={\displaystyle \frac{1}{n!}}ϵ_{ij_1\mathrm{}j_n}\mu ^id\mu ^{j_1}\mathrm{}d\mu ^{j_n},`$ (8)
$`Z_i={\displaystyle \frac{1}{(n1)!}}ϵ_{ij_1\mathrm{}j_n}\mu ^{j_1}d\mu ^{j_2}\mathrm{}d\mu ^{j_n}.`$
We find after some algebra that the dual of the field strength $`F_{\left(n\right)}`$ is given by
$$e^{a\widehat{\varphi }}\widehat{}F_{\left(n\right)}=gUϵ_{Dn}+\frac{1}{2g}X_i^1dX_id(\mu _i^2).$$
(9)
We can then substitute the Ansatz into higher dimensional equations of motion. First, we can verify that the Ansatz for $`\widehat{F}_{\left(n\right)}`$ in (6) satisfies the Bianchi identity $`d\widehat{F}_{\left(n\right)}=0`$. Next, we look at the equations of motion for the field strength $`\widehat{F}_{\left(n\right)}`$ and the dilaton $`\widehat{\varphi }`$:
$`d\left(e^{a\widehat{\varphi }}\widehat{}\widehat{F}_{\left(n\right)}\right)`$ $`=`$ $`0,`$
$`(1)^Dd\widehat{}d\widehat{\varphi }`$ $`=`$ $`ae^{a\widehat{\varphi }}\widehat{}\widehat{F}_n\widehat{F}_{\left(n\right)}.`$ (10)
After a considerable amount of algebra, we find that the Ansatz yields a consistent dimensional reduction of these $`D`$-dimensional equations to give the following $`(Dn)`$-dimensional equations for the scalar fields:
$`(1)^{Dn}d(\stackrel{~}{X}_i^1d\stackrel{~}{X}_i)=2g^2Y^{\frac{2}{n+1}}\left[2\stackrel{~}{X}_i^2\stackrel{~}{X}_i{\displaystyle \underset{j}{}}\stackrel{~}{X}_j\frac{2}{n+1}{\displaystyle \underset{j}{}}\stackrel{~}{X}_j^2+\frac{1}{n+1}({\displaystyle \underset{j}{}}\stackrel{~}{X}_j)^2\right]ϵ_{Dn},`$
$`(1)^{Dn}\frac{2(Dn2)}{(D2)a^2}d(Y^1dY)=Vϵ_{Dn}.`$ (11)
Here, we have defined the rescaled fields $`\stackrel{~}{X}_i`$ by
$$X_i=Y^{\frac{1}{n+1}}\stackrel{~}{X}_i,$$
(12)
so that $`_i\stackrel{~}{X}_i=1`$, and the potential $`V`$ is defined by
$$V\frac{1}{2}g^2\left(2\underset{i}{}X_i^2(\underset{i}{}X_i)^2\right)=\frac{1}{2}g^2Y^{\frac{2}{n+1}}\left(2\underset{i}{}\stackrel{~}{X}_i^2(\underset{i}{}\stackrel{~}{X}_i)^2\right).$$
(13)
Finally, to check the higher-dimensional Einstein equations, we need first to calculate the Ricci tensor for the metric in (6). This is most easily done by noting that it is conformally related to the metric
$$d\overline{s}_D^2=\mathrm{\Delta }^pds_{Dn}^2+\mathrm{\Delta }^q\underset{i}{}X_i^1(d\mu _i)^2,$$
(14)
with
$$d\widehat{s}_D^2=e^{2f}d\overline{s}_D^2,$$
(15)
where we have defined
$$e^{2f}=Y^{\frac{1}{D2}},p=\frac{n1}{D2},q=\frac{Dn1}{D2}.$$
(16)
It is easy to establish the standard result that the coordinate-frame components of the Ricci tensor $`\widehat{R}_{MN}`$ for the metric $`d\widehat{s}_D^2`$ are related to the coordinate-frame components $`\overline{R}_{MN}`$ for the metric $`d\overline{s}_D^2`$ by
$$\widehat{R}_{MN}=\overline{R}_{MN}+(D2)\left(_Mf_Nf\overline{}_M_Nf\overline{g}^{PQ}(_Pf)(_Qf)\overline{g}_{MN}\right)\overline{\text{ }\text{ }}f\overline{g}_{MN}.$$
(17)
Results for the Ricci tensor for certain metrics of the form (14) were derived in , and with minor modifications they can be carried over to our present case. They were obtained in a basis where one of the $`(n+1)`$ coordinates $`\mu _i`$, say $`\mu _0`$, is expressed in terms of the $`n`$ remaining ones $`\mu _\alpha `$ by using the relation $`\mu _i\mu _i=1`$. Thus the components $`g_{\alpha \beta }`$ of the distorted $`n`$-sphere metric (5), and its inverse, are given by
$`g_{\alpha \beta }`$ $`=`$ $`X_\alpha \delta _{\alpha \beta }+X_0^1\widehat{\mu }_\alpha \widehat{\mu }_\beta ,`$
$`g^{\alpha \beta }`$ $`=`$ $`X_\alpha \delta _{\alpha \beta }\mathrm{\Delta }^1X_\alpha X_\beta \mu _\alpha \mu _b,`$ (18)
where in the first line we are writing $`\widehat{\mu }_\alpha =\mu _\alpha /\mu _0`$. We refer to for many of the details of the curvature calculations. Combining these results with (17), we obtain, after extensive algebraic manipulations, the following expressions for the lower-dimensional spacetime, internal and mixed components of the $`D`$-dimensional Ricci tensor:
$`\widehat{R}_{\mu \nu }`$ $`=`$ $`R_{\mu \nu }\frac{(n1)(Dn1)}{4(D2)}\mathrm{\Delta }^2_\mu \mathrm{\Delta }_\nu \mathrm{\Delta }\frac{1}{2}p\mathrm{\Delta }^1\text{ }\text{ }\mathrm{\Delta }g_{\mu \nu }+\frac{1}{2}p\mathrm{\Delta }^2_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }g_{\mu \nu }`$
$`\frac{1}{4}X_i^2_\mu X_i_\nu X_i+\frac{1}{2}\mathrm{\Delta }^1X_i^1\mu _i^2_\mu X_i_\nu X_i+\frac{1}{4(D2)}Y^2_\mu Y_\nu Y`$
$`\frac{1}{4}q\mathrm{\Delta }^1(_\mu \mathrm{\Delta }_\nu Y+_\nu \mathrm{\Delta }_\mu Y)\frac{1}{2(D2)}\text{ }\text{ }\mathrm{log}Yg_{\mu \nu }`$
$`p\left({\displaystyle \underset{i}{}}X_i^2\mathrm{\Delta }^1X_i^2\mu _i^2{\displaystyle \underset{j}{}}X_j+2\mathrm{\Delta }^2(X_i^2\mu _i^2)^22\mathrm{\Delta }^1X_i^3\mu _i^2\right)g_{\mu \nu },`$
$`\widehat{R}_{\alpha \beta }`$ $`=`$ $`R_{\alpha \beta }+\frac{1}{2}qg_{\alpha \beta }\mathrm{\Delta }^2\text{ }\text{ }\mathrm{\Delta }\frac{1}{2}qg_{\alpha \beta }\mathrm{\Delta }^3_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }\frac{1}{2}\mathrm{\Delta }^1\text{ }\text{ }g_{\alpha \beta }`$
$`+\frac{1}{2}\mathrm{\Delta }^1g^{\gamma \delta }_\lambda g_{\alpha \gamma }^\lambda g_{\beta \delta }\frac{1}{4}\mathrm{\Delta }^2_\alpha \mathrm{\Delta }_\beta \mathrm{\Delta }\frac{1}{2}\mathrm{\Delta }^1_\alpha _\beta \mathrm{\Delta }`$
$`\frac{1}{4}qg_{\alpha \beta }\mathrm{\Delta }^2_\gamma \mathrm{\Delta }^\gamma \mathrm{\Delta }+\frac{1}{2}qg_{\alpha \beta }\mathrm{\Delta }^1_\gamma ^\gamma \mathrm{\Delta }\frac{1}{2(D2)}\mathrm{\Delta }^1\text{ }\text{ }\mathrm{log}Yg_{\alpha \beta },`$
$`\widehat{R}_{\alpha \mu }`$ $`=`$ $`\frac{1}{2}\mathrm{\Delta }^2U(X_\alpha ^1_\mu X_\alpha X_0^1_\mu X_0)\mu _\alpha +\frac{1}{8}a^2\mathrm{\Delta }^2_\mu \mathrm{\Delta }_a\mathrm{\Delta }\frac{1}{4}q\mathrm{\Delta }^1Y^1_\mu Y_\alpha \mathrm{\Delta }.`$
Note that here we are using a “generalised” summation convention in which summations over the $`i`$ index, where not otherwise indicated, are understood. The $`\text{ }\text{ }`$ operator denotes the d’Alembertian calculated in the lower-dimensional metric $`g_{\mu \nu }`$, and $`R_{\alpha \beta }`$ denotes the Ricci tensor of the internal metric (i.e. the Ricci tensor for the metric (5), with the $`X_i`$ are treated as parameters independent of the internal coordinates).
The $`D`$-dimensional Einstein equation reads $`\widehat{R}_{MN}=\widehat{S}_{MN}`$, where
$$\widehat{S}_{MN}=\frac{1}{2}_M\widehat{\varphi }_N\widehat{\varphi }+\frac{e^{a\widehat{\varphi }}}{2(Dn1)!}\left(\widehat{F}_{MN}^2\frac{Dn3}{(Dn)(Dn1)}\widehat{F}^2\widehat{g}_{MN}\right).$$
(20)
After some algebra we find that $`\widehat{S}_{MN}`$ is given by
$`\widehat{S}_{\mu \nu }`$ $`=`$ $`\frac{1}{2}\mathrm{\Delta }^1\mu _i^2X_i^1_\mu X_i_\nu X_i\frac{(n1)(Dn1)}{4(D2)}\mathrm{\Delta }^2_\mu \mathrm{\Delta }_\nu \mathrm{\Delta }+\frac{(Dn1)^2}{2a^2(D2)^2}Y^2_\mu Y_\nu Y`$
$`\frac{1}{4}q\mathrm{\Delta }^1Y^1(_\mu \mathrm{\Delta }_\nu Y+_\nu \mathrm{\Delta }_\mu Y)`$
$`\frac{1}{2}p\mathrm{\Delta }^2\left(U^2_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }+\mathrm{\Delta }\mu _i^2X_i^1_\lambda X_i^\lambda X_i\right)g_{\mu \nu },`$
$`\widehat{S}_{\alpha \beta }`$ $`=`$ $`\frac{1}{2}q\mathrm{\Delta }^3U^2g_{\alpha \beta }+\frac{1}{2}q\mathrm{\Delta }^2g_{\alpha \beta }X_i^1\mu _i^2_\lambda X_i^\lambda X_i\frac{1}{2}q\mathrm{\Delta }^3_\lambda \mathrm{\Delta }^\lambda \mathrm{\Delta }g_{\alpha \beta }`$ (21)
$`\frac{1}{2}\mathrm{\Delta }^2(X_\alpha ^1_\lambda X_\alpha X_0^1_\lambda X_0)(X_\beta ^1^\lambda X_\beta X_0^1^\lambda X_0)\mu _\alpha \mu _\beta `$
$`+\frac{1}{8}a^2\mathrm{\Delta }^2_\alpha \mathrm{\Delta }_\beta \mathrm{\Delta },`$
$`\widehat{S}_{\alpha \mu }`$ $`=`$ $`\frac{1}{2}\mathrm{\Delta }^2U(X_\alpha ^1_\mu X_\alpha X_0^1_\mu X_0)\mu _\alpha +\frac{1}{8}a^2_\mu \mathrm{\Delta }_\alpha \mathrm{\Delta }\frac{1}{4}q\mathrm{\Delta }^1Y^1_\mu Y_\alpha \mathrm{\Delta }.`$
After making use of the already-established equations of motion for the scalar fields, we eventually find after considerable further algebra that the $`\widehat{R}_{\mu \nu }=\widehat{S}_{\mu \nu }`$ components of the higher-dimensional Einstein equation imply
$$R_{\mu \nu }=\frac{1}{4}\stackrel{~}{X}_i^2_\mu \stackrel{~}{X}_i_\nu \stackrel{~}{X}_i+\frac{2(Dn2)}{(D2)(n+1)a^2}Y^2_\mu Y_\nu Y+\frac{1}{Dn2}Vg_{\mu \nu }.$$
(22)
The full system of $`(Dn)`$-dimensional equations of motion can therefore be derived from the Lagrangian
$$=R\text{1}\mathrm{l}\frac{2(Dn2)}{(n+1)(D2)a^2}Y^2dYdY\frac{1}{4}\underset{i}{}\stackrel{~}{X}_i^2d\stackrel{~}{X}_id\stackrel{~}{X}_iV\text{1}\mathrm{l}.$$
(23)
It remains to check the consistency of the other components of the $`D`$-dimensional Einstein equations. After making use of the lower-dimensional equations of motion for the scalar fields, we find that the internal components $`\widehat{R}_{\alpha \beta }`$ of the higher-dimensional Ricci tensor agree precisely with the expression for $`\widehat{S}_{\alpha \beta }`$ that follows from substituting the Ansäze for $`\widehat{F}_{\left(n\right)}`$ and $`\widehat{\varphi }`$, given in (6), into (21). Again, we have made extensive use of formulae derived in , appropriately modified to the case under consideration here. Finally, we note that the mixed components $`\widehat{R}_{\alpha \mu }`$ in (2) agree precisely with the mixed components of $`\widehat{S}_{\alpha \mu }`$ given in (21).
With these calculations we have now obtained a complete and explicit proof that the Ansatz (6) yields a consistent Kaluza-Klein $`n`$-sphere reduction of the $`D`$-dimensional theory described by (2), with the lower-dimensional fields appearing in the Ansatz satisfying the equations of motion that follow from the $`(Dn)`$-dimensional Lagrangian (23).
## 3 Domain walls as distributions of $`p`$-branes
We find that the $`d`$-dimensional gravity/scalar Lagrangian (23) admits a domain wall solution, given by
$`ds_d^2`$ $`=`$ $`(gr)^{\frac{a^2(D2)}{2(d2)}}\left((gr)^{n3}h^{{\scriptscriptstyle \frac{1}{2(d2)}}}dx^\mu dx_\mu +h^{{\scriptscriptstyle \frac{d3}{2(d2)}}}{\displaystyle \frac{dr^2}{g^2r^2}}\right),`$
$`X_i`$ $`=`$ $`(gr)^{{\scriptscriptstyle \frac{a^2(D2)}{4(d2)}}}h^{{\scriptscriptstyle \frac{(d3)}{4(d2)}}}H_i^1,`$ (24)
where
$$h\underset{i=1}{\overset{n+1}{}}H_i,H_i=1+\frac{\mathrm{}_i^2}{r^2}.$$
(25)
In fact there is a redundancy in the paramtrisation of these solutions, which can be seen as follows. We make the following transformation of the radial coordinate,
$$r^2=R^2L^2,$$
(26)
where $`L`$ is a constant, and define new quantities as follows:
$$\stackrel{~}{H}_i1+\frac{\stackrel{~}{\mathrm{}}_i^2}{R^2},\stackrel{~}{h}\underset{i=1}{\overset{n+1}{}}\stackrel{~}{H}_i,\stackrel{~}{\mathrm{}}_i^2\mathrm{}_i^2L^2.$$
(27)
After straightforward calculations, we find that the solution (24) becomes
$`ds_d^2`$ $`=`$ $`(gR)^{\frac{a^2(D2)}{2(d2)}}\left((gR)^{n3}\stackrel{~}{h}^{{\scriptscriptstyle \frac{1}{2(d2)}}}dx^\mu dx_\mu +\stackrel{~}{h}^{{\scriptscriptstyle \frac{d3}{2(d2)}}}{\displaystyle \frac{dR^2}{g^2R^2}}\right),`$
$`X_i`$ $`=`$ $`(gR)^{{\scriptscriptstyle \frac{a^2(D2)}{4(d2)}}}\stackrel{~}{h}^{{\scriptscriptstyle \frac{(d3)}{4(d2)}}}\stackrel{~}{H}_i^1,`$ (28)
This is identical in form to the original solution (24), but with the redefined functions given in (27). Let us suppose that, without loss of generality, the parameters $`\mathrm{}_i`$ are ordered so that $`\mathrm{}_1^2\mathrm{}_2^2\mathrm{}\mathrm{}_{n+1}^2`$. If we choose the constant $`L`$ in the coordinate transformation (26) to be equal to $`\mathrm{}_{n+1}`$, then we see that the original solution with $`(n+1)`$ parameters $`\mathrm{}_i`$ (with $`1in+1`$) is really nothing but a solution with only $`n`$ parameters $`\stackrel{~}{\mathrm{}}_i^2`$ (with $`1in`$).
When $`a=0`$, which occurs for the cases $`(D,n)=(11,4),(11,7)`$ and $`(10,5)`$, the resulting solutions become AdS domain walls. The metrics in these cases become asymptotically-AdS spacetimes in seven, four and five dimensions. These AdS domain-wall solutions are sphere reductions of the decoupling limits of ellipsoidal distributions of M-branes and D3-branes. These cases (and subsets) were studied previously in .
In this paper, we shall extend the previous analysis to include the cases where the dilaton-coupling constant $`a`$ is non-vanishing. For these cases, the domain-wall metric (24) is no longer asymptotically AdS, but instead is asymptotic to a vacuum domain wall as $`r\mathrm{}`$, given by
$$ds_d^2=\rho ^{\frac{4(n+1)}{a^2(D2)}}dx^\mu dx_\mu +g^2d\rho ^2.$$
(29)
where $`\rho (gr)^{\frac{a^2(D2)}{4(d2)}}`$. This metric is flat as $`\rho `$ approaches at infinity.
In the region near $`r=0`$, the metric structure depends on the number of non-vanishing parameters $`\mathrm{}_i`$. If $`k`$ of the $`\mathrm{}_i`$ are non-vanishing, we have
$$ds_d^2=\rho ^\gamma dx^\mu dx_\mu +d\rho ^2,$$
(30)
where
$$\gamma =\frac{4(n+1k)}{a^2(D2)+2(d3)k},\rho =(gr)^{\frac{a^2(D2)+2(d3)k}{4(d2)}}.$$
(31)
Thus we see that at $`r=0=\rho `$, the solution is generic singular. To see if the singularity is naked or not, we evaluate
$$\gamma 2=\frac{4(d2)(n3k)}{a^2(D2)+2(d3)k}.$$
(32)
Thus for $`n=0,1,2,3`$, the solution has a naked singularity for all values of $`k`$. For $`n4`$, the singularity is naked for $`k>n3`$, but marginal for $`kn3`$.
If we oxidise the solution back to $`D`$ dimensions, it acquires an interpretation as a continuous distribution of $`(Dn2)`$-branes, given by
$`ds_D^2`$ $`=`$ $`H^{{\scriptscriptstyle \frac{n1}{D2}}}dx^\mu dx_\mu +H^{\frac{Dn1}{D2}}dy^mdy^m,`$
$`e^{\frac{2}{a}\varphi }`$ $`=`$ $`H.F_{\left(n\right)}=e^{a\varphi }\widehat{}(d^{Dn1}xdH^1).`$ (33)
where $`H`$ and the transverse Euclidean metric are given by
$`dy^mdy^m`$ $`=`$ $`h^{{\scriptscriptstyle \frac{1}{2}}}\stackrel{~}{\mathrm{\Delta }}dr^2+r^2{\displaystyle H_id\mu _i^2},`$
$`H`$ $`=`$ $`{\displaystyle \frac{1}{(gr)^{n1}\stackrel{~}{\mathrm{\Delta }}}},\stackrel{~}{\mathrm{\Delta }}=h^{{\scriptscriptstyle \frac{1}{2}}}{\displaystyle \frac{\mu _i^2}{H_i}}.`$ (34)
The function $`H`$ is a harmonic function of the Euclidean transverse space, and it can be expressed as
$$H=g^{(n1)}\frac{\sigma (\stackrel{}{y}^{})d^{n+1}y^{}}{|\stackrel{}{y}\stackrel{}{y}^{}|^{n1}},$$
(35)
where $`\sigma (\stackrel{}{y})`$ is the distribution function. The harmonic functions in our cases here are associated with ellipsoidal distributions.
A detailed analysis is given in , where the charge-distribution functions are obtained in the non-dilatonic cases of 3-branes in $`D=10`$, and M2-branes and M5-branes in $`D=11`$. The analysis here is almost identical, and we shall not enumerate all the possibilities. It was observed in that although the results for the charge-distribution functions are distinctly different depending upon how many of the $`\mathrm{}_i`$ parameters are non-zero, by carefully taking limits in which some of the parameters are sent to zero one can view them all as being derived from a maximally-degenerate case with all $`(n+1)`$ parameters non-zero. The distribution function with all the $`\mathrm{}_i`$ non-vanishing is given by
$$\sigma _{n+1}=\frac{1}{V_n_{i=1}^{n+1}\mathrm{}_i}\delta ^{}(1\underset{i=1}{\overset{n+1}{}}\frac{y_i^2}{\mathrm{}_i^2}),$$
(36)
where $`V_n`$ is the volume of the $`n`$-sphere and refers to the derivative with respect to the $`\delta `$-function argument. This same charge distribution arises in our present cases, too.
As an example, let us consider what happens if one of the parameters, say $`\mathrm{}_{n+1}`$ is sent to zero. It is clear from (36) that the integration in (35) over the associated direction $`y_{n+1}^{}`$ will become dominated by the contribution from $`y_{n+1}^{}`$ close to zero, and so the $`(n+1)`$-parameter charge distribution $`\sigma _{n+1}`$ in the $`\mathrm{}_{n+1}0`$ limit will become the $`n`$-parameter distribution
$$\sigma _n(y_1,\mathrm{},y_{n+1})=\delta (y_{n+1})_{\mathrm{}}^{\mathrm{}}𝑑\stackrel{~}{y}_{n+1}\sigma _{n+1}(y_1,\mathrm{},y_n,\stackrel{~}{y}_{n+1}).$$
(37)
Evaluating the integral, we obtain
$$\sigma _n=\frac{1}{2V_n_{i=1}^n\mathrm{}_i}\left(2(1\underset{i=1}{\overset{n}{}}\frac{y_i^2}{\mathrm{}_i^2})^{1/2}\delta (1\underset{i=1}{\overset{n}{}}\frac{y_i^2}{\mathrm{}_i^2})(1\underset{i=1}{\overset{n}{}}\frac{y_i^2}{\mathrm{}_i^2})^{3/2}\mathrm{\Theta }(1\underset{i=1}{\overset{n}{}}\frac{y_i^2}{\mathrm{}_i^2})\right)\delta (y_{n+1}).$$
(38)
Sending another parameter, say $`\mathrm{}_n`$ to zero, we next obtain the $`(n1)`$-parameter charge distribution
$$\sigma _{n1}=\frac{\pi }{V_n_{i=1}^{n1}\mathrm{}_i}\delta (1\underset{i=1}{\overset{n1}{}}\frac{y_i^2}{\mathrm{}_i^2})\delta ^{\left(2\right)}(y_n,y_{n+1}).$$
(39)
Further details of the successive results for smaller numbers $`k`$ of non-vanishing parameters $`\mathrm{}_i`$ are given in . Note that the distributions associated with $`k=n+1`$ and $`k=n`$ non-vanishing $`\mathrm{}_i`$ parameters both have regions with negative as well as positive $`p`$-brane tensions. For $`kn1`$, on the other hand, the distributions contain only positive tensions.
When all the parameters $`\mathrm{}_i`$ vanish, corresponding to the “vacuum” domain-wall solution in $`d=Dn`$ dimensions, the $`D`$-dimensional solution describes coincident $`(Dn2)`$-branes at the origin, with the constant 1 in the harmonic function $`H`$ dropped. This can be viewed as a certain decoupling limit. The metric of the solution in the Einstein frame can then be expressed as
$$ds_\mathrm{E}^2=e^{\frac{a}{n1}\varphi }\left(r^{n3}dx^\mu dx_\mu +\frac{dr^2}{r^2}+d\mathrm{\Omega }_n^2\right).$$
(40)
One can then define a dual frame $`ds_{\mathrm{dual}}^2=e^{a\varphi /(n1)}ds_\mathrm{E}^2`$, in which the Lagrangian becomes
$$=ee^{\frac{a(D2)}{2(n1)}\varphi }\left(R+\frac{(D2)(n^2nDn+3D2)}{2(n1)^2}(\varphi )^2\frac{1}{2n!}F_{\left(n\right)}^2\right).$$
(41)
In this dual frame, the metric is AdS$`\times S^n`$ if $`n3`$, and Minkowski$`\times S^3`$ when $`n=3`$. This analysis was given in detail in for $`D=10`$, leading to the conjecture of a Domain-wall/QFT correspondence. Further studies of the Domain-wall/QFT correspondence in general dimensions were given in .
It is of interest to note that in the dual frame, the metric depends only on the dimension $`n`$ of the internal sphere, but it is independent of $`D`$; the $`D`$-dependence of the Einstein-frame metric can all extracted as a conformal factor. Note that the dual frame metric has qualitative differences in the three situations $`n>3`$, $`n=3`$ and $`n<3`$. For $`n=3`$, the dual frame is Minkowskian, whilst for $`n3`$, the spacetime is AdS. However, for $`n>3`$ we have that $`r=0`$ is the horizon, whilst for $`n<3`$ the horizon is instead at $`r=\mathrm{}`$. These qualitative differences have significance for the structure of the spectrum in the dual QFT, which we shall analyse in the next section.
When the $`\mathrm{}_i`$ parameters are non-vanishing, the metric of the distributed branes in the dual frame is given by
$$ds_{\mathrm{dual}}^2=\stackrel{~}{\mathrm{\Delta }}^{\frac{n3}{n1}}ds_d^2+g^2\stackrel{~}{\mathrm{\Delta }}^{\frac{2}{n1}}H_i^2d\mu _i^2,$$
(42)
where
$$ds_d^2=(gr)^{n3}dx^\mu dx_\mu +\frac{dr^2}{(gr)^2h^{{\scriptscriptstyle \frac{1}{2}}}}.$$
(43)
Again we see that the metric does not manifestly depend on $`D`$, but on $`n`$ instead.
## 4 Analysis of the spectrum
A minimally-coupled scalar field $`\mathrm{\Phi }`$ obeys the wave equation
$$_\mu (\sqrt{g}g^{\mu \nu }_\nu \mathrm{\Phi })=0.$$
(44)
We make the Ansatz $`\mathrm{\Phi }=e^{ipx}\chi (r)`$, where $`m^2=pp`$ determines the mass of the fluctuating mode, and so the wave equation has the following general form
$$r^1_r\left[r^1\underset{i=1}{\overset{n+1}{}}\sqrt{r^2+\mathrm{}_i^2}_r\chi \right]=𝒬\chi ,$$
(45)
where $`𝒬=m^2g^{{\scriptscriptstyle \frac{1}{2}}(n+1)}`$. Remarkably, the wave equation depends only on the dimension of the internal sphere, but otherwise is independent of details of the original higher-dimensional theory.
It is helpful to cast the wave equation into the Schrödinger form, which can be done by first writing the metric in a manifestly conformally-flat frame as
$$ds^2=e^{2A(z)}(dx^\mu dx_\mu +dz^2),$$
(46)
by means of an appropriate coordinate transformation. The coordinate $`z`$ runs from 0 to $`z^{}`$, and $`A(z)`$ has the following asymptotic behaviour:
$`e^{2A}(zz^{})^\gamma ^{},\gamma ^{}={\displaystyle \frac{2(n+1)}{(d2)(n3)}}`$ $`\mathrm{for}`$ $`zz^{},`$
$`e^{2A}z^{\stackrel{~}{\gamma }},\stackrel{~}{\gamma }={\displaystyle \frac{2\gamma }{2\gamma }}={\displaystyle \frac{2(n+1k)}{(d2)(n3k)}},`$ $`\mathrm{for}`$ $`z0.`$ (47)
Making the field redefinition $`\chi =e^{(D2)A/2}\psi `$, the wave equation assumes the form
$$(^2V)\psi =\frac{1}{4}𝒬\psi ,$$
(48)
with the Schrödinger potential given by
$$V=\frac{d2}{2}A^{\prime \prime }+\frac{(d2)^2}{4}(A^{})^2.$$
(49)
The asymptotic behavior of the potential is given by
$`V`$ $``$ $`{\displaystyle \frac{c^{}}{(zz^{})^2}},\text{for}zz^{},`$
$`V`$ $``$ $`{\displaystyle \frac{c}{(z\stackrel{~}{z})^2}},\text{for}z\stackrel{~}{z},`$ (50)
where
$$c^{}=\frac{1}{4}+\frac{(n1)^2}{(n3)^2},c=\frac{1}{4}+\frac{(n1k)^2}{(n3k)^2}\frac{1}{4}.$$
(51)
The range of the coordinate $`z`$ is determined by the values of $`z^{}`$ and $`\stackrel{~}{z}`$, which in the original coordinate $`r`$ correspond to $`r\mathrm{}`$ and $`r0`$ limit, respectively. It is understood that if $`z^{}`$ or $`\stackrel{~}{z}`$ equals $`\pm \mathrm{}`$, the potential in (50) is of the form $`\pm 1/z^2`$.
Note that for $`n3`$ \[$`n4`$\] the limit $`r\mathrm{}`$ corresponds to $`zz^{}`$ with $`z^{}=\mathrm{}`$ \[$`z^{}=\mathrm{finite}`$\]. On the other hand for $`nk3`$ \[$`nk4`$\] the limit $`r0`$ corresponds to $`z\stackrel{~}{z}`$ where $`\stackrel{~}{z}=0`$ \[$`\stackrel{~}{z}=\mathrm{}`$\]. When $`n=3`$ or $`k=n3`$, where the denominator of the above expression vanishes, the coordinate $`z`$ depends logarithmically on the original coordinate $`r`$ ($`z\mathrm{log}(r)`$) and the Schrödinger potential becomes constant: $`V=1/4`$.
Note that since the wave equation is independent of $`D`$, whilst the metric depends on $`D`$, it may be more instructive to perform a field redefinition directly on the wave equation (45). This can be done by first defining $`y=r^2`$, and then introducing a new coordinate $`z`$ defined by $`y/z=\sqrt{f(y)}`$, where $`f(y)=[_{i=1}^{n+1}(y+\mathrm{}_i^2)]^{1/2}`$. (These are the defining equations that relate $`z`$ and $`r`$ coordinates.) The Schrödinger potential is then given by
$$V=\frac{1}{4}_z^2\mathrm{log}f+\frac{1}{16}(_z\mathrm{log}f)^2.$$
(52)
and it clearly depends on $`n`$ and $`\mathrm{}_i`$ ($`i=1,\mathrm{},k`$) only.
### 4.1 Vacuum excitations
When all the parameters $`\mathrm{}_i`$ vanish, the solution (24) becomes a domain-wall vacuum solution. In the case when $`a^2=0`$, which occurs for $`(D,n)=(11,7)`$, $`(10,5)`$ and $`(11,4)`$, the solution is just the AdS spacetime in $`d=4`$, 5 and 7 respectively. For $`a^20`$, the metric of the solution is (30). The metric is flat near $`\rho =\mathrm{}`$, but becomes singular as $`\rho `$ approaches zero. Since we have
$$\frac{4(n+1)}{a^2(D2)}2=\frac{4(d2)(n3)}{a^2(D2)},$$
(53)
the singularity is marginal for $`n3`$, but naked for $`n<3`$.
The characteristics of the Schrödinger potential depend only the value of $`n`$. For $`n=0,1,2`$, the potential is given by
$$V=\frac{c^{}}{z^2},$$
(54)
where $`c^{}`$ is given in (51). The coordinate $`z`$ runs from 0 to infinity as $`r`$ runs from 0 to infinity. For $`n=3`$, the potential is a constant, $`V=1/4`$, and the coordinate $`z`$ runs from minus infinity to infinity as $`r`$ runs from 0 to infinity. For $`n4`$, the potential is of the same form as (54), but the coordinate $`z`$ now runs from minus infinity to 0 as $`r`$ runs from 0 to infinity. Thus we see that although the domain-wall vacuum can have (naked) singularities, the quantum fluctuations are nevertheless well behaved. In fact it is straightforward to solve the minimally-coupled scalar wave equation in the domain-wall vacuum, namely
$$r^1_r(r^n_r\chi )=𝒬\chi .$$
(55)
If we define a new dependent variable $`y`$ by
$$\chi (r)=y(r)r^{(n1)/2},$$
(56)
and change to the new independent variable $`z`$ defined by
$$z=\frac{2\sqrt{𝒬}}{n3}r^{(n3)/2},$$
(57)
then the wave equation (55) becomes the Bessel equation
$$z^2y^{\prime \prime }(z)+zy^{}(z)+(z^2\nu ^2)y(z)=0,$$
(58)
where
$$\nu =\frac{n1}{n3}.$$
(59)
The solutions to (55) are therefore given by
$$\chi (r)=ar^{(n1)/2}J_\nu \left(\frac{2\sqrt{𝒬}}{n3}r^{(n3)/2}\right)+br^{(n1)/2}Y_\nu \left(\frac{2\sqrt{𝒬}}{n3}r^{(n3)/2}\right).$$
(60)
A special case arises for $`n=3`$ (the Schrödinger potential is constant, $`V=1/4`$, there) for which we find
$$\chi (r)=ar^{1+\mathrm{i}\sqrt{𝒬1}}+br^{1\mathrm{i}\sqrt{𝒬1}}.$$
(61)
The requirement that $`Q1`$ corresponds to the condition that there is an energy gap.
### 4.2 Domain-wall excitations
When some of the $`\mathrm{}_i`$ parameters are non-vanishing, the wave equations cannot in general be solved explicitly. Here, for simplicity, we shall consider the case where all the non-vanishing $`\mathrm{}_i`$ are equal. There are certain examples where the wave equations can be solved exactly. Two of these ($`n=5`$, $`k=2`$ with two equal charges $`\mathrm{}_i`$, and $`n=5`$, $`k=4`$ with four equal charges) are solved in . Another solvable example is $`n=3`$, $`k=2`$, with the two non-vanishing charges equal, say $`\mathrm{}_1=\mathrm{}_2\mathrm{}`$. In this case, if we let $`x=r^2/\mathrm{}^2`$, equation (45) becomes the hypergeometric equation
$$x(1x)\chi ^{\prime \prime }+(12x)\chi ^{}\frac{1}{4}𝒬\chi =0,$$
(62)
and so one solution gives
$$\chi _1={}_{2}{}^{}F_{1}^{}[a,b;1;\frac{r^2}{\mathrm{}^2}],a=\frac{1}{2}+\frac{\mathrm{i}}{2}\sqrt{𝒬1},b=\frac{1}{2}\frac{\mathrm{i}}{2}\sqrt{𝒬1}.$$
(63)
Note again that $`Q>1`$ corresponds to the condition that the (continuous) spectrum has a gap owing to the properties of the Schrödinger potential $`V1/4`$ (figure g). Note that at small $`r`$ we therefore have $`\chi _11`$, while at large $`r`$ the asymptotic behaviour is of the same form as in (61). Since the $`c`$ argument of the hypergeometric function $`{}_{2}{}^{}F_{1}^{}[a,b;c;x]`$ in (63) is an integer, the second solution $`\chi _2`$ of (62) must be obtained by taking an appropriate limit of the standard second solution $`x_2^{1c}F_1[ac+1,bc+1;2c;x]`$. This gives a logarithmic behaviour of the form $`\chi _2\mathrm{log}r`$ at small $`r`$.
For the remaining examples, although we cannot solve the wave equation analytically we can determine the structure of the spectra for the various cases from the forms of their Schrödinger potentials. The results are summarised in Table 1.
$`n`$ $`k`$ $`z`$-range $`V`$ type Spectrum 0,1 0,1 $`(0,\mathrm{})`$ a continuous 2 0 $`(0,\mathrm{})`$ b continuous 1,2 $`(0,\mathrm{})`$ c continuous 3 0 $`(\mathrm{},\mathrm{})`$ $`V=\frac{1}{4}`$ cont. with gap 1 $`(0,\mathrm{})`$ d disc., cont. with gap 2,3 $`(0,\mathrm{})`$ e cont. with gap $`4`$ $`n4`$ $`(\mathrm{},0)`$ f continuous $`n2`$ $`(\mathrm{},0)`$ g cont. with gap $`n3`$ $`(1,0)`$ h discrete $`n,n1`$ $`(1,0)`$ i discrete
Table 1: Spectral analysis for domain-wall solutions for various $`n`$’s and $`k`$’s.
The various different types of structures of the potentials are sketched in Figure 1.
Figure 1: Sketches of the various Schrödinger potentials
## 5 Conclusions
In this paper we have studied consistent $`n`$-sphere reductions of a $`D`$-dimensional theory of gravity coupled to an $`n`$-form field-strength and a dilaton. Provided that the dilaton has a specific strength of coupling to the $`n`$-form, given by (3), we have proven the consistency of the non-linear Kaluza-Klein Ansatz for the $`n`$-sphere reduction in which there are $`n`$ scalars parameterising right-ellipsoidal inhomogeneous deformations of the sphere.<sup>2</sup><sup>2</sup>2We did not turn on the Kaluza-Klein gauge-fields in the reduction, which corresponds to a consistent truncation of the theory. However in the appendix we also discuss an $`n`$-sphere reduction of this Lagrangian that corresponds to making pair-wise identifications of the diagonal scalar fields, together with turning on the electric components of the Abelian Kaluza-Klein fields. This reduction provides a $`D`$-dimensional embedding of the $`(Dn)`$-dimensional non-extreme (large) charged-black holes as (near-extreme) spinning electric $`(Dn2)`$-branes. In the BPS limit the charged black holes become neutral BPS domain-wall solutions.
The generality of these consistent reductions provides a framework within which we can address the $`D`$-dimensional embedding of a class of solutions of the reduced gauged supergravities in $`d=Dn`$ dimensions. In general, these gauged supergravities have potentials for the scalar fields that do not admit AdS ground-states, and thus in general, the typical solutions correspond to domain walls that are asymptotic to the “dilatonic” vacuum. In particular, we found the general class of BPS domain-wall solutions that are specified by $`k=\{1,\mathrm{},n\}`$ parameters, which characterise the harmonic functions of $`k`$ non-trivial scalars.
All these solutions have explicit representations as continuous distributions of extremal $`(Dn2)`$-branes, and thus in the context of the Domain-wall/QFT correspondence they describe the Coulomb phase of the dual strongly-coupled field theory.
The universal properties of these gravity solutions manifest themselves in the properties of the wave equations in these backgrounds. For minimally-coupled scalars, the wave equations are completely universal and depend only on the dimension $`n`$ of the compactifying sphere and the number $`k`$ of parameters in the harmonic functions specifying the non-trivial scalar fields. Remarkably, the wave equations are independent of the original dimension $`D`$. Thus in the dual field theory the bound-state spectrum is completely specified by $`n`$ and $`k`$. We gave an analysis of the spectra for all these cases.
One of the interesting outcomes of our study is the generality and universality of the BPS solutions for the specific subsector of the sphere-reduced gravity theories. This provides a strong indication that the dual field theories should exhibit the same intriguing features, irrespective of the dimension.
## Acknowledgments
We should like to thank K. Skenderis for a useful discussion. C.N.P. would like to thank High Energy Theory Group at Penn for hospitality.
## Appendix A Single-charge rotating $`p`$-branes
The Lagrangian (2) also admits rotating $`p`$-brane solutions. In this appendix, we show that such a rotating $`p`$-brane associated with $`a`$ given by (3) can be dimensionally reduced on the transverse spherical space, and it then gives rise to a domain-wall black hole in the lower dimension. The Lagrangian (2) admits an electric $`(d1)`$-brane with $`d=n1`$, or a magnetic $`(d1)`$-brane with $`d=Dn1`$. We shall consider only the magnetic solution here, since the electric one can be viewed as a magnetic solution of the dual $`(Dn)`$-form field strength $`F_{\left(Dn\right)}`$. There are two cases arising, depending on whether $`\stackrel{~}{d}`$ is even or odd.
Case 1: $`n=2N1`$
In this case, there are $`N`$ angular momenta $`\mathrm{}_i`$, with $`i=1,2,\mathrm{},N`$. We find that the metric of the rotating $`(n2)`$-brane solution to the equations following from (2) is
$`ds_D^2`$ $`=`$ $`H^{{\scriptscriptstyle \frac{n1}{D2}}}((1{\displaystyle \frac{2m}{r^{n1}\mathrm{\Delta }}})dt^2+d\stackrel{}{x}d\stackrel{}{x})+H^{{\scriptscriptstyle \frac{Dn1}{D2}}}[{\displaystyle \frac{\mathrm{\Delta }dr^2}{H_1\mathrm{}H_N\frac{2m}{r^{(n1)}}}}`$ (64)
$`+r^2{\displaystyle \underset{i=1}{\overset{N}{}}}H_i(d\stackrel{~}{\mu }_i^2+\stackrel{~}{\mu }_i^2d\varphi _i^2){\displaystyle \frac{4m\mathrm{cosh}\alpha }{r^{n1}H\mathrm{\Delta }}}dt({\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{}_i\stackrel{~}{\mu }_i^2d\varphi _i)`$
$`+{\displaystyle \frac{2m}{r^{n1}H\mathrm{\Delta }}}({\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{}_i\stackrel{~}{\mu }_i^2d\varphi _i)^2],`$
where the functions $`\mathrm{\Delta }`$, $`H`$ and $`H_i`$ are given by
$`\mathrm{\Delta }=H_1\mathrm{}H_N{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{\stackrel{~}{\mu }_i^2}{H_i}},H=1+{\displaystyle \frac{2m\mathrm{sinh}^2\alpha }{r^{n1}\mathrm{\Delta }}},`$
$`H_i=1+{\displaystyle \frac{\mathrm{}_i^2}{r^2}},i=1,2,\mathrm{},N.`$ (65)
The dilaton $`\varphi `$ and the field strength $`F_{\left(n\right)}`$ are given by
$$e^{2\varphi /a}=H,e^{a\varphi }F_{\left(n\right)}=\frac{dH^1}{\mathrm{sinh}\alpha }\left(\mathrm{cosh}\alpha dt+\underset{i=1}{\overset{N}{}}\mathrm{}_i\mu _i^2d\varphi _i\right)d^{Dn2}x.$$
(66)
The $`N`$ quantities $`\stackrel{~}{\mu }_i`$ are subject to the constraint $`_i\stackrel{~}{\mu }_i^2=1`$. They are related to our previous coordinates constrained $`\mu _i`$ on the sphere as follows:
$$\mu _1+\mathrm{i}\mu _2=\stackrel{~}{\mu }_1e^{\mathrm{i}\varphi _1},\mu _3+\mathrm{i}\mu _4=\stackrel{~}{\mu }_2e^{\mathrm{i}\varphi _2},\text{etc}.$$
(67)
We now consider the decoupling limit, which is obtained by making the rescalings
$`mϵ^{n1}m,\mathrm{sinh}\alpha ϵ^{(n1)/2}\mathrm{sinh}\alpha ,`$
$`rϵr,x^\mu ϵ^2x^\mu ,\mathrm{}_iϵ\mathrm{}_i`$ (68)
and then sending $`ϵ0`$. In this limit, the additive constant 1 in the function $`H`$ in (65) can be dropped. Furthermore, the last term in (64) can also be dropped. The remaining metric can be expressed as
$$ds_D^2=Y^{\frac{2}{D2}}\left(\mathrm{\Delta }^{\frac{n1}{D2}}ds_d^2+g^2\mathrm{\Delta }^{\frac{(Dn1)}{(D2)}}\underset{i=1}{\overset{N}{}}\overline{X}_i^1(d\stackrel{~}{\mu }_i^2+\stackrel{~}{\mu }_i^2(d\varphi _i+gA_{\left(1\right)}^i)^2)\right),$$
(69)
where $`\mathrm{\Delta }=\overline{X}_i\stackrel{~}{\mu }_i^2`$, and $`g=(2m\mathrm{sinh}^2\alpha )^{1/(n1)}`$. The $`d=Dn`$ dimensional metric and the scalar fields $`X_i`$ are given by
$`ds_d^2`$ $`=`$ $`h^{{\scriptscriptstyle \frac{d3}{d2}}}fdt^2+h^{{\scriptscriptstyle \frac{1}{d2}}}\left({\displaystyle \frac{dr^2}{(gr)^{5n}f}}+(gr)^{n3}d\stackrel{}{x}d\stackrel{}{x}\right),`$
$`X_i`$ $`=`$ $`(gr)^{{\scriptscriptstyle \frac{a^2(D2)}{4(d2)}}}h^{{\scriptscriptstyle \frac{(d3)}{2(d2)}}}H_i^1,`$
$`Y`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}X_i=\left((gr)^{n+1}h^1\right)^{{\scriptscriptstyle \frac{a^2(D2)}{8(d2)}}},`$
$`A_{\left(1\right)}^i`$ $`=`$ $`{\displaystyle \frac{1H_i^1}{g\mathrm{}_i\mathrm{sinh}\alpha }}dt,h={\displaystyle \underset{i=1}{\overset{N}{}}}H_i.`$
$`f`$ $`=`$ $`(gr)^{n3}(h{\displaystyle \frac{2m}{r^{n1}}}).`$ (70)
This solution describes $`N`$ electrically-charged black holes in a $`d`$-dimensional domain-wall background.
Note that in general, the abstract metric Ansatz that we have written in (69) does not (at least as it stands) correspond to part of a consistent Kaluza-Klein reduction. It can be viewed as a modification of the general consistent metric Ansatz in (6) in which we first (consistently) partially truncate the scalars by setting them equal in pairs ($`X_1=X_2=\overline{X}_1`$, $`X_3=X_4=\overline{X}_2`$, etc)̇. Then, having also made the redefinitions (67), we introduce a $`U(1)`$ gauge field $`A_{\left(1\right)}^i`$ associated with the rotation in each of the original 2-planes $`(\mu _1,\mu _2)`$, $`(\mu _3,\mu _4)`$, etc. (Although we have not presented it here, we can also straightforwardly carry out the same steps on the original Ansätze for $`\widehat{\varphi }`$ and $`\widehat{F}_{\left(n\right)}`$ in (6) too.) This does give a consistent reduction in the case $`(10,5)`$ discussed in , but in general additional fields would have to be included too. The reason for this is that the $`U(1)`$ gauge fields, in quadratic products of the form $`F_{\left(2\right)}^iF_{\left(2\right)}^j`$, will act as sources for other fields. In the special case of $`(D,n)=(10,5)`$, they actually act as sources for themselves (corresponding to cubic Chern-Simons terms in the five-dimensional theory), but in the other cases they will act as sources for additional fields, requiring a larger set of fields in the Kaluza-Klein reduction Ansatz.
The metric (69), together with analogously-obtained expressions for $`\widehat{\varphi }`$ and $`\widehat{F}_{\left(n\right)}`$, is nevertheless still usable in appropriate circumstances. The problematic terms $`F_{\left(2\right)}^iF_{\left(2\right)}^j`$ actually vanish for our specific domain-wall black hole solutions since all the $`U(1)`$ charges are purely electric. This means that these particular lower-dimensional configurations will lift to the higher dimension without necessitating the turning-on of the additional fields that would be needed for a fully-consistent Ansatz, but which have been omitted in our discussion. Thus we still have an exact embedding of these specific solutions in the higher dimension.
Case 2: $`n=2N`$
Here, the solution has the same form as (64), but with the range of the index $`i`$ extended to include 0. However, there is no angular momentum parameter or azimuthal coordinate associated with the extra index value, and so $`\mathrm{}_0=0`$ and $`\varphi _0=0`$. The $`\stackrel{~}{\mu }_i`$ and $`\varphi _i`$ coordinates are now related to the original coordinates $`\mu _i`$ on the sphere by
$$\mu _0=\stackrel{~}{\mu }_0,\mu _1+\mathrm{i}\mu _2=\stackrel{~}{\mu }_1e^{\mathrm{i}\varphi _1},\mu _3+\mathrm{i}\mu _4=\stackrel{~}{\mu }_2e^{\mathrm{i}\varphi _2},\text{etc}.$$
(71)
Otherwise, all the formulae in Case 1 generalise to this case, simply by extending the summation to span the range $`0iN`$. Of course $`H_0=1`$ as a consequence of $`\mathrm{}_0=0`$.
Note that for $`a=0`$, we have $`(D,n)=(11,7),(11,4)`$ and $`(10,5)`$. These correspond to the rotating M-branes and D3-branes .
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# Warm Inflation: Towards a realistic COBE data power spectrum for matter and metric thermal coupled fluctuations
## I Introduction
Warm inflation takes into account separately, the matter and radiation energy fluctuations. In this scenario the matter field $`\phi `$ interacts with the particles of a thermal bath with mean temperature $`T_r`$, which is smaller than the Grand Unified Theories (GUT) critical temperature $`T_r<T_{GUT}10^{15}`$ GeV. This scenario was firstly studied by Berera. The warm inflation scenario served as a explicit demonstration that inflation can occur in the presence of a thermal component. In the formalism developed by Berera the temperature of the universe remains constant during the inflationary expansion.
Warm inflation was originally formulated in a phenomenological setting, but some attemps of a fundamental justification has also been presented. Furthermore, a dynamical system analysis showed that a smooth transition from inflationary to a radiation phase is attained for many values of the friction parameter, thereby showing that the warm inflation scenario may be a workable variant to standard inflation. During the warm iflationary era, vacuum fluctuations on scales smaller than the size of the horizon are magnified into classical perturbations on scales bigger than the Hubble radius. The classical perturbations can lead to an effective curvature of spacetime and energy density perturbations .
In an alternative formalism for warm inflation, I studied a model where the mean temperature and the amplitude of the temperature’s fluctuations decreases with time for a rapid power - law expanding universe. This is the most significative difference with the Berera’s formalism in which the warm inflation expansion is isothermal. During the warm inflationary expansion, the kinetic energy density $`\rho _{kin}`$ is smaller with respect to the vacuum energy
$$\rho (\phi )\rho _mV(\phi )\rho _{kin}.$$
The kinetic energy density is given by
$$\rho _{kin}=\rho _r(\phi )+\frac{\dot{\phi }^2}{2},$$
where
$$\rho _r(\phi )=\frac{\tau (\phi )}{8H(\phi )}\dot{\phi }^2.$$
(1)
Here, the dot denotes the derivative with respect to the time. Furthermore, $`\tau (\phi )`$ and $`H(\phi )`$ are the friction and Hubble parameters. The eq. (1) comes from the assumption that the radiation energy density remains constant during the inflationary era ($`\dot{\rho }_r0`$).
The Lagrangian that describes the warm infaltion scenario is
$$(\phi ,\phi _{,\mu })=\sqrt{g}\left[\frac{R}{16\pi }+\frac{1}{2}g^{\mu \nu }\phi _{,\mu }\phi _{,\nu }+V(\phi )\right]+_{int},$$
(2)
where $`R`$ is the scalar curvature, $`g^{\mu \nu }`$ gives the metric tensor and $`g`$ is the metric. The Lagrangian $`_{int}`$ takes into account that the particles in the thermal bath interact with the scalar field $`\phi `$. In principle, a permanent or temporary coupling of the scalar field $`\phi `$ with others fields might also lead to dissipative processes producing entropy at different eras of the cosmic evolution. It is expected that progress in nonequilibrium statistics of quantum fields will provide the necessary theoretical framework for discussing dissipation in more general cases.
The semiclassical Friedmann equation is
$$H^2(\phi )=\frac{8\pi }{3M_p^2}E\left|\rho _m(\phi )+\rho _r(\phi )\right|E,$$
(3)
where $`M_p=1.2\times 10^{19}`$ GeV is the Planckian mass.
Now I consider the semiclassical expansion for the inflaton field $`\phi `$
$$\phi (\stackrel{}{x},t)=\varphi _c(t)+\alpha (t)\varphi (\stackrel{}{x},t).$$
(4)
Here, $`\varphi _c(t)=<E|\phi |E>`$, $`<E|\varphi |E>=<E|\dot{\varphi }|E>=0`$, and $`|E>`$ is an arbitrary state. Furthermore, $`\alpha (t)`$ is a dimensionless time - dependent function that characterize the gravitational coupling between the fluctuations of the matter field and the fields in the thermal bath. A lot of work can be done on phenomenological grounds, as, for instance, by applying nonequilibrium thermodynamic techniques to the problem or even studying particular models with dissipation. An example of this latter case is the warm inflationary picture recently proposed.
The aim of this work is the study of the power spectrum in warm inflation with the semiclassical expansion (4), taking into account the COBE data coarse - grained field introduced in a previous work. This topic was studied in but with the semiclassical expansion $`\phi =\varphi _c+\varphi `$. In this work I incorporate in the formalism the backreaction of the metric for the study of the effective curvature for the now observable universe, when the fluctuations are coupled with the thermal bath.
## II Dynamics of the inflaton
### A Dynamics of the classical field
The dynamics for the classical field in warm inflation was obtained in previous works. The equation of motion for $`\varphi _c`$ is
$$\ddot{\varphi }_c+\left[3H_c+\tau _c\right]\dot{\varphi }_c+V^{}(\varphi _c)=0,$$
(5)
where $`H_cH(\varphi _c)+\frac{\dot{a}}{a}`$, and $`\tau _c\tau (\varphi _c)`$ and $`V^{}(\varphi _c)\frac{dV(\phi )}{d\phi }|_{\varphi _c}`$. The term $`\tau _c\dot{\varphi }_c`$ in eq. (5) shows as the scalar field evolves with the time in a damped regime generating an expansion which depends on the mean temperature $`T_r`$ of the thermal bath. As a consequence, the subsequent reheating mechanism is not needed and thermal fluctuations produce the primordial spectrum of density perturbations. Furthermore, $`\dot{\varphi }_c=\frac{M_p^2}{4\pi }H_c^{}\left(1+\frac{\tau _c}{3H_c}\right)^1`$ and the classical effective potential is
$`V(\varphi _c)`$ $`=`$ $`{\displaystyle \frac{M_p^2}{8\pi }}[H_c^2{\displaystyle \frac{M_p^2}{12\pi }}\left(H_c^{}\right)^2(1+{\displaystyle \frac{\tau _c}{4H_c}})`$ (6)
$`\times `$ $`(1+{\displaystyle \frac{\tau _c}{3H_c}})^2].`$ (7)
The radiation energy density of the background is
$$\rho _r\left[\varphi _c(t)\right]\frac{\tau _c}{8H_c}\left(H_c^{}\right)^2\left(\frac{M_p^2}{4\pi }\right)^2\left(1+\frac{\tau _c}{3H_c}\right)^2,$$
(8)
and the temperature of this background is
$$T_r\rho ^{1/4}\left[\varphi _c(t)\right].$$
(9)
Note that the temperature depends on time. In the warm inflation model here studied, I will suppose that it decreases with time, in agreement one expects in an expanding universe. The temporal evolution of the background temperature depends on the particular model which one considers. For example, in a power - law expanding universe $`T_rt^{1/2}`$.
### B First order $`\varphi `$ \- fluctuations
In this section I will study the first order $`\varphi `$ \- fluctuations for the matter field $`\phi `$, on a globally flat Friedmann -Robertson - Walker (FRW) metric
$$ds^2=dt^2+a^2d\stackrel{}{x}^2,$$
(10)
which describes a globally isotropic and homogeneous spacetime. The equation of motion for the quantum perturbations $`\varphi `$, is
$`\ddot{\varphi }+\left[2{\displaystyle \frac{\dot{\alpha }}{\alpha }}+(3H_c+\tau _c)\right]\dot{\varphi }{\displaystyle \frac{1}{a^2}}^2\varphi `$ (11)
$`+`$ $`\left[(3H_c+\tau _c){\displaystyle \frac{\dot{\alpha }}{\alpha }}+V^{\prime \prime }(\varphi _c)\right]\varphi =0.`$ (12)
The function $`\alpha (t)`$ depends on time. I consider $`\alpha (t)=F[T_r(t)/M]`$, where $`T_r(t)`$ is the temperature of the background and $`M10^{15}`$ GeV is the GUT mass. The structure of the equation (12) can be simplified by means of the map $`\chi =e^{3/2{\scriptscriptstyle \left(H_c+\tau _c/3+{\scriptscriptstyle \frac{2\dot{\alpha }}{3\alpha }}\right)𝑑t}}\varphi `$
$$\ddot{\chi }\frac{1}{a^2}^2\chi \mu ^2(t)\chi =0,$$
(13)
where $`\mu ^2(t)=\frac{k_o^2}{a^2}`$ is the time dependent parameter of mass and $`k_o(t)`$ is the time dependent wave number which separates the long wavelength ($`kk_o`$) and the short wave ($`kk_o`$) sectors.
The square time dependent parameter of mass is
$$\mu ^2(t)=\frac{9}{4}\left(H_c+\frac{\tau _c}{3}\right)^2V^{\prime \prime }(\varphi _c)+\frac{3}{2}\left(\dot{H}_c+\frac{\dot{\tau }_c}{3}\right).$$
(14)
Note that $`\mu (t)`$ do not depends on the function $`\alpha (t)`$.
### C Second order $`\varphi `$ \- fluctuations and Backreaction
Making a second order $`\varphi `$ \- fluctuations expansion for $`\phi `$, one obtains the following semiclassical Friedmann equation
$$H_c^2+\frac{K}{a^2}=\frac{8\pi }{3M_p^2}E\left|\rho _m+\rho _r\right|E,$$
(15)
where $`K`$ is an effective curvature produced by the backreaction of the metric with the fluctuations of the scalar field. This curvature is given by
$`{\displaystyle \frac{K}{a^2}}`$ $`=`$ $`{\displaystyle \frac{8\pi }{3M_p^2}}[(1+{\displaystyle \frac{\tau _c}{8H_c}})({\displaystyle \frac{\dot{\alpha }^2}{2}}\varphi ^2+{\displaystyle \frac{\alpha ^2}{2}}\dot{\varphi }^2`$ (16)
$`+`$ $`\alpha \dot{\alpha }\varphi \dot{\varphi })+{\displaystyle \frac{\alpha ^2}{a^2}}\left(\stackrel{}{}\varphi \right)^2+{\displaystyle \frac{V^{\prime \prime }}{2}}\alpha ^2\varphi ^2].`$ (17)
Note that $`K`$ depends on the temporal evolution of $`\alpha (t)`$ as well as the expectation values for $`\varphi ^2`$, $`\varphi \dot{\varphi }`$, $`\dot{\varphi }^2`$ and $`\left(\stackrel{}{}\varphi \right)^2`$. If $`\alpha (t)`$ is a function of the temperature, $`\alpha (t)=F[T_r(t)/M]`$, the instantaneous comoving temperature will be very important during the warm inflationary regime.
To study the backreaction of the metric with the fluctuations $`\varphi `$, I introduce the metric
$$ds^2=dt^2+a^2\left[1+h(\stackrel{}{x},t)\right]d\stackrel{}{x}^2,$$
(18)
where $`h(\stackrel{}{x},t)`$ represents the fluctuations of the metric produced by the $`\varphi `$ \- fluctuations. Making the following expansion for $`H(\phi )`$
$$H[\phi (\stackrel{}{x},t)]H_c[\varphi _c(t)]+H^{}[\varphi _c(t)]\alpha (t)\varphi (\stackrel{}{x},t),$$
(19)
one obtains the following expression for $`h(\stackrel{}{x},t)`$
$$h(\stackrel{}{x},t)2^t𝑑t^{}\alpha (t^{})\varphi (\stackrel{}{x},t^{})H^{}[\varphi _c(t^{})],$$
(20)
and the effective curvature can be represented by
$`{\displaystyle \frac{K}{a^2}}`$ $`=`$ $`E\left|\left({\displaystyle \frac{\dot{h}(\stackrel{}{x},t)}{2}}\right)^2\right|E`$ (21)
$`=`$ $`E\left|\left[\alpha (t)\varphi (\stackrel{}{x},t)H^{}(\varphi _c)\right]^2\right|E.`$ (22)
This expression shows that the temporal evolution of the effective curvature arises from the matter field fluctuations $`\varphi (\stackrel{}{x},t)`$ and the temperature of the thermal bath, due to the fact I am considering that $`\alpha (t)`$ is a function of the temperature of this bath. In order to study the evolution of the fluctuations on the infrared (long wavelength) sector, firstly one can write the fields $`\chi `$ and $`h`$ as two Fourier expanded fields
$`\chi (\stackrel{}{x},t)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle d^3k\left[a_k\chi _k(\stackrel{}{x},t)+a_k^{}\chi _k^{}(\stackrel{}{x},t)\right]},`$ (23)
$`h(\stackrel{}{x},t)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle d^3k\left[a_kh_k(\stackrel{}{x},t)+a_k^{}h_k^{}(\stackrel{}{x},t)\right]},`$ (24)
where $`\chi _k(\stackrel{}{x},t)=e^{\mathrm{i}\stackrel{}{k}.\stackrel{}{x}}\xi _k(t)`$ and $`h_k(\stackrel{}{x},t)=e^{\mathrm{i}\stackrel{}{k}.\stackrel{}{x}}\stackrel{~}{\xi }_k(t)`$. Here, $`\stackrel{~}{\xi }_k(t)=2^t𝑑t^{}\alpha (t^{})H^{}[\varphi _c(t^{})]\xi _k(t^{})`$ and the asterisk denote the complex conjugate. The operators $`a_k`$ and $`a_k^{}`$ are the well known annihilation and creation operators, which satisface $`[a_k,a_k^{}^{}]=\mathrm{i}\delta ^{(3)}(kk^{})`$ and $`[a_k^{},a_k^{}^{}]=[a_k,a_k^{}]=0`$. The commutation relations for the fields $`\chi `$ and $`h`$ are
$`[\chi (\stackrel{}{x},t),\dot{\chi }(\stackrel{}{x}^{},t)]=`$ (25)
$`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle d^3k\left(\xi _k\dot{\xi }_k^{}\dot{\xi }_k\xi _k^{}\right)e^{\mathrm{i}\stackrel{}{k}.(\stackrel{}{x}\stackrel{}{x}^{})}},`$ (26)
$`[h(\stackrel{}{x},t),\dot{h}(\stackrel{}{x}^{},t)]=`$ (27)
$`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle d^3k\left(\stackrel{~}{\xi }_k\dot{\stackrel{~}{\xi }}_k^{}\dot{\stackrel{~}{\xi }}_k\stackrel{~}{\xi }_k^{}\right)e^{\mathrm{i}\stackrel{}{k}.(\stackrel{}{x}\stackrel{}{x}^{})}}.`$ (28)
To obtain $`[\chi (\stackrel{}{x},t),\dot{\chi }(\stackrel{}{x}^{},t)]=\mathrm{i}\delta ^{(3)}(\stackrel{}{x}\stackrel{}{x}^{})`$ in eq. (26), one requires that $`\left(\xi _k\dot{\xi }_k^{}\dot{\xi }_k\xi _k^{}\right)=\mathrm{i}`$.
## III Data COBE coarse - grained fields and Stochastic representation
The data COBE coarse - grained matter field $`\chi _{Ccg}`$ were introduced in a previous work
$$\chi _{Ccg}=\frac{1}{(2\pi )^{3/2}}d^3kG(k,t)\left[a_k\chi _k+a_k^{}\chi _k^{}\right].$$
(29)
Now we can introduce the data COBE coarse - grained field $`h_{Ccg}`$ for the fluctuations of the metric
$$h_{Ccg}=\frac{1}{(2\pi )^{3/2}}d^3kG(k,t)\left[a_kh_k+a_k^{}h_k^{}\right].$$
(30)
In eqs. (29) and (30) the suppression factor $`G(k,t)`$ is given by
$$G(k,t)=\sqrt{\frac{1}{1+\left(\frac{k_o(t)}{k}\right)^N}},$$
(31)
with $`N=mn`$. Causality places a strict constraint on the suppression index: $`N4n`$. A suppression factor like (31) also was founded in a model with cosmic strings plus cold or hot dark matter Furthermore, the square fluctuations for the data COBE coarse - grained matter field is
$`E\left|\chi _{Ccg}^2\right|E`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dk}{k}}𝒫_{\chi _{Ccg}}(k)`$ (32)
$`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _0^{k_o(t)}}𝑑kk^2\left|\xi _k(t)\right|^2G^2(k,t),`$ (33)
where the power spectrum $`𝒫_{\chi _{Ccg}}(k)`$ when the horizon exit is
$$𝒫_{\chi _{Ccg}}(k)=A(t_{})\left(\frac{k}{k_o(t_{})}\right)^nf(k).$$
(34)
Here, $`t_{}`$ denotes the time when the horizon entry, for which $`k_o(t_{})\pi H_o`$ in comoving scale. The parameters in eq. (34) are the amplitude $`A(t_{})`$ on time $`t_{}`$, the spectral index $`n`$, the suppression wavenumber $`k_o`$ and the suppression index $`m`$.
The stochastic equation for $`\chi _{Ccg}`$ is
$$\ddot{\chi }_{Ccg}\left(\frac{k_o(t)}{a(t)}\right)^2\chi _{Ccg}=\frac{N}{k_o(t)}\left[\xi _1+\xi _2\right],$$
(35)
where
$`\xi _1(\stackrel{}{x},t)`$ $`=`$ $`{\displaystyle \frac{\dot{k}_ok_o^N}{(2\pi )^{3/2}}}{\displaystyle d^3kk^NG^3(k,t)}`$ (36)
$`\times `$ $`\left[a_k\dot{\chi }_k+a_k^{}\dot{\chi }_k^{}\right],`$ (37)
$`\xi _2(\stackrel{}{x},t)`$ $`=`$ $`{\displaystyle \frac{\dot{k}_ok_o^N}{(2\pi )^{3/2}}}{\displaystyle d^3kk^NG^5(k,t)}`$ (38)
$`\times `$ $`\left[\left({\displaystyle \frac{k_o}{k}}\right)^N\left(3\dot{k}_o^22k_o\ddot{k}_o\right)+2\dot{k}_o^2(1N)2k_o\ddot{k}_o\right]`$ (39)
$`\times `$ $`\left[a_k\chi _k+a_k^{}\chi _k^{}\right].`$ (40)
The noises (37) and (40) arise from the increasing number of degrees of freedom of the ifrared sector from the short-wavelength sector. For the special case considered in eq. (31), $`\xi _1`$ is a colored noise, while $`\xi _2`$ gives non - local dissipation.
Since $`h2^t𝑑t^{}\alpha (t^{})H^{}(t^{})\varphi (\stackrel{}{x},t)`$, one can rewrite it as $`h2^t𝑑t^{}\stackrel{~}{\alpha }(t^{})H^{}(t^{})\chi (\stackrel{}{x},t)`$, where $`\stackrel{~}{\alpha }(t)=e^{3/2{\scriptscriptstyle 𝑑t(H_c+\tau _c/3+{\scriptscriptstyle \frac{2\dot{\alpha }}{3\alpha }})}}\alpha (t)`$. With this representation for $`h`$, the data COBE coarse - grained metric field $`h_{Ccg}`$ becomes
$$h_{Ccg}(\stackrel{}{x},t)2^t𝑑t^{}\stackrel{~}{\alpha }(t^{})H^{}(t^{})\chi _{Ccg}(\stackrel{}{x},t^{}).$$
(41)
Replacing (41) in eq. (35), one obtains the following stochastic equation for $`h_{Ccg}`$
$`{\displaystyle \frac{^3h_{Ccg}}{t^3}}+2{\displaystyle \frac{\dot{g}(t)}{g(t)}}{\displaystyle \frac{^2h_{Ccg}}{t^2}}\left({\displaystyle \frac{k_o}{a}}\right)^2{\displaystyle \frac{h_{Ccg}}{t}}`$ (42)
$`+`$ $`{\displaystyle \frac{\ddot{g}(t)}{g(t)}}h_{Ccg}={\displaystyle \frac{N}{g(t)k_o(t)}}\left[\xi _1(\stackrel{}{x},t)+\xi _2(\stackrel{}{x},t)\right],`$ (43)
where $`g(t)=\left[2\stackrel{~}{\alpha }(t)H^{}(t)\right]^1`$. The square fluctuations for the field $`\varphi _{Ccg}`$ is
$`E\left|\varphi _{Ccg}^2\right|E`$ $`=`$ $`{\displaystyle \frac{e^{3{\scriptscriptstyle \left(H_c+{\scriptscriptstyle \frac{\tau _c}{3}}+{\scriptscriptstyle \frac{2\dot{\alpha }}{3\alpha }}\right)𝑑t}}}{2\pi ^2}}`$ (44)
$`\times `$ $`{\displaystyle _0^{k_o}}𝑑kk^2\left|\xi _k^2(t)\right|G^2(k,t).`$ (45)
Thus, the effective curvature $`K/a^2`$ for the now observable universe is \[see eq. (22)\]
$$\frac{K}{a^2}|_{\mathrm{COBE}}=\left[\alpha (t)H^{}[\varphi _c(t)]\right]^2E\left|\varphi _{Ccg}^2\right|E.$$
(46)
Hence, the power spectrum for $`\varphi _{Ccg}`$ and $`h_{Ccg}`$ when the horizon entry, are
$`𝒫_{\varphi _{Ccg}}(k)`$ $`=`$ $`B(t_{})\left({\displaystyle \frac{k}{k_o(t_{})}}\right)^nf(k),`$ (47)
$`𝒫_{h_{Ccg}}(k)`$ $`=`$ $`C(t_{})\left({\displaystyle \frac{k}{k_o(t_{})}}\right)^nf(k).`$ (48)
Here, $`B(t_{})`$ and $`C(t_{})`$ are the amplitude such that
$`B(t_{})`$ $`=`$ $`A(t_{})e^{3^t_{}\left(H_c+\frac{\tau _c}{3}+\frac{2\dot{\alpha }}{3\alpha }\right)𝑑t},`$ (49)
$`C(t_{})`$ $`=`$ $`\left[\alpha (t)H^{}[\varphi _c(t)]\right]^2|_{t=t_{}}B(t_{}).`$ (50)
Due to $`\left|\delta _k\right|^2=𝒫_{\varphi _{Ccg}}(k)`$, the spectral density become $`|\delta _k|=k^nf(k)`$. The standard choice $`n=1`$ and $`f(k)`$ as constant, was invoked by Harrison and Zel’dovich on the grounds that it is scale invariant at the epoch of the horizon entry. The constraint $`|n1|<0.3`$ was obtained from the data COBE spectrum. Note that both $`B(t_{})`$ and $`C(t_{})`$ depends on the temperature of the background when the horizon entry. This is a very important characteristic that becomes from this formulation, once one consider $`\phi =\varphi _c+\alpha (t)\varphi `$ and $`H(\phi )=H_c+\alpha (t)H^{}\varphi `$ as semiclassical expansions for $`\phi `$ and $`H(\phi )`$. From eq. (50) one obtains
$$\frac{C(t_{})}{B(t_{})}=\left[\alpha (t)H^{}[\varphi _c(t)]\right]^2|_{t=t_{}}.$$
(51)
Taking $`\rho _r(t_{})=\frac{\pi ^2}{30}N[T_r(t_{})]T_r^4(t_{})`$, where $`N[T_r(t_{})]`$ is the number of relativistic degrees of freedom at temperature $`T_r(t_{})`$ and replacing $`(H_c^{})^2`$ in eq. (8), one obtains (for $`N[T_r(t_{})]10^3`$, $`\alpha (t_{})=\left(\frac{T_r(t_{})}{M}\right)^\beta `$$`\beta 0`$ — and $`M10^4M_p`$)
$$\frac{C(t_{})}{B(t_{})}=\frac{64}{45}\pi ^410^{3+8\beta }\frac{(3H_c+\tau _c)^2}{3H_c\tau _c}\left(\frac{T_r(t_{})}{M_p}\right)^{2(\beta +2)}.$$
(52)
For the case $`\tau _c(t_{})H_c(t_{})`$, one obtains the expression
$$\frac{C(t_{})}{B(t_{})}10^{(6+8\beta )}\left(\frac{T_r(t_{})}{M_p}\right)^{2(\beta +2)}.$$
(53)
For example, for $`\beta =1`$ one obtains $`T_r(t_{})=10^5M_p`$ for $`\frac{C(t_{})}{B(t_{})}10^{15}`$. This implies that the amplitude for the fluctuations of the metric when the horizon entry, should be very small for the expected values of temperature.
## IV Final Remarks
To summarize, in this letter I considered a model for warm inflation where the fluctuations of the scalar field are coupled with the thermal bath. This coupling, depends on the temperature of the background which is a function of the temperature. The temperature decreases with the time, as well as the Hubble parameter.
By means of the COBE data coarse - grained field for the fluctuations of the scalar field I characterize these fluctuations on the scale of the now observable universe. Once on knows the stochastic equation for the field $`\chi _{Ccg}`$, it is possible to obtain the stochastic equation for $`h_{Ccg}`$ \[see eq. (43)\]. The square fluctuations for $`\varphi _{Ccg}`$ and $`h_{Ccg}`$ give the spectral density $`\delta _k`$ for both, $`\varphi _{Ccg}`$ and $`h_{Ccg}`$. The spectral density $`\delta _k`$ depends on the function $`G^2(k,t)`$, the modes $`\xi _k`$ and the index $`n`$. Finally, I find that the amplitude for the fluctuations of the metric when the horizon entry, should be very small for the expected values of temperature.
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# POLARIMETRY AT A FUTURE LINEAR COLLIDER – HOW PRECISE?
## 1 Introduction
This paper addresses the issue of how accurately one should measure the electron beam polarization, $`P_e`$, at a future linear collider. The collider performance parameters shown in Table 1 are used. I begin by considering how accurately Standard Model (SM) asymmetries may be measured, and the requirements these measurements place on the polarimetry. In addition to Compton polarimetry,<sup>?</sup> the measurement of SM physics asymmetries for polarimetry is considered. This type of polarimetry has previously been proposed when both colliding beams are polarized.<sup>?,?</sup> Here, I additionally consider the possibility to use the asymmetry in forward W-pairs in $`e^+e^{}`$ collisions, when only the electron beam is polarized.
The precision of the polarimetry can affect the discovery potential for a new physics signal. I examine the importance of precise polarimetry for accurately assessing W-pair backgrounds.
Beam-beam effects in the collision process have a significant impact on polarim-etry. First, significant depolarization can result, and one needs to determine how the measured $`P_e`$is related to the luminosity-weighted polarization, $`P_e^{lum}`$. At high luminosities with large beamsstrahlung disruption, it may no longer be possible to place a Compton polarimeter in the extraction line from the Interaction Region (IR). This will limit the ability to accurately determine the amount of depolarization in the collision process.
## 2 Standard Model Asymmetries in $`𝒆^\mathbf{+}𝒆^{\mathbf{}}`$
For an $`e^+e^{}`$collider with 500 GeV center-of-mass energy, the dependence of SM production cross sections on $`P_e`$is plotted in Figure 2. This Figure is taken from a Report for Snowmass 1996 on the Next Linear Collider<sup>?</sup> and assumes the positron beam is unpolarized. Following the study in that Report, some SM asymmetries are estimated for a detector with an acceptance of $`|cos\theta |<0.99`$, for an integrated luminosity of $`80fb^1`$. These are summarized in Table 2. The left-right asymmetry, $`A_{LR}`$, and its statistical and systematic uncertainties are given by
$`A_{LR}`$ $`=`$ $`{\displaystyle \frac{(\sigma _L\sigma _R)}{(\sigma _L+\sigma _R)}}`$
$`\delta A_{LR}`$ $`=`$ $`\sqrt{{\displaystyle \frac{[1(PA_{LR})^2]}{P^2N}}+\left({\displaystyle \frac{\delta P}{P}}\right)^2A_{LR}^2}`$
$`=`$ $`\sqrt{(\delta A_{LR}^{stat})^2+(\delta A_{LR}^{syst})^2},`$
where $`P=90\%`$ is the electron polarization;<sup>?</sup> $`N_L`$ ($`N_R`$) is the number of observed events with the left-(right-)polarized beam, and $`N=N_L+N_R`$. Equal integrated luminosities are accumulated with the left- and right-polarized electron beam.
Table 2 indicates that it is necessary to have better than $`1\%`$ polarimetry to fully exploit testing SM predictions for these asymmetry measurements. This can be achieved with a precise Compton polarimeter. It is also interesting to consider whether the asymmetry in W-pairs could be used as a polarimeter. The Feynman diagrams contributing to W-pair production are shown in Figure 2. The diagram in Figure 2b is highly suppressed. In the forward detector regions ($`\mathrm{cos}\theta >0.7`$), the exchange diagram in Figure 2c dominates and $`A_{LR}`$for W-pairs is essentially $`100\%`$. The measured left-right asymmetry in forward W-pairs can therefore be used to determine the electron beam polarization.<sup>a</sup><sup>a</sup>aOne will still be able to check for new physics processes that might affect this polarization determination, by measuring the dependence of the result on the polar scattering angle. To achieve sub-1$`\%`$ accuracy for $`P_e`$will require achieving backgrounds to the W-pair sample below $`1\%`$. This has been achieved at LEP200 for the W mass measurements, where one of the Ws is required to decay to $`\mu \nu _\mu `$ and tight cuts are placed on the reconstructed W mass.<sup>?</sup> A detailed study needs to be done to determine whether this low background could also be achieved with the forward detector regions at a linear collider. An advantage of using a detector physics asymmetry for polarimetry rather than Compton polarimetry is that $`P_e^{lum}`$is directly determined and beam-beam depolarization effects are properly accounted for.<sup>b</sup><sup>b</sup>bIf depolarization effects are significant, however, there may be a significant dependence of $`P_e`$on the electron’s effective collision energy. In this case, the energy distribution of collision electrons in the W-pair sample should be the same as the distribution of collision electrons in the physics sample being studied to ensure that $`P_e^{lum}`$is being accurately determined.
A wonderful example of the physics possibilities with precise asymmetry measurements is the linear collider $`Z^0`$-Factory option. It is very desirable to accumulate a large $`Z^0`$ sample ($`>10`$ million) with a polarized electron beam. For example, Ref. ? considers achieving a sample of $`10^9`$ $`Z^0`$ decays with $`80\%`$ electron polarization and $`60\%`$ positron polarization. This enables the determination of the weak mixing angle, by measuring the polarization-dependent cross sections for $`Z^0`$ production, with an unprecedented accuracy of $`\delta \mathrm{sin}^2(\theta _W^{eff})=1.310^5`$. In this case, the availability of a polarized positron beam allows for very precise polarimetry.<sup>?</sup> For the case where the positron beam is not polarized, it will still be very desirable to make a significant improvement on the SLD weak mixing angle measurement, $`\delta \mathrm{sin}^2(\theta _W^{eff})=2.810^4`$.<sup>?</sup> This will require better than $`0.5\%`$ polarimetry from a Compton polarimeter located in the extraction line from the IR. One can also hope to resolve whether the anomalies observed in the $`Zb\overline{b}`$ asymmetry measurements<sup>?</sup> at SLD and LEP are due to statistical fluctuations, systematic problems or new physics.
## 3 Standard Model Asymmetries in $`𝒆^{\mathbf{}}𝒆^{\mathbf{}}`$
For the linear collider operating with $`e^{}e^{}`$collisions, both beams can be polarized. One of the important measurements that will be made is an accurate measurement of the weak mixing angle away from the $`Z^0`$-pole, by measuring the polarization-dependent cross sections in Moller scattering ($`e^{}e^{}e^{}e^{}`$). With both beams polarized, one can measure three independent asymmetries which can be chosen to be
$`A_1`$ $`=`$ $`{\displaystyle \frac{N_{LL}N_{RR}}{N_{LL}+N_{RR}}}`$
$`A_2`$ $`=`$ $`{\displaystyle \frac{N_{RR}N_{LR}}{N_{RR}+N_{LR}}}`$
$`A_3`$ $`=`$ $`{\displaystyle \frac{N_{LR}N_{RL}}{N_{LR}+N_{RL}}}.`$
From these asymmetry measurements, one can determine $`\mathrm{sin}^2(\theta _W^{eff}),P_1`$, and $`P_2`$. A detailed study of this has been done in Ref. ? for a 500 GeV collider, beam polarizations of 90$`\%`$, detector acceptance with $`|\mathrm{cos}\theta |<0.995`$, and integrated luminosity of 25 $`fb^1`$. They find that the beam polarizations can be determined with an accuracy of 0.9$`\%`$, and that the weak mixing angle can be determined with an accuracy of $`\delta \mathrm{sin}^2(\theta _W^{eff})=0.00026`$. This accuracy is comparable to that achieved with SLD’s $`A_{LR}`$measurement at the $`Z^0`$-pole and will be the best measurement of the weak mixing angle away from the $`Z^0`$-pole. The running of $`\mathrm{sin}^2(\theta _W^{eff})`$<sup>?</sup> with $`Q^2`$ will be measured with excellent precision, a factor 3 better than that expected from the SLAC E158 experiment.<sup>?</sup> Excellent sensitivity to additional $`Z^0`$ bosons (up to $`m_Z^{}10`$ TeV) and to electron compositeness (up to a compositeness scale, $`\mathrm{\Lambda }100`$ TeV) will be achieved. The beam polarization uncertainty is comparable to what one can expect with a Compton polarimeter, and has the advantage that it directly measures $`P_e^{lum}`$. The determination of $`P_e^{lum}`$from the Moller scattering analysis can also be applied to other physics analyses.<sup>c</sup><sup>c</sup>cOne may have to make small corrections for the dependence of $`P_e^{lum}`$on the energy distribution of the collision electrons if depolarization effects are significant.
## 4 Background Suppression of W-pairs
For the $`e^+e^{}`$collider, W-pair background will be an obstacle for observing new physics reactions. Beam polarization will be an important tool for understanding and reducing this background. The production cross section for W-pairs may be written as
$`\sigma (P_1,P_2)={\displaystyle \frac{1}{4}}[(1P_1)(1+P_2)\sigma _{LR}+(1+P_1)(1P_2)\sigma _{RL}],`$
where $`P_1`$ is the electron beam polarization; $`P_2`$ is the positron beam polarization; and $`\sigma _{LR}`$ ($`\sigma _{RL}`$) is the W-pair production cross section for a left-(right-)handed electron colliding with a right-(left-)handed positron. As noted in Ref. ?, $`\sigma _{RL}`$ is highly suppressed, $`(\sigma _{RL}/\sigma _{LR}0.004)`$.
For assessing the utility of achieving high beam polarizations, it is useful to construct a Figure-of-Merit (FOM) defined as the ratio, $`R=\sigma (P_1,P_2)^{max}/\sigma (P_1,P_2)^{min}`$ where the maximum (minimum) W-pair production cross section is achieved with a left-(right-)polarized electron beam and a right-(left-)polarized positron beam. Table 4 summarizes this FOM for some possibilities for the beam polarizations. It is desirable to achieve a high electron beam polarization, since the FOM increases by a factor of 2 when improving $`P_1`$ from $`80\%`$ to $`90\%`$. The utility of polarizing the positron beam is also evident,<sup>?</sup> though this may be difficult to implement in a cost-effective way.
The accuracy of the polarization determination will be important for assessing the suppression of the W-pair backgrounds. To illustrate this, consider a potential experiment where the electron beam polarization is $`90\%`$ and the positron beam is unpolarized. Suppose an analysis for isolating a new physics signal yields 400 candidate events after analysis cuts, but with no cut on the polarization state. If the right-polarized electron state is chosen, suppose 40 events are observed to survive. This would be an excess of 20 events over what would be expected if the entire sample were due to W-pair backgrounds. The measured left-right asymmetry would be
$$A_{LR}^{meas}=0.80\pm 0.024(stat)\pm \frac{\delta P}{P}.$$
The uncertainty on $`A_{LR}^{meas}`$ is summarized in Table 4 for 3 possible values of the accuracy of the polarization determination. To achieve a $`4\sigma `$ signal will require better than $`1\%`$ polarimetry. More precise polarimetry is generally desirable as the beam polarization increases, to assure accurate assessment of the W-pair background.
## 5 Beam-beam Effects on Precise Polarimetry
Beam-beam effects in the collision process can cause significant depolarization due to spin precession and the Sokolov-Ternov spin flip mechanism.<sup>?</sup> For a 500 GeV NLC with a luminosity of $`610^{33}`$ cm<sup>-2</sup>s<sup>-1</sup>, the luminosity-weighted depolarization is estimated to be about $`0.15\%`$. For a 1 TeV NLC with a luminosity of $`1.410^{34}`$ cm<sup>-2</sup>s<sup>-1</sup>, it is estimated to be about $`1.5\%`$.<sup>?</sup> These calculations should be checked experimentally, however. This can be done with a Compton polarimeter in the extraction line from the IR, by comparing polarization measurements with and without collisions.<sup>?</sup>
The extraction line Compton polarimeter measures the total depolarization in the collision process. The luminosity-weighted depolarization is typically one-quarter of this.<sup>?</sup> This is easily understood for spin precession effects, where the depolarization has a quadratic dependence on the precession angle and one assumes that half the precession occurs before the hard collision.<sup>d</sup><sup>d</sup>dThis assumption may not be valid if the beams undergo significant betatron oscillations during the collision. In that case, the luminosity-weighted depolarization may be comparable to the total depolarization. For example, the depolarization due to the large disruption angles of the beams is
$`\mathrm{\Delta }P_e`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\gamma {\displaystyle \frac{g2}{2}}\right)^2\left[\sigma _x^{}^2+\sigma _y^{}^2\right]`$
$`\mathrm{\Delta }P_e^{lum}`$ $``$ $`{\displaystyle \frac{1}{2}}\left(\gamma {\displaystyle \frac{g2}{2}}\right)^2\left[\left({\displaystyle \frac{\sigma _x^{}}{2}}\right)^2+\left({\displaystyle \frac{\sigma _y^{}}{2}}\right)^2\right]`$
$`\mathrm{\Delta }P_e^{lum}`$ $``$ $`{\displaystyle \frac{1}{4}}\mathrm{\Delta }P_e,`$
where $`\sigma _x^{}`$ ($`\sigma _y^{}`$) is the disrupted x (y) angular divergence; $`\gamma (\frac{g2}{2})`$ is the spin precession factor; $`\mathrm{\Delta }P_e`$ is the total depolarization and $`\mathrm{\Delta }P_e^{lum}`$ is the luminosity-weighted depolarization.
It is important for IR physicists to include depolarization in their tabulations of beam-beam effects. Extensive tables of beam-beam effects for energy distributions, luminosity distributions and outgoing angular distributions are produced in the design reports for NLC, JLC and TESLA. Surprisingly, depolarization effects are not included. This is presumably due to the lack of depolarization calculations in the beam-beam effect simulation codes used, CAIN<sup>?</sup> and GUINEA-PIG.<sup>?</sup> These programs should be improved to include depolarization effects, and depolarization should be included as a key element in tables summarizing beam-beam effects.
At high luminosities, or for the more severe beam disruption experienced in $`e^{}e^{}`$collisions, it becomes difficult to transport the disrupted beam cleanly to the beam dumps.<sup>?</sup> This may lead to an extraction line design that precludes a Compton polarimeter and other beam diagnostics. In this scenario, one will have to rely on a Compton polarimeter before the IP, and possibly utilize SM physics asymmetries for polarimetry as well. By comparing the upstream Compton polarization measurement with that provided by a SM physics asymmetry, the beam-beam depolarization can be determined. This depolarization determination will have greater systematic uncertainty than that achievable with an extraction line Compton polarimeter, for which many sources of systematic error will cancel when determining the amount of depolarization. However, $`P_e^{lum}`$and $`\mathrm{\Delta }P_e^{lum}`$ will be more directly determined, perhaps offsetting this disadvantage. Ideally, both extraction line Compton polarimetry and SM physics asymmetry polarimetry will be achievable.
## 6 Conclusions
Precise measurements of SM asymmetries in both $`e^+e^{}`$and $`e^{}e^{}`$collider modes require better than $`1\%`$ polarimetry. Sub-$`1\%`$ polarimetry may also be required to accurately assess W-pair backgrounds in a discovery search for a new physics signal at an $`e^+e^{}`$collider.
A Compton polarimeter in the extraction line from the IP is desirable, especially for its ability to accurately measure depolarization effects. SM physics asymmetries are useful for polarimetry when both colliding beams are polarized. For the $`e^+e^{}`$collider with only the electron beam polarized, the asymmetry in forward W-pairs may also prove useful as a polarimeter. For the $`Z^0`$-factory $`e^+e^{}`$collider (below W-pair threshold) with no positron polarization, a very precise Compton polarimeter in the extraction line is required.
There is a need to include depolarization in the tables summarizing beam-beam effects. Including depolarization calculations in the simulation programs for beam-beam effects will assist this.
Acknowledgements
I would like to thank Michael Peskin for discussions and comments regarding this paper. I would also like to thank Clemens Heusch for organizing this meeting and for promoting the less conventional aspects of linear colliders.
References
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# 1 Introduction
## 1 Introduction
In recent years, there has been renewed interest in the understanding of the physics of systems with large number of partons. These studies have been essentially motivated by two large experimental programmes — low $`x`$ deep inelastic scattering (DIS) at HERA and heavy ion collision experiments at RHIC and LHC. Both physical situations involve a large number of participating gluons. In low $`x`$ DIS these gluons are generated in the proton light cone wave function by the evolution to low $`x`$, whereas in the nuclear collision this evolution is enhanced since the nuclear wave function contains many gluons already at moderate values of energy.
The growth of gluon density leads to interesting physical consequences, the physical understanding of which has been steadily improving. One universal feature now believed to be true, is the saturation of gluon densities. Apparently, the number of gluons per unit phase space volume practically saturates and at large densities grows only very slowly (logarithmically) as a function of the parameter that triggers the growth. The relevant parameter could be $`1/x`$ in the low $`x`$ regime or the atomic number of the nucleus $`A`$ in heavy ion collisions. This saturation takes place at values of transverse momentum below a certain saturation momentum $`k_s`$, which itself depends on $`1/x`$ and $`A`$. The nature of this dependence is less well understood. In the analysis based on the Balitsky-Fadin-Kuraev-Lipatov (BFKL) evolution and on the double logarithmic approximation (DLA) the dependence is powerlike $`k_s(1/x)^{\alpha _s\delta }`$, while other approaches suggest a much slower dependence. In the case of power dependence, the saturation momentum at HERA is estimated to be in the range of 1-2 GeV with similar, slightly higher, values at LHC. Optimistically, one can hope that the saturation region is itself semiperturbative, that is the value of the coupling constant is reasonably small and, therefore, weak coupling methods can be applied to the quantitative analysis of the phenomenon.
The physics of saturation must have experimental manifestations. The simplest and the most direct, in a way, is the unitarization of the total DIS cross section. This is, however, also the least interesting one. First, since the effect of unitarization is almost kinematical, one does not need high partonic density, it is enough to have a large number of partons, not necessarily in the same bin of the phase volume . Second, because the experimental status of unitarization is unclear. So far, all DIS data on the total cross section can be reasonably well described by linear Dokshitzer-Gribov-Lipatov-Altarelli-Parisi (DGLAP) evolution without the need to include nonlinear effects . Although physically it is hard to believe that the leading twist perturbative approximation can be applied at $`Q^2`$ as low as 1 GeV<sup>2</sup> and although some aspects of the gluon distribution that emerge form these fits are intuitively not satisfactory, present inclusive DIS data cannot be considered as an unambiguous confirmation of nonlinear effects.
The realm of nonlinear effects is, however, much richer than the total cross section. In particular, one expects qualitative changes in the structure of the final states as one moves into the saturation region. The study of these effects has, however, not started in earnest yet and we have a long way to go before being able to make verifiable quantitative predictions.
In particular, one needs a well defined formal framework to perform calculations. Several approaches to the problem have been developed in recent years by different groups. They all rely on the smallness of the coupling constant while resumming the effects of large number of partons/partonic density. The aim of all these approaches is essentially to derive the evolution of the hadronic scattering cross section with $`1/x`$. They, however, utilize different techniques and conceptual frameworks and the resulting evolution equations look rather different. It is the aim of this paper to explore the relation between some of these different approaches in an attempt to understand where they diverge from each other in terms of physics input.
In particular, we will concern ourselves with three recent works Ref. , Ref. and Refs. . In Ref. the evolution equation for the scattering amplitude is derived using the effective action and the eikonal approximation in the target rest frame. Ref. uses the dipole model method of Refs. . And, finally, Refs. uses the effective action in the projectile rest frame to derive the evolution of the hadron light cone wave function with $`1/x`$. We will refer to the resulting evolution equation as the Jalilian-Marian–Kovner–Leonidov–Weigert (JKLW) equation.
The outline of this paper is the following. In Sec. 2 we rederive the evolution equation of Ref. in a simple and intuitive way. This derivation makes it obvious that this approach is equivalent to the approach of Ref. up to subleading corrections in $`1/N_c`$. This is not new and was noted already in Ref. . In the following, we will refer to this evolution equation as the Balitsky-Kovchegov (BK) equation<sup>1</sup><sup>1</sup>1The $`O(1/N_c)`$ differences between the equations derived in Ref. and do not carry essential new physics. They therefore do not affect our understanding of the relationship between the generic frameworks of the BK and JKLW equations.. We discuss the physical picture of this evolution and resulting unitarization of the total cross section in both target and projectile rest frames and point out the effects due to which the approximations involved should break down at extremely small $`x`$. The breakdown of the approximation should have very little effect on the unitarization of the total cross section, since, especially for large targets like nuclei, the black disk limit should be reached while the approximation is still valid. However, one does expect the structure of the final states to be strongly affected. Our discussion here is, in large measure, parallel to that of Refs. .
In Sec. 3 we relate explicitly the calculation of Ref. to that of Refs. . In particular, we calculate the basic physical quantities appearing in the evolution equation of Refs. in the approximation of Ref. . We show that the results of Ref. are recovered from Refs. in the limit of small induced fields. We also show that the double logarithmic limit of the evolution of Ref. is trivial. That is, in the double logarithmic limit, the evolution equation for the gluon distribution function (defined operatorially as the number of gluons in the lightcone gauge in the infinite momentum frame) becomes linear and does not contain any Gribov-Levin-Ryskin (GLR) type corrections<sup>2</sup><sup>2</sup>2This is not to say that the evolution of the DIS cross section for which the equation of Ref. and has been derived is linear in the double logarithmic limit. The GLR type nonlinearity does indeed appear in the evolution equation for the virtual photon cross section due to the nonlinear relation between the cross section and the gluon distribution function.. This is in contrast with the result of Ref. where it was shown that the double logarithmic limit of the evolution Refs. results in a nonlinear equation. We point out that this is indeed a very natural result from the point of view of the dipole model approach.
In Sec. 4 we transform the full calculation of Refs. into the framework of Ref. . We show that in the approach of Ref. it corresponds to abandoning the eikonal approximation or, equivalently, to the inability to fully describe the target by a classical $`A^+`$ field. We also point out technical reasons which lead us to believe that, in fact, in the framework of the effective action of Ref. such a failure is expected when the evolution is continued to very low values of $`x`$.
Finally, we conclude with a brief discussion in Sec. 5.
## 2 A simple derivation of the BK equation
In this section we will give a simple derivation of the evolution equation first derived in Ref. and discuss the physical picture behind it. Consider the deep inelastic scattering at low $`x`$. We will work in the frame in which the photon fluctuates into an energetic quark-antiquark pair long before it reaches the target, but where most of the energy resides in the target hadron which moves very fast. The scattering of the quark-antiquark pair is dominated by its interaction with the gluons in the target. Since the target hadron moves fast, the time evolution of the gluon fields is slowed by Lorentz time dilation. Also, due to Lorentz contraction, the gluon fields are well localized in the plane perpendicular to the direction of motion, which we take to be the positive $`x_3`$ axis. The target can, therefore, be modeled by a distribution of static gluon fields localized at $`x^{}=0`$. As the scattering energy increases ($`x`$ decreases) the gluon fields of the target change due to contributions of quantum fluctuations. It is this evolution in $`x`$ of the hadronic ensemble that we intend to describe in terms of the evolution equation.
### 2.1 The BK equation.
In this section we will use the lightcone gauge $`A^{}=0`$. In this gauge, following Ref. , we take the vector potentials representing the relevant gluon field configurations to be of the form
$$b^i=0,b^+=b(x_{})\delta (x^{}).$$
(1)
Here and in the rest of this section, unless otherwise specified, we use the matrix notation for the gauge field $`b^+=b_a^+t_a`$ etc., where $`t_a`$ are the generators of the $`SU(N)`$ group in the fundamental representation. One can reasonably ask whether the vector potential of this form is the only relevant one. This turns out to be a nontrivial question. In fact, we will argue later in Sec. 4 that this is not quite the case if we want to be able to describe the evolution up to arbitrarily small values of $`x`$. At this point, however, we follow Ref. . We will return to this question in Sec. 4.
The DIS structure function can be written in the following general form
$$F_2(x,Q^2)=\frac{Q^2}{4\pi ^2\alpha _{em}}\frac{dzdx_{}dy_{}}{4\pi }\mathrm{\Phi }(x_{}y_{},z)N(x_{},y_{},y).$$
(2)
Here, $`x_{}`$ and $`y_{}`$ are the transverse coordinates of the quark and the antiquark in the pair, $`z`$ is the fraction of the pairs longitudinal momentum carried by the quark and $`y`$ is the rapidity of the slowest particle in the pair. Also, $`\mathrm{\Phi }(x_{}y_{},z)`$ is the square of the “wave function” of the photon — the probability that the virtual photon fluctuates into the pair with given coordinates and momenta — and $`N(x_{},y_{},y)`$ is the cross section for the scattering of the pair.
The wave function $`\mathrm{\Phi }`$ is well known. It is given for example in Ref. , but its explicit form will not be of interest to us. We concentrate our discussion on the scattering cross section $`N`$. If the quark-antiquark pair is energetic enough the scattering cross section is eikonal
$$N(x_{},y_{})=\mathrm{tr}V(x_{})V^{}(y_{})1_A,$$
(3)
where $`V`$ ($`V^{}`$) is the eikonal phase for the scattering of the energetic quark (antiquark)
$$V(x^+=0,x_{})=𝒫\mathrm{exp}\left[ig\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑x^{}A^+(x^+=0,x_{},x^{})\right]$$
(4)
with the vector potential in the fundamental representation. We, therefore, have to calculate the average of $`V(x^+=0,x_{})V^{}(y^+=0,y_{})_A`$ over the hadronic wave function as indicated by $`\mathrm{}_A`$. In our frame, the quark and the antiquark move with velocity of light in the negative $`x_3`$ direction. All the fields in Eq. (4), therefore, have vanishing $`x^+`$ coordinate. This will also be the case for all the fields in the rest of this section. For simplicity, we suppress the $`x^+`$ coordinate in the following.
In the leading approximation the vector potential is given by Eq. (1) and the scattering amplitude is
$$N(x_{},y_{})=\mathrm{tr}[U(x_{})U^{}(y_{})1]_b$$
(5)
with
$$U(x_{})=𝒫\mathrm{exp}\left[i\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑x^{}b^+(x_{},x^{})\right],$$
(6)
To calculate the order $`\alpha _s`$ correction to this expression we write the vector potential as
$$A^+=\frac{1}{g}b^++a^+$$
(7)
with $`a^+`$ being a small fluctuation and expand the eikonal factors to second order in $`a^+`$.
Recalling that the classical background vector potential is a delta function in $`x^{}`$, we have
$`V(x_{})=`$ $`𝒫\mathrm{exp}\left[ig{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}𝑑x^{}a^+(x_{},x^{})\right]U(x_{})𝒫\mathrm{exp}\left[ig{\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}𝑑x^{}a^+(x_{},x^{})\right]`$ (8a)
$`=`$ $`U(x_{})ig\left\{{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}𝑑x^{}a^+(x_{},x^{})U(x_{})+U(x_{}){\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}𝑑x^{}a^+(x_{},x^{})\right\}`$
$`g^2\{{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}dx^{}dy^{}\theta (y^{}x^{})a^+(x_{},x^{})a^+(x_{},y^{})U(x_{})`$
$`+{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}𝑑x^{}a^+(x_{},x^{})U(x_{}){\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}𝑑x^{}a^+(x_{},x^{})`$
$`+U(x_{}){\displaystyle \underset{0}{\overset{+\mathrm{}}{}}}dx^{}dy^{}\theta (y^{}x^{})a^+(x_{},x^{})a^+(x_{},y^{})\}.`$ (8b)
All contributions break down into $`x^{}`$ ordered pieces because of the $`x^{}`$ structure in Eq. (1).
Now, together with the analogous expansion for $`V^{}`$, we insert this into Eq. (3) and obtain
$`\mathrm{tr}V(x_{})V^{}(y_{}`$ $`)_A\mathrm{tr}U(x_{})U^{}(y_{})_b=`$
$`=g^2\mathrm{tr}`$ $`{\displaystyle _{\mathrm{}}^0}𝑑w^{}a^+(x_{},w^{})U(x_{})U^{}(y_{}){\displaystyle _{\mathrm{}}^0}𝑑z^{}a^+(y_{},z^{})`$
$`+{\displaystyle _{\mathrm{}}^0}𝑑w^{}a^+(x_{},w^{})U(x_{}){\displaystyle _0^+\mathrm{}}𝑑z^{}a^+(y_{},z^{})U^{}(y_{})`$
$`+U(x_{}){\displaystyle _0^+\mathrm{}}𝑑w^{}a^+(x_{},w^{})U^{}(y_{}){\displaystyle _{\mathrm{}}^0}𝑑z^{}a^+(y_{},z^{})`$
$`+U(x_{}){\displaystyle _0^+\mathrm{}}𝑑w^{}a^+(x_{},w^{}){\displaystyle _0^+\mathrm{}}𝑑z^{}a^+(y_{},z^{})U^{}(y_{})`$
$`{\displaystyle _{\mathrm{}}^0}𝑑w^{}𝑑z^{}\theta (z^{}w^{})a^+(x_{},w^{})a^+(x_{},z^{})U(x_{})U^{}(y_{})`$
$`{\displaystyle _{\mathrm{}}^0}𝑑w^{}a^+(x_{},w^{})U(x_{}){\displaystyle _0^+\mathrm{}}𝑑z^{}a^+(x_{},z^{})U^{}(y_{})`$
$`U(x_{}){\displaystyle _0^+\mathrm{}}𝑑w^{}𝑑z^{}\theta (z^{}w^{})a^+(x_{},w^{})a^+(x_{},z^{})U^{}(y_{})`$
$`U(x_{})U^{}(y_{}){\displaystyle _{\mathrm{}}^0}𝑑w^{}𝑑z^{}\theta (w^{}z^{})a^+(y_{},w^{})a^+(y_{},z^{})`$
$`U(x_{}){\displaystyle _0^+\mathrm{}}𝑑w^{}a^+(y_{},w^{})U^{}(y_{}){\displaystyle _{\mathrm{}}^0}𝑑z^{}a^+(y_{},z^{})`$
$`U(x_{}){\displaystyle _0^+\mathrm{}}dw^{}dz^{}\theta (w^{}z^{})a^+(y_{},w^{})a^+(y_{},z^{})U^{}(y_{})_{b,a}.`$ (9)
In writing Eq. (9) we have anticipated that $`a_u^+_a=0`$ as in the free case. This can easily be shown using the explicit expression for the fluctuation propagator given below.
Although this expression is a little cumbersome, the physical meaning of the various terms is very clear. Put more compactly the structure of the above is determined by
$$\begin{array}{cc}\hfill \mathrm{tr}V(x_{})V^{}(y_{}& )_A\mathrm{tr}U(x_{})U^{}(y_{})_b=g^2a_u^+a_v^+_a\frac{1}{2}(2\frac{\delta }{\delta b_u^+}U(x_{})\frac{\delta }{\delta b_v^+}U^{}(y_{})\hfill \\ & +\left(\frac{\delta }{\delta b_u^+}\frac{\delta }{\delta b_v^+}U(x_{})\right)U^{}(y_{})+U(x_{})\left(\frac{\delta }{\delta b_u^+}\frac{\delta }{\delta b_v^+}U^{}(y_{})\right))_b.\hfill \end{array}$$
(10)
Diagrammatically the r.h.s. can be represented as follows
$$2\times \text{}+\text{}+\text{}_b$$
(11)
The straight lines represent the eikonal factors $`U`$, while the curly lines denote the gluon fluctuation propagator $`a_u^+a_v^+_a`$, evaluated in the fixed background $`b^+(x)`$. The terms in Eq. (10) with first order derivatives correspond to processes where the gluon is emitted by the quark and absorbed by the antiquark (or vice versa). Those will be hereafter referred to as “exchange contributions”.
The terms in Eq. (10) with second derivatives acting on $`U`$ (or $`U^{}`$) correspond to the diagramms where where the quark (or antiquark) emits the gluon and then reabsorbs it at a later time; typical self energy corrections. To contrast them against the exchange contributions we will also refer to them as “non-exchange contributions”.
In Eq.(9) we are looking at a $`x^{}`$ ordered breakdown of the diagrams in Eq. (11) with the first four terms summing up to the exchange contribution and the rest to the non-exchange contributions.
With the vertices known, we only lack an explicit expression for $`a_u^+a_v^+_a`$. The QCD action expanded to second order in the fluctuation field $`a^+`$, in the presence of the classical background $`b^+`$, in our lightcone gauge is
$$\begin{array}{cc}\hfill S=\frac{1}{2}\{& a_a^+\left[(^{})^2\right]a_a^+2(^ia_a^+)(^{}a_a^i)\hfill \\ & a_a^i[(2D_{ab}^+[b]^{}(_{})^2\delta _{ab})\delta ^{ij}+^i^j\delta _{ab}]a_b^j\}.\hfill \end{array}$$
(12)
Recall that we are interested only in the propagator of the fields at equal $`x^+`$. Consequently, it is only the on-shell part of the propagator that is relevant for our purposes. We can, therefore, use the classical equation of motion for $`a^+`$
$$a^+=\frac{^i}{^{}}a^i.$$
(13)
Substituting this in Eq. (12) we get
$$S=\frac{1}{2}a_a^i(D^2)_{ab}\delta ^{ij}a_b^j,$$
(14)
where
$$(D^2)_{ab}=2D_{ab}^+[b]^{}(_{})^2\delta _{ab}.$$
(15)
The propagator of the spatial components of the vector potential is, therefore, simply given by $`i/D^2`$. The explicit form is very simple and can be found, for example, in Ref. .
$$\begin{array}{cc}\hfill i\left[\frac{1}{(D^2)}\right]_{ab}=\frac{dp^{}}{2p^{}(2\pi )^3}& [\theta (x^{}y^{})\theta (p^{})\theta (y^{}x^{})\theta (p^{})]\times \hfill \\ \hfill \times d^2& p_{}d^2q_{}e^{ipx+iqy}\frac{d^2z_{}}{(2\pi )^2}e^{i(p_{}q_{})z_{}}\stackrel{~}{U}_{ab}^1(x^{},y^{},z_{})\hfill \end{array}$$
(16)
with $`p^+=\frac{p_{}^2}{2p^{}}`$, $`q^+=\frac{q_{}^2}{2p^{}}`$ and $`q^{}=p^{}`$. The adjoint color matrix<sup>3</sup><sup>3</sup>3More rigorously, the structure of $`\stackrel{~}{U}^1`$ is given by $`\stackrel{~}{U}_{ab}^1(z_{})=e^{i(\theta (x^{})\theta (y^{}))b(z_{})}`$. However, the difference between this expression and that given in Eq.(17) only shows up if it is multiplied by $`^+`$ derivatives or $`\delta (x^{})`$ factors. Since we encounter no such factors in our calculation we will be using Eq.(17) throughout.
$$\begin{array}{cc}\hfill \stackrel{~}{U}_{ab}^1(x^{},y^{},z_{})=(\theta (x^{})\theta (& y^{})+\theta (x^{})\theta (y^{}))\delta _{ab}\hfill \\ & +\theta (x^{})\theta (y^{})\stackrel{~}{U}_{ab}(z_{})+\theta (x^{})\theta (y^{})\stackrel{~}{U}_{ab}^{}(z_{}).\hfill \end{array}$$
(17)
represents a phase factor one picks up when crossing the $`x^{}=0`$ plane due to interaction with a field of type Eq. (1). Here, $`\stackrel{~}{U}_{ab}(z_{})`$ is the adjoint version of the fundamental $`U`$ in Eq. (6). If the $`x^{}=0`$ plane is not crossed the propagation remains free.
We can now write the on-shell correlator of the “$`+`$” component of the vector potential as
$$\begin{array}{c}a_a^+(x^+=0,x_{},x^{})a_b^+(y^+=0,y_{},y^{})_a\hfill \\ \hfill =\frac{_x^i}{_x^{}}a_a^i(x^+=0,x_{},x^{})a_b^i(y^+=0,y_{},y^{})\frac{_y^i}{_y^{}}\\ \hfill =_x^i_y^i\frac{dp^{}}{(p^{})^3}\frac{1}{4\pi }\left[\theta (x^{}y^{})\theta (p^{})\theta (y^{}x^{})\theta (p^{})\right]\\ \hfill \times d^2z_{}\frac{d^2p_{}}{(2\pi )^2}\frac{d^2q_{}}{(2\pi )^2}e^{ip_{}(x_{}z_{})+iq_{}(z_{}y_{})}e^{i\frac{p_{}^{}{}_{}{}^{2}}{2p^{}}x^{}+i\frac{q_{}^{}{}_{}{}^{2}}{2p^{}}y^{}}\\ \hfill \times \{\theta (x^{})\theta (y^{})+\theta (x^{})\theta (y^{})\\ \hfill +\theta (x^{})\theta (y^{})\stackrel{~}{U}(z_{})+\theta (x^{})\theta (y^{})\stackrel{~}{U}^{}(z_{})\}.\end{array}$$
(18)
This expression displays a separation into $`x^{}`$ ordered contributions that seamlessly matches up with what we have already seen for the vertices in Eq. (9). Diagrammatically, the $`x^{}`$ ordered exchange contributions are given by
$`=\text{}+\text{}+\text{}+\text{},`$ (19)
while the (connected parts of the) non-exchange ones are represented by
$`=\text{}+\text{}+\text{},`$ (20a)
$`=\text{}+\text{}+\text{}.`$ (20b)
Whenever the $`x^{}=0`$ plane cuts the fluctuation propagator, the diagram contains a factor $`U_{ab}^{()}`$, otherwise the fluctuation propagator is free.
Algebraically, the corrections to the scattering cross section involve integrals of this propagator with respect to $`x^{}`$ and $`y^{}`$ from zero to either $`+`$ or $``$ infinity. This is straightforward to do. For example, for the second term in Eq. (9) we need
$$\begin{array}{cc}\hfill \underset{\mathrm{}}{\overset{0}{}}dw^{}\underset{0}{\overset{+\mathrm{}}{}}dz^{}a_a^+& (x_{},w^{})a_b^+(y_{},z^{})=\hfill \\ & =_x^i_y^i\underset{0}{\overset{+\mathrm{}}{}}\frac{dp^{}}{p^{}}\frac{1}{\pi }d^2z_{}\frac{d^2p_{}}{(2\pi )^2}\frac{d^2q_{}}{(2\pi )^2}\frac{e^{ip_{}(x_{}z_{})+iq_{}(z_{}y_{})}}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}\stackrel{~}{U}_{ab}(z_{})\hfill \\ & =\frac{1}{\pi }\underset{0}{\overset{+\mathrm{}}{}}\frac{dp^{}}{p^{}}x_{}|\frac{^i}{_{}^2}\stackrel{~}{U}_{ab}\frac{^i}{_{}^2}|y_{}.\hfill \end{array}$$
(21)
Here, $`x_{}|O|y_{}`$ means the matrix element of the operator $`O`$ in the coordinate basis in the usual sense. We treat $`\stackrel{~}{U}`$ as an operator in the coordinate space with matrix elements $`x_{}|\stackrel{~}{U}|y_{}=\stackrel{~}{U}(x_{})\delta (x_{}y_{})`$ and the products in the last line are understood in the operatorial sense. Explicitly,
$$\underset{\mathrm{}}{\overset{0}{}}𝑑w^{}\underset{0}{\overset{+\mathrm{}}{}}𝑑z^{}a_a^+(x_{},w^{})a_b^+(y_{},z^{})=\frac{1}{4\pi ^3}\underset{0}{\overset{+\mathrm{}}{}}\frac{dp^{}}{p^{}}d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}\stackrel{~}{U}_{ab}(z_{}).$$
(22)
In the same way we obtain for the other contributions with interaction with the background
$$\begin{array}{c}\underset{0}{\overset{+\mathrm{}}{}}𝑑w^{}\underset{\mathrm{}}{\overset{0}{}}𝑑z^{}a_a^+(x_{},w^{})a_b^+(y_{},z^{})=\hfill \\ \hfill =\frac{1}{\pi }_0^+\mathrm{}\frac{dp^{}}{p^{}}x_{}|\frac{^i}{_{}^2}\stackrel{~}{U}_{}^{}{}_{ab}{}^{}\frac{^i}{_{}^2}|y_{}=\frac{1}{4\pi ^3}_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}\stackrel{~}{U}_{}^{}{}_{ab}{}^{}(z_{})\end{array}$$
(23)
$$\begin{array}{c}\underset{\mathrm{}}{\overset{0}{}}𝑑w^{}\underset{0}{\overset{+\mathrm{}}{}}𝑑z^{}a_a^+(x_{},w^{})a_b^+(x_{},z^{})=\hfill \\ \hfill =\frac{1}{\pi }_0^+\mathrm{}\frac{dp^{}}{p^{}}x_{}|\frac{^i}{_{}^2}\stackrel{~}{U}_{ab}\frac{^i}{_{}^2}|x_{}=\frac{1}{4\pi ^3}_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{1}{(xz)_{}^2}\stackrel{~}{U}_{ab}(z_{})\end{array}$$
(24)
$$\begin{array}{c}\underset{0}{\overset{+\mathrm{}}{}}𝑑w^{}\underset{\mathrm{}}{\overset{0}{}}𝑑z^{}a_a^+(y_{},w^{})a_b^+(y_{},z^{})=\hfill \\ \hfill =\frac{1}{\pi }_0^+\mathrm{}\frac{dp^{}}{p^{}}y_{}|\frac{^i}{_{}^2}\stackrel{~}{U}_{}^{}{}_{ab}{}^{}\frac{^i}{_{}^2}|y_{}=\frac{1}{4\pi ^3}_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{1}{(yz)_{}^2}\stackrel{~}{U}_{}^{}{}_{ab}{}^{}(z_{})\end{array}$$
(25)
To simplify the color structure of these expressions we use the identity
$`\stackrel{~}{U}_{ab}(z_{})(t_aU(x_{}))^{\alpha \beta }(t_bU^{}(y_{}))^{\gamma \delta }=2\mathrm{t}\mathrm{r}\left[t_aU(z_{})t_bU^{}(z_{})\right](t_aU(x_{}))^{\alpha \beta }(t_bU^{}(y_{}))^{\gamma \delta }`$
$`={\displaystyle \frac{1}{2N_c}}\left[N_c\left(U(z_{})U^{}(y_{})\right)^{\alpha \delta }\left(U^{}(z_{})U(x_{})\right)^{\gamma \beta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }\right].`$ (26)
Note that the integral over the frequency $`p^{}`$ logarithmically diverges. In fact, we have to integrate only over a finite interval of frequencies. The gluon field modes of very low frequency have been already included in the background field $`b^+`$ and, therefore, the fluctuation fields at these low frequencies should not be considered. The lower cutoff on the frequency of the modes that are being integrated is inversely proportional to the initial value of $`x_0`$ at which we start the evolution. The upper limit on the high frequency side is furnished by the maximal rapidity of the quark (or antiquark) in the virtual photon which is of the order of $`1/x`$. The ratio of these two cutoffs is of order $`x_0/x`$. Thus, in the leading logarithmic approximation we identify
$$\frac{dp^{}}{p^{}}=\mathrm{ln}\frac{x_0}{x}.$$
(27)
The calculation of the remaining contributions (the ones with no interaction with the background) proceeds along similar lines and is given in the appendix.
Collecting all contributions together we obtain
$`\mathrm{tr}V(x_{})V^{}(y_{})_A\mathrm{tr}U(x_{})U^{}(y_{})_b={\displaystyle \frac{g^2}{8\pi ^3}}\mathrm{ln}({\displaystyle \frac{x_0}{x}}){\displaystyle d^2z_{}}`$
$`\times [2\mathrm{t}\mathrm{r}\left(U(x_{})U^{}(z_{})\right)\mathrm{tr}\left(U^{}(y_{})U(z_{})\right)2N_c\mathrm{tr}\left(U(x_{})U^{}(y_{})\right)]{\displaystyle \frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}}`$
$`\left[\mathrm{tr}\left(U(x_{})U^{}(z_{})\right)\mathrm{tr}\left(U(z_{})U^{}(y_{})\right)N_c\mathrm{tr}(U(x_{})U^{}(y_{}))\right]{\displaystyle \frac{1}{(xz)_{}^2}}`$
$`[\mathrm{tr}\left(U^{}(y_{})U(z_{})\right)\mathrm{tr}\left(U^{}(z_{})U(x_{})\right)N_c\mathrm{tr}(U^{}(y_{})U(x_{}))]{\displaystyle \frac{1}{(yz)_{}^2}}_b.`$ (28)
The eikonal factors $`V`$ themselves should be considered as functions of $`x`$, so that $`U=V(x_0)`$. Differentiating this equation with respect to $`\mathrm{ln}1/x`$ we recover the evolution step for $`\mathrm{tr}V(x_{})V^{}(y_{})_A`$. This is precisely what was found in Ref. .<sup>4</sup><sup>4</sup>4As it stands, Eq. (2.1) does not provide a closed equation — it has to be supplemented by evolution equations for arbitrary products $`V_1^{()}\mathrm{}V_n^{()}_A`$. The evolution of these higher correlators are derived following the same procedure as described above and leads to the full set of operator equations derived in Ref. . We will give a compact representation of the whole set of the evolution equations in Eq. (74), Sec. 4.
At large $`N_c`$ the products of traces in Eq. (2.1) factorize:
$$\mathrm{tr}(U(x_{})U^{}(z_{}))\mathrm{tr}(U^{}(y_{})U(z_{}))\stackrel{N_c\mathrm{}}{}\mathrm{tr}(U(x_{})U^{}(z_{}))\mathrm{tr}(U^{}(y_{})U(z_{}).$$
(29)
Equation (2.1) then becomes a closed equation for the evolution of $`N(x_{},y_{})=\mathrm{tr}(U(x)U^{}(y)1)`$. It is identical to the nonlinear evolution equation of Ref. .
### 2.2 The physical interpretation.
Now, let us discuss the physical picture of this evolution. As always with DIS, the physical picture depends on the frame in which one chooses to view the process. We have specified the frame to some extent by declaring that the photon fluctuates into a $`q\overline{q}`$ pair long before the target. However, we are still free to put the subsequent evolution in $`x`$ either into the evolution of the photon wave function or into the evolution of the gluon field distribution in the target. We will refer to the former picture as the “projectile evolution picture” and to the latter as the “target evolution picture”. The calculation is, fortunately, noncommittal on this point and we will consider both pictures in turn.
In the projectile evolution picture, the higher energy of the scattering is achieved by boosting the $`q\overline{q}`$ pair. In this picture, the quark and antiquark have very high energy and consequently their wave function develops extra gluon components. The growth of the cross section with $`1/x`$ then is interpreted as due to the scattering of extra gluons in the projectiles wavefunction. This is precisely how the low $`x`$ evolution is viewed in the dipole model of Mueller . Our calculation of this section has a simple interpretation from this point of view. The quark and the antiquark is the pair of pointlike color charges moving with velocity of light and located at $`x^+=0`$ and transverse coordinates $`x_{}`$ and $`y_{}`$. These color charges carry with them a fluctuating gluon field. When the pair is boosted to higher rapidity the gluon fields “freeze” due to time the dilation and become static. In the approximation when the gluon fields are frozen, they are given by the Weiszäcker-Williams static ($`p^+=0`$) fields created by the $`q\overline{q}`$ pair. In the leading order in $`\alpha _s`$ the Weiszäcker-Williams (WW) fields are small and are emitted independently by the quark and the antiquark. The total WW field is
$$A^i=g\frac{1}{^{}}\frac{^i}{_{}^2}[j_q^{}+j_{\overline{q}}^{}],$$
(30)
where $`j_q^{}`$ ($`j_{\overline{q}}^{}`$) is the color current due to the quark (antiquark) which in our frame has only a “$``$” component. For pointlike quark and antiquark the charge densities are delta functions in the transverse coordinates and in $`x^+`$. The WW field is therefore
$$A^i(p^{},z_{})=g\frac{1}{p^{}}[\tau \frac{x^iz^i}{(xz)^2}+\tau ^{}\frac{y^iz^i}{(yz)^2}].$$
(31)
Here, $`\tau `$ and $`\tau ^{}`$ are fundamental color matrices corresponding to the orientation of the quark and antiquark wavefunctions in the color space. Their exact form does not matter for our purposes. The WW field, if written in the particle basis, can be thought of as representing equivalent gluons. The number of gluons at a given transverse position is given by the familiar expression
$`n_{WW}(z_{})`$ $`{\displaystyle _0^{\mathrm{}}}𝑑p^{}p^{}\mathrm{tr}F^i(p^{},z_{})F^i(p^{},z_{})`$
$`=\alpha _s{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dp^{}}{p^{}}}\mathrm{tr}[\tau {\displaystyle \frac{x^iz^i}{(xz)^2}}+\tau ^{}{\displaystyle \frac{y^iz^i}{(yz)^2}}][\tau {\displaystyle \frac{x^iz^i}{(xz)^2}}+\tau ^{}{\displaystyle \frac{y^iz^i}{(yz)^2}}].`$ (32)
If we do not take the trace over the color indices, this expression gives the probability to have one extra WW gluon in the wave function of the $`q\overline{q}`$ pair (at the transverse position $`z_{}`$ with a particular color orientation). These WW gluons scatter on the gluon field of the target eikonally just like the quark and the antiquark, apart from the fact that they carry adjoint charge, and so their eikonal amplitude is given by $`\stackrel{~}{U}`$ rather than $`U`$. The terms in this expression are in one-to-one correspondence with the real contributions in Eq. (9), that is the terms in which the gluons interact with the target (background). The rest of the terms in Eq.(9) - the virtual terms - as usual serve to restore the correct normalization of the wave function.
To summarize, in the projectile evolution picture our calculation describes emission of the WW gluons into the wave function of $`q\overline{q}`$ long before the scattering. The transverse coordinates of these gluons are frozen due to the Lorentz time dilation. Subsequently, both the $`q`$ and $`\overline{q}`$, and also the gluons, scatter eikonally and independently of each other on the target gluon field. Clearly, this picture is identical to the dipole evolution picture of Mueller which was used in Ref. to derive a nonlinear evolution equation. The only difference is that the dipole model uses the simplifications in the color algebra which arise in the large $`N_c`$ limit.
The calculation presented above also has a simple interpretation in the target evolution picture. In this picture, it is the target rather than the projectile that is boosted when going to lower $`x`$. As already made explicit above by writing $`\mathrm{}_b`$, one should think about the target as being represented by an ensemble of the configurations of $`b^+`$. The corresponding statistical weight $`Z[b]`$ is determined, of course, by the structure of the target at the relevant resolution scale. We will have more to say about it in Sec. 4. The boost of the target freezes the gluon field fluctuations around the target background $`b^+`$ and, consequently, some field modes which were not important at higher $`x`$ are now capable of inducing scattering. Thus, the ensemble of the relevant field configurations which characterizes the target changes. In fact, every $`b^+`$ now forks into a “subensemble” $`b^+^{}=b^++a^+`$. In the weak coupling regime $`a^+`$ have Gaussian distribution with the width determined by the inverse of their correlation function Eq. (18). One can work back from here and calculate the modification of the distribution of the background fields. We will do this in the following sections.
The fluctuations of $`a^+`$ are, therefore, considered in the target evolution picture as modifying the ensemble of the target background fields very much like in the approach of Refs. .
### 2.3 Unitarization in different approximations.
From what has been said so far, it is clear that although the calculation presented in this section includes into the evolution some nonlinear effects, it is not the end of the story. At very low $`x`$, this approximation should break down. There are clear reasons why this should happen in both pictures. In the projectile evolution picture, it is not true indefinitely that the WW fields are emitted independently from the partons in the projectile. Due to the evolution, more and more gluons are emitted into the wave function of the projectile and so the density of partons grows. At some point, the approximation of independent emissions as well as of independent scattering of the partons on the target must break down. This is the point at which, in the parlance of Ref. , the pomeron loop diagrams must come into play. In the target evolution picture, the problematic point is the eikonal approximation for the scattering of the $`q\overline{q}`$ pair. At not very low $`x`$ the target fields are not too strong. Since the $`q\overline{q}`$ pair is very energetic, the eikonal approximation is perfectly valid. However, with the evolution the strength of the target fields grows. The energy of the $`q\overline{q}`$ pair, on the other hand, stays fixed. When the fields are strong enough the quark and antiquark will start losing a finite fraction of their energy and, therefore, the no recoil eikonal approximation cannot stay valid indefinitely.
The effect of the nonlinear evolution Eq. (2.1) on the behavior of the total cross section was studied in Refs. . It was concluded that the nonlinearities slow down the BFKL type rise of the cross section and lead to its unitarization so that the cross section approaches the black disk limit. We want to conclude this section with a comment on the nature of the unitarization in this approximation. Essentially, the unitarization is brought about by purely kinematical effects. This is especially clear in the projectile evolution picture. At the initial value of $`x=x_0`$ one starts with the $`q\overline{q}`$ pair as the only relevant component of the photon wave function, which has a certain probability $`P_{q\overline{q}}`$ to scatter on the target. So, initially, the total scattering probability $`P_{x_0}`$ is
$$P_{x_0}=P_{q\overline{q}}.$$
(33)
At lower $`x=x_0\delta x`$ the wave function also contains a component with an extra gluon. Let the probability to have an extra gluon in the wave function be $`\mathrm{\Delta }`$ and the probability for this gluon to scatter on the target $`P_g`$. In the linear approximation (the BFKL limit) the total probability of scattering is additive
$$P_x=(1\mathrm{\Delta })P_{q\overline{q}}+\mathrm{\Delta }(P_{q\overline{q}}+P_g)=P_{q\overline{q}}+\mathrm{\Delta }P_g.$$
(34)
However, this is, in fact, overcounting, since there are events where both the gluon and the $`q\overline{q}`$ pair undergo scattering and those events are counted twice in the linear approximation. One should, therefore, subtract the probability of these double scattering events from the total probability. This deficiency is corrected by writing
$$P_x=(1\mathrm{\Delta })P_{q\overline{q}}+\mathrm{\Delta }(P_{q\overline{q}}+P_gP_{q\overline{q}}P_g)=P_{q\overline{q}}+\mathrm{\Delta }(P_gP_{q\overline{q}}P_g)=P_{x+\delta x}+\mathrm{\Delta }(1P_{x+\delta x})P_g.$$
(35)
At arbitrary low $`x`$, the same argument leads to a similar expression where $`P_{x+\delta x}`$ denotes the total scattering probability of the projectile (which itself contains the $`q\overline{q}`$ pair and some number of gluons) at a slightly higher value of $`x`$. This is precisely the nonlinear term in the evolution equation Eq. (2.1) with the only difference that the extra gluon in the wave function can have arbitrary transverse coordinate and one should, of course, integrate over this extra degree of freedom. It is clear that this negative nonlinear correction leads to the unitarization of the cross section since as $`P_x`$ tends to unity the emission of the extra gluon does not increase the total scattering probability. This effect is somewhat similar to the Glauber mechanism, not in the sense that each parton undergoes multiple scattering, but that the unitarization is of a purely geometrical nature. A similar discussion, in the framework of the dipole model, is given in Refs. .
In the next section we will show how to relate the approach just discussed with that of Refs. .
## 3 The JKLW equation and the small induced field limit.
We start this section by recalling the framework and results of Refs. .
### 3.1 The JKLW equation.
In this approach, following Refs. , the averages of gluonic observables in a hadron are calculated via the following path integral<sup>5</sup><sup>5</sup>5An alternative form of the effective action was suggested in where it was also shown that it leads to the same evolution equation.
$$\begin{array}{c}O(A)=D\alpha ^iDA^\mu O(A)Z[\alpha ]\times \hfill \\ \hfill \times \mathrm{exp}\left\{id^4x\frac{1}{4}\mathrm{tr}F^{\mu \nu }F_{\mu \nu }\frac{1}{N_c}d^2x_{}𝑑x^{}\delta (x^{})^i\alpha _a^i(x_{})\mathrm{tr}T_a𝒲_\mathrm{},\mathrm{}[A^{}](x^{},x_{})\right\},\end{array}$$
(36)
where the gluon field strength tensor is given by
$$F_a^{\mu \nu }=^\mu A_a^\nu ^\nu A_a^\mu gf_{abc}A_b^\mu A_c^\nu $$
(37)
and $`𝒲`$ is the Wilson line in the adjoint representation along the $`x^+`$ axis
$$𝒲_{\mathrm{},+\mathrm{}}[A^{}](x^{},x_{})=𝒫\mathrm{exp}\left[+ig𝑑x^+A_a^{}(x^+,x^{},x_{})T_a\right].$$
(38)
The hadron is represented by an ensemble of chromoelectric fields, localized in the plane $`x^{}=0`$, of the form
$$f^{+i}=\frac{1}{g}\delta (x^{})\alpha ^i(x_{}),$$
(39)
where the two dimensional vector potential $`\alpha ^i(x_{})`$ is “pure gauge”
$`^i\alpha _a^j`$ $`^j\alpha _a^if_{abc}\alpha _b^i\alpha _c^j=0.`$ (40)
In Eq. (36), $`Z[\alpha ]`$ is the statistical weight of a configuration $`\alpha _i(x_{})`$ in the hadronic ensemble.
The evolution in Refs. is derived in the target evolution picture where decreasing $`x`$ corresponds to boosting the hadronic target. This leads to freezing of part of the gluonic degrees of freedom. Integrating out these slow modes of the vector potential generates the renormalization group equation, which has the form of the evolution equation for the statistical weight $`Z`$
$$\frac{d}{d\mathrm{ln}\frac{1}{x}}Z=\alpha _s\left\{\frac{1}{2}\frac{\delta ^2}{\delta \alpha (u)\delta \alpha (v)}\left[Z\chi (u,v)\right]\frac{\delta }{\delta \alpha (u)}\left[Z\sigma (u)\right]\right\}.$$
(41)
In the compact notation used in Eq. (41), both $`u`$ and $`v`$ stand for color and rotational index and transverse coordinates, with summation and integration over repeated occurrences implied. This evolution equation for the statistical weight can be rewritten as the set of the evolution equations for the correlation functions of the chromoelectric field
$$\begin{array}{c}\frac{d}{d\mathrm{ln}\frac{1}{x}}\alpha _{a_1}^{i_1}(x_1)\mathrm{}\alpha _{a_n}^{i_n}(x_n)=\hfill \\ \hfill =\alpha _s[\underset{0<l<n+1}{}<\alpha _{a_1}^{i_1}(x_1)\mathrm{}\alpha _{a_{l1}}^{i_{l1}}(x_{l1})\alpha _{a_{l+1}}^{i_{l+1}}(x_{l+1})\mathrm{}\alpha _{a_n}^{i_n}(x_n)\sigma _{a_l}^{i_l}(x_l)\\ \hfill +\underset{0<m<k<n+1}{}<\alpha _{a_1}^{i_1}(x_1)\mathrm{}\alpha _{a_{m1}}^{i_{m1}}(x_{m1})\alpha _{a_{m+1}}^{i_{m+1}}(x_{m+1})\mathrm{}\\ \hfill \times \alpha _{a_{k1}}^{i_{k1}}(x_{k1})\alpha _{a_{k+1}}^{i_{k+1}}(x_{k+1})\mathrm{}\alpha _{a_n}^{i_n}(x_n)\chi _{a_ma_k}^{i_mi_k}(x_m,x_k)].\end{array}$$
(42)
The quantities $`\chi [\alpha ]`$ and $`\sigma [\alpha ]`$ have the meaning of the mean fluctuation and the average value of the induced vector potential which arises from the field modes which become frozen due to extra boost of the hadronic target. In the leading logarithmic approximation of Refs. the two quantities $`\chi `$ and $`\sigma `$ completely specify the low $`x`$ evolution. We give here an explicit expression for the mean fluctuation $`\chi `$ which will be the focus of our interest throughout this section
$$\chi _{ab}^{ij}(x_{},y_{})=2x_{}|\{\frac{D^i}{D_{}^2}[D_{}^2S^1]\frac{D^j}{D_{}^2}\}_{ab}|y_{}.$$
(43)
For convenience, we have defined
$`\alpha _{ab}^i`$ $`=f_{abc}\alpha _c^i,`$
$`D_{ab}^i`$ $`=^i\delta _{ab}+\alpha _{ab}^i.`$ (44)
The operator $`S`$ in Eq. (43) is given by
$$S=\frac{1}{D_{}^2}+2[\frac{^i}{_{}^2}\frac{D^i}{D_{}^2}][\frac{^i}{_{}^2}\frac{D^i}{D_{}^2}]=\frac{1}{D_{}^2}2\frac{1}{_{}^2}_{}\alpha \frac{1}{D_{}^2}+2\frac{1}{D_{}^2}D_{}\alpha \frac{1}{_{}^2}.$$
(45)
### 3.2 Where does it come from?
Technically, these results are derived as follows. One considers the quantum corrections in the classical background field Eq. (39). The calculation is performed in the lightcone gauge $`A^+=0`$ with the residual gauge fixing $`^iA^i(x^{}\mathrm{})=0`$ which fixes the gauge completely. In this gauge the chromoelectric field Eq. (39) corresponds to the background vector potential
$$b^i=\theta (x^{})\alpha ^i(x_{}).$$
(46)
Note that, as opposed to the previous section, here we are using a different lightcone gauge: $`A^+=0`$. As a consequence, the background vector potential has a different form.
The complete set of on-shell small fluctuation solutions of the classical equations is
$$\begin{array}{c}a_{p^{},r}^i=e^{ip^{}x^+}d^2p_{}[\theta (x^{})\mathrm{exp}(i\frac{p_{}^2}{2p^{}}x^{}ip_{}x_{})v_{,r}^i(p_{})\hfill \\ \hfill +\theta (x^{})U(x_{})\mathrm{exp}(i\frac{p_{}^2}{2p^{}}x^{}ip_{}x_{})\left[U^{}v_{+,r}^i\right](p_{})+\theta (x^{})\gamma _{+,r}^i].\end{array}$$
(47)
Here, $`r`$ is the degeneracy label, which labels independent solutions with the frequency $`p^{}`$. In the free case it is conventionally chosen as the transverse momentum, $`\{r\}=\{p_{}\}`$. The matrix $`U(x_{})`$ is the $`SU(N)`$ matrix that parameterizes the two dimensional “pure gauge” vector potential $`\alpha ^i(x_{})`$
$$\alpha ^i(x_{})=iU(x_{})^iU^{}(x_{}).$$
The auxiliary functions $`\gamma _+^i,v_\pm ^i`$ are all determined in terms of one vector function. Choosing this independent function as $`v_{}^i`$ we have
$`v_{+,r}^i`$ $`=\left[T^{ij}L^{ij}\right]\left[t^{jk}l^{jk}\right]v_{,r}^k,`$ (48)
$`\gamma _{+,r}^i`$ $`=2D^i\left[{\displaystyle \frac{D^j}{D_{}^2}}{\displaystyle \frac{^j}{_{}^2}}\right]\left[t^{jk}l^{jk}\right]v_{,r}^k,`$ (49)
where we have defined the projection operators
$`T^{ij}`$ $`\delta ^{ij}{\displaystyle \frac{D^iD^j}{D_{}^2}},`$ $`L^{ij}`$ $`{\displaystyle \frac{D^iD^j}{D_{}^2}},`$
$`t^{ij}`$ $`\delta ^{ij}{\displaystyle \frac{^i^j}{_{}^2}},`$ $`l^{ij}`$ $`{\displaystyle \frac{^i^j}{_{}^2}}.`$ (50)
The $`\gamma _+`$ piece of the eigenfunction Eq. (47) is responsible for the induced vector potential since this is the only contribution that does not vanish at $`x^{}\mathrm{}`$, so that
$$\chi _{ab}^{ij}(x_{},y_{})=4\pi 𝑑p^{}\gamma _{+,a}^i(x_{},p^{})\gamma _{+,b}^j(y_{},p^{}).$$
(51)
Note that the essential nonlinearity of the expression Eq. (43) is due to the denominator in the operator $`S^1`$ Eq. (45). The reason this arises is due to the nontrivial normalization of the small fluctuation eigenfunctions. As discussed in detail in Refs. the proper normalization of the eigenfunctions requires $`v_{}^i`$ to be chosen as a complete set of eigenfunctions of the two dimensional Hermitian operator $`O^1`$
$$[(tl)O^1(tl)]_{ab}^{ij}(x_{},y_{})=x_{}|\delta _{ab}^{ij}2\left[[^i\frac{1}{_{}^2}D^i\frac{1}{D_{}^2}]S^1[\frac{1}{_{}^2}^j\frac{1}{D_{}^2}D^j]\right]_{ab}|y_{},$$
(52)
such that
$$d^2r_{}v_{,r,a}^i(x_{})v_{,r,b}^j(y_{})=\frac{1}{4\pi |p^{}|}[O^1]_{ab}^{ij}(x_{},y_{}).$$
(53)
This nontrivial normalization is the consequence of the presence of the $`\gamma _+`$ piece in the solution Eq. (47). Equation (51), supplemented by Eq. (49) and the normalization Eq. (53), leads to the final expression Eq. (43).
If the contribution of $`\gamma _+`$ could be neglected in the normalization condition, the normalization of the eigenfunctions would be trivial and we would have $`O=1`$ in Eq. (53). One can consider the limit in which $`\gamma _+`$, or equivalently $`\chi `$, is small. In the leading order in the expansion in $`\gamma _+`$ we have a very simple expression for $`\chi `$ <sup>6</sup><sup>6</sup>6If so desired this expression can be written in a simple form in terms of the unitary matrix $`U`$, since operatorially $`D^i=U^iU^{}`$. In Fourier space this gives convolutions of $`U(p)`$ and powers of transverse momentum.
$$\stackrel{~}{\chi }_{ab}^{ij}(x_{},y_{})=4x_{}|\left[D^i\{\frac{1}{_{}^2}+\frac{1}{D_{}^2}\frac{1}{_{}^2}_{}D_{}\frac{1}{D_{}^2}\frac{1}{D_{}^2}D_{}_{}\frac{1}{_{}^2}\}D^j\right]_{ab}|y_{}.$$
(54)
Note that this is a different limit than the one in which the JKLW evolution reduces to the BFKL equation . The BFKL limit corresponds to the expansion in powers of the background field $`\alpha ^i`$. Now, we are not assuming that $`\alpha ^i`$ is small, but rather that the correction induced by the evolution, $`\gamma ^i`$ is small.
### 3.3 BK to JKLW: transforming between the gauges.
We will now see that Eq. (54) is reproduced precisely by translating the calculation of the previous section into the language of the JKLW evolution.
In the previous section, following Ref. , we used the gauge $`A^{}=0`$. This is a very convenient gauge from the point of view of the projectile evolution since the eikonal amplitudes in this gauge are given by simple Wilson line factors. We will refer to this gauge as the “projectile lightcone gauge”, or the “projectile gauge” for short. The JKLW approach on the other hand uses the $`A^+=0`$ gauge, which is convenient for the target evolution picture since it simplifies the relation between the distribution functions and the correlators of the gluon fields. We will call this gauge the “target lightcone gauge” or, simply, the “target gauge”. Our immediate aim is, therefore, to calculate $`\chi `$ using the results of the calculation in the projectile gauge.
To do this note that the relation between the fields in the target and projectile gauges is given by
$$\frac{1}{g}B^\mu +A^\mu =V(\frac{1}{g}b^\mu +a^\mu )V^{}+\frac{i}{g}V^\mu V^{}.$$
(55)
To simplify the notation, from now on we will denote the fields in the target ($`A^+=0`$) gauge by capital letters and fields in the projectile ($`a^{}=0`$) gauge by lower case letters. This we do for both the background part of the field and for the small fluctuation part. The field dependent matrix $`V`$ is given by
$$V=P\mathrm{exp}\left[i_{\mathrm{}}^x^{}𝑑x^{}(b^++ga^+)\right].$$
(56)
The condition $`A^+=0`$ does not by itself specify the lower limit of the integration over $`x^{}`$ in the exponential. However, choosing this limit to be at minus infinity ensures that $`V(x^{}\mathrm{})=\mathrm{𝟣}`$ and, as a consequence, $`A^i(x^{}\mathrm{})=a^i(x^{}\mathrm{})`$. The projectile gauge fields satisfy the standard vanishing boundary conditions at infinity. This choice of the lower limit of the integration, therefore, guarantees that the target gauge fields also vanish at $`x^{}\mathrm{}`$ and, further, satisfy the residual gauge condition $`^iA^i(x^{}\mathrm{})=0`$ that was imposed in Refs. . To calculate $`\chi `$ we only need to consider the linearized relation between the small fluctuations of the fields in the two gauges. To do this we need to expand $`V`$ to first order in $`a^+`$. This has been done in the previous section. Taking only linear terms in $`a^+`$ in Eq. (8a) and substituting them into Eq. (55) we find for the transverse components of the field
$`A_a^i(x)`$ $`=\theta (x^{})\left[a_a^i(x){\displaystyle _{\mathrm{}}^x^{}}𝑑x^{}^ia^+\right]+`$
$`+\theta (x^{})\left[\stackrel{~}{U}^{ab}a_b^i(x)D_{ab}^i\left({\displaystyle _{\mathrm{}}^0}𝑑x^{}a_b^++\stackrel{~}{U}_{bc}{\displaystyle _0^x^{}}𝑑x^{}a_c^+\right)\right]`$
$`=\left[\delta ^{ij}W^iW^{}{\displaystyle \frac{1}{^+^{}}}W^jW^{}\right]_{ab}(Wa^j)_b.`$ (57)
Here, the matrix $`\stackrel{~}{U}`$ is the same as in the previous section and is related to the classical background by
$`\stackrel{~}{U}(x_{})=𝒫\mathrm{exp}\{i{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x^{}b^+\},`$
$`B^i=\theta (x^{})i\stackrel{~}{U}^i\stackrel{~}{U}^{}`$ (58)
as per Eq. (46). We have also defined
$$W(x^{})=𝒫\mathrm{exp}\left\{i_{\mathrm{}}^x^{}𝑑x^{}b^+(x^{})\right\}=\theta (x^{})1+\theta (x^{})\stackrel{~}{U},$$
(59)
which is essentially the classical part of $`V`$. The operator $`\frac{1}{^+}`$ in the last line of Eq. (57) is defined as the integral from $`\mathrm{}`$<sup>7</sup><sup>7</sup>7 Eq. (57) has been derived also in Refs. . The only difference in our derivation is that the meaning of the $`\frac{1}{p^\pm }`$ pole is entirely unambiguous and, as discussed above, is dictated by the residual gauge condition.. We will further simplify this expression by using the on-shellness conditions
$`a^+={\displaystyle \frac{^i}{^{}}}a^i,`$
$`(2^{}D^+[b^+]_{}^2)a^i=0.`$ (60)
The resulting relation between the on-shell transverse fields in the two gauges is
$$\begin{array}{c}A_a^i=\theta (x^{})(tl)^{ij}a_a^j(x^{})+\hfill \\ \hfill +\theta (x^{})\left\{\stackrel{~}{U}_{ab}(tl)^{ij}a_b^j(x^{})2D_{ab}^i\left[\frac{^j}{_{}^2}a_b^j(x^{}0^{})\stackrel{~}{U}_{bc}\frac{^j}{_{}^2}a^j(x^{}0^+)\right]\right\}.\end{array}$$
(61)
Note that we have to specify on which side of $`x^{}=0`$ the fields are taken since the solutions of the small fluctuation equations in the projectile gauge are discontinuous at zero. Now, recall that $`a^i`$ satisfies, everywhere except at $`x^{}=0`$, the free equations of motion. With this in mind we can compare this equation with Eq. (47). We see that Eq. (61) is indeed precisely of the form Eq. (47) with
$`\gamma _{+,a}^i`$ $`=2D_{ab}^i\left[{\displaystyle \frac{^j}{_{}^2}}a_b^j(x^{}0^{})\stackrel{~}{U}_{bc}{\displaystyle \frac{^j}{_{}^2}}a_c^j(x^{}0^+)\right],`$
$`v_{,a}^i`$ $`=(tl)^{ij}a_a^j(x^{}0^{}),`$
$`v_{+,a}^i`$ $`=\stackrel{~}{U}_{ab}(tl)^{ij}a_b^j(x^{}0^+).`$ (62)
Remembering that (see for example Ref. )
$$a^i(x^{}0^+)=\stackrel{~}{U}^{}a^i(x^{}0^{}),$$
(63)
we see that the functions $`v_+^i,v_{}^i`$ and $`\gamma _+^i`$ are related precisely by the same relations as in Eq. (49).
We have established, therefore, that if $`a^i`$ satisfies the equation of motion in the projectile gauge, then the transformed field $`A^i`$ of Eq. (55) satisfies the equations of motion in the target gauge. The only remaining question is that of the normalization of the eigenfunctions. Recall that the functions $`a^i`$ in the calculation of Ref. , which was reproduced in the previous section, were normalized in the same way as the eigenfunctions of the free theory. That is to say, the full set of on-shell eigenfunctions is obtained by choosing $`a_p_{}^i(x^{}0^{})`$ as a complete set of normalized eigenfunctions of the unit operator in the transverse space
$$\frac{d^2p_{}}{4\pi ^2}a_p_{}^i(x_{},x^{}0^{})a_p_{}^j(y_{},x^{}0^{})=\delta ^{ij}\delta (x_{}y_{}).$$
(64)
Since $`tl`$ is a unitary operator, Eq. (62) tells us that $`v_{}^i`$ is also normalized to unity rather than to a nontrivial operator $`O`$ as in the JKLW calculation Eq. (53). Using this, as well as the relations Eq. (62) and Eq. (51), we find that when translated into the language of JKLW, the results of Ref. give Eq. (54) as the mean fluctuation of the induced chromoelectric field. The essential nonlinearity of Eq. (43) is, therefore, absent in this calculation.
So far, we have only considered the real part of the JKLW kernel, $`\chi `$. Of course, the same method can be applied to find what is the form of the virtual part $`\sigma `$, Eq. (41), that arises from the calculation of Ref. . To reproduce the virtual part it is clearly necessary to keep the quadratic terms in the relation between $`A^i`$ and $`a^i`$. Thus the quadratic terms in Eq. (8a) will be important in this calculation. Other than that, the calculation is straightforward. Again the gauge invariance ensures that all the “kinematical” factors of $`\sigma `$ of Refs. are reproduced in the projectile gauge calculation and the only difference comes from the difference in the normalization of the eigenfunctions. It is clear, therefore, that the result of such a calculation is again the lowest order expansion of $`\sigma `$ in powers of $`\gamma _+`$
$$\begin{array}{c}\stackrel{~}{\sigma }_a^i=\left[\frac{D^i}{D^2}\right]_{ab}(\frac{N_c}{2}(^j\alpha _b^j)x_{}|\frac{1}{^2}|x_{}\hfill \\ \hfill f_{bcd}x_{}|[4D^j\frac{1}{^2}D\frac{D^j}{D^2}+2\frac{1}{^2}\alpha 2\alpha \frac{1}{^2}+4\alpha ^j\frac{1}{^2}\alpha ^j]_{cd}|x_{})\\ \hfill 2ϵ^{ij}\left[\frac{D^i}{D^2}\right]_{ab}(x,y)f_{bcd}ϵ^{kl}y_{}|\left[D^k\left(\frac{1}{^2}+\frac{1}{D^2}\frac{1}{^2}D\frac{1}{D^2}\frac{1}{D^2}D\frac{1}{^2}\right)D^l\right]_{cd}|y_{}.\end{array}$$
(65)
### 3.4 The Doubly Logarithmic Limit.
Before exploring the relationship between the two approaches further in the next section, we want to make a comment about the form Eq. (54). Although this equation certainly gives $`\stackrel{~}{\chi }`$ in general as a nonlinear function of the background field $`\alpha ^i`$, this nonlinearity disappears in the double logarithmic limit. Following Ref. we take the double logarithmic limit as the limit when the background field $`\alpha ^i`$ does not depend on $`x_{}`$. In this limit, the covariant and the simple derivatives commute and it is easy to see that Eq. (54) reduces to
$$\stackrel{~}{\chi }^{ij}=4\frac{\alpha ^2}{^2}\frac{D^iD^j}{D^2}$$
(66)
or
$$\mathrm{tr}\stackrel{~}{\chi }=4\mathrm{t}\mathrm{r}\frac{\alpha ^2}{^2}.$$
(67)
When substituted into the evolution equation Eq. (42) this gives the simple linear double logarithmic DGLAP evolution for the gluon distribution function $`G\mathrm{tr}\alpha ^2`$ (see Ref. for a detailed derivation). This is in contrast with the situation discussed in Ref. where the double logarithmic limit of Eq. (43) was studied. It was shown there that the nonlinearities in Eq. (43) survive in the doubly logarithmic limit and, in fact, lead even in this limit to the “almost saturation” of the gluon distribution.
The absence of the nonlinearities is in contradiction with the explicit calculation of Mueller and Qiu who showed that the QCD evolution of the gluon distribution in the doubly logarithmic approximation does indeed contain contributions from higher twist operators. This again underscores our observation that the nonlinearities included in the evolution of Ref. are not the whole story. Those are the “kinematical” nonlinearities in the sense discussed in the previous section and do not include interesting dynamical effects which come into play when the parton density becomes large.
In fact, the triviality of the doubly logarithmic limit of the calculation of the previous section is easy to understand using the intuition of the dipole model approach. In the projectile evolution picture, the doubly logarithmic limit is achieved by assuming that, in every step in the evolution, the extra gluon that is emitted into the virtual photon wave function has the smallest transverse momentum or, in the coordinate space, has the largest transverse coordinate . The part of the wave function that contains this gluon, therefore, essentially describes one adjoint dipole of large transverse size. One leg of this dipole is the newly emitted gluon, while the other leg is the remainder partons which are closely bunched together in the coordinate space<sup>8</sup><sup>8</sup>8These partons are in the adjoint representation of the color group, since together with the extra emitted gluon the state must be an overall singlet.. The cross section for the scattering of the adjoint dipole in the large $`N_c`$ limit is simply related to the cross section for the fundamental dipole $`\sigma _{\mathrm{adj}}=2\sigma _{\mathrm{fund}}\sigma _{\mathrm{fund}}^2`$. Thus, if initially one starts (like in Ref. or Ref. ) from a fundamental dipole, the appearance of the large adjoint dipole in the wave function leads to a nonlinear GLR type quadratic term in the evolution equation for the scattering cross section. This result indeed has been derived in Ref. . However, if we want to consider the evolution of the gluon distribution itself, the initial state should contain an adjoint rather than a fundamental dipole. This can be achieved by considering “DIS” of a virtual particle that couples to $`\mathrm{tr}F^2`$ . In this case, in any step in the doubly logarithmic evolution, the state contains only one adjoint dipole. The probability for the appearance of a larger dipole in the approximation of independent emissions is itself proportional to the number of gluons. The evolution of the gluon distribution in this approximation is, therefore, necessarily linear and is merely the simple DGLAP DLA.
## 4 More on the target versus projectile gauge
The discussion of the previous section may seem a little paradoxical on the purely technical level. Indeed, we have been calculating the same physical quantity in two different ways. The quantity in question is the equal (lightcone) time propagator of the transverse components of the vector potential $`A^i`$ in the target lightcone gauge. The first way of performing the calculation is to work entirely in the target gauge as was done in Refs. . This gives the result Eq. (43). The second way to calculate the same quantity is to first calculate the propagator of $`a^i`$ in the projectile gauge and then gauge transform the result into the target gauge using Eqs. (57,61). This results in an inequivalent expression Eq. (54).
### 4.1 The $`iϵ`$ complication.
Our first aim in this section is to resolve this technical paradox. To do this let us consider in more detail the calculation of Refs. and its transformation into the projectile gauge. The equal time propagator of the transverse components of vector potential is calculated in the following way . One starts with the quadratic part of the action for the small fluctuations of $`A^\mu `$. Integrating $`A^{}`$ it is reduced to a quadratic action for the small fluctuations of the transverse components of the vector potential
$$S=d^4xd^4yA^i(x)G^{1ij}(x,y)A^j(y).$$
(68)
For the purpose of this discussion we use somewhat simplified notations and omit the color indices on the fields. The explicit form of $`G^1`$ is given in Refs. . One then finds properly normalized eigenfunctions of $`G^1`$
$`G^{1ij}(x,y)A_{\lambda ,p^{},r}^j(y)=\lambda A_{\lambda ,p^{},r}^i(x),`$
$`{\displaystyle d^4xA_{\lambda ,p^{},r}^i(x)A_{\lambda ^{},p^{^{}},r^{}}^i(x)}=\delta (\lambda \lambda ^{})\delta (p^{}p^{^{}})\delta ^2(rr^{}).`$ (69)
Using the complete set of eigenfunctions one constructs the propagator with the standard $`iϵ`$ prescription as
$$G^{ij}(x,y)=\frac{d\lambda }{\lambda +iϵ}𝑑p^{}d^2rA_{\lambda ,p^{},r}^i(x)A_{\lambda ,p^{},r}^j(y).$$
(70)
The limit $`x^+=y^+`$, and $`x^{},y^{}\mathrm{}`$ is then taken to calculate $`\chi ^{ij}`$. Clearly, the equal time limit selects the on-shell eigenfunctions $`\lambda =0`$ and, therefore, when transforming into the projectile gauge it is important to keep track of the $`iϵ`$ prescription. The simplest way to do this is to include the $`iϵ`$ term directly in the action
$$S_t=d^4xd^4yA^i(x)\left[G^{1ij}(x,y)+iϵ\delta ^{ij}\delta (xy)\right]A^j(y).$$
(71)
The propagator Eq. (70) is then just the inverse of the quadratic form in Eq. (71) without any additional regulators.
To transform this expression into the projectile gauge one has to use Eqs. (57,61). The gauge invariance of the QCD action ensures that the first term in Eq. (71) under this transformation transforms into
$$a^i(x)D^2(x,y)a^i(y),$$
(72)
which is exactly the action used in Sec. 2 to calculate the projectile gauge propagator. However, the $`iϵ`$ term is not so simple. If the transformation Eq. (61) was unitary, the norm of the field $`A^i`$ would be preserved and the $`iϵ`$ term in Eq. (71) would transform into the standard $`iϵd^4xa_i(x)a_i(x)`$ term in the projectile gauge. The problem is that, as we saw in the previous section, the transformation Eq. (61) is not unitary. A normalized function $`a^i`$ is transformed into a function $`A^i`$ normalized not to unity but rather to an eigenvalue of the operator $`O`$ in Eq. (52).<sup>9</sup><sup>9</sup>9The fact that the transformation between the two gauges is non-unitary is not unusual. Even though it is a gauge transformation and therefore formally unitary, the gauge parameter itself depends on the dynamical field. Such transformations are generically non-unitary and do not preserve the scalar product. This was precisely the root for the discrepancy between $`\chi `$ and $`\stackrel{~}{\chi }`$. The resulting projectile gauge action can be written as<sup>10</sup><sup>10</sup>10In writing this expression we have made use of the fact that the $`iϵ`$ term is important only for functions $`a^i`$ that satisfy $`D^2a^i=0`$.
$$\begin{array}{c}S_p=d^4xa^i(x)D^2a^i(x)+\hfill \\ \hfill +iϵd^4x\left[a^i(x)a^i(x)+2a^i(0^{})[^i\frac{1}{_{}^2}D^i\frac{1}{D_{}^2}]D_{}^2[^i\frac{1}{_{}^2}D^i\frac{1}{D_{}^2}]a^i(0^{})\right].\end{array}$$
(73)
Thus, the standard $`iϵ`$ prescription in the target gauge is equivalent to a fairly complicated momentum dependent prescription in the projectile gauge. Since the calculation of Sec. 2, following Ref. , was performed using the standard $`iϵ`$ prescription in the projectile gauge the result is, indeed, expected to differ from that of Refs. .
While the technical reason for the difference between the results of Refs. and Ref. is clear, the physics behind it is not so obvious. In the rest of this section, we will make an attempt to understand the physical reason for this difference.
As we have just explained, the calculation of Refs. is equivalent to a calculation in the projectile gauge with a nonstandard momentum dependent $`iϵ`$ prescription. It is well known that such a change of prescription is equivalent to a calculation not in the vacuum state but rather in a state which contains gluons . We, therefore, ask ourselves why the projectile gauge calculation should be performed in a state which, on top of the background field $`b^+`$, also contains additional gluons.
### 4.2 Evolution as Renormalization Group in the projectile gauge.
To answer this question let us first try to reformulate the projectile gauge calculation of Sec. 2 in terms of the Wilson renormalization group akin to the approach of Refs. . The hadron is represented as a statistical ensemble of the static $`b^+`$ fields of the form Eq. (1) with a statistical weight $`Z[b]`$. Evolution in $`x`$ generates induced vector potential which changes the statistical weight. Strictly speaking the induced vector potential is not static. It has components in the frequency range $`p^{}<\mathrm{\Lambda }1/x`$. However, as long as the frequency of the components of the projectile wave function are large enough, one can treat the induced potential as static during the interaction with the $`q\overline{q}`$ pair. Also, as long as the wavelength of the projectile in the $`x^{}`$ direction is large enough ($`p^+`$ is small) the induced vector potential can be approximated by a $`\delta (x^{})`$ shaped function. One can, equivalently, describe the hadron by a statistical ensemble of $`V`$ and $`V^{}`$ with some statistical weight $`Z[V,V^{}]`$. It is clear that both descriptions encode exactly the same information. For this purpose one has to define analogs of $`\chi `$ and $`\sigma `$, that is the (connected) fluctuation correlation functions of order $`\alpha _s`$. The resulting evolution equation
$$\begin{array}{c}\frac{d}{d\mathrm{ln}\frac{1}{x}}Z[U,U^{}]=\hfill \\ \hfill \alpha _s[\frac{1}{2}(\frac{\delta ^2}{\delta U(u)\delta U^{}(v)}\left[Z\chi _{q\overline{q}}(u,v)\right]+\frac{\delta ^2}{\delta U^{}(u)\delta U(v)}\left[Z\chi _{\overline{q}q}(u,v)\right]\\ \hfill +\frac{\delta ^2}{\delta U(u)\delta U(v)}\left[Z\chi _{qq}(u,v)\right]+\frac{\delta ^2}{\delta U^{}(u)\delta U^{}(v)}\left[Z\chi _{\overline{q}\overline{q}}(u,v)\right])\\ \hfill \frac{\delta }{\delta U(u)}\left[Z\sigma _q(u)\right]\frac{\delta }{\delta U^{}(u)}\left[Z\sigma _{\overline{q}}(u)\right]]\end{array}$$
(74)
is the analog of Eq. (41).
Using the formulae of Sec. (2) and the appendix we find
$$\sigma _q^{\alpha \beta }(x_{})=\frac{1}{2\pi ^2}d^2z_{}\frac{1}{(xz)_{}^2}\left[\mathrm{tr}\left(U(x_{})U^{}(z_{})\right)\left(U(z_{})\right)^{\alpha \beta }N_c\left(U(x_{})\right)^{\alpha \beta }\right],$$
(75)
$`\chi _{q\overline{q}}^{\alpha \beta ,\gamma \delta }(x_{},y_{})`$ $`={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}}`$
$`[(U(z_{})U^{}(y_{}))^{\alpha \delta }(U^{}(z_{})U(x_{}))^{\gamma \beta }`$
$`+\left(U(x_{})U^{}(z_{})\right)^{\alpha \delta }\left(U^{}(y_{})U(z_{})\right)^{\gamma \beta }`$
$`\delta ^{\alpha \delta }(U^{}(y_{})U(x_{}))^{\gamma \beta }(U(x_{})U^{}(y_{}))^{\alpha \delta }\delta ^{\gamma \beta }],`$ (76)
$`\chi _{qq}^{\alpha \beta ,\gamma \delta }(x_{},y_{})`$ $`={\displaystyle \frac{1}{2\pi ^2}}{\displaystyle d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}}`$
$`[\left(U(z_{})\right)^{\alpha \delta }(U(y_{})U^{}(z_{})U(x_{}))^{\gamma \beta }`$
$`+\left(U(x_{})U^{}(z_{})U(y_{})\right)^{\alpha \delta }\left(U(z_{})\right)^{\gamma \beta }`$
$`\left(U(x_{})\right)^{\alpha \delta }\left(U(y_{})\right)^{\gamma \beta }\left(U(y_{})\right)^{\alpha \delta }\left(U(x_{})\right)^{\gamma \beta }]`$ (77)
$`\sigma _{\overline{q}}`$ is obtained from $`\sigma _q`$ by replacing $`U`$s by $`U^{}`$s. The same is true for $`\chi _{\overline{q}\overline{q}}`$ and $`\chi _{qq}`$. $`\chi _{\overline{q}q}`$ is obtained from $`\chi _{qq}`$ by swapping $`(x,\alpha \beta )`$ with $`(y,\gamma \delta )`$ as this exchanges the $`q`$ and $`\overline{q}`$ lines there. We stress that Eqs.(74-77) contain all the information that is contained in the BK equation as well as in the equations for higher correlation functions of $`U`$ that appear in Ref..
### 4.3 The snag.
There is one implicit assumption in this procedure: namely that the $`a^+`$ component of the vector potential is the only relevant one. One assumes that if, for example, an $`a^i`$ component is generated in the evolution it does not affect the subsequent evolution of the physical cross section. This is, however, not quite right. What is important for the interaction with the projectile is not merely $`b^+`$ but rather the $`F^{+i}`$ component of the color electric field. The interaction between the projectile and the target is due to the term $`F^iF^{+i}`$ in the QCD Lagrangian. The $`F^i`$ component is the Weiszäker-Williams field of the $`q\overline{q}`$ pair, while the $`F^{+i}`$ component is generated by the color charges in the target. In the eikonal approximation it is true that $`F^{+i}=^ib^+`$ and, therefore, the coupling can be written as $`b^+J^{}`$, where $`J^{}=^iF^i`$. However, if there is a contribution to $`F^{+i}`$ coming from the transverse component of the vector potential, it should be taken into account.
It is easy to see that such a contribution is indeed generated by the low $`x`$ evolution. Suppose one starts the evolution initially with the background field configuration as in Eq. (1). In the first step of the evolution one generates both the increment in $`a^+`$ and the increment in $`a^i`$. The two are related by the condition Eq. (60)
$$a^+=\frac{^i}{^{}}a^i.$$
(78)
Naively, one would expect that since all the fluctuation fields in this step have small frequencies, it should be true that $`a^+a^i`$ and, therefore, it should be safe to forget about $`a^i`$. The reason this is incorrect is that the on-shell solutions for $`a^i`$ are discontinuous at $`x^{}=0`$. Therefore, even though the field $`a^i`$ is indeed small, it has a large derivative with respect to $`x^{}`$ which contributes to the field strength. In fact, the induced chromoelectric field is
$$\delta F^{+i}=^ia^+D^+a^i.$$
(79)
Recalling that on-shell $`a^i`$ satisfy the second of the equations Eq. (60), we see that the second term in this expression is $`\frac{_{}^2}{2^{}}a^i`$ and is of the same order as the first term $`\frac{^i^j}{^{}}a^j`$. Clearly, even if one starts initially from a background which only contains $`b^+`$ after long enough evolution a large transverse component of the vector potential is generated. When the contribution of the transverse component to the field strength is comparable to the contribution of the “$`+`$” component, the eikonal approximation breaks down and the evolution discussed in Sec. 2 ceases to be valid. It looks indeed very natural that in order to take into account the presence of the (potentially large) transverse field, the calculation in the projectile gauge should be performed around a state that contains transverse gluons apart from the $`b^+`$ background.
One could try to argue that the transverse part of the vector potential can be somehow gauged away and the calculation could still be performed consistently around a pure $`b^+`$ background. Even if this is possible the evolution of the background defined by such a procedure will be different from the evolution of Sec. 2. In any case, we do not see how such “regauging” is possible.
It is instructive to see in more detail how the gauge fixing works in both the projectile and the target gauges and why the two seem to have different status as far as the renormalization group structure is concerned. As we mentioned above, the chromoelectric field is created by the color charges in the target. In fact, the whole renormalization group procedure can be formulated in terms of the color charge density $`j^+`$ rather than the vector potentials themselves, which was in fact originally done in Refs. . The background vector potentials are found as static solutions of classical equations of motion in the presence of the color charge density $`j^+=\rho \delta (x^{})`$
$`F^{ij}=0,`$
$`D^iF^{+i}=D^i[D^ib^+^+b^i]=j^+.`$ (80)
An important property of these equations is that for a given $`\rho `$ they have infinitely many solutions. By considering an arbitrary unitary matrix $`V(x_{},x^{})`$ it is straightforward to see that all the following are solutions
$`b^i=iV^{}^iV,`$
$`b^+={\displaystyle \frac{1}{D_{}^2}}[j^++D^i^+b^i].`$ (81)
The difference between the target and the projectile gauges at this point becomes important. In the target gauge $`B^+=0`$, the equations reduce to
$`B^i=iV^{}^iV,`$
$`D^i^+B^i=j^+.`$ (82)
We, therefore, get rid of almost all the solutions, the only residual degeneracy being the value of the matrix $`V`$ at $`x^{}\mathrm{}`$. The imposition of the residual gauge condition $`^iA^i(x^{}\mathrm{})`$ then removes all solutions except one. In the projectile gauge the situation is very different. The condition $`a^{}=0`$ does not eliminate any of the infinite number of solutions Eq. (81). The choice of the residual gauge fixing is thus crucial to eliminate the redundant solutions. If those are not eliminated, the perturbative calculation will be plagued with zero mode problems. The calculation in Sec. 2 was in fact performed with the residual gauge fixing $`^ia^i(p^{}=0)=0`$. This gauge fixing does indeed eliminate all the solutions except the one which has vanishing $`a^i`$ and has therefore precisely the form of Eq. (1).
Now, consider the renormalization group calculation. Here, we have to integrate out modes which have higher frequency $`p^{}`$. In the target gauge this is straightforward: the residual gauge condition does not care about frequency. It therefore eliminates nonzero frequency fluctuation modes which do not vanish at $`x^{}\mathrm{}`$ in the same way as it eliminated the static background solutions with this behavior. As a result, the fluctuation modes have a very similar structure to the background field and the induced field is similar to the background. It is, therefore, straightforward to formulate a selfsimilar renormalization group transformation in this gauge. The situation is quite different in the projectile gauge. The residual gauge condition, although it fixes unambiguously the background, has nothing to say about the fluctuations — it only fixes the static modes! It is impossible therefore to ensure that the fluctuations will have the same form as the static background. In fact, as we have seen above, it will not be the case. The first equation of Eq. (60) in the projectile gauge is just one of the equations of motion (with or without the external source). This means that it is possible to have nonvanishing $`a^+`$ with vanishing $`a^i`$ only at exactly zero frequency. At any finite frequency nonvanishing, $`a^i`$ is required. As we have seen, this $`a^i`$ contributes to the induced chromoelectric field, or equivalently to the induced color charge density. The calculation thus explicitly lacks a selfsimilar structure and proper renormalization group setup does not seem possible<sup>11</sup><sup>11</sup>11It may be possible to reformulate RG so that it would include also transverse background fields or equivalently finite number of gluons in addition to $`b^+`$. This seems, however, to be quite a complicated problem and is far beyond the scope of our discussion here. unless extra eikonal approximation is invoked.
The discussion of this section leads us to conclude that the projectile gauge calculation, as formulated in Ref. and Sec. 2 is only valid as long as the eikonal approximation is applicable. When the evolution is continued for a large span of $`1/x`$, the eikonal approximation breaks down and the higher nonlinear corrections of Refs. should become important. This is not to say that one cannot learn much from this simplified evolution. Quite to the contrary — clearly there is a range of $`x`$ values where this evolution captures the relevant physics. This is particularly true when the target is large — the case of a large nucleus discussed in Ref. . In this case, the eikonal cross section is significantly different from the simple perturbative one which assumes single scattering. The nonlinearity of the evolution becomes important much faster than for a small hadron. One, therefore, expects unitarization to appear already within the eikonal regime. Subsequent appearance of other nonlinear corrections will not change the fact that the total cross section has unitarized. It must, however, affect other more exclusive properties of the process such as the structure of final states. The spectrum of the target gauge fields is presumably directly related to the spectrum of the emitted gluons . Thus, when the evolution of these fields changes even locally one expects this change to be visible in the spectrum of final state gluons. Assuming the local parton-hadron duality this then has to be mirrored in the spectrum of final state hadrons.
## 5 Conclusions
In recent years several approaches to the evolution of dense gluonic systems in the saturation regime were developed. The approaches differ from one another in many technical respects and relationship between the physics is also not always clear. In this paper, our aim was to relate two of these approaches and thereby to try and reduce the entropy in the field. We have shown that the nonlinear JKLW equation of Refs. coincides with the BK evolution equation derived in Ref. and as long as the gluon field induced by the evolution is small. We have argued that the approach of Ref. should break down when the field is large enough so that the eikonal approximation intrinsic in the derivation of Ref. ceases to be valid. We have also argued that the evolution of Refs. when translated into the language of Ref. corresponds to taking into account some non-eikonal contributions.
We should note that our discussion puts into perspective the discrepancy between the double logarithmic limit (DLL) of the evolution of Refs. and the evolution suggested in Refs. . The AGL equation has been shown to arise from the BK equation in the regime where the evolution on the projectile side is dominated by production of small size dipoles, or equivalently large transverse momentum gluons . This is a natural regime when the target is a small object rather than a large object of the typical hadronic size. It was suggested in Ref. that these configurations also dominate in usual DIS in the saturation regime. It seems to us that this point warrants further study. In any case this is not the standard DLL, where the evolution on the projectile side is dominated by large dipoles. There is, therefore, no reason to expect that the DLL of Refs. has much to do with the AGL equation. In fact, as we have shown the DLL of the BK evolution is itself extremely simple when considered as the evolution of the gluon distribution operator rather than the physical DIS cross section. It turns out to be entirely devoid of nonlinear corrections and coincides with the standard DGLAP double logarithmic equation. The DIS cross section still evolves nonlinearly and in fact saturates in this limit due to the nonlinear Glauber type relation between the cross section and the gluon distribution . On the other hand, the DLL of Refs. is also nonlinear for the gluon distribution and as a result the evolution is slowed down already on the level of the gluon distribution.
We hope that this paper clarifies to some extent the relationship between the different approaches to the nonlinear low $`x`$ evolution. There are still many questions to be answered. In particular, it is desirable to find a more explicit relation between the nonlinearities of the JKLW equation and the breakdown of the eikonal approximation and to better understand the physics of these nonlinearities. Perhaps the most interesting question concerns the effect of these nonlinearities on the structure of the final states. Some work on the analytic understanding of quantities less inclusive than the total cross section has appeared recently . There is also an ongoing numerical effort in connection with the heavy ion physics in the framework of the McLerran-Venugopalan model . Further progress in this direction is extremely important both for our understanding of the nonlinear physics and for disentangling linear and nonlinear effects in the existing data.
Acknowledgments The work of A.K. is supported by PPARC. The work of J.G.M. is supported by PRAXIS XXI/BD/11277/97 grant (Subprograma Ciência e Tecnologia do 2$`^{\underset{¯}{o}}`$ Quadro Comunitário de Apoio — Portugal). He is grateful to the CERN Theory Division, J. Barbosa, J. Seixas and P. Sonderegger for hospitality during June-August 1999. H.W. was supported by the EC TMR Program, contract ERB FMRX-CT96-0008. He wants to thank the members of the Nuclear Theory Group at Brookhaven National Lab for their hospitality during October and November 1999.
## Appendix A Appendix
In this appendix, we give a more detailed derivation of the evolution equation of Sec. 2 including the calculation of the contributions where the glue does not interact with the background (target).
We start with the quadratic action for small fluctuations in the projectile gauge
$$\begin{array}{cc}\hfill S=\frac{1}{2}\{& a_a^+\left[(^{})^2\right]a_a^+2(^ia_a^+)(^{}a_a^i)\hfill \\ & a_a^i[(2D_{ab}^+[b]^{}(_{})^2\delta _{ab})\delta ^{ij}+^i^j\delta _{ab}]a_b^j\}.\hfill \end{array}$$
(83)
The equation of motion for $`a^+`$ is
$$a^+=\frac{^i}{^{}}a^i.$$
(84)
Substituting this in Eq. (83) we get
$$S=\frac{1}{2}a_a^i(D^2)_{ab}\delta ^{ij}a_b^j,$$
(85)
where
$$(D^2)_{ab}=2D_{ab}^+[b]^{}(_{})^2\delta _{ab}$$
(86)
is the same as the inverse propagator for a charged scalar field in the presence of a background field $`b^+`$. It has been calculated, for example, in Ref. . The most useful form for us is
$$\begin{array}{c}i\left[\frac{1}{(D^2)}\right]_{ab}=\frac{dp^{}}{2p^{}(2\pi )^3}[\theta (x^{}y^{})\theta (p^{})\theta (y^{}x^{})\theta (p^{})]\times \hfill \\ \hfill \times d^2p_{}d^2q_{}e^{ipx+iqy}\frac{d^2z_{}}{(2\pi )^2}e^{i(p_{}q_{})z_{}}\stackrel{~}{U}_{ab}^1(x^{},y^{},z_{})\end{array}$$
(87)
with $`p^+=\frac{p_{}^2}{2p^{}}`$, $`q^+=\frac{q_{}^2}{2p^{}}`$ and $`q^{}=p^{}`$. The color matrix $`\stackrel{~}{U}_{ab}^1(x^{},y^{},z_{})`$ is<sup>12</sup><sup>12</sup>12More rigorously, the structure of $`\stackrel{~}{U}^1`$ is given by $`\stackrel{~}{U}_{ab}^1(z_{})=e^{i(\theta (x^{})\theta (y^{}))b(z_{})}`$. However, the difference between this expression and that given in Eq.(17) only shows up if it is multiplied by $`^+`$ derivatives or $`\delta (x^{})`$ factors. Since we encounter no such factors in our calculation we will be using Eq.(17) throughout.
$$\begin{array}{c}\stackrel{~}{U}_{ab}^1(x^{},y^{},z_{})=\left(\theta (x^{})\theta (y^{})+\theta (x^{})\theta (y^{})\right)\delta _{ab}\hfill \\ \hfill +\theta (x^{})\theta (y^{})\stackrel{~}{U}_{ab}(z_{})+\theta (x^{})\theta (y^{})\stackrel{~}{U}_{ab}^{}(z_{}).\end{array}$$
(88)
The on-shell two-point correlator of $`a^+`$ can be written as
$$\begin{array}{c}a_a^+(x^+=0,x_{},x^{})a_b^+(y^+=0,y_{},y^{})=\frac{_x^i}{_x^{}}a_a^i(x^+=0,x_{},x^{})a_b^i(y^+=0,y_{},y^{})\frac{_y^i}{_y^{}}\hfill \\ \hfill =_x^i_y^i\frac{dp^{}}{p^{}}\frac{1}{2(2\pi )}\frac{1}{(p^{})^2}[\theta (x^{}y^{})\theta (p^{})\theta (y^{}x^{})\theta (p^{})]\times \\ \hfill d^2z_{}\frac{d^2p_{}}{(2\pi )^2}e^{+ip_{}(x_{}z_{})}\frac{d^2q_{}}{(2\pi )^2}e^{+iq_{}(z_{}y_{})}e^{i\frac{p_{}^{}{}_{}{}^{2}}{2p^{}}x^{}}e^{+i\frac{q_{}^{}{}_{}{}^{2}}{2p^{}}y^{}}\stackrel{~}{U}_{ab}^1(x^{},y^{},z_{}).\end{array}$$
(89)
We now need to expand the eikonal factors
$$V(x_{})=𝒫\mathrm{exp}\left[i_{\mathrm{}}^+\mathrm{}𝑑x^{}(b^++ga^+)(x_{},x^{})\right]$$
(90)
to second order in the fields. Recalling that the background part of the field $`b^+\delta (x^{})`$ and that the fluctuation field $`a^+`$ is nonsingular at $`x^{}=0`$, this becomes
$`V(x^+=0,x_{})`$ $`=𝒫\mathrm{exp}[ig{\displaystyle _{\mathrm{}}^0}dx^{}a^+(x^+=0,x_{},x^{})]\times `$
$`\times U(x_{})𝒫\mathrm{exp}\left[ig{\displaystyle _0^+\mathrm{}}𝑑x^{}a^+(x^+=0,x_{},x^{})\right],`$ (91)
where
$$U(x_{})=𝒫\mathrm{exp}\left[i_{\mathrm{}}^+\mathrm{}𝑑x^{}b^+(x^+=0,x_{},x^{})\right]$$
(92)
is the classical part of the eikonal factor.
To second order in $`a^+`$ we have
$$\begin{array}{cc}\hfill V(x^+=0,& x_{})=\hfill \\ & \{1ig_{\mathrm{}}^0dx^{}a^+(x^+=0,x_{},x^{})\hfill \\ & g^2_{\mathrm{}}^0dx^{}dy^{}\theta (y^{}x^{})a^+(x^+=0,x_{},x^{})a^+(x^+=0,x_{},y^{})\}\hfill \\ & \times U(x_{})\{1ig_0^+\mathrm{}dx^{}a^+(x^+=0,x_{},x^{})\hfill \\ & g^2_0^+\mathrm{}dx^{}dy^{}\theta (y^{}x^{})a^+(x^+=0,x_{},x^{})a^+(x^+=0,x_{},y^{})\},\hfill \end{array}$$
(93)
or
$$\begin{array}{c}V(x_{})=U(x_{})\hfill \\ \hfill ig\left\{_{\mathrm{}}^0𝑑w^{}a^+(x^+=0,x_{},w^{})U(x_{})+U(x_{})_0^+\mathrm{}𝑑w^{}a^+(x^+=0,x_{},w^{})\right\}\\ \hfill g^2\{_{\mathrm{}}^0dw^{}dz^{}\theta (z^{}w^{})a^+(x^+=0,x_{},w^{})a^+(x^+=0,x_{},z^{})U(x_{})\\ \hfill +U(x_{})_0^+\mathrm{}𝑑w^{}𝑑z^{}\theta (z^{}w^{})a^+(x^+=0,x_{},w^{})a^+(x^+=0,x_{},z^{})\\ \hfill +_{\mathrm{}}^0dw^{}a^+(x^+=0,x_{},w^{})U(x_{})_0^+\mathrm{}dz^{}a^+(x^+=0,x_{},z^{})\}.\end{array}$$
(94)
Similarly,
$$\begin{array}{c}V^{}(y_{})=U^{}(y_{})+\hfill \\ \hfill +ig\left\{U^{}(y_{})_{\mathrm{}}^0𝑑z^{}a^+(y^+=0,y_{},z^{})+_0^+\mathrm{}𝑑z^{}a^+(y^+=0,y_{},z^{})U^{}(y_{})\right\}\\ \hfill g^2\{U^{}(y_{})_{\mathrm{}}^0dw^{}dz^{}\theta (w^{}z^{})a^+(y^+=0,y_{},w^{})a^+(y^+=0,y_{},z^{})\\ \hfill +_0^+\mathrm{}𝑑w^{}𝑑z^{}\theta (w^{}z^{})a^+(y^+=0,y_{},w^{})a^+(y^+=0,y_{},z^{})U^{}(y_{})\\ \hfill +_0^+\mathrm{}dw^{}a^+(y^+=0,y_{},w^{})U^{}(y_{})_{\mathrm{}}^0dz^{}a^+(y^+=0,y_{},z^{})\}.\end{array}$$
(95)
Rather than directly calculating the eikonal cross section, we will first calculate the tensor product of two eikonal factors and later take the trace over the color indices. We use the notation $`AB=A^{\alpha \beta }B^{\gamma \delta }`$. To second order in the fluctuation field we have
$$\begin{array}{cc}& V(x_{})V^{}(y_{})U(x_{})U^{}(y_{})=\hfill \\ & =g^2_{\mathrm{}}^0dw^{}a^+(x^+=0,x_{},w^{})U(x_{})U^{}(y_{})_{\mathrm{}}^0dz^{}a^+(y^+=0,y_{},z^{})\hfill \\ & +_{\mathrm{}}^0𝑑w^{}a^+(x^+=0,x_{},w^{})U(x_{})_0^+\mathrm{}𝑑z^{}a^+(y^+=0,y_{},z^{})U^{}(y_{})\hfill \\ & +U(x_{})_0^+\mathrm{}𝑑w^{}a^+(x^+=0,x_{},w^{})U^{}(y_{})_{\mathrm{}}^0𝑑z^{}a^+(y^+=0,y_{},z^{})\hfill \\ & +U(x_{})_0^+\mathrm{}𝑑w^{}a^+(x^+=0,x_{},w^{})_0^+\mathrm{}𝑑z^{}a^+(y^+=0,y_{},z^{})U^{}(y_{})\hfill \\ & _{\mathrm{}}^0𝑑w^{}𝑑z^{}\theta (z^{}w^{})a^+(x^+=0,x_{},w^{})a^+(x^+=0,x_{},z^{})U(x_{})U^{}(y_{})\hfill \\ & _{\mathrm{}}^0𝑑w^{}a^+(x^+=0,x_{},w^{})U(x_{})_0^+\mathrm{}𝑑z^{}a^+(x^+=0,x_{},z^{})U^{}(y_{})\hfill \\ & U(x_{})_0^+\mathrm{}𝑑w^{}𝑑z^{}\theta (z^{}w^{})a^+(x^+=0,x_{},w^{})a^+(x^+=0,x_{},z^{})U^{}(y_{})\hfill \\ & U(x_{})U^{}(y_{})_{\mathrm{}}^0𝑑w^{}𝑑z^{}\theta (w^{}z^{})a^+(y^+=0,y_{},w^{})a^+(y^+=0,y_{},z^{})\hfill \\ & U(x_{})_0^+\mathrm{}𝑑w^{}a^+(y^+=0,y_{},w^{})U^{}(y_{})_{\mathrm{}}^0𝑑z^{}a^+(y^+=0,y_{},z^{})\hfill \\ & U(x_{})_0^+\mathrm{}dw^{}dz^{}\theta (w^{}z^{})a^+(y^+=0,y_{},w^{})a^+(y^+=0,y_{},z^{})U^{}(y_{}).\hfill \end{array}$$
(96)
For calculational purposes it is fruitful to separate the different contributions into the set terms where the glue interacts with the background field (the first four terms in Eq. (96) and the set of contributions where there is no such interaction (the other terms in Eq. (96)).
Let us first consider the contribution where the glue is exchanged between the quark in the negative half plane ($`x^{}<0`$) in the amplitude and an antiquark in the positive halfplane in the complex conjugate amplitude:
$$_{\mathrm{}}^0𝑑w^{}_0^+\mathrm{}𝑑z^{}a_a^+(x^+=0,x_{},w^{})a_b^+(y^+=0,y_{},z^{})\left[(t_aU(x_{}))(t_bU^{}(y_{}))\right].$$
(97)
One notices that only the $`\theta (p^{})`$ term in Eq. (89) survives and that, Eq. (88), the color matrix $`\stackrel{~}{U}_{ab}^1(w^{},z^{},z_{})`$ reduces to $`\stackrel{~}{U}_{ab}(z_{})`$.
$$\begin{array}{c}_{\mathrm{}}^0𝑑w^{}_0^+\mathrm{}𝑑z^{}a_a^+(x^+=0,x_{},w^{})a_b^+(y^+=0,y_{},z^{})=\hfill \\ \hfill =_x^i_y^i\frac{dp^{}}{p^{}}[\theta (p^{})]\frac{1}{2(2\pi )}\frac{1}{(p^{})^2}d^2z_{}\frac{d^2p_{}}{(2\pi )^2}e^{+ip_{}(x_{}z_{})}\frac{d^2q_{}}{(2\pi )^2}e^{+iq_{}(z_{}y_{})}\times \\ \hfill \times \left(_{\mathrm{}}^0𝑑w^{}e^{i\frac{p_{}^{}{}_{}{}^{2}}{2p^{}}w^{}}_0^+\mathrm{}𝑑z^{}e^{+i\frac{q_{}^{}{}_{}{}^{2}}{2p^{}}z^{}}\right)\stackrel{~}{U}_{ab}(z_{}).\end{array}$$
(98)
The $`w^{}`$ and $`z^{}`$ integrations are easily performed
$$_{\mathrm{}}^0𝑑w^{}e^{i\frac{p_{}^{}{}_{}{}^{2}}{2p^{}}w^{}}_0^+\mathrm{}𝑑z^{}e^{+i\frac{q_{}^{}{}_{}{}^{2}}{2p^{}}z^{}}=\frac{(2p^{})^2}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}.$$
(99)
Noting that
$$\frac{dp^{}}{p^{}}\theta (p^{})=_{\mathrm{}}^0\frac{dp^{}}{p^{}}=_0^+\mathrm{}\frac{dp^{}}{p^{}}=\frac{dp^{}}{p^{}}\theta (p^{}),$$
(100)
we have
$$\frac{1}{\pi }_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}_x^i_y^i\frac{d^2p_{}}{(2\pi )^2}\frac{1}{p_{}^{}{}_{}{}^{2}}e^{+ip_{}(x_{}z_{})}\frac{d^2q_{}}{(2\pi )^2}\frac{1}{q_{}^{}{}_{}{}^{2}}e^{+iq_{}(z_{}y_{})}\stackrel{~}{U}_{ab}(z_{}).$$
(101)
The transverse momenta integrals yield
$$\frac{d^2p_{}}{(2\pi )^2}\frac{1}{p_{}^{}{}_{}{}^{2}}e^{+ip_{}(x_{}z_{})}=\frac{1}{4\pi }\mathrm{log}(x_{}z_{})^2.$$
(102)
Taking the derivatives
$$\frac{1}{4\pi ^3}_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}\stackrel{~}{U}_{ab}(z_{}).$$
(103)
We use the following identity valid for $`SU(N)`$ group
$`\stackrel{~}{U}_{ab}(z_{})(t_aU(x_{}))(t_bU^{}(y_{}))=`$
$`=2\mathrm{t}\mathrm{r}\left[t_aU(z_{})t_bU^{}(z_{})\right](t_aU(x_{}))^{\alpha \beta }(t_bU^{}(y_{}))^{\gamma \delta }`$
$`={\displaystyle \frac{1}{2N_c}}\left[N_c\left(U(z_{})U^{}(y_{})\right)^{\alpha \delta }\left(U^{}(z_{})U(x_{})\right)^{\gamma \beta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }\right].`$ (104)
This contribution, the second term in Eq. (96), therefore is
$$\begin{array}{c}\frac{1}{4\pi ^3}_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}\times \hfill \\ \hfill \times \frac{1}{2N_c}\left[N_c\left(U(z_{})U^{}(y_{})\right)^{\alpha \delta }\left(U^{}(z_{})U(x_{})\right)^{\gamma \beta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }\right].\end{array}$$
(105)
The third contribution, due to the exchange of the gluon between the quark in the positive half plane and the antiquark in the negative halfplane, is calculated similarly with the only difference that now $`\stackrel{~}{U}_{ab}^1(w^{},z^{},z_{})=\stackrel{~}{U}_{ab}^{}(z_{})`$ and we pick up the $`\theta (p^{})`$ term in the propagator.
$`{\displaystyle _0^+\mathrm{}}𝑑w^{}{\displaystyle _{\mathrm{}}^0}𝑑z^{}a_a^+(x^+=0,x_{},w^{})a_b^+(y^+=0,y_{},z^{})\left[(U(x_{})t_a)(U^{}(y_{})t_b)\right]=`$
$`={\displaystyle \frac{1}{4\pi ^3}}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dp^{}}{p^{}}}{\displaystyle }d^2z_{}{\displaystyle \frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}}\times `$
$`\times {\displaystyle \frac{1}{2N_c}}\left[N_c\left(U(x_{})U^{}(z_{})\right)^{\alpha \delta }\left(U^{}(y_{})U(z_{})\right)^{\gamma \beta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }\right].`$ (106)
For the quark self energy correction everything is the same as in Eq. (105) except that the transverse coordinates of the two fields coincide ($`y_{}\stackrel{~}{x}_{}=x_{}`$)
$`{\displaystyle _{\mathrm{}}^0}𝑑w^{}{\displaystyle _0^+\mathrm{}}𝑑z^{}a_a^+(x^+=0,x_{},w^{})a_b^+(y^+=0,\stackrel{~}{x}_{}=x_{},z^{})\left[(t_aU(x_{})t_b)U^{}(y_{})\right]=`$
$`={\displaystyle \frac{1}{4\pi ^3}}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dp^{}}{p^{}}}{\displaystyle }d^2z_{}{\displaystyle \frac{1}{(xz)_{}^2}}\times `$
$`\times {\displaystyle \frac{1}{2N_c}}[N_c\mathrm{tr}((U(x_{})U^{}(z_{}))\left(U(z_{})\right)^{\alpha \beta }\left(U(x_{})\right)^{\alpha \beta }]U^{}(y_{})^{\gamma \delta }.`$ (107)
Finally, the antiquark self energy correction
$`{\displaystyle _0^+\mathrm{}}𝑑w^{}{\displaystyle _{\mathrm{}}^0}𝑑z^{}a_a^+(x^+=0,y_{},w^{})a_b^+(y^+=0,\stackrel{~}{y}_{}=y_{},z^{})\left[U(x_{})(t_at_bU^{}(y_{}))\right]=`$
$`={\displaystyle \frac{1}{4\pi ^3}}{\displaystyle _{\mathrm{}}^0}{\displaystyle \frac{dp^{}}{p^{}}}{\displaystyle }d^2z_{}{\displaystyle \frac{1}{(yz)_{}^2}}\times `$
$`\times {\displaystyle \frac{1}{2N_c}}\left[N_c\mathrm{tr}\left(U^{}(y_{})U(z_{})\right)\left(U^{}(z_{})\right)^{\gamma \delta }\left(U^{}(y_{})\right)^{\gamma \delta }\right]U(x_{})^{\alpha \beta }.`$ (108)
Now let us look at the other set of contributions. Since for all of these there is no interaction with the background we have, Eq. (88), $`\stackrel{~}{U}_{ab}^1(w^{},z^{},z_{})=\delta _{ab}`$. The $`z_{}`$ integral is trivial for all terms in this set (it yields a $`\delta `$ function for the transverse momenta). However, as we want to combine both sets of contributions in the final result, this integral will not be performed.
Consider the first term in Eq. (96) — the exchange between the quark and the antiquark in the negative half-plane —
$$_{\mathrm{}}^0𝑑w^{}𝑑z^{}a_a^+(x^+=0,y_{},w^{})a_b^+(y^+=0,\stackrel{~}{y}_{}=y_{},z^{})\left[(t_aU(x_{}))(U^{}(y_{})t_b)\right].$$
(109)
In the correlator part now both $`\theta `$ functions survive, corresponding to the two possible orderings of $`w^{}`$ and $`z^{}`$
$$\begin{array}{c}_{\mathrm{}}^0𝑑w^{}𝑑z^{}a_a^+(x^+=0,x_{},w^{})a_b^+(y^+=0,y_{},z^{})=\hfill \\ \hfill =_x^i_y^i\frac{dp^{}}{p^{}}\frac{1}{2(2\pi )}\frac{1}{(p^{})^2}d^2z_{}\frac{d^2p_{}}{(2\pi )^2}e^{+ip_{}(x_{}z_{})}\frac{d^2q_{}}{(2\pi )^2}e^{+iq_{}(z_{}y_{})}\times \\ \hfill \times \left(_{\mathrm{}}^0𝑑w^{}𝑑z^{}\left[\theta (w^{}z^{})\theta (p^{})\theta (z^{}w^{})\theta (p^{})\right]e^{i\frac{p_{}^2}{2p^{}}w^{}}e^{+i\frac{q_{}^{}{}_{}{}^{2}}{2p^{}}z^{}}\right)\delta _{ab}.\end{array}$$
(110)
The $`w^{}`$ and $`z^{}`$ integrals now have to be performed with more care. Doing this we get
$$\begin{array}{c}\frac{1}{\pi }_x^i_y^i\frac{dp^{}}{p^{}}d^2z_{}\frac{d^2p_{}}{(2\pi )^2}e^{+ip_{}(x_{}z_{})}\frac{d^2q_{}}{(2\pi )^2}e^{+iq_{}(z_{}y_{})}\times \hfill \\ \hfill \times \left[\theta (p^{})\left(\frac{1}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}\frac{1}{p_{}^{}{}_{}{}^{2}(q_{}^{}{}_{}{}^{2}p_{}^{}{}_{}{}^{2})}\right)+\theta (p^{})\left(\frac{1}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}\frac{1}{q_{}^{}{}_{}{}^{2}(q_{}^{}{}_{}{}^{2}p_{}^{}{}_{}{}^{2})}\right)\right]\delta _{ab}.\end{array}$$
(111)
Recalling Eq. (100) we obtain for the momentum denominators
$`{\displaystyle \frac{dp^{}}{p^{}}\left[\theta (p^{})\left(\frac{1}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}\frac{1}{p_{}^{}{}_{}{}^{2}(q_{}^{}{}_{}{}^{2}p_{}^{}{}_{}{}^{2})}\right)+\theta (p^{})\left(\frac{1}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}\frac{1}{q_{}^{}{}_{}{}^{2}(q_{}^{}{}_{}{}^{2}p_{}^{}{}_{}{}^{2})}\right)\right]}=`$
$`={\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dp^{}}{p^{}}}[2{\displaystyle \frac{1}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}}+({\displaystyle \frac{1}{q_{}^{}{}_{}{}^{2}(q_{}^{}{}_{}{}^{2}p_{}^{}{}_{}{}^{2})}}{\displaystyle \frac{1}{p_{}^{}{}_{}{}^{2}(q_{}^{}{}_{}{}^{2}p_{}^{}{}_{}{}^{2})}})`$
$`={\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dp^{}}{p^{}}}{\displaystyle \frac{1}{p_{}^{}{}_{}{}^{2}q_{}^{}{}_{}{}^{2}}}.`$ (112)
The correlator part is then
$$\begin{array}{c}_{\mathrm{}}^0𝑑w^{}𝑑z^{}a_a^+(x^+=0,x_{},w^{})a_b^+(y^+=0,y_{},z^{})=\hfill \\ \hfill =\frac{1}{4\pi ^3}_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}\delta _{ab}.\end{array}$$
(113)
The color algebra raises no problems
$`\delta _{ab}(t_aU(x_{}))(U^{}(y_{})t_b)=`$ (114)
$`=\delta _{ab}(t_aU(x_{}))^{\alpha \beta }(U^{}(y_{})t_b)^{\gamma \delta }`$
$`={\displaystyle \frac{1}{2N_c}}\left[N_c\delta ^{\alpha \delta }\left(U^{}(y_{})U(x_{})\right)^{\gamma \beta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }\right]`$
and so the first term in Eq. (96) is
$`{\displaystyle \frac{1}{4\pi ^3}}{\displaystyle _0^+\mathrm{}}{\displaystyle \frac{dp^{}}{p^{}}}{\displaystyle d^2z_{}}`$ $`{\displaystyle \frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}}\times `$
$`\times {\displaystyle \frac{1}{2N_c}}\left[N_c\delta ^{\alpha \delta }\left(U^{}(y_{})U(x_{})\right)^{\gamma \beta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }\right].`$ (115)
The quark to antiquark exchange in the positive half-plane gives
$$\begin{array}{c}_0^+\mathrm{}𝑑w^{}𝑑z^{}a_a^+(x^+=0,y_{},w^{})a_b^+(y^+=0,\stackrel{~}{y}_{}y_{},z^{})\left[(U(x_{})t_a)(t_bU^{}(y_{}))\right]=\hfill \\ \hfill =\frac{1}{4\pi ^3}\underset{0}{\overset{+\mathrm{}}{}}\frac{dp^{}}{p^{}}d^2z_{}\frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}\times \\ \hfill \times \frac{1}{2N_c}\left[N_c\delta ^{jk}\left(U(x_{})U^{}(y_{})\right)^{\alpha \delta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }\right].\end{array}$$
(116)
Now we combine the two terms that give corrections to the quark line — the fifth and the seventh terms in Eq. (96). It is easy to see that they have the same color structure and will also yield the same transverse structure.
The color algebra is trivial
$$\delta _{ab}(t_at_bU(x_{}))U^{}(y_{})=\delta _{ab}(t_at_bU(x_{}))^{\alpha \beta }(U^{}(y_{}))^{\gamma \delta }=\frac{(N_c^21)}{2N_c}U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }.$$
(117)
The fifth term is the same as Eq. (115), but with $`y_{}\stackrel{~}{x}_{}=x_{}`$ and only the $`\theta (p^{})`$ term surviving
$$\begin{array}{c}_{\mathrm{}}^0𝑑w^{}𝑑z^{}\theta (z^{}w^{})a_a^+(x^+=0,x_{},w^{})a_b^+(x^+=0,x_{},z^{})\delta _{ab}(t_at_bU(x_{}))U^{}(y_{})=\hfill \\ \hfill =\frac{1}{\pi }_x^i_y^i_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{d^2p_{}}{(2\pi )^2}e^{+ip_{}(x_{}z_{})}\frac{d^2q_{}}{(2\pi )^2}e^{+iq_{}(z_{}\stackrel{~}{x}_{})}\times \\ \hfill \times \left[\frac{1}{p_{}^2q_{}^2}+\frac{1}{q_{}^2(q_{}^2p_{}^2)}\right]\frac{(N_c^21)}{2N_c}U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }.\end{array}$$
(118)
And the seventh term is the same as Eq. (116)
$$\begin{array}{c}_0^+\mathrm{}𝑑w^{}𝑑z^{}\theta (z^{}w^{})a_a^+(x^+=0,x_{},w^{})a_b^+(x^+=0,x_{},z^{})\delta _{ab}(U(x_{})t_at_b)U^{}(y_{})=\hfill \\ \hfill =\frac{1}{\pi }_x^i_y^i_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{d^2p_{}}{(2\pi )^2}e^{+ip_{}(x_{}z_{})}\frac{d^2q_{}}{(2\pi )^2}e^{+iq_{}(z_{}\stackrel{~}{x}_{})}\times \\ \hfill \times \left[\frac{1}{p_{}^2q_{}^2}\frac{1}{p_{}^2(q_{}^2p_{}^2)}\right]\frac{(N_c^21)}{2N_c}U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }.\end{array}$$
(119)
Adding the two terms Eqs. (118,119) and the performing the transverse integrations we get
$$\frac{1}{4\pi ^3}_0^+\mathrm{}\frac{dp^{}}{p^{}}d^2z_{}\frac{1}{(xz)_{}^2}\frac{(N_c^21)}{2N_c}U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }.$$
(120)
The correction to the antiquark line — the eighth and the tenth terms in Eq. (96) give similarly
$$\frac{1}{4\pi ^3}_0^{\mathrm{}}\frac{dp^{}}{p^{}}d^2z_{}\frac{1}{(yz)_{}^2}\frac{(N_c^21)}{2N_c}U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }.$$
(121)
Finally, combining all the terms together we get
$`V(x_{})^{\alpha \beta }V^{}(y_{})^{\gamma \delta }U(x_{})^{\alpha \beta }U^{}(y_{})^{\gamma \delta }=`$
$`={\displaystyle \frac{1}{8\pi ^3}}\mathrm{log}({\displaystyle \frac{x_0}{x}}){\displaystyle }d^2z_{}`$
$`\{[(U(z_{})U^{}(y_{}))^{\alpha \delta }(U^{}(z_{})U(x_{}))^{\gamma \beta }+(U(x_{})U^{}(z_{}))^{\alpha \delta }(U^{}(y_{})U(z_{}))^{\gamma \beta }`$
$`\delta ^{\alpha \delta }(U^{}(y_{})U(x_{}))^{\gamma \beta }(U(x_{})U^{}(y_{}))^{\alpha \delta }\delta ^{\beta \gamma }]{\displaystyle \frac{(xz)_{}(yz)_{}}{(xz)_{}^2(yz)_{}^2}}`$
$`\left[\mathrm{tr}\left(U(x_{})U^{}(z_{})\right)\left(U(z_{})\right)^{\alpha \beta }N_cU(x_{})^{\alpha \beta }\right]U^{}(y_{})^{\gamma \delta }{\displaystyle \frac{1}{(xz)_{}^2}}`$
$`U(x_{})^{\alpha \beta }[\mathrm{tr}(U^{}(y_{})U(z_{}))\left(U^{}(z_{})\right)^{\gamma \delta }N_cU^{}(y_{})^{\gamma \delta }]{\displaystyle \frac{1}{(zy)_{}^2}}\}.`$ (122)
This coincides with the result of Ref. . When comparing this evolution equation with the results of Ref. , one should keep in mind that there the evolution is considered with respect to the variable $`\zeta `$. The relation between the two evolution equations is given by $`\frac{d}{d\mathrm{ln}\frac{1}{x}}\mathrm{}=2\frac{d}{d\mathrm{ln}\zeta }\mathrm{}`$.
Now taking trace over the color indices we obtain the evolution equation for the scattering cross section given in Sec. 2.
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# Quantum Background Independence and Witten Geometric Quantization of the Moduli of CY Threefolds.
## 1 Introduction
By definition a Calabi-Yau (CY) manifold is a compact complex $`n`$dimensional Kähler manifold M with a holomorphic $`n`$form $`\mathrm{\Omega }_\text{M}`$ which has no zeroes and $`H^0(\text{M},\mathrm{\Omega }_\text{M}^k)=0`$ for $`0<k<n`$. Calabi-Yau manifolds are playing important role in string theory. The powerful ideas from string theory played a very important role in the recent developments in some branches of mathematics and especially in the study of moduli of CY manifolds. In this paper we will study moduli space of CY threefolds based on the ideas introduced in , and .
In and it was proved that there are no obstructions to the deformations of the complex structures on CY manifolds. This means that the local moduli space of CY manifolds is smooth of dimension
$$h^{n1,1}=dim_{}H^1(\text{M},\mathrm{\Omega }_\text{M}^{n1}).$$
From the theory of moduli of polarized algebraic manifolds developed by Viehweg in it follows that the moduli space of polarized CY manifolds is a quasi-projective variety.
The moduli space $``$ of three dimensional CY manifolds has a very rich structure. According to the theory of variations of Hodge structures there exists is a well defined map from the moduli space of polarized CY manifolds to $`\left(H^n(\text{M},)\right)`$ which is called the period map. It assigns to each point $`\tau `$ of the moduli space the line in $`H^n(\text{M},))`$ spanned by the cohomology class represented by the non-zero holomorphic n-form. According to local Torelli Theorem the period map is a local isomorphism. Local Torelli Theorem implies that locally the moduli space of CY manifolds can be embedded in $`\left(H^n(\text{M},)\right).`$ When the dimension $`n`$ of the CY manifold is odd, Griffiths and Bryant noticed that the intersection form on $`H^n(`$M,$`)`$ defines on the standard charts $`U_i=^{n1}`$ of $`\left(H^n(\text{M},)\right)`$ holomorphic one forms $`\alpha _i`$ such that $`d\alpha _i`$ is a skew symmetric form of maximal rank on $`^{n1}.`$ This means that on $`U_i=^{n1}`$ a natural contact structure is defined. In the case of three dimensional CY manifolds Griffiths and Bryant proved that the restrictions of $`d\alpha _i`$ on the tangent space of the image of the local moduli space of CY manifolds is zero. Thus the image of the local moduli space is a Legandre submanifold. See . Arnold described the local structure of the Legandre submanifolds in a contact manifold in . This description implies the existence of a generating holomorphic function for the local moduli space.
Based on the work A. Strominger noticed that the potential of the Weil-Petersson metric on the local moduli space can be expressed through the generating holomorphic function. See . By using this observation, Strominger introduced the notion of special Kähler geometry. V. Cortes showed that on the tangent bundle of the special Kähler manifold one can introduce a Hyper-Kähler structure. See and . From here it follows that on the tangent bundle of $`\left(\text{M}\right)`$ one can introduce a Hyper-Kähler structure. Earlier R. Donagi and E. Markman constructed in an analytically completely integrable Hamiltonian system which is canonically associated with the family of CY manifolds over the relative dualizing line bundle over the moduli space $`\left(\text{M}\right)`$. They showed that the space of the Griffiths intermediate Jacobians, associated with the family of three dimensional CY manifolds on $``$ carries a Hyper-Kähler structure. B. Dubrovin introduced the notion of Frobenius manifolds in . The relations of the structure of Frobenius manifolds and Gromov-Witten invariants were studied by Yu. I. Manin and M. Kontsevich in .
The importance of all these structures is justified by the work of Candelas and coauthors in their seminal paper . In this paper Candelas and his coauthors gave an explicit formula for the number of rational curves on the quintic hypersurface in the four dimensional projective space. M. Kontsevich defined the correct compactification of the stable maps and realized that one can use the localization formula for the computations of the rational curves. See . Recently B. Lian, K. Liu and Yau gave a rigorous mathematical proof of the Candelas formula in . See also the important paper of Givental .
In this paper we study two different topics. The first topic is the applications of the geometric quantization scheme of Witten introduced in to the problem of the quantum background independence in string theory. The second topic is the introduction of a $``$ structure on the tangent space of the moduli space of polarized CY threefolds $`\left(\text{M}\right)`$ and thus we associate an algebraic integrable structure on the tangent bundle of $`\left(\text{M}\right)`$. For both topics it is crucial to construct a flat $`𝕊p(2h^{2,1},)`$ connection on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ of polarized CY threefolds.
The problem of the quantum background independence was addressed in . In the paper Witten wrote:
”Finding the right framework for intrinsic, background independent formulation of string theory is one of the main problems in the subject, and so far has remained out of reach…”
In fact in a program was outlined how one can solve the problem of the background independence in the topological field theory:
”Though the interpretation of the holomorphic anomaly as an obstruction to background independence eliminates some thorny puzzles, it is not satisfactory to simply leave matters as this. Is there some sophisticated sense in which background independence does hold? In thinking about this question, it is natural to examine the all orders generalization of the holomorphic anomaly equation, which in the final equation of their paper ) Besrshadsky et. al. write the following form. Let $`F_g`$ be the genus $`g`$ free energy. Then
$$\overline{_i^{^{}}}F_g=\overline{C}_{i^{^{}}j^{^{}}k^{^{}}}e^{2K}G^{jj^{^{}}}G^{kk^{^{}}}\left(D_jD_kF_{g1}+\frac{1}{2}\underset{r}{}D_jF_rD_kF_{gr}\right).$$
(1)
This equation can be written as a linear equation for
$$Z=\mathrm{exp}\left(\frac{1}{2}\underset{g=0}{\overset{\mathrm{}}{}}\lambda ^{2g2}F_g\right),$$
(2)
namely
$$\left(\overline{_i^{^{}}}\lambda ^2\overline{C}_{i^{^{}}j^{^{}}k^{^{}}}e^{2K}G^{jj^{^{}}}G^{kk^{^{}}}D_jD_k\right)Z=0$$
(3)
This linear equation is called a master equation by Bershadsky et. al.; it is similar to the structure of the heat equations obeyed by theta functions…
It would be nice to interpret $`\left(\text{3}\right)`$ as a statement of some sophisticated version of background independence. In thinking about this equation, a natural analogy arises with Chern-Simon gauge theory in $`2+1`$ dimensions. In this theory, an initial value surface is a Riemann surface $`\mathrm{\Sigma }.`$ In the Hamiltonian formulation of the theory, one construct a Hilbert space $``$ upon quantization on $`\mathrm{\Sigma }.`$ should be obtained from some physical space $`𝐖`$ (a moduli space of flat connections on $`\mathrm{\Sigma }`$). Because the underlying Chern-Simon Lagrangian does not depend on the choice of the metric, one would like to construct $``$ in a natural, background independent way. In practice, however, quantization of $`𝐖`$ requires a choice of polarization, and there is no natural way or background independent choice of polarization.
The best that one can do is to pick a complex structure $`J`$ on $`\mathrm{\Sigma },`$ whereupon $`𝐖`$ gets a complex structure. Then a Hilbert space $`_J`$ is constructed as a suitable space of holomorphic functions (really sections of a line bundle) over $`𝐖`$. We denote such function as $`\psi (a^i,t^a)`$ where $`a^i`$ are complex coordinates on $`𝐖`$ and $`t^a`$ are coordinates parametrizing the choice of $`J.`$ Now background independence does not hold in a naive sense; $`\psi `$ can not be independent of $`t^^i`$ (given that it is to be holomorphic on $`𝐖`$ in a complex structure dependent on $`t^a).`$ But there is a more sophisticated sense in which background independence can be formulated. See and . The $`_J`$ can be identified with each other (projectively) using a (projectively) flat connection over the space of $`J^{}`$s. This connection $``$ is such that a covariant constant wave function should have the following property: as $`J`$ changes, $`\psi `$ should change by Bogoliubov transformation, representing the effect of a change in the representation used for the canonical commutation relations. Using parallel transport by $``$ to identify the various $`_J^{}`$ s are realizations determined by a $`J`$-dependent choice of the representation of the canonical commutators. Background independence of $`\psi (a^i,t^a)`$ should be interpreted to mean that the quantum state represented by $`\psi `$ is independent of $`t^a,`$ or equivalently that $`\psi `$ is invariant under parallel transport by $`.`$ Concretely, this can be written as an equation:
$$\left(\frac{}{t^a}\frac{1}{4}\left(\frac{J}{t^a}\omega ^1\right)^{ij}\frac{D}{Da^i}\frac{D}{Da^j}\right)\psi =0.$$
(4)
that is analogous to the heat equation for theta functions…”
In the above program is realized on the space $`𝐖=H^3(\text{M},)`$. Bershadsky, Cecotti, Ooguri and Vafa work on $`H^1(`$M,$`T_\text{M}^{1,0}),`$ i.e. the tangent space to the moduli of CY. The space $`𝐖=H^3(\text{M},)`$ has a natural symplectic form structure given by the intersection pairing
$$\omega (\alpha ,\beta ):=\underset{\text{M}}{}\alpha \beta .$$
The complex structure on M defines a complex structure on $`H^3(\text{M},).`$ On the vector bundle $`R^3\pi _{}`$ over the moduli space with a fibre $`𝐖=H^3(\text{M},)`$ we have a natural flat $`𝕊p(2h^{2,1}+2,)`$ connection. The tangent bundle to $`\omega _{𝒳\text{/}\text{ }\left(\text{M}\right)\text{ }}`$ is naturally isomorphic to $`\pi ^{}\left(R^3\pi _{}\right)`$. Thus on it we have a natural flat $`𝕊p(2h^{2,1}+2,)`$ connection.
In the present paper the program of Witten is realized for the tangent bundle of the moduli space of polarized CY threefolds. One of the most important ingredient in the realization of the Witten program is the construction of a flat $`𝕊p(2h^{2,1},)`$ connection on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ of polarized CY threefolds. In the present article we constructed such flat $`𝕊p(2h^{2,1},)`$ connection.
The idea of the construction of the flat $`𝕊p(2h^{2,1},)`$ connection on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ is to modify the unitary connection of the Weil-Petersson metric on $`\left(\text{M}\right)`$ with a Higgs field to a $`𝕊p(2h^{2,1},)`$ connection and then prove that the $`𝕊p(2h^{2,1},)`$ connection is flat. The construction of the Higgs field defined on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ is done by using the cup product $`\varphi _1\varphi _2`$ $`H^2(`$M,$`^2T^{1,0})`$ for $`\varphi _iH^1(`$M,$`T^{1,0}),`$ the identifications of $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$ with $`H^1(`$M,$`T^{1,0})`$, $`H^2(\text{M},\mathrm{\Omega }_\text{M}^1)`$ with $`H^2(`$M,$`^2T^{1,0})`$ and the identification of $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$ with $`H^2(\text{M},\mathrm{\Omega }_\text{M}^1)`$ by the Poincare duality.
The construction of a flat $`𝕊p(2h^{2,1},)`$ connection on the tangent bundle of $`\left(\text{M}\right)`$ is related to the important example of a special Hyper-Kähler manifold which occurs in four dimensional gauge theories with $`N=2`$ supersymmetry: the scalars in the vector multiplet lie in a special Kähler manifolds. The moduli space of such theories was studied by Cecotti and Vafa in . They introduced the tt\* equations. One of the observation in this paper is that the analogue of the tt\* equations in case of the moduli of polarized CY threefolds is the same as the Yang-Mills equations coupled with Higgs fields that were studied by Hitchin in case of Riemann surfaces in and by C. Simpson in general in .
The flat $`𝕊p(2h^{2,1},)`$ connection is crucial to apply the geometric quantization method of Witten to the tangent space of the moduli space $`\left(\text{M}\right)`$ of polarized CY threefolds to solve the problem of the ground quantum independence in the topological field theory. On the basis of the geometric quantization of the tangent bundle of $`\left(\text{M}\right)`$ we are able to modify the beautiful computations of E. Witten in to obtain a projective connection on some infinite dimensional Hilbert space bundle. We prove that the holomorphic anomaly equations $`\left(\text{1}\right)`$ of Bershadsky, Cecotti, Ooguri and Vafa imply that the free energy obtained from the ”counting functions” F<sub>g</sub> of curves of genus g on a CY manifold M is a parallel section of a projective flat connection. Our computations are based on the technique developed in .
The projective connection constructed in is different from ours since we work on different spaces. The difference appeared in the computation of the formula for $`\left(dJ\omega ^1\right).`$ On the space $`𝐖=H^3(\text{M},)`$ Witten obtained that
$$\left(dJ\right)_a^{\overline{b}}=2\underset{c,d}{}\overline{C}_{acd}g^{d,\overline{b}}$$
where $`\left(g_{a,\overline{b}}\right)`$ defines the symplectic structure on $`𝐖`$ coming from the cup product and $`C_{acd}`$ is the Yukawa coupling. Our formula on $`𝐖=H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$ is
$$\left(dJ\right)_a^{\overline{b}}=\underset{c,d}{}\overline{C}_{acd}g^{d,\overline{b}},$$
where $`\left(g_{a,\overline{b}}\right)`$ is the symplectic form obtained from the restriction of the cup product on $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2)H^3(\text{M},).`$ At the end we obtained exactly the formula $`\left(\text{4}\right)`$ suggested by E. Witten.
In two equations are derived. One of them is $`\left(\text{1}\right)`$. It gives a recurrent relation between $`F_g`$ ’s. The other equation in is $`\left(\text{98}\right).`$ These two equations are marked as $`(3.6)`$ and $`(3.8)`$ in . According to the free energy $`Z`$ satisfy the equation$`\left(\text{98}\right)`$. One can notice that there is a difference between our equation and the equation $`\left(\text{98}\right)`$ for the free energy $`Z`$ in . The holomorphic anomaly equation $`\left(\text{98}\right)`$ in involves the term $`F_1`$ while ours do not.
It was pointed out in that the anomaly equations are the analogue of the heat equations for the classical theta functions. Thus they are of second order. From here one can deduce that if we know the functions F<sub>0</sub> and F<sub>1</sub> that count the rational and elliptic curves on M we will know the functions F<sub>g</sub> that count all curves of given genus $`g>1.`$ It was Welters who first noticed that the heat equation of theta functions can be interpreted as a projective connection. See . Later N. Hitchin used the results of to construct a projectively flat connection on a vector bundle over the Teichmüller space constructed from the symmetric tensors of stable bundle over a Riemann surface. See . For other useful applications of the geometric approach to quantization see .
The second problem discussed in this paper is about the existence of $``$ structure on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ of polarized CY threefolds. This problem is suggested by the mirror symmetry conjecture since it suggests that
$$H^{3,0}H^{2,1}H^{1,2}H^{0,3}$$
can be ”identified” on the mirror side with
$$H^0H^2H^4H^6.$$
Thus since by the mirror conjecture $`H^{2,1}`$ can be identified with $`H^2`$ one should expect some natural $``$ structure on $`H^{2,1}`$ invariant under the flat $`𝕊p(2h^{2,1},)`$ connection. Thus we need to define at a fixed point of the moduli space $`\tau _0\left(\text{M}\right)`$ a $``$ structure on the tangent space $`T_{\tau _0,\left(\text{M}\right)}=H^1(`$M$`{}_{\tau _0}{}^{},\mathrm{\Omega }_{\text{M}_{\tau _0}}^2)`$. One way to obtain a natural $``$ structure on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ is the following one. Suppose that there exists a point $`\tau _0\left(\text{M}\right)`$ such that
$$H^{3,0}\left(\text{M}_{\tau _0}\right)H^{0,3}\left(\text{M}_{\tau _0}\right)=\mathrm{\Lambda }_0,$$
(5)
where $`\mathrm{\Lambda }_0`$ is a rank two sublattice in $`H^3(`$M$`{}_{\tau _0}{}^{},).`$ Once we construct such $``$ structure on $`T_{\tau _0,\left(\text{M}\right)},`$ we can use the parallel transport to define a $``$ structure on each tangent space of $`\left(\text{M}\right)`$. Unfortunately the existence of points that $`\left(\text{5}\right)`$ is satisfied is a very rare phenomenon for CY manifolds. There is a conjecture due to B. Mazur and Y. Andre which states that if the moduli space of CY manifold is a Shimura variety then such points are everywhere dense subset. The moduli space of polarized CY manifolds that are not locally symmetric spaces probably will not contain everywhere dense subset of points that correspond to CY manifolds for which
$$H^{3,0}\left(\text{M}_{\tau _0}\right)\overline{H^{3,0}\left(\text{M}_{\tau _0}\right)}=\mathrm{\Lambda }_0,$$
where $`\mathrm{\Lambda }_0`$ is a rank two sublattice in $`H^3(\text{M}_{\tau _0},).`$
The idea of the introduction of the $``$ structure on the tangent space of the moduli space $`\left(\text{M}\right)`$ is to consider the deformation space of M$`\times \overline{\text{M}}.`$ We relate the local deformation space on M$`\times \overline{\text{M}}`$ to the variation of Hodge structure of weight two with $`p_g=1`$. Such Variations of Hodge structures for the products M$`\times \overline{\text{M}}`$ are defined by the two dimensional real subspace $`H^{3,0}\left(\text{M}\right)\overline{H^{3,0}\left(\text{M}\right)}`$ in $`H^3(\text{M},)`$ is generated by $`\mathrm{Re}\mathrm{\Omega }_\tau `$ and $`\mathrm{Im}\mathrm{\Omega }_\tau `$ and they are parametrized by the symmetric space
$$𝕊𝕆_0(2,2h^{2,1})/𝕊𝕆(2)\times 𝕊𝕆(2h^{2,1})$$
where the set of points for which $`\left(\text{5}\right)`$ holds is an everywhere dense subset. Thus we are in situation similar to the moduli of algebraic polarized K3 surfaces. For the $`(\tau ,\overline{\nu })`$ in the local moduli space of M$`\times \overline{\text{M}}`$ that corresponds to M$`{}_{\tau }{}^{}\times \overline{\text{M}_\nu }`$ we construct a Hodge structure of weight two
$$H_{\tau ,\nu }^{2,0}H_{\tau ,\nu }^{1,1}\overline{H_{\tau ,\nu }^{2,0}}$$
where $`H_{\tau ,\nu }^{2,0}\overline{H_{\tau ,\nu }^{2,0}}`$ is the two dimensional subspace in $`H^3(\text{M},)`$ generated by
$$\mathrm{Re}\left(\mathrm{\Omega }_{\tau _1}+\overline{\mathrm{\Omega }_{\tau _2}}\right)\text{ and }\mathrm{Im}\left(\mathrm{\Omega }_{\tau _2}\overline{\mathrm{\Omega }_{\tau _2}}\right).$$
It is not difficult to show that the points $`(\tau ,\nu )`$ in the local moduli space of M$`\times \overline{\text{M}}`$ such that
$$H_{\tau ,\nu }^{2,0}\overline{H_{\tau ,\nu }^{2,0}}=\mathrm{\Lambda }_1,$$
where $`\mathrm{\Lambda }_1`$ is a rank four sublattice in $`H^3(\text{M},)`$ is an everywhere dense subset. Each point $`(\tau ,\nu )`$ of this everywhere dense subset defines a natural $``$ structure on $`H_{\tau ,\nu }^{1,1}.`$ Then by using the flat $`𝕊p(2h^{2,1},)`$ connection on the tangent space of $`\left(\text{M}\right)`$ we define a $``$ structure on $`H_{\tau ,\tau }^{2,0}`$ and thus on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2.`$ By using this $``$ structure we introduce an algebraic integrable structure on the tangent bundle of $`\left(\text{M}\right)`$. In the authors introduced algebraic integrable structure on the tangent bundle of the relative dualizing line bundle of $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ of the moduli space $`\left(\text{M}\right)`$ of polarized CY threefolds.
All the results in the Sections 3, 4, 5, 6 and 7 are new. Next we will describe the ideas and the content of each section.
In Section 2 we review the results of and . The Teichmüller space of the CY manifolds is constructed too.
In Section 3 we show that the analogue of the tt\* equations on $`\left(\text{M}\right)`$ are the same self dual equations that were studied by N. Hitchin and C. Simpson’s in and . Thus tt\* equations define a flat $`𝕊p(2h^{2,1},)`$ connection on the bundle $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)\text{ }}^2.`$ On the other hand we know that the tangent bundle $`𝒯_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ is isomorphic to $`𝔏^{}`$ $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$, where $`𝔏`$ is isomorphic to $`\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^3`$. We constructed by using the theory of determinant bundles a holomorphic non-vanishing section $`\eta _\tau \mathrm{\Gamma }((\text{M}),(𝔏))`$ in . Thus $`\eta _\tau `$ defines a flat structure on the tangent bundle $`𝒯_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ of the moduli space of three dimensional CY manifolds $`\left(\text{M}\right)`$. Using the flat structure defined by tt\* equations on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)\text{ }}^2`$ and the flat structure defined by the section $`\eta _\tau `$ on $`𝔏`$, we define a flat $`𝕊p(2h^{2,1},)`$ connection on the tangent bundle $`𝒯_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ of the moduli space $`\left(\text{M}\right)`$ of three dimensional CY manifolds. We will call this connection the Cecotti-Hitchin-Simpson-Vafa connection and will refer to it as the CHSV connection.
A beautiful theorem of Simpson proved in shows when a quasi-projective variety is covered by a symmetric domain. One can show that the tt\* equations can be interpreted in the same way. This will be done in .
In Section 3 we interpreted the holomorphic connection which is defined by the Frobenius Algebra structure on the bundle $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ as a Higgs field. It seems that the paper by Deligne is suggesting that the Higgs field that we constructed is related to variation of mixed Hodge Structure of CY threefolds, when there exists a maximal unipotent element in the mapping class group. See .
In Section 4 we review some basic constructions from .
In Section 5 we used the ideas from and some modifications of the beautiful computations done by E. Witten in to quantize the tangent bundle of $`\left(\text{M}\right)`$. This can be done since we can identify the tangent spaces at each point of the moduli space of CY manifolds by using the parallel transport defined by the flat connection $`𝕊p(2h^{2,1},)`$ defined by Cecotti-Hitchin-Simpson-Vafa and the existence of the non-zero section $`\eta _\tau `$ of the relative dualizing line bundle over $`\left(\text{M}\right)`$. We will show that the symplectic structure defined by the imaginary part of the Weil-Petersson metric is parallel with respect to the CHSV connection. In this section we construct a projective flat connection on some Hilbert vector bundle associated with the tangent bundle on the moduli space $`\left(\text{M}\right)`$. Based on these results, the method from and the technique developed in , we show that holomorphic anomaly equations $`\left(\text{1}\right)`$ of Bershadsky, Cecotti, Ooguri and Vafa of the genus $`g2`$ imply that the free energy $`Z`$ defined by $`\left(\text{97}\right)`$ is a parallel with respect to a flat projective connection constructed in Section 6.
In Section 6 we will introduce a natural $``$ structure on the tangent space of each point of the moduli space of CY threefolds by using the flat $`𝕊p(2h^{2,1},)`$ connection constructed in Section 3. In order to do that we introduce the notion of the extended period space of CY threefolds which is similar to the period domain of marked algebraic polarized K3 surfaces. We know from the moduli theory of algebraic polarized K3 surfaces that the points that define K3 surfaces with CM structure form an everywhere dense subset. This follows from the fact that the period domain is an open set on a quadric defined over $``$ in the projective space $`(^{20})`$. This fact together with the existence of a flat $`𝕊p(4h^{2,1},)`$ connection on the extended period domain will define in a natural way a lattice of maximal rank in the tangent space at each point of the moduli space of CY threefolds. Using the existence of the $`𝕊p(2h^{2,1},)`$ connection defined by the tt\* equation on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)\text{ }}^2`$, we define by the parallel translations a $``$ structure on the fibres of the bundle $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$at each point of the moduli space of M.
The mirror symmetry suggests that we can identify the second cohomology group of the mirror CY M’ with $`H^{2,1}`$ of the original CY manifold. Since the second cohomology group of a CY manifold has a natural $``$ structure, then $`H^{2,1}`$ of the original CY manifold should also carry a natural $``$ structure. This construction suggests that the existence of the natural $``$ structure on $`H^{2,1}`$ is equivalent to the tt\* equations.
In Section 7 we obtain an algebraic integrable system in the sense of R. Donagi and E. Markman using the flat Cecotti-Hitchin-Simpson-Vafa connection. From that we obtain a map from the moduli space of CY manifold M to the moduli space of principally polarized abelian varieties and the CHSV connection is the pull back of the connection defined by R. Donagi and E. Markman on the moduli space of principally polarized abelian varieties. See . We also construct a Hyper-Kähler structure on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ of polarized CY threefolds. D. Freed constructed Hyper-Kähler structure on the tangent bundle of the relative dualizing sheaf of the moduli space $`\left(\text{M}\right)`$ of polarized CY threefolds in .
Acknowledgements I want to express my special thanks to S. Shatashvili, who pointed out the important paper of Witten, where the anomaly equations of BCOV were interpreted as a projective connection. The author want to thank Y. Eliashberg, J. Li and K. Liu for stimulating conversations and their patience with me. Special thanks to N. Nekrasov whose remarks and suggestions helped me enormously.
## 2 Deformation Theory for CY manifolds.
### 2.1 Review of
###### Definition 1
Let M be an even dimensional C manifold. We will say that M has an almost complex structure if there exists a section $`IC^{\mathrm{}}(`$M$`,Hom(T^{},T^{}))`$ such that $`I^2=id.`$ $`T`$ is the tangent bundle and $`T^{}`$ is the cotangent bundle on M.
This definition is equivalent to the following one:
###### Definition 2
Let M be an even dimensional C manifold. Suppose that there exists a global splitting of the complexified cotangent bundle $`T^{}𝐂=\mathrm{\Omega }^{1,0}\mathrm{\Omega }^{0,1},`$ where $`\mathrm{\Omega }^{0,1}=\overline{\mathrm{\Omega }^{1,0}}.`$ Then we will say that M has an almost complex structure.
We are going to define the almost integrable complex structure.
###### Definition 3
We will say that an almost complex structure is an integrable one if for each point $`x`$M there exists an open set $`U`$M such that we can find local coordinates $`z^1,..,z^n`$ such that $`dz^1,..,dz^n`$ are linearly independent in each point $`mU`$ and they generate $`\mathrm{\Omega }^{1,0}|{}_{U}{}^{}.`$
It is easy to see that any complex manifold has an almost integrable complex structure.
###### Definition 4
Let M be a complex manifold. $`\varphi \mathrm{\Gamma }(`$M$`,Hom(\mathrm{\Omega }^{1,0},\mathrm{\Omega }^{0,1}))`$ is called a Beltrami differential.
Since $`\mathrm{\Gamma }(`$M$`,Hom(\mathrm{\Omega }^{1,0},\mathrm{\Omega }^{0,1}))\mathrm{\Gamma }(`$M$`,T^{1,0}\mathrm{\Omega }^{0,1}),`$ we deduce that locally $`\varphi `$ can be written as follows:
$$\varphi |{}_{U}{}^{}=\varphi _{\overline{\alpha }}^\beta \overline{dz}^\alpha \frac{}{z^\beta }.$$
From now on we will denote by
$$A_\varphi =\left(\begin{array}{cc}id& \varphi (\tau )\\ \overline{\varphi (\tau )}& id\end{array}\right):T^{}T^{}.$$
We will consider only those Beltrami differentials $`\varphi `$ such that $`det(A_\varphi )0.`$
###### Definition 5
It is easy to see that the Beltrami differential $`\varphi `$ defines a new almost complex structure operator $`I_\varphi =A_\varphi ^1IA_\varphi .`$
With respect to this new almost complex structure the space $`\mathrm{\Omega }_\varphi ^{1,0}`$ is defined as follows; if $`dz^1,..,dz^n`$ generate $`\mathrm{\Omega }^{1,0}|_U`$, then
$$dz^1+\varphi (dz^1),..,dz^n+\varphi (dz^n)$$
generate $`\mathrm{\Omega }_\varphi ^{1,0}|_U`$ and, moreover we have: $`\overline{\mathrm{\Omega }_\varphi ^{1,0}}\mathrm{\Omega }_\varphi ^{1,0}=0.`$ The Beltrami differential $`\varphi `$ defines an integrable complex structure on M if and only if the following equation holds:
$$\overline{}\varphi +\frac{1}{2}[\varphi ,\varphi ]=0.$$
where
$$[\varphi ,\varphi ]|{}_{U}{}^{}:=\underset{\nu =1}{\overset{n}{}}\underset{1\alpha ,\beta n}{}\left(\underset{\mu =1}{\overset{n}{}}(\varphi _{\overline{\alpha }}^\mu \left(_\mu \varphi _{\overline{\beta }}^\nu \right)\varphi _{\overline{\beta }}^\mu \left(_\mu \varphi _{\overline{\alpha }}^\nu \right))\right)\overline{dz}^\alpha \overline{dz}^\beta \frac{}{dz^\nu }.$$
(See .)
The main results in are the two theorems stated bellow:
###### Theorem 6
Let M be a CY manifold and let $`\left\{\varphi _i\right\}`$ be harmonic (with respect to the CY metric $`g`$) representative of the basis in $`^1(`$M$`,T^{1,0}),`$ then the equation: $`\overline{}\varphi +\frac{1}{2}[\varphi ,\varphi ]=0`$ has a solution in the form:
$$\varphi (\tau _1,..,\tau _N)=\underset{i=1}{\overset{N}{}}\varphi _i\tau ^i+\underset{|I_N|2}{}\varphi _{I_N}\tau ^{I_N}$$
where $`I_N=(i_1,..,i_N)`$ is a multi-index,
$$\varphi _{I_N}C^{\mathrm{}}(M,\mathrm{\Omega }^{0,1}T^{1,0}),$$
$`\tau ^{I_N}=(\tau ^i)^{i_1}\mathrm{}(\tau ^N)^{i_N}`$ and there exists $`\epsilon >0`$ such that
$$\varphi (\tau )C^{\mathrm{}}(M,\mathrm{\Omega }^{0,1}T^{1,0})$$
for $`|\tau ^i|<\epsilon `$ for $`i=1,..,N.`$ $`See`$ $`\text{[32]}.`$
###### Theorem 7
Let $`\mathrm{\Omega }_0`$ be a holomorphic n-form on the n dimensional CY manifold M. Let $`\left\{U_i\right\}`$be a covering of M and let $`\{z_1^i,..,z_n^i\}`$ be local coordinates in $`U_i`$ such that $`\mathrm{\Omega }_0|{}_{U_i}{}^{}=dz_1^i\mathrm{}dz_n^i.`$ Then for each $`\tau =(\tau ^1,..,\tau ^N)`$ such that $`|\tau _i|<\epsilon `$ the forms on M defined as:
$$\mathrm{\Omega }_t|{}_{U_i}{}^{}:=(dz_1^i+\varphi (\tau )(dz_1^i))..(dz_n^i+\varphi (\tau )(dz_n^i))$$
are globally defined complex n forms $`\mathrm{\Omega }_\tau `$ on M and, moreover, $`\mathrm{\Omega }_\tau `$ are closed holomorphic n forms with respect to the complex structure on M defined by $`\varphi (\tau ).`$
###### Corollary 8
We have the following Taylor expansion for
$$\mathrm{\Omega }_\tau |{}_{U}{}^{}=\mathrm{\Omega }_0+\underset{k=1}{\overset{n}{}}(1)^{\frac{k(k1)}{2}}\left(^k\varphi \right)\mathrm{}\mathrm{\Omega }_0.$$
(6)
(See .)
From here we deduce the following Taylor expansion for the cohomology class $`\left[\mathrm{\Omega }_\tau \right]`$ $`H^n(`$M,$`𝐂):`$
###### Corollary 9
$$\left[\mathrm{\Omega }_\tau \right]=[\mathrm{\Omega }_0]\underset{i=1}{\overset{N}{}}[\left(\varphi _i\mathrm{}\mathrm{\Omega }_0\right)]\tau ^i+\frac{1}{2}\underset{i,j=1}{\overset{N}{}}[(\left(\varphi _i\varphi _j\right)\mathrm{}\mathrm{\Omega }_0)]\tau ^i\tau ^j+O(\tau ^3)$$
(7)
(See .)
We are going to define the Kuranishi family for CY manifolds of any dimension.
###### Definition 10
Let $`𝒦^N`$ be the polydisk defined by $`|\tau ^i|<\epsilon `$ for every $`i=1,..,N`$, where $`\epsilon `$ is chosen such that for every $`\tau 𝒦`$ , $`\varphi (\tau )C^{\mathrm{}}(`$M,$`\mathrm{\Omega }^{0,1}T^{1,0})`$, where $`\varphi (\tau )`$ is defined as in Definition 4. On the trivial $`C^{\mathrm{}}`$ family M$`\times 𝒦`$ we will define for each $`\tau 𝒦`$ an integrable complex structure I<sub>ϕ(τ)</sub> on the fibre over $`\tau `$ of the family M$`\times 𝒦`$ ,where I<sub>ϕ(τ)</sub> was defined in Definition 5. Thus we will obtain a complex analytic family $`\pi :𝒳𝒦`$ of CY manifolds. We will call this family the Kuranishi family. Thus we introduce also a coordinate system in $`𝒦`$. We call this coordinate system a flat coordinate system.
### 2.2 Construction of the Teichmüller Space of CY Manifolds
###### Definition 11
We will define the Teichmüller space $`𝒯\left(\text{M}\right)`$ of M as follows:
$$𝒯\left(\text{M}\right):=\{\text{all integrable complex structures on M}\}/\mathrm{𝐃𝐢𝐟𝐟}_0\left(\text{M}\right),$$
where $`\mathrm{𝐃𝐢𝐟𝐟}_0\left(\text{M}\right)`$ is the group of diffeomorphisms of M isotopic to identity.
$`\mathrm{𝐃𝐢𝐟𝐟}_0\left(\text{M}\right)`$ acts on the complex structures as follows: let $`\psi \mathrm{𝐃𝐢𝐟𝐟}_0\left(\text{M}\right)`$ and let
$$IC^{\mathrm{}}(Hom(T^{}\left(\text{M}\right),T^{}\left(\text{M}\right)),$$
such that $`I^2=id,`$ then clearly $`\psi ^{}(I)`$ is such that $`(\psi ^{}(I))^2=id.`$ Moreover, if $`I`$ is an integrable complex structure, then $`\psi ^{}(I)`$ is integrable too.
We will call a pair $`(`$M, $`\{\gamma _1,\mathrm{},\gamma _{b_n}\})`$ a marked CY manifold, if M is a Calabi-Yau manifold and $`\{\gamma _1,..,\gamma _{b_n}\}`$ is a basis in $`H_n(M,𝐙)/Tor.`$ Over the Kuranishi space we have a universal family of marked Calabi-Yau manifolds $`𝒳𝒦`$ defined up to an action of a group that acts trivially on the middle homology and preserves the polarizations class. And, moreover, as a C manifold $`𝒳`$ is diffeomorphic to $`𝒦\times `$M.
###### Theorem 12
The Teichmüller space$`𝒯`$ $`\left(\text{M}\right)`$ of a Calabi Yau manifold M exists as a complex manifold of dimension $`h^{2,1}.`$
Proof: For the proof of Theorem 12 see . $`\mathrm{}`$
### 2.3 Construction of the Moduli Space
###### Definition 13
We will define the mapping class group $`\mathrm{\Gamma }^{}\left(\text{M}\right)`$ as follows:
$$\mathrm{\Gamma }^{}\left(\text{M}\right):=\mathrm{𝐃𝐢𝐟𝐟}^+\left(\text{M}\right)/\mathrm{𝐃𝐢𝐟𝐟}_0\left(\text{M}\right),$$
where $`\mathrm{𝐃𝐢𝐟𝐟}^+\left(\text{M}\right)`$ is the group of diffeomorphisms preserving the orientation of M and $`\mathrm{𝐃𝐢𝐟𝐟}_0\left(\text{M}\right)`$ is the group of diffeomorphisms isotopic to identity.
D. Sullivan proved that the mapping class group of any C manifold of dimension greater or equal to 5 is an arithmetic group. (See .) It is easy to prove that the mapping class group $`\mathrm{\Gamma }^{}\left(\text{M}\right)`$ acts discretely on the Teichmüller space $`𝒯\left(\text{M}\right)`$ of the CY manifold M.
We will consider from now on polarized CY manifolds, i.e. a pair $`(M,\omega (1,1)),`$ where
$$[\omega (1,1)]H^2(\text{M},)H^{1,1}(\text{M},)$$
is a fixed class of cohomology and it corresponds to the imaginary part of a CY metric. We will define $`\mathrm{\Gamma }\left(\text{M}\right)`$ as follows:
$$\mathrm{\Gamma }_{\omega (1,1)}\left(\text{M}\right):=\{\varphi \mathrm{\Gamma }^{}\left(\text{M}\right)|\varphi ([\omega (1,1)])=[\omega (1,1)]\}.$$
From now on we will work with this family.
###### Theorem 14
There exists a subgroup $`\mathrm{\Gamma }\left(\text{M}\right)`$ in $`\mathrm{\Gamma }_{\omega (1,1)}`$ of finite index such that $`\mathrm{\Gamma }`$ acts without fixed points on the Teichmüller space $`𝒯\left(\text{M}\right)`$. The moduli space $`𝔐\left(\text{M}\right)=𝒯\left(\text{M}\right)/\mathrm{\Gamma }\left(\text{M}\right)`$ is a smooth quasi-projective variety. There exists a family of polarized CY manifolds $`𝒴\left(\text{M}\right)\left(\text{M}\right)=𝒯\left(\text{M}\right)/\mathrm{\Gamma }\left(\text{M}\right)`$. The relative dualizing sheaf $`\omega _{𝒴/\left(\text{M}\right)}`$ is a trivial line bundle.
Proof: Viehweg proved in that the moduli space $`\left(\text{M}\right)`$ is a quasi projective variety. In it was proved that we can find a subgroup $`\mathrm{\Gamma }\left(\text{M}\right)`$ in $`\mathrm{\Gamma }_{\omega (1,1)}\left(\text{M}\right)`$ such that the space $`𝒯\left(\text{M}\right)/\mathrm{\Gamma }\left(\text{M}\right)`$ is a smooth complex manifold. We also proved that over $`𝒯\left(\text{M}\right)/\mathrm{\Gamma }\left(\text{M}\right)=\left(\text{M}\right)`$ there exists a family of CY manifolds $`𝒴\left(\text{M}\right)\left(\text{M}\right).`$ In we proved the following Theorem:
###### Theorem 15
Let $`\left(\text{M}\right)=𝒯\left(\text{M}\right)/\mathrm{\Gamma }\left(\text{M}\right)`$. Then there exists a global non vanishing holomorphic section $`\eta _\tau `$ of the line bundle $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ whose $`L^2`$ norm $`\eta _\tau _{𝐋^2}^2`$ is equal to $`\left(det_{(0,1)}\right),`$ where $`det_{(0,1)}`$ is the regularized determinant of the Laplacian of a CY acting on $`\mathrm{\Omega }_\text{M}^{0,1}`$ of the CY metric with imaginary class equal to the polarization class and $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ is a trivial holomorphic line bundle.
Theorem 14 follows from Theorem 15. $`\mathrm{}`$
###### Corollary 16
$`\eta _\tau `$ defines a flat structure on $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}.`$
### 2.4 Weil-Petersson Geometry
In our paper we define a metric on the Kuranishi space $`𝒦`$ and called this metric, the Weil-Petersson metric. We will review the basic properties of the Weil-Petersson metric which were established in . In we proved the following theorem:
###### Theorem 17
Let M be a CY manifold of dimension n and let $`\mathrm{\Omega }_\text{M}`$ be a non zero holomorphic n form on M such that
$$\left(1\right)^{\frac{n(n1)}{2}}\left(\sqrt{1}\right)^n\underset{\text{M}}{}\mathrm{\Omega }_\text{M}\overline{\mathrm{\Omega }_\text{M}}=1.$$
Let g be a Ricci flat (CY) metric on M. Then the map:
$$\psi L^2(\text{M},\mathrm{\Omega }_\text{M}^{0,k}^mT_\text{M}^{1,0})\psi \mathrm{}\mathrm{\Omega }_\text{M}L^2(\text{M},\mathrm{\Omega }_\text{M}^{nm,k})$$
gives an isomorphism between Hilbert spaces and this map preserves the Hodge decomposition..
###### Corollary 18
We can identify the tangent space $`T_\tau =H^1(`$M$`{}_{\tau }{}^{},T_\tau ^{1,0})`$ at each point $`\tau 𝒯\left(\text{M}\right)`$ with$`H^1(`$M$`{}_{\tau }{}^{},\mathrm{\Omega }_\tau ^{n1})`$, by using the map $`\psi \psi \mathrm{}\mathrm{\Omega }_\text{M}.`$
###### Notation 19
We will denote by
$$\omega _1,\omega _2:=\underset{\text{M}_\tau }{}\omega _1\overline{\omega _2}.$$
(8)
###### Definition 20
Let $`\psi _1,`$ $`\psi _2T_\tau =𝐇^1(\text{M}_\tau ,T_{\text{M}_\tau }^{1,0})`$ $`(`$the space of harmonic forms with respect to the CY metric g.). We will define the Weil-Petersson metric as follows:
$$\psi _1,\psi _2_{WP}:=\sqrt{1}\underset{\text{M}_\tau }{}\left(\psi _1\mathrm{}\mathrm{\Omega }_\tau \right)\left(\overline{\psi _2\mathrm{}\mathrm{\Omega }_\tau }\right)=\sqrt{1}\psi _1\mathrm{}\mathrm{\Omega }_\tau ,\overline{\psi _2\mathrm{}\mathrm{\Omega }_\tau }$$
and $`\mathrm{\Omega }_\tau ^2=1.`$ Thus $`\psi ,\psi _{WP}>0`$.
The Weil-Petersson metric is a Kähler metric on the Teichmüller space $`𝒯\left(\text{M}\right)`$. It defines a natural connection, namely the Levi-Civita connection$``$ We will denote the covariant derivatives in direction $`\frac{}{\tau ^i}`$ at the tangent space of a point $`\tau 𝒯\left(\text{M}\right)`$ defined by $`\varphi _i`$ by $`_i.`$ In we proved the following theorem:
###### Theorem 21
In the flat coordinate system introduced in Definition 10 the following formulas hold for the curvature operator:
$$R_{i\overline{j},k\overline{l}}=\delta _{i\overline{j}}\delta _{k\overline{l}}+\delta _{i\overline{l}}\delta _{k\overline{j}}\sqrt{1}\underset{\text{M}}{}((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_\text{M}))(\overline{(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\text{M}})$$
$$=\delta _{i\overline{j}}\delta _{k\overline{l}}+\delta _{i\overline{l}}\delta _{k\overline{j}}\sqrt{1}(\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_\text{M})),(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\text{M}.$$
(9)
## 3 Flat $`𝕊`$p(2h$`{}_{}{}^{2,1},)`$ Structure on the Moduli Space of CY Threefolds
### 3.1 A Flat Structure on the Line Bundle $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$
The flat structure on the line bundle $`\omega _{𝒴\text{ }\left(\text{M}\right)\text{ /}\left(\text{M}\right)}`$ is defined by Corollary 16.
### 3.2 Gauss-Manin Connection
###### Definition 22
On the Teichmüller space $`𝒯\left(\text{M}\right)`$ we have a trivial bundle namely
$$^n=H^n(\text{M},)\times 𝒯\left(\text{M}\right)𝒯\left(\text{M}\right).$$
Theorem 14 implies that we constructed the moduli space $`\left(\text{M}\right)`$ as $`𝒯\left(\text{M}\right)`$/$`\mathrm{\Gamma }\left(\text{M}\right)`$. Thus we obtain a natural representation of the group $`\mathrm{\Gamma }\left(\text{M}\right)`$ into $`H^n(\text{M},)`$ and a flat connection on the flat bundle
$$^n/\mathrm{\Gamma }\left(\text{M}\right)𝒯\left(\text{M}\right)/\mathrm{\Gamma }\left(\text{M}\right)=\left(\text{M}\right).$$
This connection is called the Gauss-Manin connection. The covariant derivative in direction $`\varphi _i`$ of the tangent space T$`_{\tau ,\left(\text{M}\right)}`$with respect to the Gauss-Manin connection will be denoted by $`𝒟_i`$.
The Gauss-Manin connection $`𝒟`$ is defined in a much more general situation and it is defined on the moduli space of CY manifolds of dimension $`n3`$. We will state explicit formulas for the covariant differentiation $`𝒟_i`$ defined by the Gauss-Manin connection. We will fix a holomorphic three form $`\mathrm{\Omega }_0`$ such that
$$\sqrt{1}\mathrm{\Omega }_0,\mathrm{\Omega }_0=\sqrt{1}\underset{\text{M}}{}\mathrm{\Omega }_0\overline{\mathrm{\Omega }_0}=\mathrm{\Omega }_0^2=1.$$
Using the form $`\mathrm{\Omega }_0`$ and theorem 17, we can identify the cohomology groups $`H^1(\text{M},T_\text{M}^{1,0})`$ and $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$ on M:
###### Proposition 23
The map
$$𝔦:\psi \psi \mathrm{}\mathrm{\Omega }_\text{M}$$
(10)
is an isomorphism between the groups $`H^1(\text{M},T_\text{M}^{1,0})`$ and $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$.
Proof: Our proposition follows directly from Theorem 17. $`\mathrm{}`$
###### Remark 24
From now on in the map $`\left(\text{10}\right)`$ we will use for $`\mathrm{\Omega }_\text{M}`$ the restriction of the holomorphic form $`\eta _\tau `$ on M defined by Theorem 15.
###### Remark 25
Suppose that M is a three dimensional CY manifold. Then the Poincare map identifies $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$ with $`H^2(\text{M},\mathrm{\Omega }_\text{M}^1).`$ This identification will be denoted by $`\mathrm{\Pi }`$, i.e. $`\mathrm{\Pi }:H^2(\text{M},\mathrm{\Omega }_\text{M}^1)H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$ and it is defined by identifying some basis $`\left\{\mathrm{\Omega }_i\right\}`$ of $`H^2(\text{M},\mathrm{\Omega }_\text{M}^1)`$ with the basis $`\left\{\overline{\mathrm{\Omega }_i}\right\}`$ of $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2).`$ So:
$$\mathrm{\Pi }(\mathrm{\Omega }_i):=\overline{\mathrm{\Omega }_i}.$$
(11)
###### Notation 26
Using Proposition 23 and Remark 25 one can identify the spaces $`H^1(\text{M},T_\text{M}^{1,0})`$ and $`H^2(\text{M},\mathrm{\Omega }_\text{M}^1)`$ for three dimensional CY by using the map $`F`$, where
$$F(\varphi ):=\mathrm{\Pi }(\iota (\varphi ))\text{ }$$
(12)
for $`\varphi H^1(`$M$`,T_\text{M}^{1,0}).`$
###### Lemma 27
Let $`𝔦^1:H^1(\text{M},\mathrm{\Omega }_\text{M}^2)\stackrel{\mathrm{}\mathrm{\Omega }_\text{M}^{}}{}H^1(\text{M},T_\text{M}^{1,0})`$ be the inverse identification defined by $`\left(\text{10}\right)`$. Then
$$𝒟_i(𝔦(\varphi ))=\iota \left(\varphi _i\right)\mathrm{}\varphi H^2(\text{M}_0,\mathrm{\Omega }_{\text{M}_0}^1).$$
(13)
Proof: The proof of the lemma follows directly from formula$`\left(\text{7}\right)`$. Indeed since $`\{\varphi _i\}`$ is a basis of $`H^1(`$M$`,T_\text{M}^{1,0})`$, then $`\{\varphi _i\mathrm{}\mathrm{\Omega }_\text{M}\}`$ will be a basis of $`H^1(\text{M},\mathrm{\Omega }_\text{M}^2)`$. In the flat coordinates $`(\tau ^1,\mathrm{},\tau ^N)`$ introduced in Definition 10 and according to $`\left(\text{7}\right)`$ we have
$$\left[\mathrm{\Omega }_\tau \right]=[\mathrm{\Omega }_\text{M}]\underset{i=1}{\overset{N}{}}\left[\left(\varphi _i\mathrm{}\mathrm{\Omega }_\text{M}\right)\right]\tau ^i+\underset{i,j=1}{\overset{N}{}}\left[\left(\varphi _i\varphi _j\right)\mathrm{}\mathrm{\Omega }_\text{M}\right]\tau ^i\tau ^j+O(\tau ^3).$$
From Definition 22 and the expression of $`[\mathrm{\Omega }_\tau ]`$ given by the above formula, we deduce formula $`\left(\text{13}\right)`$. Lemma 22 is proved. $`\mathrm{}`$
### 3.3 Higgs Fields and the Tangent Space of $`\left(\text{M}\right)`$
Let $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ be the relative dualizing sheaf on $`\left(\text{M}\right).`$
###### Definition 28
We define the Higgs field of a holomorphic bundle $``$ over a complex manifold M as a globally defined holomorphic map $`\mathrm{\Phi }:\mathrm{\Omega }_\text{M}^1`$ such that
$$\mathrm{\Phi }\mathrm{\Phi }=0.$$
(14)
###### Lemma 29
The tangent bundle $`T_{𝒯\left(\text{M}\right)}`$ of the Teichmüller space $`𝒯\left(\text{M}\right)`$ of any CY manifold M is canonically isomorphic to the bundles
$$T_{𝒯\left(\text{M}\right)}Hom(\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)},R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2)$$
$$\left(\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}\right)^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2.$$
Proof: The proof of this lemma is standard and follows from local Torelli Theorem. $`\mathrm{}`$
### 3.4 Construction of a Higgs Field on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$
From now on we will consider only three dimensional CY manifolds.
###### Definition 30
We define a Higgs field $`\stackrel{~}{}`$ on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ by using the Gauss-Manin connection and Poincare duality $`\mathrm{\Pi }`$ in the following manner:
$$\stackrel{~}{}_{\varphi _i}\mathrm{\Omega }=\stackrel{~}{}_i\mathrm{\Omega }:=\mathrm{\Pi }(𝒟_i(\mathrm{\Omega }))=\mathrm{\Pi }(\mathrm{\Omega }\mathrm{}\varphi _i),$$
(15)
where $`\mathrm{\Omega }H^1(\text{M},\mathrm{\Omega }^2)`$ and$`\{\varphi _i\}H^1(\text{M},T_\text{M}^{1,0})`$ is an orthonormal basis with respect to the Weil-Petersson metric on $`T_0^1(\text{M},T_\text{M}^{1,0})`$.
###### Lemma 31
$`\stackrel{~}{}`$ as defined in Definition 30 is a Higgs field.
Proof: We need to check that $`\stackrel{~}{}`$ satisfies the conditions in the definition 28 of a Higgs field. The definition $`\left(\text{15}\right)`$ of $`\stackrel{~}{}`$ implies that
$$\stackrel{~}{}Hom(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2T(\left(\text{M}\right)),R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2).$$
(16)
On the other hand standard facts from commutative algebra imply
$$Hom(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2T(\left(\text{M}\right),R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2)$$
$$\left(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2\right)^{}\left(T(\left(\text{M}\right))\right)^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2$$
$$\left(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2\right)^{}\left(\mathrm{\Omega }_{\left(\text{M}\right)}^1R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2\right)$$
$$Hom(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2,\mathrm{\Omega }_{\left(\text{M}\right)}^1R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2).$$
(17)
So we obtain
$$Hom(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2T(\left(\text{M}\right)),R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2)$$
$$Hom(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2,\mathrm{\Omega }_{\left(\text{M}\right)}^1R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2).$$
(18)
So the condition that
$$\stackrel{~}{}Hom(R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2,R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2\mathrm{\Omega }_{\left(\text{M}\right)}^1)$$
follows directly from $`\left(\text{16}\right)`$ and $`\left(\text{18}\right)`$. Next we need to prove that $`\stackrel{~}{}\stackrel{~}{}=0.`$ It is a standard fact that the relation $`\stackrel{~}{}\stackrel{~}{}=0`$ is equivalent to the relations $`[\stackrel{~}{}_i,\stackrel{~}{}_j]=0.`$ So in order to finish the proof of Lemma 31 we need to prove the following Proposition:
###### Proposition 32
The commutator of $`\stackrel{~}{}_i`$ and $`\stackrel{~}{}_j`$ is equal to zero, i.e.
$$[\stackrel{~}{}_i,\stackrel{~}{}_j]=0.$$
(19)
Proof: Lemma 32 follows directly from the definition of $`\stackrel{~}{}.`$ Indeed since the Gauss-Manin connection is a flat, i.e. $`[𝒟_i,𝒟_j]=0`$ and the fact that the covariant derivative of Poincare duality is zero we deduce that:
$$[\stackrel{~}{}_i,\stackrel{~}{}_j]=[\mathrm{\Pi }(𝒟_i),\mathrm{\Pi }\left(𝒟_j\right)]=\mathrm{\Pi }[𝒟_i,𝒟_j]=0.$$
Proposition 32 is proved. $`\mathrm{}`$ Lemma 31 is proved. $`\mathrm{}`$
### 3.5 Higgs Field on the Tangent Bundle of $`\left(\text{M}\right)`$
Let $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ be the relative dualizing sheaf on $`\left(\text{M}\right).`$ We will define bellow a Higgs Field on the tangent vector bundle
$$T_{\left(\text{M}\right)}\left(\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}\right)^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2$$
of the moduli space $`\left(\text{M}\right)`$. We denote by $`\theta =id\stackrel{~}{},`$ where $`\stackrel{~}{}`$ is defined by $`\left(\text{15}\right).`$
###### Lemma 33
Let $`\{\varphi _i\}`$ be a basis of the tangent bundle $`𝒯_{\left(\text{M}\right)}`$ restricted on some open polydisk $`U\left(\text{M}\right).`$ We can identify the fibre $`T_{\text{M}_\tau }^{1,0}`$ of $`𝒯_{\left(\text{M}\right)}`$ with the harmonic forme $`^1(\text{M}_\tau ,T_{\text{M}_\tau }^{1,0})`$ with respect to the CY metric corresponding to the \[p;arization class $`L.`$ Let us define
$$\vartheta :𝒯_{\left(\text{M}\right)}𝒯_{\left(\text{M}\right)}\left(\left(\text{M}\right)\right)\mathrm{\Omega }_{\left(\text{M}\right)}^1\left(\left(\text{M}\right)\right)$$
by
$$\vartheta \left(\eta _\tau ^1\varphi _j\right):=\eta _\tau ^1\underset{i=1}{\overset{N}{}}\left(ϝ^1\left(𝒟_i(\varphi _j\mathrm{}\mathrm{\Omega }_\tau )\right)(\varphi _i)^{}\right)$$
where $`𝒟`$ is the Gauss-Manin connection defined by $`\left(\text{15}\right)`$, $`ϝ`$ is defined by $`\left(\text{12}\right)`$ and $`\eta _\tau `$ is the holomorphic three form on M defined by Theorem 15. Then
$$\vartheta _i\left(\eta _\tau ^1\varphi _j\right)=\eta _\tau ^1\left(\left(\left(\varphi _j\mathrm{}\mathrm{\Omega }_\tau \right)\mathrm{}\varphi _i\right)\mathrm{}\left(\mathrm{\Omega }_\tau \right)^{}\right),$$
(20)
and $`\vartheta `$ defines a Higgs field on the tangent bundle of $`\left(\text{M}\right).`$
Proof: Lemma 33 follows directly from Lemma 31, the definition of $`\vartheta `$ and the fact that
$$T_{\left(\text{M}\right)}\left(\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}\right)^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2.$$
Lemma 33 is proved. $`\mathrm{}`$
We know that $`h^{2,1}=dim_{}H^1\left(\text{M\hspace{0.17em},}\mathrm{\Omega }_\text{M}^2\right)=h^{1,2}=dim_{}H^2\left(\text{M\hspace{0.17em},}\mathrm{\Omega }_\text{M}^1\right)`$ are constants. Therefore $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ and $`R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1`$ are holomorphic vector bundles over the moduli space $`\left(\text{M}\right)`$ of polarized CY threefold. We have a non degenerate pairing:
$$R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2\times R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1R^3\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^3$$
(21)
given by
$$\sqrt{1}\left(\left(\varphi _j(\tau )\mathrm{}\mathrm{\Omega }_\tau \right)\right)\omega _j(\tau )=h_{ij}(\tau )\mathrm{\Omega }_\tau \overline{\mathrm{\Omega }_\tau },$$
(22)
where $`\left(\varphi _j(\tau )\mathrm{}\mathrm{\Omega }_\tau \right)`$ and $`\omega _j(\tau )`$ are holomorphic sections of the holomorphic vector bundles $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ and $`R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1.`$ Clearly $`h_{ij}(\tau )`$ depends holomorphically on $`\tau \left(\text{M}\right).`$ By using the non-degenerate pairing $`\left(\text{21}\right)`$ we can identify $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ with $`R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1`$ as follows: to the basis of holomorphic sections
$$\left\{\varphi _i\mathrm{}\mathrm{\Omega }_{\tau _0}\right\}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2$$
we will assign
$$\left(\varphi _i\mathrm{}\mathrm{\Omega }_{\tau _0}\right)^{}Hom(R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1,R^3\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^3)$$
such that for the pairing defined by $`\left(\text{22}\right)`$ satisfies:
$$\varphi _i\mathrm{}\mathrm{\Omega }_\tau \left(\varphi _i\mathrm{}\mathrm{\Omega }_{\tau _0}\right)^{}=$$
$$\varphi _i\mathrm{}\mathrm{\Omega }_\tau \left(\varphi _i\mathrm{}\mathrm{\Omega }_{\tau _0}\right)^{}=\delta _{i,j}\mathrm{\Omega }_\tau \overline{\mathrm{\Omega }_\tau }.$$
(23)
So $`\left(\varphi _i\mathrm{}\mathrm{\Omega }_{\tau _0}\right)`$ and $`\omega _l^{}`$ are given by the formula:
$$\omega _l^{}=\underset{k}{}h^{lk}\left(\varphi _k\mathrm{}\mathrm{\Omega }_\tau \right)$$
(24)
$$\left(\varphi _l\mathrm{}\mathrm{\Omega }_\tau \right)^{}=\underset{k}{}h^{kl}\omega _k.$$
(25)
###### Lemma 34
Let us a fix a point $`\tau _0\left(\text{M}\right).`$ Suppose that
$$\{\varphi _i,\text{ }i=1,\mathrm{},N\}\text{ and }\{\omega _i,\text{ }i=1,\mathrm{},N\}$$
are some bases of $`T_{\tau ,\left(\text{M}\right)}(U)=\mathrm{\Omega }_\tau R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ and $`R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1`$ in some open polydisk of $`\tau _0.`$ Let us define
$$C_{ijl}=\sqrt{1}\underset{\text{M}}{}(\mathrm{\Omega }_\tau \left((\varphi _i\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\tau \right)).$$
(26)
Let $`h_{ij}:=\varphi _i\mathrm{}\mathrm{\Omega }_\tau ,\omega _j`$ be the pairing between the holomorphic vector bundles $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ and $`R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1`$ defined by $`\left(\text{22}\right).`$ Then
$$\vartheta _i\varphi _j=\underset{k,l=1}{\overset{N}{}}C_{ij}^k\left(\omega _k^{}\mathrm{}\mathrm{\Omega }_\tau ^{}\right),$$
(27)
$$\vartheta _i\varphi _j=\underset{k=1}{\overset{N}{}}C_{ijk}\varphi _k,$$
where $`C_{ij}^k`$ and $`C_{ijk}`$ are holomorphic functions in $`U`$. The relations between $`C_{ij}^k`$ and $`C_{ijl}`$ are given by
$$C_{ij}^k=\underset{m=1}{\overset{n}{}}C_{ijm}h^{mk},$$
(28)
and,
$$C_{ijl}=C_{jil}=C_{ilj}=C_{jli}=C_{lij}=C_{lij}.$$
(29)
Proof: $`\left(\text{20}\right)`$ implies that
$$\vartheta _i\varphi _j=\left(\left(\varphi _j\mathrm{}\mathrm{\Omega }_\tau \right)\mathrm{}\varphi _i\right)\mathrm{}\left(\mathrm{\Omega }_\tau \right)^{}=\left(\varphi _i\varphi _j\mathrm{}\mathrm{\Omega }_\tau \right)\mathrm{}\mathrm{\Omega }_\tau ,$$
where $`(\varphi _i\varphi _j\mathrm{}\mathrm{\Omega }_\tau )R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1|{}_{U}{}^{}.`$ Since $`\left\{\omega _k\right\}`$ is a basis of the holomorphic vector bundle $`R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1|_U`$ we have
$$\left(\varphi _i\varphi _j\mathrm{}\mathrm{\Omega }_\tau \right)=\underset{k}{}C_{ij}^k\omega _k.$$
(30)
Poincare duality and $`\left(\text{30}\right)`$ imply that we can identify
$$\left(\varphi _i\varphi _j\mathrm{}\mathrm{\Omega }_\tau \right)R^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1$$
with the holomorphic section
$$\underset{k}{}C_{ij}^k\omega _k^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2,$$
where $`\omega _k^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2|_U`$ are defined by $`\left(\text{24}\right)`$ and are the Poincare dual of $`\omega _kR^2\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^1|_U`$. Thus we have
$$\varphi _l\mathrm{}\mathrm{\Omega }_\tau ,\varphi _i\varphi _j\mathrm{}\mathrm{\Omega }_\tau =\varphi _l\mathrm{}\mathrm{\Omega }_\tau ,\underset{k}{}C_{ij}^k\omega _k^{}=$$
$$\varphi _l\mathrm{}\mathrm{\Omega }_\tau ,\underset{k}{}C_{ij}^k\omega _k=\underset{k}{}C_{ij}^kh_{lk}.$$
(31)
Combining $`\left(\text{26}\right)`$ and $`\left(\text{31}\right)`$ we get:
$$\varphi _l\mathrm{}\mathrm{\Omega }_\tau ,\underset{k}{}C_{ij}^k\omega _k=\varphi _l\mathrm{}\mathrm{\Omega }_\tau ,\varphi _i\varphi _j\mathrm{}\mathrm{\Omega }_\tau =$$
$$\sqrt{1}\underset{\text{M}_{\tau _0}}{}(\mathrm{\Omega }_\tau \left((\varphi _i\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\tau \right)=C_{ijl}.$$
(32)
So we can conclude from $`\left(\text{22}\right)`$ and $`\left(\text{25}\right)`$ that
$$C_{ij}^k=\underset{l=1}{\overset{N}{}}C_{ijl}h^{lk}.$$
(33)
Thus $`\left(\text{32}\right)`$ and $`\left(\text{33}\right)`$ imply $`\left(\text{27}\right)`$ and $`\left(\text{28}\right).`$
Next we will prove $`\left(\text{29}\right).`$ We can multiply the global section $`\eta _\tau `$ of $`\omega _{𝒴\left(\text{M}\right)\text{/}\left(\text{M}\right)}`$ defined in by a constant and so we can assume that at the point $`\tau _0\left(\text{M}\right)`$ we have
$$\eta _{\tau _0}=\mathrm{\Omega }_0\text{ and }\eta _\tau =\lambda (\tau )\mathrm{\Omega }_\tau .$$
(34)
From the definition of $`\iota `$ given by $`\left(\text{10}\right)`$ and formula $`\left(\text{26}\right)`$, we conclude that $`\left(\text{29}\right)`$ holds. Thus Lemma 34 is proved. $`\mathrm{}`$
###### Lemma 35
We have
$$[\vartheta _i,\vartheta _j]=0.$$
(35)
Proof: The formula for $`\vartheta _i`$ given by $`\left(\text{20}\right)`$ implies
$$[\vartheta _i,\vartheta _j]\varphi _k=ϝ^1([𝒟_i,𝒟_j](\varphi _k\mathrm{}\mathrm{\Omega }_\tau )).$$
(36)
Since Gauss-Manin connection $`𝒟`$ is a flat connection then
$$[𝒟_i,𝒟_j]=0.$$
(37)
Combining formula $`\left(\text{37}\right)`$ with $`\left(\text{36}\right)`$ we get $`[\vartheta _i,\vartheta _j]\varphi _k=0.`$ Thus Lemma 35 is proved. $`\mathrm{}`$
### 3.6 Relations with Frobenius Algebras
One can use Lemma 22 to define an associative product on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ of three dimensional CY manifolds as follows: Let $`\{\varphi _i\}`$ be a basis of T$`{}_{\tau _0,\left(\text{M}\right)}{}^{}=H^1(\text{M},T_\text{M}^{1,0}),`$ then we define the product as:
$$\varphi _i\times \varphi _j=𝔦^1\left(\mathrm{\Pi }\left(𝒟_i(𝔦(\varphi _j))\right)\right)=F_{ijk}\varphi _k.$$
(38)
###### Lemma 36
Let $`C_{ij}^k`$ be defined by $`\left(\text{26}\right),`$ then
$$F_{ijk}=\sqrt{1}C_{ijk}.$$
(39)
Proof: Lemma 36 follows directly from the formulas for $`F_{ijk}`$ and $`C_{ijk}.`$ $`\mathrm{}`$
###### Corollary 37
The relations $`\left(\text{29}\right)`$ and $`\left(\text{35}\right)`$ shows that $`F_{ijk}`$ define a structure of a commutative algebra on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ of three dimensional CY manifolds.
### 3.7 The Analogue of Cecotti-Vafa tt\* Equations on $`\left(\text{M}\right)`$
###### Definition 38
Let $`\vartheta `$ be the Higgs field defined by $`\left(\text{20}\right).`$ Theorem 14 implies the existence of a global holomorphic non vanishing section $`\eta _\tau `$ of the line bundle $`\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}.`$ Then $`\eta _\tau `$ defines a metric on the flat line bundle $`\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}`$ with curvature zero. Let us define the Weil-Petersson metric on
$$𝒯_{\left(\text{M}\right)}Hom(\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)},R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2)$$
$$\left(\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}\right)^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2$$
as follows: Let
$$\varphi _\tau =\left(\eta _\tau \right)^{}\omega _\tau (1,1)𝒯_{\left(\text{M}\right)}\left(\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}\right)^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2.$$
Let us define the function
$$\lambda (\tau ):=\frac{\eta _\tau }{\mathrm{\Omega }_\tau }.$$
(40)
Then
$$\varphi _\tau ^2:=\left|\lambda (\tau )\right|^2\sqrt{1}\underset{\text{M}}{}\omega _\tau (1,2)\overline{\omega _\tau (1,2)}=$$
$$\frac{\sqrt{1}\omega _\tau (1,2),\overline{\omega _\tau (1,2)}}{\left|\lambda (\tau )\right|^2}.$$
(41)
Let $``$ be the standard connection of the metric defined by $`\left(\text{41}\right).`$ We will define the Cecotti-Hitchin-Vafa-Simpson (CHVS) connection $`D=D+\overline{D}`$ on the tangent bundle of the moduli space $`\left(\text{M}\right)`$ of three dimensional CY manifolds as follows:
$$D_i:_i+t\vartheta _i\text{ and }D_{\overline{j}}=_{\overline{j}}+t^1\overline{\vartheta }_j,$$
(42)
where $`t^{}.`$
###### Theorem 39
The curvature of the metric defined by $`\left(\text{41}\right)`$ in the flat coordinates defined by Definition 10 is given by:
$$R_{i\overline{j},k\overline{l}}=\sqrt{1}\underset{\text{M}}{}\left((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_\text{M}\right)\left(\overline{(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\text{M}}\right)=$$
$$=\sqrt{1}(\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_\text{M},(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\text{M}.$$
(43)
Proof: Since the metric on $`(\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)})^{}`$ is flat, then it has a zero curvature. The connection of the metric defined by $`\left(\text{41}\right)`$ will be
$$:=\stackrel{~}{}idid_1,$$
where $`\stackrel{~}{}`$ is the flat connection on $`\left(\pi _{}\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}\right)^{}`$ and $`_1`$ is the connection on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2.`$ Thus we get that for the curvature $`[,]`$ we have that $`[,]=[_1,_1].`$ So the curvature of the metric defined by $`\left(\text{41}\right)`$ is equal equal to the curvature on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2.`$ Formula $`\left(\text{7}\right)`$ implies
$$\sqrt{1}_i\overline{}_{\overline{j}}\left(\mathrm{\Omega }_\tau ,\mathrm{\Omega }_\tau \right)=$$
$$\sqrt{1}\mathrm{\Omega }_0\mathrm{}\varphi _i,\mathrm{\Omega }_0\mathrm{}\varphi _j\frac{\sqrt{1}}{2}\underset{k,l}{}(\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0)),(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0\tau ^k\overline{\tau }^l+O(|\tau |^3).$$
(44)
Thus $`\left(\text{44}\right)`$ implies $`\left(\text{43}\right).`$ Theorem 39 is proved. $`\mathrm{}`$
We will show that the Cecotti-Hitchin-Vafa-Simpson connection is flat.
###### Theorem 40
The connection $`D`$ defined by $`\left(\text{42}\right)`$ is a flat one, i.e.:
$$[D_i,D_j]=[D_i,D_{\overline{j}}]=[D_{\overline{i}},D_{\overline{j}}]=0$$
(45)
for all $`0i,jN.`$
Proof: First we will prove that $`[D_i,D_j]=0.`$ The definition of $`D_i=_i+t\vartheta _i`$ implies that
$$[D_i,D_j]=[_i,_j]+t[_i,\vartheta _j]+t^2[\vartheta _i,\vartheta _j].$$
(46)
We know that $`\mathrm{}_i`$ is a Hermitian connection of the Weil-Petersson metric defined by $`\left(\text{41}\right)`$, which is a Kähler metric and thus the (2,0) part of its curvature is zero. This implies that $`[_i,_j]=0.`$ From Lemma 34 we know that $`[\vartheta _i,\vartheta _j]=0.`$
###### Lemma 41
We have
$$[_i,\vartheta _j]=0.$$
(47)
Proof: In the flat coordinates $`(\tau ^1,\mathrm{},\tau ^N)`$ at a fixed point $`\tau _0=0\left(\text{M}\right)`$ of the moduli space we have that $`_i=_i`$. The definition of $`\vartheta _i`$ given by the formula $`\left(\text{20}\right)`$ and applied to
$$\varphi _k(\tau ):=\left(\frac{}{\tau ^k}\mathrm{\Omega }_\tau \right)\mathrm{}\left(\mathrm{\Omega }_\tau \right)^{},\text{ }k=1,..,N,$$
(48)
where $`\mathrm{\Omega }_\tau `$ is defined by $`\left(\text{6}\right)`$ gives that we have at the point $`\tau _0=0:`$
$$_i\left(\vartheta _j\left(\varphi _k(\tau )\right)\right)|{}_{\tau =\tau _0}{}^{}=$$
$$\vartheta _i\left(_j\left(\varphi _k(\tau )\right)\right)=\left(\left(\frac{^2}{\tau ^i\tau ^j}\left(\frac{}{\tau ^k}\mathrm{\Omega }_\tau \right)\right)\mathrm{}\left(\mathrm{\Omega }_\tau \right)^{}\right).$$
(49)
So $`\left(\text{49}\right)`$ implies $`\left(\text{47}\right).`$ Lemma 41 is proved. $`\mathrm{}`$
###### Corollary 42
$`[D_i,D_j]=[D_{\overline{i}},D_{\overline{j}}]=0.`$
###### Lemma 43
We have $`[_i+t\vartheta _i,_{\overline{j}}+t^1\overline{\vartheta _j}]=0.`$
Proof: We will identify $`T_{0,\left(\text{M}\right)}`$ with $`H^1(\text{M}_0,\mathrm{\Omega }_{\text{M}_0}^2)`$ as in Proposition 23. We assumed that $`\{\varphi _i\}`$ given by $`\left(\text{48}\right)`$ is an orthonormal basis of the tangent space $`T_{0,𝒦}=H^1(\text{M},T_\text{M}^{1,0})`$ at at point $`0\left(\text{M}\right)`$. We have in the flat coordinates $`(\tau ^1,\mathrm{},\tau ^N)`$ that $`\varphi _k(\tau )=(_k\mathrm{\Omega }_\tau )\mathrm{}\eta _\tau ^1).`$ We will need the following Propositions
###### Proposition 44
We have
$$_i\varphi _k,_j\varphi _l=\frac{1}{\left|\lambda \left(\tau \right)\right|^2}\underset{\text{M}}{}\left(\left(\varphi _i\varphi _k\right)\mathrm{}\mathrm{\Omega }_\tau \right)\overline{((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\tau )},$$
(50)
where $`\left(\varphi _i\varphi _k\right)\mathrm{}\mathrm{\Omega }_\tau `$ and $`\left(\varphi _j\varphi _l\right)\mathrm{}\mathrm{\Omega }_\tau `$ are forms of type $`(1,2)`$ on M$`_\tau .`$
Proof: Since $`\varphi _i(\tau )=\left(\left(_i\mathrm{\Omega }_\tau \right)\mathrm{}\eta _\tau ^1\right)`$ we get that
$$_i\varphi _k,_j\varphi _l\left|{}_{\tau =0}{}^{}=\left(_i_k\mathrm{\Omega }_\tau \right)\mathrm{}\eta _\tau ^1,(_j_l\mathrm{\Omega }_\tau )\mathrm{}\eta _\tau ^1\right|{}_{\tau =0}{}^{}.$$
From $`\left(\text{13}\right)`$ we derive that
$$_i_k\mathrm{\Omega }_\tau \left|{}_{\tau =0}{}^{}=(\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0\text{ and }_j_l\mathrm{\Omega }_\tau \right|{}_{\tau =0}{}^{}=(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0$$
are form of type $`(1,2).`$ Thus we get
$$_i\varphi _k,_j\varphi _l\left|{}_{\tau =0}{}^{}=(\frac{1}{\left|\lambda \left(\tau \right)\right|^2}\underset{\text{M}}{}\left((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_\tau \right)\overline{((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_\tau )})\right|{}_{\tau =0}{}^{},$$
where $`\left(\varphi _i\varphi _k\right)\mathrm{}\mathrm{\Omega }_0`$ and $`\left(\varphi _j\varphi _l\right)\mathrm{}\mathrm{\Omega }_0`$ are forms of type $`(1,2).`$ Proposition 44 is proved. $`\mathrm{}`$
###### Proposition 45
We have
$$\vartheta _i(\varphi _k(\tau ))\left|{}_{\tau =0}{}^{}=\left(\left(\mathrm{\Pi }\left((\varphi _i(\tau )\varphi _k(\tau ))\mathrm{}\mathrm{\Omega }_\tau \right)\right)\mathrm{}\left(\mathrm{\Omega }_\tau \right)^1\right)\right|_{\tau =0}$$
and
$$\vartheta _i\varphi _k,\vartheta _j\varphi _l|{}_{\tau =0}{}^{}=$$
$$(\frac{1}{\left|\lambda \left(\tau \right)\right|^2}\underset{\text{M}}{}\left(\mathrm{\Pi }\left((\varphi _i(\tau )\varphi _k(\tau ))\mathrm{}\mathrm{\Omega }_\tau \right)\right)\overline{(\mathrm{\Pi }((\varphi _j(\tau )\varphi _l(\tau ))\mathrm{}\mathrm{\Omega }_\tau ))})|{}_{\tau =0}{}^{},$$
(51)
where $`\mathrm{\Pi }\left(\left(\varphi _i\varphi _k\right)\mathrm{}\mathrm{\Omega }_0\right)`$ and $`\mathrm{\Pi }\left(\left(\varphi _j\varphi _l\right)\mathrm{}\mathrm{\Omega }_0\right)`$ are forms of type $`(2,1).`$
Proof: It follows from the definition of $`\vartheta `$
$$\vartheta _i((_k\mathrm{\Omega }_\tau )\mathrm{}\mathrm{\Omega }_\tau ^1)|{}_{\tau =0}{}^{}=\left(\left(\mathrm{\Pi }(\varphi _i\varphi _k\mathrm{}\mathrm{\Omega }_0)\right)\mathrm{}\mathrm{\Omega }_0^1\right),$$
(52)
where $`\mathrm{\Pi }\left(\varphi _i\varphi _k\mathrm{}\mathrm{\Omega }_\tau \right)`$ is the Poincare dual of the form $`\left(\varphi _i\varphi _k\mathrm{}\mathrm{\Omega }_0\right)`$ of type $`(1,2).`$ Thus $`\mathrm{\Pi }\left(\varphi _i\varphi _k\mathrm{}\mathrm{\Omega }_0\right)`$ is a form of type $`(2,1)`$ and $`\left(\text{52}\right)`$ implies $`\left(\text{51}\right).`$ Proposition 45 is proved. $`\mathrm{}`$
###### Proposition 46
We have the following formula:
$$((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0),\overline{((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0)}=\underset{\text{M}}{}((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0)\overline{((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0)}=$$
$$\mathrm{\Pi }((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0),\overline{\mathrm{\Pi }((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0)}=$$
$$=\underset{\text{M}}{}\mathrm{\Pi }((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0)\overline{\mathrm{\Pi }((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0)}.$$
(53)
Proof: Since $`\{\varphi _i\}`$ is an orthonormal basis in $`T_{0,\left(\text{M}_0\right)}`$ with respect to the W.-P. metric we get that $`\{\varphi _i\mathrm{}\mathrm{\Omega }_0\}`$ is orthonormal basis in $`H^1\left(\text{M}\right)`$ and $`\{\overline{\varphi _i\mathrm{}\mathrm{\Omega }_0}\}`$ is an orthonormal basis in $`H^2(\text{M}_0,\mathrm{\Omega }_{\text{M}_0}^1).`$ Thus
$$\sqrt{1}\varphi _i\mathrm{}\mathrm{\Omega }_0,(\varphi _j\mathrm{}\mathrm{\Omega }_0)=\sqrt{1}\underset{\text{M}}{}(\varphi _i\mathrm{}\mathrm{\Omega }_0)\overline{(\varphi _j\mathrm{}\mathrm{\Omega }_0)}=\delta _{i\overline{j}}$$
(54)
and
$$\sqrt{1}\overline{\varphi _i\mathrm{}\mathrm{\Omega }_0},\overline{\varphi _j\mathrm{}\mathrm{\Omega }_0}=\sqrt{1}\underset{\text{M}}{}\overline{\varphi _i\mathrm{}\mathrm{\Omega }_0}\varphi _j\mathrm{}\mathrm{\Omega }_0=\delta _{i\overline{j}}.$$
(55)
Let
$$(\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0=\underset{\nu =1}{\overset{N}{}}\alpha _\nu \overline{\varphi _\nu \mathrm{}\mathrm{\Omega }_0},\text{ }(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0=\underset{i=1}{\overset{N}{}}\beta _\mu \left(\overline{\varphi _\mu \mathrm{}\mathrm{\Omega }_0}\right),$$
(56)
then
$$\mathrm{\Pi }\left((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0\right)=\underset{\nu =1}{\overset{N}{}}\alpha _\nu \left(\varphi _\nu \mathrm{}\mathrm{\Omega }_0\right)\text{ and }\mathrm{\Pi }\left((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0\right)=\underset{i=1}{\overset{N}{}}\beta _\mu \left(\varphi _\mu \mathrm{}\mathrm{\Omega }_0\right).$$
(57)
Combining $`\left(\text{54}\right),`$ $`\left(\text{55}\right),`$ $`\left(\text{56}\right)`$ and $`\left(\text{57}\right)`$ we get that
$$\sqrt{1}((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0),((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0)=\underset{\nu =1}{\overset{N}{}}\alpha _\nu \beta _\nu $$
(58)
and
$$\sqrt{1}\mathrm{\Pi }((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0),(\mathrm{\Pi }(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0)=\underset{\nu =1}{\overset{N}{}}\alpha _\nu \beta _\nu .$$
(59)
Thus $`\left(\text{58}\right)`$ and $`\left(\text{59}\right)`$ imply Proposition 46. $`\mathrm{}`$
We have
$$[_i+t\vartheta _i,_{\overline{j}}+t^1\overline{\vartheta _j}]\varphi _k,\varphi _l=[_i,_{\overline{j}}]\varphi _k,\varphi _l+[\vartheta _i,\overline{\vartheta _j}]\varphi _k,\varphi _l.$$
(60)
From $`\left(\text{43}\right)`$ we derive that
$$R_{i\overline{j}.k\overline{l}}=[_i,_{\overline{j}}]\varphi _k,\varphi _l=\frac{\sqrt{1}}{\left|\lambda (0)\right|^2}(\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0)),(\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0.$$
(61)
$`\left(\text{51}\right)`$ implies
$$[\vartheta _i,\overline{\vartheta _j}]\varphi _k,\varphi _l=\frac{1}{\left|\lambda (0)\right|^2}\left(\mathrm{\Pi }\left((\varphi _i\varphi _k)\mathrm{}\mathrm{\Omega }_0\right)\right),\overline{(\mathrm{\Pi }((\varphi _j\varphi _l)\mathrm{}\mathrm{\Omega }_0))}.$$
(62)
Combining $`\left(\text{60}\right),`$ $`\left(\text{61}\right),`$ $`\left(\text{62}\right)`$ with $`\left(\text{53}\right)`$ we conclude that $`[_i+t\vartheta _i,_{\overline{j}}+t^1\overline{\vartheta _j}]=0.`$ Theorem 40 is proved. $`\mathrm{}`$
###### Corollary 47
The connection constructed in Theorem 40 is a flat $`𝕊p(2h^{2,1},)`$ connection on the tangent space of the moduli space $`\left(\text{M}\right)`$ of three dimensional CY manifolds. The imaginary form of the Weil-Petersson metric is a parallel form with respect to the CHSV connection.
Proof: It is an well known fact that the imaginary part of the Weil-Petersson metric $`\omega _\tau (1,1)`$ is parallel with respect to the CHSV connection since it is the imaginary part of a Kähler metric. On the other hand $`\omega _\tau (1,1)`$ is just the restriction of the intersection form on $`H^3\left(\text{M}\right)`$ and so it is parallel with respect to the Gauss-Manin connection and so with respect to the connection $`\theta .`$ From here Corollary 47 follows directly. $`\mathrm{}`$
###### Remark 48
It is easy to see that Cecotti-Vafa tt\* equations are exactly the Hitchin-Simpson self duality equations studied in , and .
We will call the connection that we constructed a Cecotti-Hitchin-Simpson-Vafa connection and will denote it as a CHSV connection.
## 4 Review of the Geometric Quantization
### 4.1 Basic Notions of ADW Geometric Quantization
In this paragraph we are going to review the method of the geometric quantization described in and . We will consider a linear space $`𝐖^{2n}`$ with a constant symplectic structure
$$\omega =\frac{1}{2}\omega _{ij}dt^idt^j,$$
(63)
where $`\omega _{ij}`$ is a constant invertible matrix and the $`x^i`$ linear coordinates on $`^{2n}=𝐖.`$ We will denote by $`\omega ^1`$ the matrix inverse to $`\omega `$ and obeying $`\omega _{ij}\left(\omega ^1\right)^{jk}=\delta _i^k.`$
###### Definition 49
The ”prequantum line bundle” is a unitary line bundle $``$ over $`𝐖`$ with a connection whose curvature is $`\sqrt{1}\omega `$. Up to an isomorphism, there is only one such choice of $``$. One can take $``$ to be the trivial unitary line bundle, with a connection given by the covariant derivatives
$$\frac{D}{Dt^i}=\frac{}{t^i}+\frac{\sqrt{1}}{2}\omega _{ij}t^j.$$
(64)
###### Definition 50
We define the $`𝐋^2`$ norm of the sections of $``$ as follows; Let $`h`$ be a positive function on $`𝐖`$ which define a metric on $``$ and $`dd^c\mathrm{log}h=\omega .`$ Then we will say that the $`𝐋^2`$ norm of a section $`\varphi `$ of $``$ is defined as
$$\varphi _{𝐋^2}^2=\left(1\right)^{n(n1)/2}\left(\frac{\sqrt{1}}{2}\right)^n\underset{𝐖}{}\mathrm{exp}\left(h\right)\left|\varphi \right|^2𝑑z^1\mathrm{}dz^n\overline{dz^1\mathrm{}dz^n}.$$
(65)
Then we define the”prequantum Hilbert space” $`_0`$ as the Hilbert space that consists of sections of $``$ with a finite $`𝐋^2`$ norm.
In order to define the quantum Hilbert space, we will introduce the notion of polarization.
###### Definition 51
We will define the polarization as a choice of a complex structure $`J`$ on $`𝐖`$ with the following properties: a. $`J`$ is a translation invariant, so it is defined by a constant matrix $`J_j^i`$ with $`J^2=1`$. b. The two-form $`\omega `$ is of type (1,1) with respect to the complex structure $`J.`$ c. $`J`$ is positive in the sense that the bilinear form g defined by g$`(u,v)=\omega (u,Jv)`$ is strictly positive.
###### Definition 52
Given such a complex structure $`J`$, we define the quantum Hilbert space $`H_J`$ to be space of all holomorphic functions $`\varphi _J(z^1,\mathrm{},z^n)`$ on $`𝐖`$ with respect to the complex structure $`J`$ with a finite $`𝐋^2`$ norm.
It is well known that the Heisenberg group of $`𝐖`$ has an irreducible projective representation in $`_J.`$ Thus each such representation of the Heisenberg group of $`𝐖`$ depends on the choice of the complex structure $`J`$ of $`𝐖.`$ We want to construct a projectively flat connection on the infinite dimensional dimensional vector space over the parameter space of the complex structures of $`𝐖.`$ Construction of such a connection enables one to identify all the irreducible projective representation of the Heisenberg group in $`_J.`$
### 4.2 Siegel Space
We will introduce some notations following . First of all, one has the projection operators $`\frac{1}{2}(1\sqrt{1}J)`$ on $`W`$<sup>1,0</sup> and on $`\overline{𝐖^{1,0}}=`$$`W`$$`^{0,1},`$ where
$$𝐖^{1,0}:=\{u𝐖|Ju=\sqrt{1}u\}\text{ .}$$
We know that the vector space $`W`$ with a complex structure $`J`$ can be identified as a complex vector space canonically with the spaces $`W`$<sup>1,0</sup> and $`W`$<sup>0,1</sup> by the maps
$$u\frac{1}{2}(1\sqrt{1}J)u.$$
We need to write down explicitly in fixed coordinates these two identifications. We will follow the notations from in the above identifications. For any vector $`v=(\mathrm{},v^i,\mathrm{})`$, we denote by
$$v^{\underset{¯}{i}}=\frac{1}{2}(1\sqrt{1}J)_j^iv^j\text{ and }v^{\overline{i}}=\frac{1}{2}(1+\sqrt{1}J)_j^iv^j.$$
For one forms $`w=(\mathrm{},w_i,\mathrm{})`$ we have
$$w_{\underset{¯}{j}}=\frac{1}{2}(1\sqrt{1}J)_j^iw_i\text{ and }w_{\overline{j}}=\frac{1}{2}(1+\sqrt{1}J)_j^iw_j.$$
Thus $`J_{\underset{¯}{j}}^{\underset{¯}{i}}=\sqrt{1}\delta _{\underset{¯}{j}}^{\underset{¯}{i}}`$ and $`J_{\overline{j}}^{\overline{i}}=\sqrt{1}\delta _{\overline{j}}^{\overline{i}}.`$ This means that the projections of $`J_j^i`$ and $`\delta _j^i`$ on $`W`$<sup>1,0</sup> and $`W`$<sup>0,1</sup> are proportional.
Let $``$ be the space of all $`J`$ obeying the conditions in Definition 51. Then $``$ is a symmetric space, i.e. $`=𝕊p(2n,)/𝕌(n).`$ $``$ is called Siegel space. It is a well known fact that we have the following realization of $``$ $`=𝕊p(2n,)/𝕌(n)`$ as a tube domain;
$$=\{Z|Z\text{ }is\text{ }a\text{ }(n\times n)\text{ complex matrix such that }Z^t=Z\text{ }and\text{ }\mathrm{Im}Z>0\}.$$
$``$ has a natural complex structure, defined as follows. The condition $`J^2=1`$ implies that for a first order variation $`\delta J`$ of $`J`$, one must have
$$J\delta J+\delta JJ=0$$
This means that the non-zero projections of $`\delta J`$ on (1,0) vectors and (0,1) vectors are $`\delta J_{\overline{j}}^{\underset{¯}{i}}`$ and $`\delta J_{\underset{¯}{j}}^{\overline{i}}`$. We define the complex structure on $``$ by declaring $`\delta J_{\overline{j}}^{\underset{¯}{i}}`$ to be of type $`(1,0)`$ and $`\delta J_{\underset{¯}{j}}^{\overline{i}}`$ to be of type $`(0,1)`$. Notice that the projection of $`\delta J`$ on $`(0,1)`$ vectors is a map from $`(1,0)`$ vectors to $`(0,1)`$ vectors. So the $`(1,0)`$ part of $`\delta J`$ is the Beltrami differential as it was defined in Section 2 of this article.
### 4.3 Construction of Witten Projective Connection
Over $``$ we introduce two Hilbert space bundles. One of them, say $`^0`$, is the trivial bundle $`\times _0`$, where $`_0`$ is the Hilbert space of all function $`\psi (t^i;J)`$ with finite $`𝐋^2`$ norm. The definition of $`_0`$ is independent of $`J`$. The second is the bundle $`^Q`$, whose fibre over a point $`J`$ are functions $`\psi (t^i;J)`$ of $`t^i`$, holomorphic in the complex structure defined by $`J,`$ i.e. the following equation holds:
$$\frac{D}{Dt^{\overline{i}}}\psi (t^i;J)=0.$$
This equation has a dependence on $`J`$ coming from the projection operators used in defining $`t^{\overline{i}}.`$
A connection on the bundle $`^0`$ restricts to a connection on $`^Q`$ if and only if its commutator with $`D_{\overline{i}}`$ is a linear combination of the $`D_{\overline{j}}`$. Since $`^0`$ is defined as a product bundle $`\times _0`$, there is a trivial connection $`\delta `$ on this bundle
$$\delta :=\underset{i,j}{}dJ_j^i\frac{}{J_j^i}.$$
(66)
We can expand $`\delta `$ in $`(1,0)`$ and $`(0,1)`$ pieces, $`\delta =\delta ^{(1,0)}+\delta ^{(0,1)}`$, with
$$\delta ^{(1,0)}=\underset{i,j}{}dJ_{\overline{j}}^{\underset{¯}{i}}\frac{}{J_{\overline{j}}^{\underset{¯}{i}}}.$$
(67)
Unfortunately, as it was shown in the commutator of $`\delta ^{(1,0)}`$ with $`D_{\overline{i}}`$ is not a linear combination of the $`D_{\overline{j}}`$. So one needs to modify the trivial connection so that its commutator with $`D_{\overline{i}}`$ will be a linear combination of the $`D_{\overline{j}}`$
###### Definition 53
Witten defined in the following connection $``$ on the bundle $`^0=\times _0:`$
$$^{(1,0)}:=\delta ^{(1,0)}\frac{1}{4}(dJ\omega ^1)^{\underset{¯}{i}\underset{¯}{j}}\frac{D}{Dt^{\underset{¯}{i}}}\frac{D}{Dt^{\underset{¯}{j}}}\text{ and }^{(0,1)}:=\delta ^{(0,1)}.$$
(68)
Witten proved the following theorem in :
###### Theorem 54
A. The connection $``$ descends to a connection on $`^\mathrm{Q}`$. B. The curvature of the connection $``$ on $`^\mathrm{Q}`$ is of type $`(1,1)`$ and it is equal to
$$[^{(0,1)},^{(1,0)}]=[\delta ^{(0,1)},\frac{1}{4}(dJ\omega ^1)^{\underset{¯}{i}\underset{¯}{j}}\frac{D}{Dt^{\underset{¯}{i}}}\frac{D}{Dt^{\underset{¯}{j}}}]=\frac{1}{8}dJ_{\overline{k}}^{\underset{¯}{i}}dJ_{\underset{¯}{j}}^{\overline{k}}.$$
(69)
This theorem shows that the curvature of $``$ is not zero, even when it is restricted to $`^Q`$. The curvature is a c-number, that is, it depends only on $`J`$, and not on the variables $`t^i`$ that are being quantized. The fact that the curvature of $``$ is a c-number means that parallel transport by $``$ is unique up to a scalar factor which, moreover, it is of modulus 1 since the curvature is real or more fundamentally since $``$ is unitary. So up to this factor one can identify the various $`_J`$ ’s, and regard them as a different realization of the quantum Hilbert space $``$.
## 5 Geometric Quantization of Moduli Space of CY Manifolds and Quantum Background Independence
### 5.1 Symplectic Structures and the Flat Coordinates
Our goal is to quantize geometrically the cotangent bundle of the moduli space $`\left(\text{M}\right)`$. This means to define a flat $`𝕊p(2h^{2,1},)`$ connection on the cotangent bundle of $`\left(\text{M}\right)`$. We have done this in the Section 3. Then will construct the prequantum line bundle. Once the prequatum line bundle is constructed we will compute explicitly the projective connection defined in and . After the geometric quantization is done we will derive BCOV holomorphic anomaly equations established in as the projective flat connections on the Hilbert bundle $`^Q`$. To perform these computations, it is important to fix first the flat symplectic structure on the tangent bundle $`𝒯_{\left(\text{M}\right)}`$, then the local coordinates on $`\left(\text{M}\right)`$ that describe the change of the complex structures on the CY manifold M<sub>τ</sub> and the linear coordinates in $`W`$$`{}_{\tau }{}^{}=\mathrm{\Omega }_{\tau ,\left(\text{M}\right)}^1.`$
###### Remark 55
We will define the symplectic form $`\omega _1(\tau )`$ on the vector bundle
$$R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2\left(\text{M}\right)$$
as the imaginary part of the metric
$$g_{i,\overline{j}}:=\frac{}{\tau ^i}\mathrm{\Omega }_\tau ,\frac{}{\tau ^j}\mathrm{\Omega }_\tau =\sqrt{1}\underset{\text{M}}{}\frac{}{\tau ^i}\mathrm{\Omega }_\tau \overline{\frac{}{\tau ^j}\mathrm{\Omega }_\tau }\text{ .}$$
(70)
The symplectic structure on the tangent space $`𝒯_{\left(\text{M}\right)}`$ is defined by the imaginary part of the Weil-Petersson metric, i.e.:
$$\omega (\tau )=\mathrm{Im}G|{}_{T_{\tau ,𝒦}}{}^{},$$
(71)
where
$$G_{i,\overline{j}}:=\left(\frac{}{\tau ^i}\mathrm{\Omega }_\tau \right)\mathrm{}\eta _\tau ^1,\left(\frac{}{\tau ^j}\mathrm{\Omega }_\tau \right)\mathrm{}\eta _\tau ^1_{WP}=\frac{\sqrt{1}}{\eta _\tau ^2}\underset{\text{M}}{}\frac{}{\tau ^i}\mathrm{\Omega }_\tau \overline{\frac{}{\tau ^j}\mathrm{\Omega }_\tau }\text{ .}$$
(72)
Thus $`\left(\text{72}\right)`$ implies the following relations between the two symplectic structures:
$$g_{i,\overline{j}}=e^KG_{i,\overline{j}},$$
(73)
where
$$\eta _\tau ^2=\sqrt{1}\underset{\text{M}}{}\eta _\tau \overline{\eta _\tau }=\mathrm{exp}(K).$$
(74)
We know from Remark 48 that the forms $`\eta _\tau `$ are parallel with respect to the Cecotti-Hitchin-Simpson-Vafa connection. Thus we can identified the symplectic structures of the tangent bundle by the flat $`𝕊p(2h^2,)`$ Cecotti-Hitchin-Simpson-Vafa connection.
###### Remark 56
We know that if we fix a basis of the orthogonal vectors $`\{\varphi _{i,\tau }\}`$ in the holomorphic tangent space
$$T_{\tau ,𝒦}=H^1(\text{M}_\tau ,T_{\text{M}_{gt}}^{1,0})$$
then we obtain a linear coordinate system $`(\tau ^1,\mathrm{},\tau ^N)`$ in the dual $`\mathrm{\Omega }_{\tau ,𝒦}^1`$ of $`T_{\tau ,𝒦}.`$ $`\mathrm{\Omega }_{\tau ,𝒦}^1`$ can be canonically identified with $`T_{\tau ,𝒦}`$ by using the parallel symplectic form. According to the results obtained in then linear coordinate system $`(\tau ^1,\mathrm{},\tau ^N)`$ defined by the choice of the orthogonal vectors $`\{\varphi _{i,\tau _0}\}`$ in the holomorphic tangent space
$$T_{\tau _\tau ,𝒦}=H^1(\text{M}_{\tau _0},T_{\text{M}_{\tau _0}}^{1,0})$$
defines the same local coordinate system $`(\tau ^1,\mathrm{},\tau ^N)`$ in $`𝒦`$ since
$$𝒦H^1\left(\text{M,}T_{\text{M}_{\tau _0}}^{1,0}\right).$$
See Theorem 6.
###### Remark 57
In we construct in a canonical way a family of holomorphic form $`\mathrm{\Omega }_\tau `$ given by $`\left(\text{6}\right)`$. Then the family of holomorphic form $`\mathrm{\Omega }_\tau `$ defines a choice of a basis
$$\{\varphi _{i,\tau }\mathrm{}\mathrm{\Omega }_\tau \}\text{ and }\{\overline{(\varphi _{i,\tau }\mathrm{}\mathrm{\Omega }_\tau )}\}$$
(75)
of the complexified tangent space: $`H^1\left(\text{M,}\mathrm{\Omega }_\text{M}^2\right)H^2\left(\text{M,}\mathrm{\Omega }_\text{M}^1\right).`$ Thus the flat local coordinate system $`(\tau ^1,\mathrm{},\tau ^N)`$ is defined by a choice of the basis in the holomorphic tangent space $`T_{\tau ,𝒦}`$ defines in a canonical way a coordinate system on $`H^1\left(\text{M}\right)`$ which is the same as $`(\tau ^1,\mathrm{},\tau ^N)`$ in case we choose the basis $`\left(\text{75}\right).`$ We identified the dual of $`H^1\left(\text{M,}\right)`$ with $`H^1\left(\text{M}\right)`$ by using the parallel symplectic form.
###### Remark 58
The complex structure $`J_\tau `$ on $`T_{\tau ,\left(\text{M}\right)}`$ is defined by the complex structure on $`H^1\left(\text{M,}T_\text{M}^{1,0}\right).`$ Indeed by local deformation theory we have
$$T_{\tau ,\left(\text{M}\right)}^{1,0}=H^1\left(\text{M,}T_\text{M}^{1,0}\right).$$
On CY manifold $`H^1(`$M$`{}_{\tau }{}^{},T_\tau ^{1,0})`$ can be identified with $`H^1\left(\text{M,}\mathrm{\Omega }_\text{M}^2\right)`$ by contraction with the non zero holomorphic form $`\mathrm{\Omega }_\tau .`$ Thus the complex structure operator $`J_\tau `$ on the CY manifolds acts on
$$T_{\tau ,\left(\text{M}\right)}^{1,0}=H^1\left(\text{M,}T_\text{M}^{1,0}\right)$$
in a natural way as follows
$$J((dz^i)(J\left(dz^j\right)(J\left(\overline{dz}^k\right))=\sqrt{1}(dz^idz^j\overline{dz}^k)$$
since $`J(dz^i)=\sqrt{1}dz^i.`$
###### Remark 59
We are using two different identifications of $`H^1(`$M$`{}_{\tau }{}^{},T_{\text{M}_\tau }^{1,0})`$ with $`H^1\left(\text{M,}\mathrm{\Omega }_\text{M}^2\right)`$ by using using the contraction with the families of the two holomorphic $`3`$forms $`\eta _\tau `$ and $`\mathrm{\Omega }_\tau .`$ Recall that $`\eta _\tau `$ was defined globally by Theorem 15. We obtained two coordinate systems $`(\tau ^1,\mathrm{},\tau ^N)`$ and $`(t^1,\mathrm{},t^N)`$ on $`H^1\left(\text{M,}\mathrm{\Omega }_\text{M}^2\right).`$ The relation between them is given by the relations between $`\mathrm{\Omega }_\tau `$ and $`\eta _\tau ,`$ i.e. the two coordinate systems are proportional with the coefficients of proportionality $`\lambda (\tau ),`$ where $`\mathrm{\Omega }_\tau =\lambda (\tau )\eta _\tau .`$
Next we will define the prequantum line bundles over $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)\text{/}\left(\text{M}\right)}^2`$ and $`𝒯_{\left(\text{M}\right)}`$.
###### Theorem 60
The prequantum line bundles on the tangent bundle $`𝒯_{\left(\text{M}\right)}`$ is $`p^{}\left(\pi _{}\left(\omega _{𝒴(\text{M)/}(\text{M)}}^{}\right)\right),`$ where $`p:𝒯_{\left(\text{M}\right)}\left(\text{M}\right),`$ $`\omega _{𝒴(\text{M)/}(\text{M)}}`$ is the relative dualizing line bundle and the metric on its fibre is defined by the $`L^2`$ norm of the holomorphic three form. The Chern class of $`\pi ^{}\left(\omega _{𝒴(\text{M) /}(\text{M)}}^{}\right)`$ is given by the restriction of
$$\omega (\tau ):=\{\text{Imaginary Part of the }WP\text{ metric}\}$$
on $`T_{\tau ,\left(\text{M}\right)}.`$
Proof: We proved in that the natural metric
$$\mathrm{\Omega }_\tau ^2=\sqrt{1}\underset{\text{M}}{}\mathrm{\Omega }_\tau \overline{\mathrm{\Omega }_\tau }$$
on $`\omega _{𝒴\text{ }\left(\text{M}\right)\text{ /}\left(\text{M}\right)}`$ is such that its Chern form is
$$\sqrt{1}\frac{^2}{\tau ^i\overline{\tau ^j}}\mathrm{log}(\mathrm{\Omega }_\tau ^2)=\omega (\tau ).$$
This shows that $`=\pi ^{}(\omega _{𝒴(\text{M) /}(\text{M)}}^{})`$ is the prequantum line bundle on the tangent bundle $`𝒯_{\left(\text{M}\right)}`$. This proves Theorem 60. $`\mathrm{}`$
###### Corollary 61
The prequantum line bundles on the bundle $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ is $`p^{}\left(\pi _{}\left(\omega _{𝒴\left(\text{M}\right)\text{ /}\left(\text{M}\right)}^{}\right)\right),`$ where
$$p:R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2\left(\text{M}\right),\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}$$
is the relative dualizing line bundle and the metric on its fibre is defined by the $`L^2`$ norm of the holomorphic three form. The Chern class of $`\pi ^{}\left(\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^{}\right)`$ is given by the restriction of $`\omega _1(\tau )`$ on $`W`$$`{}_{\tau }{}^{}=H^1\left(\text{M,}T_\text{M}^{1.0}\right).`$
### 5.2 The Geometric Quantization of $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$
We will compute the projective flat connection on the Hilbert space bundle $`^Q`$ over the vector bundle $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)\text{/}\left(\text{M}\right)}^2.`$ The prequantum line bundle is the line bundle $`p^{}\left(\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^{}\right)`$ over $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)\left(\text{M}\right)}^2.`$
###### Theorem 62
The matrix of the operator $`(dJ)`$ in the basis
$$\left\{\frac{}{\tau ^i}\mathrm{\Omega }_\tau ,\text{ }i=1,\mathrm{},N\right\}$$
is given on each fibre of $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ by,
$$(dJ)_{\underset{¯}{a}}^{\overline{b}}=\underset{c,d}{}C_{\underset{¯}{a}\underset{¯}{c}\underset{¯}{d}}g^{d,\overline{b}}.$$
(76)
Proof: The proof of Theorem 62 follows directly from the following Lemma:
###### Lemma 63
We have
$$\left(\frac{}{\tau ^i}J\right)\left(\frac{}{\tau ^j}\mathrm{\Omega }_\tau \right)=\underset{k,l=1}{\overset{N}{}}C_{ijk}g^{k,\overline{l}}\overline{\frac{}{\tau ^l}\mathrm{\Omega }_\tau },\text{ }$$
(77)
where
$$C_{ijk}=\sqrt{1}\underset{\text{M}}{}\left(\frac{^2}{\tau ^i\tau ^j}\mathrm{\Omega }_\tau \right)\left(\frac{}{\tau ^k}\mathrm{\Omega }_\tau \right)|{}_{\tau =0}{}^{}=$$
$$\sqrt{1}\underset{\text{M}}{}\left(\mathrm{\Omega }_\tau \right)\left(\frac{^3}{\tau ^i\tau ^j\tau ^k}\mathrm{\Omega }_\tau \right)|{}_{\tau =0}{}^{}.$$
(78)
The idea of the proof of formula $`\left(\text{77}\right)`$ is the following one; we know that in the basis
$$\left\{\frac{}{\tau ^i}\mathrm{\Omega }_\tau ,\text{ }i=1,\mathrm{},N\right\}$$
of $`H^1(`$M$`{}_{\tau }{}^{},\mathrm{\Omega }_\tau ^2)`$ the complex structure operator is given by the matrix:
$$\left(\begin{array}{cc}\sqrt{1}I_{h^{1,2}}\hfill & 0\hfill \\ 0\hfill & \sqrt{1}I_{h^{1,2}}\hfill \end{array}\right).$$
(79)
Thus $`\left(\text{79}\right)`$ implies
$$\left(\frac{}{\tau ^i}J\right)\left(\frac{}{\tau ^j}\mathrm{\Omega }_\tau \right)\left|{}_{\tau =0}{}^{}=J\left(\frac{^2\mathrm{\Omega }_\tau }{\tau ^i\tau ^j}\right)\right|{}_{\tau =0}{}^{}.$$
So we need to compute the expression of the vectors $`\left\{\frac{^2}{\tau ^i\tau ^j}\mathrm{\Omega }_\tau ,\text{ }i,j=1,\mathrm{},N\right\}`$ as a linear combinations of
$$\left\{\overline{\frac{}{\tau ^i}\mathrm{\Omega }_\tau },\text{ }i=1,\mathrm{},N\right\}.$$
(80)
We know from that$`\left\{\frac{^2\mathrm{\Omega }_\tau }{\tau ^i\tau ^j}\right\}H^2\left(\text{M}\right)`$ where $`\tau =(\tau ^1,\mathrm{},\tau ^N)`$ are the flat coordinates. Therefore if we express the vectors
$$\left\{\frac{^2\mathrm{\Omega }_\tau }{\tau ^i\tau ^j}\right\}H^2(\text{M}_\tau ,\mathrm{\Omega }_{\text{M}_\tau }^1)$$
as linear combination of the basis $`\left(\text{80}\right),`$ we will get explicitly the matrices of the operators $`\frac{}{\tau ^i}J_{\tau =0},`$ $`i=0,\mathrm{},N.`$ From the explicit formulas of the operators $`\frac{}{\tau ^i}J|_{\tau =0}`$ we will get the formula $`\left(\text{77}\right)`$.
Proof: We know that
$$\sqrt{1}\underset{\text{M}}{}\left(\frac{}{\tau ^i}\mathrm{\Omega }_\tau \right)\left(\overline{\frac{}{\tau ^j}\mathrm{\Omega }_\tau }\right)|{}_{\tau =0}{}^{}=\delta _{i,\overline{j}}.$$
(81)
By using the expression $`\left(\text{6}\right)`$ for $`\mathrm{\Omega }_\tau `$ and the natural $`L^2`$ metric on $`H^2\left(\text{M}\right)`$ we get the following expression of the vector $`\left(\frac{^2}{\tau ^i\tau ^j}\mathrm{\Omega }_\tau \right)|_{\tau =0}`$ in the orthogonal basis
$$\left\{\overline{\frac{}{\tau ^k}\mathrm{\Omega }_\tau }\right|{}_{\tau =0}{}^{},\text{ }k=1,\mathrm{},N\}$$
of $`H^{2,1}\left(\text{M}\right):`$
$$\left(\frac{^2}{\tau ^i\tau ^j}\mathrm{\Omega }_\tau \right)|_{\tau =0}=\sqrt{1}\underset{k=1}{\overset{N}{}}\left(\frac{^2}{\tau ^i\tau ^j}\mathrm{\Omega }_\tau ,\overline{\frac{}{\tau ^k}\mathrm{\Omega }_\tau }\right)\left(\overline{\frac{}{\tau ^k}\mathrm{\Omega }_\tau }\right)|{}_{\tau =0}{}^{}=$$
$$\sqrt{1}\underset{k=1}{\overset{N}{}}(\underset{\text{M}}{}\left(\frac{^2}{\tau ^i\tau ^j}\mathrm{\Omega }_\tau \right)\left(\frac{}{\tau ^k}\mathrm{\Omega }_\tau \right))\left(\overline{\frac{}{\tau ^k}\mathrm{\Omega }_\tau }\right)|{}_{\tau =0}{}^{}=$$
$$\sqrt{1}\underset{k=1}{\overset{N}{}}C_{ijk}g^{k,\overline{l}}\left(\overline{\frac{}{\tau ^l}\mathrm{\Omega }_\tau }\right)|_{\tau =0}$$
(82)
Formula $`\left(\text{82}\right)`$ implies that for any $`\tau 𝒦`$ we have
$$\frac{^2}{\tau ^i\tau ^j}\mathrm{\Omega }_\tau =\sqrt{1}\underset{k=1}{\overset{N}{}}C_{ijk}g^{k,\overline{l}}\left(\overline{\frac{}{\tau ^l}\mathrm{\Omega }_\tau }\right).$$
(83)
We know that
$$J\left(\frac{}{\tau ^i}\mathrm{\Omega }_\tau \right)=\sqrt{1}\left(\frac{}{\tau ^i}\mathrm{\Omega }_\tau \right).$$
(84)
Combining $`\left(\text{83}\right)`$ and $`\left(\text{84}\right)`$ we get that
$$\frac{}{\tau ^j}\left(J\left(\frac{}{\tau ^i}\mathrm{\Omega }_\tau \right)\right)=\frac{}{\tau ^j}\left(\sqrt{1}\left(\frac{}{z^i}\mathrm{\Omega }_\tau \right)\right)=$$
$$\sqrt{1}\left(\frac{^2}{\tau ^jz^i}\mathrm{\Omega }_\tau \right)=\underset{k=1}{\overset{N}{}}C_{ijk}g^{k,\overline{l}}\left(\overline{\frac{}{\tau ^l}\mathrm{\Omega }_\tau }\right).$$
(85)
Lemma 63 is proved. $`\mathrm{}`$
Lemma 63 implies directly Theorem 62. $`\mathrm{}`$
###### Corollary 64
The projective flat connection on the Hilbert space bundle $`^Q`$ over the vector bundle $`R^1\pi _{}\mathrm{\Omega }_{𝒳\text{/}𝒦}^2`$ is given by
$$\frac{}{\tau ^a}+\frac{\sqrt{1}}{4}\overline{C_{abc}}g^{\overline{b},i}g^{\overline{c},j}\frac{D}{Dt^i}\frac{D}{Dt^j}.$$
(86)
Proof: According to Theorem 62 on $`R^1\pi _{}\mathrm{\Omega }_{𝒳\text{/}𝒦}^2`$ we have
$$\left(dJ\omega _1^1\right)^{ij}D_iD_j=\sqrt{1}\underset{a=1}{\overset{N}{}}\overline{C_{abc}}g^{\overline{b},i}g^{\overline{c},j}D_iD_j.$$
(87)
Thus the projective connection on $`R^1\pi _{}\mathrm{\Omega }_{𝒳\text{/}𝒦}^2`$ is given by $`\left(\text{86}\right)`$. $`\mathrm{}`$
### 5.3 Computations on the Tangent Bundle of the Moduli Space
To quantize geometrically the tangent bundle $`𝒯_{\left(\text{M}\right)}`$ on the moduli space $`\left(\text{M}\right)`$ means to compute explicitly the prequantum line bundle and then projective connection on the Hilbert space bundle $`^Q`$ related to the prequantum line bundle on the tangent bundle $`𝒯_{\left(\text{M}\right)}`$ of the moduli space $`\left(\text{M}\right).`$
###### Theorem 65
We have on $`𝒯_{\left(\text{M}\right)}`$
$$(dJ)_i^{\overline{l}}=e^K\underset{j,k,l=1}{\overset{N}{}}C_{ijk}G^{k,\overline{l}}.$$
(88)
Proof: Formula $`\left(\text{88}\right)`$ follows directly from formula $`\left(\text{87}\right)`$, the globally defined isomorphism $`\iota _\tau :H^1\left(\text{M,}T_\text{M}^{1,0}\right)H^1\left(\text{M,}\mathrm{\Omega }_\text{M}^2\right),`$ given by $`\varphi \varphi \mathrm{}(\eta _\tau ),`$ and the relation $`g_{i,\overline{j}}=e^KG_{i,\overline{j}},`$ where $`e^K`$ is given by $`\left(\text{73}\right).`$ Theorem 65 is proved. $`\mathrm{}`$
### 5.4 BCOV Anomaly Equations
###### Definition 66
The following data, a. The moduli space $`\left(\text{M}\right)`$ of CY threefolds, b. The CHSV $`𝕊p(2h^{2,1},)`$ flat connection on the tangent bundle constructed in , c. The ”prequantized line bundle” $`\pi ^{}(\omega _{𝒴\left(\text{M}\right)\text{/}\left(\text{M}\right)}^{})`$ on the tangent bundle $`𝒯_{\left(\text{M}\right)}`$, d. The imaginary part of the Weil-Petersson metric $`\omega `$ and e. the bundle of the Hilbert spaces $`^Q`$ over the tangent bundle $`𝒯_{\left(\text{M}\right)}`$ with a projective flat connection on it, will be called a CY quantum system.
###### Theorem 67
The expression of Witten projective flat connection as defined in Definition 53 in the flat coordinates for the CY quantum system defined in Definition 66 coincides with the BCOV anomaly equations $`\left(\text{1}\right)`$ in .
Proof: In order to prove Theorem 67 we need to compute explicitly the Witten projective connection constructed in on the Hilbert vector bundle $`^Q`$ over the tangent bundle of the moduli space $`\left(\text{M}\right)`$. The explicit formula $`\left(\text{68}\right)`$ and since BCOV anomaly equations were established in the flat coordinate system $`(\tau ^1,\mathrm{},\tau ^N)`$ imply that we need to compute the expression of
$$(dJ\omega ^1)^{ij}D_iD_j$$
on $`𝒯_{\left(\text{M}\right)}`$ in the same coordinate system. As we pointed out the flat coordinate system $`(\tau ^1,..,\tau ^N)𝒦`$ introduced on the basis of Theorem 6 is same coordinate system used in .
We already established in Theorem 62 the explicit expression of $`(dJ\omega _1^1)^{ij}D_iD_j`$ on $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ in the coordinates $`(\tau ^1,\mathrm{},\tau ^N).`$ We will establish first the local expression of $`(dJ\omega ^1)^{ij}D_iD_j`$ on $`𝒯_{\left(\text{M}\right)}`$ in the coordinate system $`(t^1,\mathrm{},t^N)`$ defined by the identification $`𝒯_{\left(\text{M}\right)}`$ with $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ that uses the parallel section $`\eta _\tau ^1`$ of $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^{}.`$ Then we will compute $`(dJ\omega ^1)^{ij}D_iD_j`$ on $`𝒯_{\left(\text{M}\right)}`$ in the flat coordinate $`(\tau ^1,\mathrm{},\tau ^N)`$ defined by the identification of $`𝒯_{\left(\text{M}\right)}`$ with $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ by the holomorphic tensor $`\mathrm{\Omega }_\tau ^1.`$ Thus we will establish BCOV anomaly equations $`\left(\text{1}\right)`$ in as the projective Witten connection. Theorem 67 will be proved.
In order to compute $`\left(dJ\omega _{1,\tau }^1D_iD_j\right)`$ we need to established the relations between $`g^{\overline{i},j}`$ and $`G^{\overline{i},j}.`$ $`\left(\text{73}\right)`$ implies that these relations given by
$$g^{\overline{i},j}=e^KG^{\overline{i},j},$$
(89)
where $`e^K=\eta _\tau ^2.`$
Let $`(t^1,\mathrm{},t^N)`$ be the complex linear coordinates on $`T_{\tau ,\left(\text{M}\right)}`$ defined by the identification of $`𝒯_{\left(\text{M}\right)}`$ with $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^{}R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ that uses that parallel section $`\eta `$ of $`\omega _{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^{}.`$ Remark 59 and the relation $`\mathrm{\Omega }_\tau =\lambda (\tau )\eta _\tau `$ imply that we have
$$(\tau ^1,..,\tau ^N)=\lambda (t^1,..,t^N).$$
(90)
We know that the expression of $`\left(dJ\omega _1^1\right)^{ij}D_iD_j`$ on the fibre of $`R^1\pi _{}\mathrm{\Omega }_{𝒳\text{/}𝒦}^2`$ is given by $`\left(\text{86}\right)`$ in the coordinates $`(\tau ^1,..,\tau ^N).`$ Thus combining $`\left(\text{89}\right)`$ with $`\left(\text{86}\right)`$ we get that on $`T_{\tau ,\left(\text{M}\right)}`$ in the coordinates $`(t^1,..,t^N)`$ we have:
$$\left(dJ\omega ^1\right)^{ij}D_iD_j=\sqrt{1}e^{2K}\underset{a=1}{\overset{N}{}}\overline{C_{abc}}G^{\overline{b},i}G^{\overline{c},j}D_iD_j.$$
(91)
From $`\left(\text{91}\right)`$ and $`\left(\text{68}\right)`$ we get that the projective connection on the tangent bundle of $`\left(\text{M}\right)`$ is given by
$$\frac{}{\tau ^a}+\frac{\sqrt{1}}{4}e^{2K}\overline{C_{abc}}G^{\overline{b},i}G^{\overline{c},j}\frac{D}{Dt^i}\frac{D}{Dt^j}$$
(92)
By using the symplectic identifications of the fibres of $`𝒯_{\left(\text{M}\right)}`$ by using the flat $`𝕊p(2h^{2,1},)`$ and then by using the projective flat connection on the Hilbert space fibration $`^Q`$ we obtain that the quantum state represented by a parallel vector $`\mathrm{\Psi }(\tau ,t)`$ is independent of $`t`$. This means that $`\left(\text{92}\right)`$ implies that if $`\mathrm{\Psi }(\tau ,t)`$ is independent of $`t`$ then $`\mathrm{\Psi }(\tau ,t)`$ satisfies the following equations:
$$\left(\frac{}{\tau ^a}+\frac{\sqrt{1}}{4}e^{2K}\overline{C_{abc}}G^{\overline{b},i}G^{\overline{c},j}\frac{D}{Dt^i}\frac{D}{Dt^j}\right)\mathrm{\Psi }(\tau ,t)=0$$
(93)
and
$$\frac{}{\overline{t^k}}\mathrm{\Psi }(\tau ,t)=0.$$
(94)
Based on $`\left(\text{90}\right)`$ we have
$$\frac{D}{Dt^i}=\lambda \frac{D}{D\tau ^i}.$$
Thus formulas $`\left(\text{93}\right)`$ and $`\left(\text{94}\right)`$ can be written as follows:
$$\left(\frac{}{\tau ^a}+\frac{\sqrt{1}}{4}\lambda ^2e^{2K}\overline{C_{abc}}G^{\overline{b},i}G^{\overline{c},j}\frac{D}{D\tau ^i}\frac{D}{D\tau ^j}\right)\mathrm{\Psi }(\tau )=0$$
(95)
and
$$\frac{}{\overline{\tau ^k}}\mathrm{\Psi }(\tau )=0.$$
(96)
As we pointed out the flat coordinate system $`(\tau ^1,..,\tau ^N)𝒦`$ introduced on the basis of Theorem 6 is same coordinates are use in . Thus the BCOV anomaly equations $`\left(\text{1}\right)`$ in are the same as the equations $`\left(\text{95}\right)`$ and $`\left(\text{96}\right).`$
Direct computations show that if $`F_g`$ satisfy the equations $`\left(\text{1}\right)`$ then
$$\mathrm{\Psi }(\tau )=\mathrm{exp}\left(\underset{g=1}{\overset{\mathrm{}}{}}\lambda ^{2g2}F_g\right)$$
satisfy $`\left(\text{95}\right)`$. Theorem 67 is proved. $`\mathrm{}`$
### 5.5 Comments
* Bershadsky, Cecotti, Ooguri and Vafa used for the free energy $`Z(\lambda ;t,\overline{t})`$ the following expression:
$$𝔉(\lambda ;t,\overline{t}):=\underset{g=1}{\overset{\mathrm{}}{}}\lambda ^{2g2}F_g\text{ and }Z(\lambda ;t,\overline{t})=\mathrm{exp}\left(𝔉(\lambda ;t,\overline{t})\right).$$
(97)
Compare this with the expression $`\left(\text{2}\right)`$ for $`Z,`$ i.e. $`Z=\mathrm{exp}\left(\frac{1}{2}𝔉(\lambda ;t,\overline{t})\right)`$ used by Witten in .
* It is proved in by using physical arguments that $`g`$-genus partition function $`F_g`$ satisfy the equation $`\left(\text{1}\right)`$. The holomorphic anomaly equation $`(\mathbf{3.8})`$ derived in is
$$\left(\overline{\frac{}{t^i}}\overline{\frac{}{t^i}}F_1\right)\mathrm{exp}𝔉(\lambda ;t,\overline{t})=\frac{\lambda ^2}{2}\overline{C}_{ijk}e^{2K}G^{i,\overline{j}}G^{k,\overline{k}}\widehat{D}_j\widehat{D}_k\mathrm{exp}𝔉(\lambda ;t,\overline{t}),$$
(98)
where
$$\widehat{D}_j𝔉(\lambda ;t,\overline{t})=\underset{g}{}\lambda ^{2g2}D_jF_g=$$
$$\underset{g}{}\lambda ^{2g2}\left(_j\left(2g2\right)_jK\right)F_g=$$
$$\left(_j_jK\lambda _\lambda \right)𝔉(\lambda ;t,\overline{t}).$$
Thus the equation $`\left(\text{98}\right)`$ is different by term involving $`F_1`$ from the equation $`\left(\text{95}\right)`$.
* In the authors used the normalized holomorphic form, namely they normalized $`\mathrm{\Omega }_\tau `$ in such a way that $`{\displaystyle \underset{\gamma }{}}\mathrm{\Omega }_\tau =1,`$ where $`\gamma `$ is the invariant vanishing cycle. This normalized form is the same as the form defined in 9. They showed by using string theory that the functions $`F_g`$ ”count” curves of genus g. It seems to me that it is a very deep mathematical fact.
## 6 The Extended Period Space of CY Manifolds.
### 6.1 Definition of the Extended Period Space and Basic Properties
In this paragraph we will study the extended period space $`𝔥_{2,2h^{2,1}}(H^3(`$M$`,),`$ which parametrizes all possible filtrations of the type:
$$F^0=H^{3,0}F^1=H^{3,0}+H^{2,1}+H^{1,2}F^2=H^3(\text{M},𝐂),$$
where $`dimF^0=1`$ plus some extra properties which are motivated from the above filtration and Variations of Hodge Structures on K3 surfaces.
We will use the following notation for the cup product for $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ $`H^3(`$M,$`)`$, i.e. $`\mathrm{\Omega }_1,\mathrm{\Omega }_2={\displaystyle \underset{\text{M}}{}}\mathrm{\Omega }_1\mathrm{\Omega }_2.`$
###### Definition 68
$`𝔥_{2,2h^{2,1}}`$ by definition is the set of lines in $`H^3(`$M,$`)`$ spanned by the cohomology class $`[\mathrm{\Omega }]`$ of the holomorphic three form in $`H^3(`$M,$`)`$ such that
$$\sqrt{1}\underset{\text{M}}{}\mathrm{\Omega }\overline{\mathrm{\Omega }}=\sqrt{1}\mathrm{\Omega },\overline{\mathrm{\Omega }}>0.$$
(99)
###### Theorem 69
There is a one to one map between the points $`\tau `$ of $`𝔥_{2,2h^{1,2}}`$ and the two dimensional oriented planes $`E_\tau `$ in $`H^3`$(M,$`)`$, where $`E_\tau `$ is generated by $`\gamma _1`$ and $`\mu _1`$ such that $`\gamma _1,\mu _1=1.`$
Proof: Let $`[\mathrm{\Omega }_\tau ]H^3`$(M,$``$) be a non-zero vector corresponding to a point $`\tau 𝔥_{2,2h^{1,2}}.`$ Since the class of cohomology $`[\mathrm{\Omega }_\tau ]`$ satisfies $`\left(\text{99}\right)`$ we may chose $`\mathrm{\Omega }_\tau `$ such that
$$\sqrt{1}\mathrm{\Omega }_\tau ,\overline{\mathrm{\Omega }_\tau }=2.$$
(100)
Then $`\left(\text{99}\right)`$ and $`\left(\text{100}\right)`$ imply$`{\displaystyle \underset{\text{M}}{}}\mathrm{Im}\mathrm{\Omega }_\tau \mathrm{Re}\mathrm{\Omega }_\tau =\mathrm{Re}\mathrm{\Omega }_\tau ,\mathrm{Im}\mathrm{\Omega }_\tau =1.`$ We will define $`E_\tau `$ to be the the two dimensional oriented subspace in $`H^3`$(M,$``$) spanned by $`\mathrm{Re}\mathrm{\Omega }_\tau `$ and $`\mathrm{Im}\mathrm{\Omega }_\tau .`$ So to each point $`\tau 𝔥_{2,2h^{1,2}}`$ we have assigned an oriented two plane $`E_\tau `$ in $`H^3`$(M,$`).`$
Suppose that $`E`$ is a two dimensional oriented plane in $`H^3`$(M,$``$) spanned by $`\gamma `$ and $`\mu H^3`$(M,$``$) be such that $`\gamma ,\mu =1.`$ Let $`\mathrm{\Omega }_E:\mu +\sqrt{1}\gamma .`$ Then we have
$$\sqrt{1}\underset{\text{M}}{}\mathrm{\Omega }_E\overline{\mathrm{\Omega }_E}=\sqrt{1}\mathrm{\Omega }_E,\overline{\mathrm{\Omega }_E}=2\gamma ,\mu =2.$$
So to the plane $`E`$ we assign the line in $`H^3`$(M,$``$) spanned by $`\mathrm{\Omega }_E.`$ This proves Theorem 69. $`\mathrm{}`$
###### Corollary 70
$`𝔥_{2,2h^{1,2}}`$ is an open set in $`\mathrm{𝐆𝐫𝐚𝐬𝐬}(2,2h^{1,2}+2).`$ So it has a complex dimension $`2h^{1,2}.`$
###### Remark 71
It is easy to see that each point $`\tau 𝔥_{2,2h^{1,2}}`$ defines a natural filtration of length two in $`H^3`$(M,$``$). Indeed, let $`H_\tau ^{3,0}`$ be the subspace in $`H^3`$(M,$``$) spanned by a non-zero element $`\mathrm{\Omega }_\tau F^0.`$ Let $`\gamma _0=\mathrm{Re}\mathrm{\Omega }_\tau `$ and $`\mathrm{Im}\mathrm{\Omega }_\tau =\mu _0.`$ From Theorem 69 we know that $`\gamma _0,\mu _0`$ is a positive number. Let
$$\{\gamma _0,\mu _0,..,\gamma _{h^{1,2}},\mu _{h^{1,2}}\}$$
be a symplectic basis of $`H^3`$(M,$``$) such that $`\gamma _i,\mu _j=\delta _{ij}.`$ Let $`\mathrm{\Omega }_i:=\mu _i+\sqrt{1}\gamma _i.`$ We will define $`H_\tau ^{2,1}`$ to be the subspace in $`H^3(`$M,$`)`$ spanned by the vectors $`\mathrm{\Omega }_i`$ for $`i=1,..,h^{1,2}.`$ Then we define $`H_\tau ^{1,2}:=\overline{H_\tau ^{2,1}}.`$ It is easy to see that
$$\left(H_\tau ^{2,1}+H_\tau ^{1,2}\right)^{}=H_\tau ^{3,0}+H_\tau ^{0,3},$$
where $`H_\tau ^{0,3}:=\overline{H_\tau ^{3,0}}`$ and the orthogonality is with respect to
$$\omega _1,\overline{\omega _2}=\underset{\text{M}}{}\omega _1\overline{\omega _2}.$$
The natural filtration in $`H^3`$(M,$``$) that corresponds to $`\tau 𝔥_{2,2h^{1,2}}`$ is defined as follows:
$$F_\tau ^0=H_\tau ^{3,0}F_\tau ^1=H_\tau ^{3,0}+H_\tau ^{2,1}+H_\tau ^{1,2}H^3(M,).$$
(101)
Next we will introduce a metric G on $`H^3(`$M,$`𝐂).`$ We will use the metric G to show that the filtration defined by 101 is a Hodge filtration of weight two.
###### Definition 72
Let M be a fixed CY manifold and let
$$\mathrm{\Omega }=\mathrm{\Omega }^{3,0}+\mathrm{\Omega }^{2,1}+\mathrm{\Omega }^{1,2}+\mathrm{\Omega }^{0,3}H^3(\text{M},𝐂)$$
be the Hodge decomposition of $`\mathrm{\Omega }`$, then $`G(\mathrm{\Omega },\overline{\mathrm{\Omega }})`$ is defined as follows:
$$G(\mathrm{\Omega },\overline{\mathrm{\Omega }}):=\sqrt{1}\left(\underset{\text{M}}{}\mathrm{\Omega }^{3,0}\overline{\mathrm{\Omega }^{3,0}}+\underset{\text{M}}{}\mathrm{\Omega }^{2,1}\overline{\mathrm{\Omega }^{2,1}}\right)+$$
$$\sqrt{1}\left(\underset{\text{M}}{}\overline{\mathrm{\Omega }^{1,2}}\mathrm{\Omega }^{1,2}+\underset{\text{M}}{}\overline{\mathrm{\Omega }^{0,3}}\mathrm{\Omega }^{0,3}\right).$$
(102)
From the definition of the metric, it follows that it has a signature $`(2,2h^{2,1})`$ on H<sup>3</sup>(M,$``$). We will denote the quadratic form of this metric by Q.
###### Lemma 73
The metric defined by $`\left(\text{102}\right)`$ does not depend on the choice of the complex structure on M.
Proof: Let M<sub>0</sub> and M<sub>τ</sub> be two different complex structures on M. Let $`\mathrm{\Omega }_0`$ and $`\mathrm{\Omega }_\tau `$ be two non zero holomorphic three forms on M<sub>0</sub> and M<sub>τ</sub> respectively. Let $`\{\mathrm{\Omega }_{0,i}\}`$ and $`\{\mathrm{\Omega }_{\tau ,i}\}`$ be two bases of $`H^{2,1}(`$M$`{}_{0}{}^{})`$ and $`H^{2,1}(`$M$`{}_{\tau }{}^{})`$ respectively, where $`i=1,..,h^{2,1}`$ such that $`\mathrm{\Omega }_{0,i},\overline{\mathrm{\Omega }_{0,j}}=\mathrm{\Omega }_{\tau ,i},\overline{\mathrm{\Omega }_{\tau ,j}}=\delta _{ij}.`$ Then one see immediately that
$$\{\mathrm{Re}\mathrm{\Omega }_0,\mathrm{Im}\mathrm{\Omega }_0,\mathrm{},\mathrm{Re}\mathrm{\Omega }_{0,i},\mathrm{Im}\mathrm{\Omega }_{0,i},..\}and\text{ }\{\mathrm{Re}\mathrm{\Omega }_\tau ,\mathrm{Im}\mathrm{\Omega }_\tau ,..,\mathrm{Re}\mathrm{\Omega }_{\tau ,i},\mathrm{Im}\mathrm{\Omega }_{\tau ,i},..\}$$
are two different symplectic bases of $`H^3(`$M,$``$). So there exists an element g$`𝕊p(2h^{2,1}+2)`$ such that
$$\text{g}(\mathrm{\Omega }_0)=\mathrm{\Omega }_\tau \text{ and g}(\mathrm{\Omega }_{0,i})=\mathrm{\Omega }_{\tau ,i}.$$
(103)
So $`\left(\text{103}\right)`$ implies that
$$\text{g}(H^{2,1}(\text{M}_0))=H^{2,1}(\text{M}_\tau )$$
(104)
and
$$\text{g}(H^{3,0}(\text{M}_0))=H^{3,0}(\text{M}_\tau )$$
(105)
Lemma 73 follows directly from $`\left(\text{104}\right)`$, $`\left(\text{105}\right)`$ and the definition of the metric G. $`\mathrm{}`$
###### Theorem 74
Let $`𝔥_{2,2h^{1,2}}`$ be the space that parametrizes all filtrations in $`H^3(`$M$`,)`$ defined in Remark 71. Then these filtration are Hodge filtrations of weight two and their moduli space is isomorphic to the symmetric space
$$𝔥_{2,2h^{1,2}}:=𝕊𝕆_0(2,2h^{1,2})/𝕊𝕆(2)\times 𝕊𝕆(2h^{1,2}).$$
Proof: Since the signature of the metric G is $`(2,2h^{2,1})`$ the proof of Theorem 74 is standard and follows directly from definition of the variations of Hodge structures of weight two. See for example . $`\mathrm{}`$
We will use the fact that the space $`𝔥_{2,2h^{1,2}}=𝕊𝕆_0(2,2h^{1,2})/𝕊𝕆(2)\times 𝕊𝕆(2h^{1,2})`$ is as an open set in the Grassmannian $`Grass`$$`(2,b_3)`$ identified with all two dimensional oriented subspaces in $`H^3(\text{M},)`$ such that the restriction of the quadratic form Q on them is positive, i.e.
$$𝔥_{2,h^{1,2}}:=\{EH^3(\text{M},)dimE=2,Q|{}_{E}{}^{}>0+orientation\}.$$
###### Definition 75
We will define a canonical map from $`𝔥_{2,2h^{1,2}}Gr(2,2h^{1.2}+2)`$ to $`(H^3(\text{M},))`$ as follows; let $`E_\tau 𝔥_{2,2h^{1,2}},`$ i.e. $`E_\tau `$ is an oriented two dimensional subspace in $`H^3(\text{M},)`$ on which the restriction of Q is positive. Let $`e_1`$ and $`e_2`$ be an orthonormal basis of $`E_\tau `$. Let $`\mathrm{\Omega }_\tau :=e_1+\sqrt{1}e_2.`$ Then $`\mathrm{\Omega }_\tau `$ defines a point $`\tau (H^3(`$M$`,)))`$ that corresponds to the line in $`H^3(\text{M},))`$ spanned by $`\mathrm{\Omega }_\tau .`$ It is a standard fact that the points $`\tau `$ in $`(H^3(\text{M},)))`$ is such that Q($`\tau ,\tau `$)$`=0`$ and Q($`\tau ,\overline{\tau })>0`$ are in one to one corresponds with the points in $`𝔥_{2,h^{1,2}}.`$ See . This follows from the arguments used in the proof of Theorem 69 or see .
It is a well known fact that $`𝔥_{2,h^{1,2}}`$ is isomorphic to one of the irreducible component of the open set of the quadric in $`(H^3(\text{M},))`$ defined as follows:
$$𝔥_{2,2h^{1,2}}\{\tau (H^3(\text{M},))|Q(\tau ,\tau )=0\text{ }and\text{ }Q(\tau ,\overline{\tau })>0\}.$$
(See .)
We will consider the family $`𝒳\times \overline{𝒳}𝒦\times \overline{𝒦},`$ where the family $`\overline{𝒳}\overline{𝒦}`$ is the family that corresponds to the conjugate complex structures, i.e. the point $`(\tau _1,\overline{\tau _2})𝒦\times 𝒦`$ $`\tau _1`$ corresponds to the complex structure M$`_{\tau _1}`$ and $`\overline{\tau _2}`$ corresponds to the $`\overline{\text{M}_{\tau _2}},`$ where $`\overline{\text{M}_{\tau _2}}`$ is the conjugate complex structure on M$`_{\tau _2}.`$
We will define the period map
$$p:𝒦\times \overline{𝒦}(H^3(\text{M},))$$
(106)
as follows; to each point $`(\tau ,\upsilon )𝒦\times \overline{𝒦}`$ we will assign the complex line $`p(\tau ,\upsilon )`$ in $`H^3(\text{M},)`$ defined by the oriented two plane $`E_{\tau ,\upsilon }H^3(`$M$`,)`$ spanned by $`\mathrm{Re}(\mathrm{\Omega }_\tau +\overline{\mathrm{\Omega }_\upsilon })`$ and $`\mathrm{Im}(\mathrm{\Omega }_\tau \mathrm{\Omega }_\upsilon )`$, where $`\mathrm{\Omega }_\tau `$ and $`\mathrm{\Omega }_\nu `$ are defined as in Theorem 7. $`𝒦`$ is the Kuranishi space defined in Definition 10. We will show that the analogue of local Torelli Theorem holds, i.e. we will show that the period map $`p`$ is a local embedding $`p:𝒦\times \overline{𝒦}𝔥_{2,2h^{1,2}}.`$
###### Remark 76
We will define an embedding of the Kuranishi family $`𝒦`$ into $`𝒦\times \overline{𝒦}`$ as follows; to each $`\tau 𝒦`$ we will associate the complex structure ($`I_\tau ,I_\tau )`$ on M$`\times `$M.
###### Theorem 77
The period map $`p`$ defined by $`\left(\text{106}\right)`$ is a local isomorphism. Moreover the image $`p(𝒦\times \overline{𝒦})`$ is contained in $`𝔥_{2,2h^{1,2}}\left(H^3(\text{M,})\right)`$ for small enough $`\epsilon .`$
Proof: The fact that the period map $`p`$ is a local isomorphism follows directly from the local Torelli theorem for CY manifolds proved in . Let $`(\tau ,\overline{\tau })𝒦\times \overline{𝒦},`$ then clearly the point $`p(\tau ,\overline{\tau })𝔥_{2,2h^{1,2}}(H^3(\text{M},)),`$ i.e. Q$`|`$$`{}_{E_{\tau ,\overline{\tau }}}{}^{}>0.`$ This follows directly from the definition of Q and the fact that $`E_{\tau ,\overline{\tau }}`$ is the subspace in $`H^3(`$M$`{}_{\tau }{}^{},)`$ spanned by $`\mathrm{Re}\mathrm{\Omega }_\tau `$ and $`\mathrm{Im}\mathrm{\Omega }_\tau .`$
Let $`(\tau _1,\overline{\tau _2})𝒦\times \overline{𝒦}`$ be a point ”close” to the diagonal $`\mathrm{\Delta }`$ in $`𝒦\times \overline{𝒦}`$ then Q$`|{}_{E_{\tau ,\upsilon }}{}^{}>0.`$ Indeed this follows from the fact that the condition Q$`|{}_{E_{\tau ,\upsilon }}{}^{}>0`$ is an open one on $`Gr(2,2h^{1,2}+2)`$. Then the two dimensional oriented space $`E_{(\tau _1,\overline{\tau _2})}H^3(\text{M},)`$ spanned by $`\{\mathrm{Re}\mathrm{\Omega }_{\tau _1}+\mathrm{Re}\mathrm{\Omega }_{\tau _2},\mathrm{Im}\mathrm{\Omega }_{\tau _1}+\mathrm{Im}\mathrm{\Omega }_{\tau _2}\},`$ where $`\mathrm{\Omega }_{\tau _1}`$ and $`\mathrm{\Omega }_{\tau _2}`$ are defined by formula $`\left(\text{6}\right)`$ in Theorem 7 will be such that Q$`|{}_{E_{\tau ,\upsilon }}{}^{}>0`$. From here we deduce that $`p(\tau _1,\overline{\tau _2})𝔥_{2,2h^{1,2}}.`$ Theorem 77 is proved. $`\mathrm{}`$
We defined the Kuranishi space $`𝒦`$ to be a open polydisk $`|\tau ^i|<\epsilon `$ for $`i=1,\mathrm{},N`$ in $`H^1\left(\text{M}\right).`$ Since $`p`$ is a local isomorphism we may assume that $`𝒦\times \overline{𝒦}`$ is contained in $`𝔥_{2,2h^{1,2}}`$ for small enouph $`\epsilon .`$
### 6.2 Construction of a $``$ Structure on the Tangent Space of $`\left(\text{M}\right)`$
###### Definition 78
To define a $``$ structure on a complex vector space $`V`$ means the construction of a free abelian group $`AV`$ such that the rank of $`A`$ is equal to the dimension of V, i.e. $`A=V.`$
###### Definition 79
We define $`𝔥_{2,2h^{1,2}}()`$ as follows; A point $`\tau 𝔥_{2,2h^{1,2}}()`$ if the two dimensional oriented subspace $`E_\tau =H_\tau ^{3,0}+H_\tau ^{0,3}`$ that corresponds to $`\tau `$ constructed in Theorem 69 is such that $`E_\tau H^3(`$M,$`).`$
###### Theorem 80
$`𝔥_{2,2h^{1,2}}()`$ is an everywhere dense subset in $`𝔥_{2,2h^{1,2}}.`$
Proof: Our claim follows directly from two facts. The first one is that the set of the points in $`Gr(2,2+2h^{2,1})`$ that corresponds to two dimensional subspaces in $`H^3(`$M,$`)`$ form an everywhere dense subset in $`Gr(2,2+2h^{2,1})`$ and the second one is that $`𝔥_{2,2h^{2,1}}`$ is an open set in $`Gr(2,2+2h^{2,1}).`$ Theorem 80 is proved. $`\mathrm{}`$
###### Theorem 81
For each $`\tau 𝔥_{2,2h^{2,1}}()`$ a natural $``$ structure is defined on the tangent space T$`_{\tau ,𝔥_{2,2h^{2,1}}}`$ at the point $`\tau 𝔥_{2,2h^{2,1}}.`$ This means that there exists a subspace $`^{2h^{2,1}}H^3(`$M$`,)`$ such that $`H_\tau ^{2,1}+H_\tau ^{1,2}T_{\tau ,𝔥_{2,2h^{2,1}}}=^{2h^{2,1}}.`$
Proof: From the theory of Grassmannians we know that $`T_{\tau ,𝔥_{2,2h^{2,1}}}`$ can be identified with $`H_\tau ^{2,1}+H_\tau ^{1,2}.`$ Our corollary follows directly from the construction of $`H_\tau ^{2,1}+H_\tau ^{1,2}`$ described in Remark 71. Indeed the point $`\tau 𝔥_{2,2h^{2,1}}()`$ defines two vectors $`\gamma _0`$ and $`\mu _0H^3(`$M$`,)/Tor`$ that form a basis of $`H^{3,0}\left(\text{M}\right)H^{0,3}\left(\text{M}\right)`$ such that $`\mu _0,\gamma _0`$ and $`\mu _0,\gamma _0>0.`$ We choose the vectors
$$\{\gamma _0,\mu _0,\gamma _1,\mu _1,..,\gamma _{h^{2,1}},\mu _{h^{2,1}}\}$$
to be in $`H^3(`$M,$``$)/Tor and we require that $`\mu _i,\gamma _j=\delta _{ij}.`$ Then, from the way we defined $`H_\tau ^{2,1}`$ and $`H_\tau ^{1,2},`$ it follows that
$$H_\tau ^{2,1}+H_\tau ^{1,2}=(\gamma _1\mu _1\mathrm{}\gamma _{h^{2,1}}\mu _{h^{2,1}}).$$
Theorem 81 is proved. $`\mathrm{}`$
We will consider the embedding of $`𝒦`$ in $`𝒦\times \overline{𝒦}`$ defined in Remark 76. Next we choose a point $`\kappa 𝒦\times \overline{𝒦}𝔥_{2,2h^{2,1}}(H^3(`$M$`,))`$ such that $`\kappa (𝒦\times \overline{𝒦})𝔥_{2,2h^{2,1}}().`$ We know that $`\kappa `$ corresponds to a two dimensional space $`E_\kappa H^3(`$M,$`)`$, with the additional condition, that there exists vectors $`\gamma _0`$ and $`\mu _0H^3(`$M,$``$)/Tor that span $`E_\kappa `$ and $`\mu _0,\gamma _0=1`$. The existence of such points follows from the fact that the set of all two dimensional space, $`E_\kappa H^3(`$M$`,)`$ such that there exists vectors $`\gamma _0`$ and $`\mu _0H^3(`$M$`,)/Tor`$ and $`\mu _0,\gamma _0=1`$ is an everywhere dense subset in $`𝔥_{2,2h^{2,1}}`$. Let
$$\{\gamma _0,\mu _0,\gamma _1,\mu _1,..,\gamma _{h^{2,1}},\mu _{h^{2,1}}\}H^3(\text{M},)/Tor$$
be such that $`\gamma _0`$ and $`\mu _0`$ span $`E_\kappa `$ and $`\mu _i,\gamma _j=\delta _{ij}.`$ It follows from the construction in Remark 71 that the vectors
$$\gamma _1,\mu _1,..,\gamma _{h^{2,1}},\mu _{h^{2,1}}H^3(\text{M},)/Tor$$
span $`H_\kappa ^{2,1}+H_\kappa ^{1,2},`$ i.e. they span the tangent space $`T_{\kappa ,𝔥_{2,2h^{2,1}}}=H_\kappa ^{2,1}+H_\kappa ^{1,2}.`$ We know from Corollary 47 that there exists an $`𝕊p(2h^{2,1},)`$ flat connection on $`𝒦`$ and so we define a flat connection on the product $`𝒦\times \overline{𝒦}`$ as the direct sum of the two connections. Using this $`𝕊p(4h^{2,1},)`$ flat connection we can perform a parallel transport of the vectors $`\gamma _{1,}\mu _1,..,\gamma _{h^{2,1}},\mu _{h^{2,1}}T_\kappa `$ to a basis
$$\gamma _{1,\tau },\mu _{1,\tau },..,\gamma _{h^{2,1},\tau },\mu _{h^{2,1},\tau }$$
in the tangent space$`H_\tau ^{2,1}+H_\tau ^{1,2}T_{\tau ,𝔥_{2,2h^{2,1}}}`$ to each point $`\tau 𝒦𝒦\times \overline{𝒦}.`$ Thus we can conclude that $`\mu _{i,\tau },\gamma _{j,\tau }=\delta _{ij}`$ and the free abelian group
$$A_\tau :=\gamma _{1,\tau }\mu _{1,\tau }\mathrm{}\gamma _{h^{2,1},\tau }\mu _{h^{2,1},\tau }$$
in $`H_\tau ^{2,1}+H_\tau ^{1,2}`$ is such that
$$H_\tau ^{2,1}+H_\tau ^{1,2}=(\gamma _{1,\tau }\mu _{1,\tau }\mathrm{}\gamma _{h^{2,1},\tau }\mu _{h^{2,1},\tau }).$$
So we defined for each $`\tau 𝒦`$ an abelian subgroup $`A_\tau `$ T$`{}_{\tau ,𝒦}{}^{}=(H_\tau ^{2,1}+H_\tau ^{1,2})`$ such that
$$A_\tau =T_{\tau ,𝒦}=(H_\tau ^{2,1}+H_\tau ^{1,2})$$
and $`\gamma ,\mu `$ for $`\gamma `$ and $`\mu A_\tau `$.
###### Definition 82
The image of the projection of the abelian subgroup $`A_\tau `$ of T$`{}_{\tau ,𝒦}{}^{}=(H_\tau ^{2,1}+H_\tau ^{1,2})`$ to $`H_\tau ^{2,1}`$ will be denoted by $`\mathrm{\Lambda }_\tau `$ for each $`\tau \left(\text{M}\right)`$.
###### Theorem 83
There exists a holomorphic map $`\varphi `$ from the moduli space $`\left(\text{M}\right)`$ of CY manifolds to the moduli space of principally polarized abelian varieties $`𝕊p(2h^{2,1},)\backslash _{h^{2,1}}`$, where $`_{h^{2,1}}:=𝕊p(2h^{2,1},)/U(h^{2,1}).`$
Proof: From 82 we know that there exists a lattice $`\mathrm{\Lambda }_\tau `$ T<sub>τ,K</sub> such that the restriction of the imaginary part of the Weil-Petersson metric
$$\mathrm{Im}(g)(u,v)=u,v$$
on $`\mathrm{\Lambda }_\tau `$ is such that $`u,v`$ and $`|det\gamma _i,\gamma _j|=1`$ for any symplectic basis of $`\mathrm{\Lambda }_\tau .`$ Thus over $`𝒦`$ we can construct a family of principally polarized abelian varieties
$$𝒜_K𝒦.$$
(107)
In fact we constructed a family of principally polarized abelian varieties
$$𝒜\left(\text{M}\right)$$
(108)
over the moduli space $`\left(\text{M}\right)`$ since CHSV connection is a flat connection globally defined over $`\left(\text{M}\right)`$. This means that we defined the holomorphic map $`\varphi `$ between the quasi-projective varieties
$$\varphi :\left(\text{M}\right)𝕊p(2h^{2,1},)\backslash _{h^{2,1}}.$$
The existence of $`\varphi `$ follows from the fact that there exists a versal family of principally polarized abelian varieties
$$𝔄𝕊p(2h^{2,1},)\backslash _{h^{2,1}}.$$
Theorem 83 is proved. $`\mathrm{}`$
Notice that the family of principally polarized varieties $`𝒜\left(\text{M}\right)`$ is constructed by using the vector bundle $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2.`$ Using the identification between $`𝒯_{\left(\text{M}\right)}`$ and $`R^1\pi _{}\mathrm{\Omega }_{𝒴\left(\text{M}\right)/\left(\text{M}\right)}^2`$ given by $`\varphi \varphi \mathrm{}\eta _\tau ,`$ we define a family of principally polarized varieties isomorphic to the family $`𝒜\left(\text{M}\right).`$
### 6.3 Holomorphic Symplectic Structure on the Extended Period Domain
###### Definition 84
Let us fix a symplectic basis
$$\{\gamma _0,\gamma _1,..,\gamma _{h^{1,2}};\upsilon _0,\mathrm{},\upsilon _{h^{1,2}}\}$$
in $`H^3(`$M,$``$)/Tor, i.e. $`\gamma _i,\upsilon _j=\delta _{ij}.`$ This basis defines a coordinate system in $`(H^3(`$M,$`))`$ which we will denote by $`(z_0:..:z_{2h^{1,2}}).`$On the open set:
$$U_0:=\{(z_0:..:z_{2h^{1,2}+1})|z_00\},$$
we define a holomorphic one forms:
$$\alpha _0:=dt_{h^{1,2}+1}+t_1dt_{h^{1,2}+2}+..+t_{h^{1,2}}dt_{2h^{1,2}+1},$$
(109)
where $`t_i=\frac{z_i}{z_0},`$ for $`i=1,..,h^{1,2}.`$ Let us restrict $`\alpha _0`$ on $`𝔥_{2,2h^{1,2}}U_0`$ and denote this restriction by $`\alpha _0.`$ In the same way we can define the forms $`\alpha _i`$ on the open set $`U_i:=\{(z_0:..:z_{2h^{1,2}})|`$ $`z_i0\}.`$
###### Theorem 85
There exists a closed holomorphic non degenerate two form $`\psi `$ on $`𝔥_{2,2h^{1,2}}`$ such that
$$\psi \left|{}_{U_i𝔥_{2,2h^{1,2}}}{}^{}=d\alpha _i\right|{}_{U_i𝔥_{2,2h^{1,2}}}{}^{}.$$
(110)
Proof: The proof of Theorem 85 is based on the following Proposition:
###### Proposition 86
We have
$$d\alpha _i=d\alpha _j$$
(111)
on $`𝔥_{2,2h^{1,2}}(U_iU_j).`$ Thus there exists a holomorphic non degenerate form $`\psi `$ such that $`\psi |{}_{U_i}{}^{}=d\alpha _i.`$
Proof: It is easy to see that since the extended period domain $`𝔥_{2,2h^{1,2}}`$ is an open set of a quadric in $`(H^3(`$M,$`))`$ then that the tangent space $`T_{\tau ,𝔥_{2,2h^{1,2}}}`$ to any point $`\tau 𝔥_{2,2h^{1,2}}`$ can be identified with the orthonormal complement $`(H_\tau ^{3,0}+\overline{H^{3,0}})^{}H^3(`$M$`,)`$ with respect to the metric G defined in Definition 72. Let $`\upsilon `$ and $`\mu T_{\tau ,𝔥_{2,2h^{1,2}}}H^3(`$M,$`)`$. From the definition of the form $`d\alpha _i,`$ it follows that
$$d\alpha _i(\upsilon ,\mu )=\upsilon ,\mu ,$$
(112)
where $`\upsilon ,\mu `$ is the symplectic form defined by the intersection form on $`H^3(`$M,$`)`$. From here Proposition 86 follows directly. $`\mathrm{}`$
Theorem 85 follows directly from Proposition 86. $`\mathrm{}`$
###### Corollary 87
The holomorphic two form $`\psi `$ is a parallel form when restricted to $`𝒦\times 𝒦𝔥_{2,2h^{1,2}}`$ with respect to the CHSV connection. (See Definition 38.)
Proof: The corollary follows directly from Definition 38 and the fact that the imaginary part of the Weil-Petersson metric when restricted to the tangent space $`T_{\tau ,𝒦}=H^{2,1}H^3(`$M$`,)`$ is just the restriction of intersection form $`\upsilon ,\mu `$ on $`H^3(`$M,$`)`$. $`\mathrm{}`$
## 7 Algebraic Integrable System on the Moduli Space of CY Manifolds.
###### Definition 88
Let N be an algebraic variety. An algebraic integrable system is a holomorphic map $`\pi :X`$N where a. X is a complex symplectic manifold with holomorphic symplectic form $`\psi \mathrm{\Omega }^{2,0}(X)`$; b. The fibres of $`\pi `$ are compact Lagrangian submanifolds, hence affine tori; c. There is a family of smoothly varying cohomology classes $`[\rho _n]H^{1,1}(X_n)H^2(X_n,),nN,`$ such that \[$`\rho _n`$\] is a positive polarization of the fibre $`X_n`$. Hence $`X_n`$ is an abelian tersor, i.e. on $`X_n`$ we do not have a point which represents zero to define a structure of a group on $`X_n`$. See .
This notion is the complex analogue of completely integrable (finite dimensional) systems in classical mechanics was introduced by R. Donagi and E. Markman in . We will show that the family $`𝒜\left(\text{M}\right)`$ as defined in Definition 82 is an algebraic integrable system in the sense of Donagi-Markrman.
###### Theorem 89
The holomorphic family $`𝒜\left(\text{M}\right)`$ defines an algebraic integrable system on the moduli space of three dimensional CY manifolds $`\left(\text{M}\right)`$ in the sense of R. Donagi and Markman.
Proof: We must check properties a, b and c stated in Definition 88. In order to check property a and b, we need to construct a non-degenerate closed holomorphic two form $`\mathrm{\Omega }_1`$ on the cotangent space $`T^{}𝒦\left(\text{M}\right)`$. The cotangent space $`T_\tau ^{}\left(\text{M}\right)`$ at a point $`\tau 𝒦\left(\text{M}\right)`$ can be identified with $`H^{1,2}(`$M$`{}_{\tau }{}^{})`$ by contraction with $`\overline{\mathrm{\Omega }_\tau },`$ where $`\mathrm{\Omega }_\tau `$ is a holomorphic three form such that
$$\sqrt{1}\underset{\text{M}}{}\mathrm{\Omega }_\tau \overline{\mathrm{\Omega }_\tau }=1.$$
Then the local Torelli theorem for CY manifolds shows that the restriction of the symplectic form $`\upsilon ,\mu `$ defined by the intersection form on $`H^3(`$M,$`)`$/Tor by formulas $`\left(\text{110}\right)`$ and $`\left(\text{112}\right)`$ will give a globally defined holomorphic two form $`\mathrm{\Omega }_1`$ on the cotangent bundle of $`𝒦\left(\text{M}\right)`$. Then the properties a and b as stated in are obvious.
Next we will construct the smoothly varying cohomology classes $`[\rho _\tau ]`$ which fulfill property c. Let $`\rho (1,1)`$ be the imaginary form of the Weil-Petersson metric on $`\left(\text{M}\right)`$. Since we proved in that the potential of the Weil-Petersson metric is defined from a metric on the relative dualizing line bundle of the family $`𝒳\left(\text{M}\right)`$, we deduce that $`\rho (1,1)`$ is a smoothly varying cohomology class of type $`(1,1)`$. From here we deduce that for each $`\tau \left(\text{M}\right)`$, $`[\rho _\tau ]H^{1,1}(𝒜_\tau ,)`$ and $`[\rho _\tau ]`$ varies smoothly. We need to show that $`[\rho _\tau ]H^2(𝒜_\tau ,).`$ This statement is equivalent to saying that if $`\nu `$ and $`\mu `$ are any two vectors in the lattice
$$\mathrm{\Lambda }_\tau H^{2,1}(\text{M}_\tau )H^{1,2}(\text{M}_\tau ),$$
then $`\rho _\tau (\nu ,\mu )`$ and if $`\gamma _1,..,\gamma _{2h^{1,2}}`$ is a $``$-basis of the lattice $`\mathrm{\Lambda }_\tau `$, then $`det(\rho _\tau (\gamma _i,\gamma _j))=1.`$ We proved that the imaginary part of the Weil-Petersson metric is a parallel with respect to the Cecotti-Hitchin-Simpson-Vafa connection. (See Remark 48.) We used the Cecotti-Hitchin-Simpson-Vafa parallel transport to define the $``$ structure on
$$T_{\tau ,𝒦𝒦\times 𝒦}=H^{2,1}\left(\text{M}\right)+H^{1,2}\left(\text{M}\right).$$
From here it follows that the number $`[\rho _\tau ](\nu ,\mu )`$ is equal to the cup product of the parallel transport of the vectors $`\nu `$ and $`\mu `$ at a point $`(\tau ,\upsilon )𝔥_{2,2h^{1,2}}(),`$ which is an integer. Exactly the same arguments show that
$$det(\rho _\tau (\gamma _i,\gamma _j))=1.$$
So the family $`𝒜\left(\text{M}\right)`$ fulfills properties b and c. Our Theorem is proved. $`\mathrm{}`$
###### Corollary 90
On the tangent bundle of the moduli space of polarized CY threefolds there exists a canonical Hyper-Kähler metric.
Proof: Cor. 90 follows directly from . $`\mathrm{}`$
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# Chiral symmetry and properties of hadron correlators in matter
## Acknowledgments
The author wants to thank Mike Birse for enlightening discussions concerning the role of chiral symmetry in nuclei. Useful discussions with T.Cohen, L.Kisslinger, T.Hatsuda, M. Johnson and A.Thomas are gratefully acknowledged. The author is grateful for the support from Center for Subatomic Structure of Matter at the University of Adelaide where the final part of this work was done.
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# hep-ph/0004217RIKEN-BNL preprint Double transverse spin asymmetries in vector boson production
## I Introduction
Double transverse spin asymmetries in high energy hadron-hadron collisions have attracted much theoretical attention (starting with the early investigations ), but no experimental studies have been performed so far. At the BNL Relativistic Heavy Ion Collider (RHIC) transversely polarized protons will be collided for the first time. Therefore it is important to make an analysis of the double transverse spin asymmetries. The proposed DESY HERA-$`\stackrel{}{N}`$ experiment also prompts such a study (see e.g. ). For the transversity double spin asymmetry in the Drell-Yan process such studies have been performed in considerable detail . In this paper we will mainly investigate another type of double transverse spin asymmetry, one which does not involve transversity functions.
In general, helicity non-flip quark and gluon states in transversely polarized hadron-hadron scattering will lead to power suppression \[$`𝒪(1/Q^2)`$, where $`Q^2`$ is the vector boson virtuality\]. In the present study we will exploit the fact that if the transverse momentum of the produced vector boson is observed, this no longer holds true. By observing the transverse momentum certain azimuthal asymmetries can occur in the cross section without explicit power suppression. The possibility of such an unsuppressed helicity non-flip double transverse spin asymmetry has been noted before in the literature and a tree level expression for the case of a virtual photon has been given . Here we investigate specifically the case of weak vector boson production and the effect of inclusion of Sudakov form factors.
The helicity non-flip nature will allow for double transverse spin asymmetries even in $`W`$ production (for which the helicity flip contribution is absent ). Unfortunately, such transverse momentum dependent azimuthal asymmetries will turn out to suffer from suppression due to Sudakov factors, which increases with energy, as was forseen in Ref. . We will explore this issue quantitatively in detail and we will show that for the azimuthal double transverse spin asymmetry of interest, the inclusion of Sudakov factors causes suppression by at least an order of magnitude compared to the tree level result and effectively produces a power behavior of $`1/Q^\alpha `$, with $`\alpha 0.6`$. Moreover, the asymmetry will be shown to be approximately proportional to $`x_1g_1(x_1)x_2\overline{g}_1(x_2)`$, which gives rise to additional suppression at small values of the light cone momentum fractions. The conclusion will be that this asymmetry is of interest mainly at lower energies, i.e. for $`\gamma ^{}`$ production. This also leaves the option of studying possible contributions to double transverse spin asymmetries in $`W`$ production from physics beyond the standard model.
The outline of this paper is as follows. In Sec. II we will repeat the essentials of the transversity double spin asymmetry in the Drell-Yan process in order to contrast it to the helicity non-flip asymmetry (Sec. III). We study the latter asymmetry in the neutral (Sec. IIIA) and in the charged (Sec. IIIB) vector boson case. In order to obtain estimates we will assume Gaussian transverse momentum dependence of the quarks (discussed in Sec. IIIC). We will then include Sudakov form factors in the asymmetry (Sec. IIID) and estimate its quantitative effects (Sec. IIIE). In Sec. IIIF we will comment on the possible study of physics beyond the standard model via double transverse spin asymmetries in $`W`$ production.
## II Transversity double spin asymmetry
The main characteristic of the transversity double transverse spin asymmetry of vector boson production in hadron-hadron collisions is that the gluon distribution does not contribute. Hence, at leading order in an expansion in inverse powers of the hard scale(s) only the quark transversity distribution function (denoted as $`h_1,\delta q`$ or $`\mathrm{\Delta }_Tq`$) contributes. This leads to the well-known expression for the double transverse spin asymmetry in the Drell-Yan process:
$$A_{TT}=\frac{\sigma (p^{}p^{}\mathrm{}\mathrm{}^{}X)\sigma (p^{}p^{}\mathrm{}\mathrm{}^{}X)}{\sigma (p^{}p^{}\mathrm{}\mathrm{}^{}X)+\sigma (p^{}p^{}\mathrm{}\mathrm{}^{}X)}=\frac{\mathrm{sin}^2\theta \mathrm{cos}2\varphi _S^{\mathrm{}}}{1+\mathrm{cos}^2\theta }\frac{_{a,\overline{a}}e_a^2h_1^a(x_1)\overline{h}{}_{1}{}^{a}(x_2)}{_{a,\overline{a}}e_a^2f_1^a\overline{f}_1^a}.$$
(1)
We would like to note that the sum is over flavors including anti-flavors, otherwise one should add a term in both numerator and denominator for the exchange $`\left(x_1x_2\right)`$, since one can use that $`h_1^{\overline{a}}=\overline{h}_1^a`$. The above asymmetry comes from the following azimuthal dependence in the cross section
$$\frac{d\sigma (p^{}p^{}\mathrm{}\mathrm{}^{}X)}{d\mathrm{\Omega }dx_1dx_2}=\frac{\alpha ^2}{3Q^2}\underset{a,\overline{a}}{}e_a^2\{y(1y)|𝑺_{1T}||𝑺_{2T}|\mathrm{cos}(\varphi _{S_1}^{\mathrm{}}+\varphi _{S_2}^{\mathrm{}})h_1^a(x_1)\overline{h}{}_{1}{}^{a}(x_2)+\mathrm{}\}.$$
(2)
The above is expressed in the so-called Collins-Soper frame (for details see e.g. ):
$`\widehat{t}`$ $``$ $`q/Q,`$ (3)
$`\widehat{z}`$ $``$ $`{\displaystyle \frac{x_1}{Q}}\stackrel{~}{P_1}{\displaystyle \frac{x_2}{Q}}\stackrel{~}{P_2},`$ (4)
$`\widehat{h}`$ $``$ $`q_T/Q_T=(qx_1P_1x_2P_2)/Q_T,`$ (5)
where $`q`$ is the vector boson momentum, $`P_i`$ are the hadron momenta and $`\stackrel{~}{P_i}P_iq/(2x_i)`$. The azimuthal angles lie inside the plane orthogonal to $`t`$ and $`z`$. In particular, $`\varphi ^{\mathrm{}}`$ gives the orientation of $`\widehat{l}_{}^\mu \left(g^{\mu \nu }\widehat{t}^\mu \widehat{t}^\nu +\widehat{z}^\mu \widehat{z}^\nu \right)l_\nu `$, the perpendicular part of the lepton momentum $`l`$; $`\varphi _{S_i}^{\mathrm{}}`$ is the angle between $`\widehat{𝒍}_{}`$ and $`𝑺_{iT}`$. In the cross sections we also encounter a dependence on $`y=l^{}/q^{}`$, which in the lepton center of mass frame equals $`y=(1+\mathrm{cos}\theta )/2`$, where $`\theta `$ is the angle of $`\widehat{z}`$ with respect to the momentum of the outgoing lepton $`l`$.
The perturbative corrections to the double transverse spin asymmetry, Eq. (1), have been calculated in and using the assumption that at low energies the transversity distribution function $`h_1`$ equals the helicity distribution function $`g_1`$ or by saturating the Soffer bound, it has been shown in Ref. that $`A_{TT}`$ is expected to be of the order of a percent at RHIC energies.
In Ref. it was discussed that in the process $`p^{}p^{}WX`$ the transversity distribution cannot contribute and this is a general feature of chiral-odd functions and charged current exchange processes (see also ). This means that only the suppressed contributions from the twist-3 distribution function $`g_T`$ and its gluon analogue $`\mathrm{\Delta }_Tg`$ contribute (these are chiral-even functions; they mix under evolution). Of course there are contributions of the transversity functions via quark mass terms or via production of other particles that can compensate for the helicity flip, but these are all of higher order in the strong and/or weak coupling constants (e.g. one can think of $`p^{}p^{}ZXW^+W^{}X`$, but this is negligible at RHIC energies). Since neither quarks nor gluons contribute without suppression to the asymmetry $`A_{TT}`$ in the process $`p^{}p^{}WX`$, it might make this asymmetry a good place to look for contributions from physics beyond the standard model. For instance, scalar or tensor couplings of the quarks to the $`W`$ could in principle produce an asymmetry. We will return to this issue at the end of this article. First one has to investigate and estimate another standard model mechanism, namely, there is the possibility that the quarks (and gluons) are not exactly collinear to the initial proton, leading to a helicity non-flip asymmetry without explicit suppression factors of $`1/Q`$. In other words, if one measures the cross section differential in the transverse momentum of the vector boson, either in its angle compared to the other particles or in its magnitude, the helicity non-flip double transverse spin asymmetry can receive contributions at leading order, even for $`W`$ production. If one averages over this transverse momentum, then the asymmetry will vanish, but an (inadvertent) incomplete averaging, for instance due to imposed cuts, might still have observable consequences, cf. for instance Ref. . Even though we will show that for $`W`$ production this will not be a problem since the asymmetry turns out to be negligible, for $`\gamma ^{}`$ at lower energies this is important to take into consideration.
## III Helicity non-flip double transverse spin asymmetry
If one can measure the cross section differential in the transverse momentum of the vector boson, either in its angle compared to the other particles or in its magnitude, then there is the possibility to have a double transverse spin asymmetry at leading order, in principle even for $`W`$ production. To illustrate this we will make use of the formalism pioneered by Ralston and Soper , which will be applicable in the region where the observed transverse momentum is small compared to the hard scale(s). Since this is a tree level formalism, we will later on include the effects of resummation of soft gluons by combining it with the approach of Ref. . We focus on the Drell-Yan process, first on neutral vector boson production and later on charged vector boson production. However, the expressions also apply to $`p^{}p^{}2\text{jets}`$ (same formulas with minor trivial replacements) and a similar analysis can be applied to $`p^{}p^{}\pi (\mathrm{\Lambda },\mathrm{})X`$, where the observed hadron is part of a high-$`p_T`$ jet. But in those cases there are background processes, which should be considered also and this requires more detailed study.
### A Neutral vector boson production
In Eq. (2) we have given the contribution to the cross section that depends on the sum of the two transverse spin angles with respect to the lepton pair production plane, i.e. $`\mathrm{cos}(\varphi _{S_1}^{\mathrm{}}+\varphi _{S_2}^{\mathrm{}})`$. This means that if one integrates over the lepton pair orientation, then this azimuthal dependence will average to zero. At order $`1/Q^2`$ both quarks and gluons can contribute to $`A_{TT}`$ via a term in the cross section which does not depend on the lepton scattering plane
$$A_{TT}\mathrm{cos}(\varphi _{S_1}^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})\frac{M_1M_2}{Q^2}g_T\overline{g}_T,$$
(6)
but this is expected to be negligible at $`Q^2=M_Z^2`$. Moreover, it is not at all clear that such a factorized description of the asymmetry holds at this level of next-to-next-to-leading twist, since it is well known that for the unpolarized case this order \[$`𝒪(1/Q^4)`$\] does not factorize.
On the other hand, if one were to observe the transverse momentum of the lepton pair compared to the protons, there will be a double transverse spin asymmetry as a function of this transverse momentum $`𝒒_T`$, which appears at the leading order in $`1/Q`$. It will involve one more angle ($`\varphi _h^{\mathrm{}}`$), but even if one would integrate over this angle (keeping only the magnitude of $`𝒒_T`$) and over the lepton pair orientation ($`\varphi ^{\mathrm{}}`$), then there will remain an azimuthal dependence in the cross section that depends on the orientations of the two transverse polarization vectors only.
To make this explicit we will look at Eq. (A1) of Ref. , which gives the cross section for the polarized Drell-Yan process $`p^{}p^{}\gamma (Z)X`$ in the formalism of Ralston and Soper using transverse momentum dependent distribution functions. From the expressions for the production of the $`Z`$ boson it is easy to obtain the expressions for the production of the $`W`$ boson. We will not repeat all the details of the calculation of the cross section expressions, rather we will focus on the expression as given in the Appendix in Ref. . For the sake of argument it is unimportant to include contributions from the (formally difficult) $`T`$-odd distribution functions, hence we neglect them, but they can be easily included. This leaves
$`{\displaystyle \frac{d\sigma (h_1h_2\mathrm{}\overline{\mathrm{}}X)}{d\mathrm{\Omega }dx_1dx_2d^2𝒒_T}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^2}{3Q^2}}{\displaystyle \underset{a,\overline{a}}{}}\{K_1^a(y)|𝑺_{1T}||𝑺_{2T}|\mathrm{sin}(\varphi _h^{\mathrm{}}\varphi _{S_1}^{\mathrm{}})\mathrm{sin}(\varphi _h^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})\left[{\displaystyle \frac{𝒑_T𝒌_T}{M_1M_2}}g_{1T}\overline{g}_{1T}\right]`$ (8)
$`K_1^a(y)|𝑺_{1T}||𝑺_{2T}|\mathrm{cos}(2\varphi _h^{\mathrm{}}\varphi _{S_1}^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})\left[{\displaystyle \frac{\widehat{𝒉}𝒑_T\widehat{𝒉}𝒌_T}{M_1M_2}}g_{1T}\overline{g}_{1T}\right]+\mathrm{}\},`$
where
$$K_1(y)=\left(\frac{1}{2}y+y^2\right)\left[e_a^2+2g_V^le_ag_V^a\chi _1+c_1^lc_1^a\chi _2\right]\frac{12y}{2}\left[2g_A^le_ag_A^a\chi _1+c_3^lc_3^a\chi _2\right],$$
(9)
which contain the combinations of the couplings
$`c_1^j`$ $`=`$ $`(g_V^j{}_{}{}^{2}+g_A^j{}_{}{}^{2}),`$ (10)
$`c_3^j`$ $`=`$ $`2g_V^jg_A^j.`$ (11)
The $`Z`$-boson propagator factors are given by
$`\chi _1`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{sin}^2(2\theta _W)}}{\displaystyle \frac{Q^2(Q^2M_Z^2)}{(Q^2M_Z^2)^2+\mathrm{\Gamma }_Z^2M_Z^2}},`$ (12)
$`\chi _2`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{sin}^2(2\theta _W)}}{\displaystyle \frac{Q^2}{Q^2M_Z^2}}\chi _1,`$ (13)
and $`g_V`$ and $`g_A`$ are the vector and axial-vector couplings to the $`Z`$ boson. We have summed over the polarization of the outgoing leptons. Furthermore, we use the convolution notation (Ralston and Soper use $`I[\mathrm{}]`$)
$$\left[f\overline{f}\right]d^2𝒑_Td^2𝒌_T\delta ^2(𝒑_T+𝒌_T𝒒_T)f^a(x_1,𝒑_T^2)\overline{f}{}_{}{}^{a}(x_2,𝒌_T^2),$$
(14)
where $`a`$ is the flavor index.
The function $`g_{1T}`$ is the function $`h^{LT}`$ of Ralston and Soper and has as interpretation the distribution of longitudinally polarized quarks ($`\gamma ^+\gamma _5`$ projection) inside a transversely polarized hadron. It enters into the calculation compared to the unpolarized distribution function as follows:
$`\mathrm{\Phi }(x_1,𝒑_T)`$ $`=`$ $`{\displaystyle \frac{M_1}{2P_1^+}}\{f_1(x_1,𝒑_T^2){\displaystyle \frac{\overline{)}P_1}{M_1}}{\displaystyle \frac{(𝒑_T𝑺_{1T})}{M_1}}g_{1T}(x_1,𝒑_T^2){\displaystyle \frac{\overline{)}P_1\gamma _5}{M_1}}+\mathrm{}\}.`$ (15)
The details of the momenta are: the momenta of the quarks which annihilate into the photon with momentum $`q`$, are predominantly along the direction of the parent hadrons. One hadron momentum ($`P_1`$) is chosen to be along the lightlike direction given by the vector $`n_+`$ (apart from mass corrections). The second hadron with momentum $`P_2`$ is predominantly in the $`n_{}`$ direction which satisfies $`n_+n_{}=1`$, such that $`P_1P_2=𝒪(Q^2)`$. We decompose the momenta in $`+,`$ and transverse components, defined through $`p^\pm =pn_{}`$, where we note that \[cf. Eqs. (3)–(5)\]
$`n_+^\mu `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[\widehat{t}^\mu +\widehat{z}^\mu {\displaystyle \frac{Q_T}{Q}}\widehat{h}^\mu \right],`$ (16)
$`n_{}^\mu `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left[\widehat{t}^\mu \widehat{z}^\mu {\displaystyle \frac{Q_T}{Q}}\widehat{h}^\mu \right].`$ (17)
The four-momentum conservation delta function at the vector boson vertex is written as (neglecting $`1/Q^2`$ contributions)
$$\delta ^4(qkp)=\delta (q^+p^+)\delta (q^{}k^{})\delta ^2(𝒑_T+𝒌_T𝒒_T),$$
(18)
and allows for integration over $`p^{}`$ and $`k^+`$. However, the transverse momentum integrations cannot be separated, unless one integrates over the transverse momentum of the vector boson or –equivalently– of the lepton pair.
The Drell-Yan cross section is obtained by contracting the lepton tensor with the hadron tensor. At the tree level we find, for the hadron tensor,
$$𝒲^{\mu \nu }=\frac{1}{3}d^2𝒑_Td^2𝒌_T\delta ^2(𝒑_T+𝒌_T𝒒_T)\text{Tr}\left(\mathrm{\Phi }(x_1,𝒑_T)V_1^\mu \overline{\mathrm{\Phi }}(x_2,𝒌_T)V_2^\nu \right)|_{p^+,k^{}}+\left(\begin{array}{c}qq\\ \mu \nu \end{array}\right).$$
(19)
The vertices $`V_i^\mu `$ can be either the photon vertex $`V^\mu =e\gamma ^\mu `$ or the $`Z`$-boson vertex $`V^\mu =g_V\gamma ^\mu +g_A\gamma _5\gamma ^\mu `$.
The above given azimuthal dependence in the cross section, Eq. (8), means that if one observes the transverse momentum of the $`\gamma `$ or $`Z`$ boson, one can consider the cross section differential in the magnitude of the transverse momentum only and integrate over the orientations of the leptons and of $`𝒒_T`$ itself. This results in the following double transverse spin asymmetry:
$$A_{TT}(Q_T)=\frac{d\sigma \left[p^{}p^{}\gamma (Z)X\right]d\sigma \left[p^{}p^{}\gamma (Z)X\right]}{d\sigma \left[p^{}p^{}\gamma (Z)X\right]+d\sigma \left[p^{}p^{}\gamma (Z)X\right]}=\frac{_{a,\overline{a}}K_1^a(y)\left[𝒑_T𝒌_Tg_{1T}\overline{g}_{1T}\right]}{2M_1M_2_{a,\overline{a}}K_1^a(y)\left[f_1\overline{f}_1\right]},$$
(20)
where $`Q_T^2q_T^2𝒒_T^2Q^2`$. This can be seen from the following considerations. The angular dependence $`\mathrm{sin}(\varphi _h^{\mathrm{}}\varphi _{S_1}^{\mathrm{}})\mathrm{sin}(\varphi _h^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})`$ can be rewritten after integration over the angle $`\varphi _h^{\mathrm{}}`$ as $`\mathrm{cos}(\varphi _{S_1}^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})/2=\mathrm{cos}(\varphi _{S_2}^{S_1})/2`$. Since this does not depend on the orientation of $`\mathrm{}`$ itself, one can integrate over it also. The angular dependence $`\mathrm{cos}(2\varphi _h^{\mathrm{}}\varphi _{S_1}^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})`$ averages out.
If on the other hand one only observes the angle of the transverse momentum and averages over its magnitude, one can also obtain a nonvanishing asymmetry (which can still be differentiated from the transversity asymmetry). The whole point is to prevent the averaging: $`d^2𝒒_T\left[𝒑_T𝒌_Tg_{1T}\overline{g}_{1T}\right]=0`$. As said before, an incomplete averaging due to imposed experimental cuts might also result in a nonvanishing asymmetry, cf. for instance Ref. .
Before we continue we would like to point out that unlike for $`h_1`$ there is a leading twist gluon analogue of the function $`g_{1T}`$. The function arises with $`iϵ_T^{\alpha \beta }`$ in the correlation function $`\mathrm{\Phi }^{\alpha \beta }PS|F^{+\alpha }F^{+\beta }|PS`$, which means that it is a $`\mathrm{\Delta }g`$ type of function with transverse momentum dependence. But since the transition $`gg\gamma (Z)`$ is only possible via quarks, we implicitly include the gluon in the sum over flavors.
### B Charged vector boson production
In order to arrive at the expressions for the cross sections of the charged current process (cf. Ref. ), one can take $`e_a=0`$ and replace
$$\chi _2^Z\chi _2^W=\left(\frac{1}{8\mathrm{sin}^2\theta _W}\right)^2\frac{Q^4}{(Q^2M_W^2)^2+\mathrm{\Gamma }_W^2M_W^2},$$
(21)
in the above given coupling $`K_1^a`$. In addition, one replaces $`c_1=\pm c_3=1`$, depending on the chirality of the quark or lepton, since $`c_1=(g_R^2+g_L^2)/2`$ and $`c_3=(g_L^2g_R^2)/2`$. Hence, for a left-handed quark one finds $`c_1^a=c_3^a=1`$ and for a right-handed quark one finds $`c_1^a=c_3^a=1`$; similarly for the leptons. We also note that a left- or right-handed quark or lepton has helicity $`\lambda _{q/e}=1`$. This results in
$`K_1^{ab}(y)`$ $`=`$ $`4\chi _2^W|V_{ab}|^2\left({\displaystyle \frac{1}{2}}y+y^2\lambda _q\lambda _e{\displaystyle \frac{12y}{2}}\right)`$ (22)
$`=`$ $`4\chi _2^W|V_{ab}|^2\times \{\begin{array}{cc}y^2& \text{for equal quark and lepton chiralities},\hfill \\ (1y)^2& \text{for opposite quark and lepton chiralities},\hfill \end{array}`$ (25)
where $`a,b`$ are the incoming quark and antiquark flavor indices, respectively, and $`V_{ab}`$ stands for the appropriate Cabibbo-Kobayashi-Maskawa (CKM) matrix element. We illustrate the above by assuming that only $`u`$ and $`d`$ quark distribution functions contribute. This leaves two elementary subprocesses: $`u\overline{d}W^+e^+\nu `$ ($`u`$ and $`\nu `$ have equal chiralities) and $`d\overline{u}W^{}e^{}\overline{\nu }`$ ($`d`$ and $`e^{}`$ have equal chiralities) for which one finds the couplings $`K_1^{u\overline{d}}=K_1^{\overline{u}d}=4\chi _2^W|V_{ud}|^2y^2`$. For the cross section one has to take into account that in the sum over final state polarizations there is now only one state that contributes, but for the asymmetry this is not relevant.
In case of $`p^{}p^{}WX`$ and subsequent leptonic decay of the $`W`$, we encounter the problem that the produced neutrino will prevent a determination of the transverse momentum of the $`W`$ boson. Hence, in the case of a produced neutrino one cannot define a lepton scattering plane (one does not observe $`l^{}`$), hence azimuthal angles cannot be defined compared to the $`W`$ boson direction. This holds unless one can reconstruct the direction of the neutrino by the momentum imbalance .
Another possible way of observing the transverse momentum of the $`W`$ boson is looking at the $`W`$ decay into 2 jets. The expressions for lepton pair production stay essentially the same for the 2 jets case after the obvious replacement of the coupling constants. By measuring the direction of the 2 jets, the transverse momentum of the $`W`$ can be determined, but one problem is that it has $`\gamma ^{}/Z2\text{jets}`$ as a background. Separation of $`\gamma ^{}/Z`$ and $`W`$ might only be possible with a very high transverse momentum cut , but then the given expressions are not applicable anymore. Another problem is that it also receives contributions from quark-quark scattering next to the quark-antiquark scattering, but such contributions have rather large color factor suppression . In any case, the contribution coming from the transversity functions to the $`2\text{jets}`$ asymmetry can always be eliminated by averaging over the orientation of the 2 jets, as explained above.
Here we will focus only on lepton pair production and assume that in case of $`W`$ production the direction of the neutrino can be reconstructed to obtain the transverse momentum of the $`W`$.
### C Gaussian transverse momentum dependence
In order to obtain an estimate of the above asymmetry, we will consider a Gaussian transverse momentum dependence of the functions. Instead of using Gaussians, another way of obtaining an estimate of the asymmetry would be to use the spectator model for the function $`g_{1T}(x,𝒑_T^2)`$ . But for simplicity we will assume a Gaussian transverse momentum dependence, e.g.
$$g_{1T}(x,𝒑_T^2)=g_{1T}(x)\frac{R^2}{\pi }\mathrm{exp}\left(R^2𝒑_T^2\right)g_{1T}(x)𝒢(𝒑_T^2).$$
(26)
We would like to relate the function $`g_{1T}(x)`$ to a well-known function in order to be able to make some predictions in the end. This can be achieved by using $`𝒑_T^2`$ weighted functions
$$f^{(1)}(x)=d^2𝒑_T\frac{𝒑_T^2}{2M^2}f(x,𝒑_T^2).$$
(27)
For a Gaussian transverse momentum dependence we find that $`g_{1T}^{(1)}=g_{1T}(x)/(2M^2R^2)`$. In the Wandzura-Wilczek approximation the function $`g_{1T}^{(1)}`$ is a well-known quantity: it equals (upon neglecting quark masses) $`xg_T^{WW}(x)`$, where $`g_T^{WW}(x)=g_1+g_2^{WW}`$ is the Wandzura-Wilczek part of the function $`g_T`$ (see also Ref. for a discussion on this topic). This can be shown by using the equations of motion. The function $`g_T`$ has been studied by SLAC and the Spin Muon Collaboration (SMC) and the data are (still) consistent with $`g_T=g_T^{WW}`$. Also, the data show that $`g_2`$ is small compared to $`g_1`$, therefore, up to a few percent one can take $`xg_T^{WW}(x)xg_1(x)`$. Thus, we find $`g_{1T}(x)=g_{1T}^{(1)}\mathrm{\hspace{0.17em}2}M^2R^2xg_1(x)\mathrm{\hspace{0.17em}2}M^2R^2`$.
Furthermore, will assume that the Gaussians are the same for $`f_1`$ and $`g_{1T}`$ and for both protons, i.e., we take $`R_1=R_2=R`$ and $`M_1=M_2=M`$. In this way we find for example
$$[𝒑_T𝒌_Tg_{1T}\overline{g}_{1T}]\frac{M^4R^4}{\pi }(1\frac{Q_T^2R^2}{2})\mathrm{exp}(\frac{Q_T^2R^2}{2})x_1g_1(x_1)x_2\overline{g}{}_{1}{}^{}(x_2).$$
(28)
This results in the following tree level double transverse spin asymmetry at $`Q_T=0`$
$$A_{TT}(Q_T=0)=M^2R^2\frac{_{a,\overline{a}}K_1^a(y)x_1g_1^a(x_1)x_2\overline{g}{}_{1}{}^{a}(x_2)}{_{a,\overline{a}}K_1^a(y)f_1^a(x_1)\overline{f}{}_{1}{}^{a}(x_2)}.$$
(29)
We will also encounter the Fourier transforms of these functions. The function $`\stackrel{~}{f}`$ will denote the Fourier transform of $`f`$, and since we use the notation $`f(x)=d^2𝒑_Tf(x,𝒑_T)`$, we see that $`f(x)=\stackrel{~}{f}(x,b=0)`$. Taking the Fourier transform of Eq. (26) yields
$$\stackrel{~}{g}{}_{1T}{}^{}(x,b^2)=g_{1T}(x)\mathrm{exp}(\frac{b^2}{4R^2}).$$
(30)
### D Beyond the range of intrinsic transverse momentum
Monte Carlo studies including soft gluon resummation show that the largest contribution to the unpolarized cross section arises when the transverse momentum of the $`W`$ is several GeV. This transverse momentum is too high to trust tree level expressions which involve only intrinsic transverse momenta. In order to go beyond this region, we will also include the Sudakov factor arising from resummed perturbative corrections to the transverse momentum distribution.
Resummation of soft gluons into Sudakov form factors results in a replacement in Eqs. (14) and (19) of
$$\delta ^2(𝒑_T+𝒌_T𝒒_T)\frac{d^2𝒃}{(2\pi )^2}e^{i𝒃(𝒑_T+𝒌_T𝒒_T)}e^{S(b)},$$
(31)
where $`e^{S(b)}`$ is the Sudakov form factor and $`b^2=𝒃^2`$. This has been shown in Refs. for the leading twist and is discussed for the present context in more detail in Ref. . The Sudakov form factor is found to be
$$S(b,Q)=_{b_0^2/b^2}^{Q^2}\frac{d\mu ^2}{\mu ^2}\left[A(\alpha _s(\mu ))\mathrm{ln}\frac{Q^2}{\mu ^2}+B(\alpha _s(\mu ))\right].$$
(32)
One can expand the functions $`A`$ and $`B`$ in $`\alpha _s`$ and the first few coefficients are known for unpolarized scattering and for longitudinally polarized scattering . The latter result is needed here since the function $`g_{1T}`$ is a distribution of longitudinally polarized quarks; the asymmetry on the parton level is $`a_{LL}`$. In order to obtain a first estimate of the effect of including the Sudakov factor we will take into account only the first term in the expansion of $`A`$: $`A^{(1)}=C_F/\pi `$. This leads to the expression
$$S(b,Q)=\frac{16}{332n_f}\left[\mathrm{log}\left(\frac{b^2Q^2}{b_0^2}\right)+\mathrm{log}\left(\frac{Q^2}{\mathrm{\Lambda }^2}\right)\mathrm{log}\left[1\frac{\mathrm{log}\left(b^2Q^2/b_0^2\right)}{\mathrm{log}\left(Q^2/\mathrm{\Lambda }^2\right)}\right]\right],$$
(33)
with $`b_0=2\mathrm{exp}(\gamma _E)1.123`$. We will take for the number of flavors $`n_f=5`$ and also $`\mathrm{\Lambda }_{QCD}=200\text{MeV}`$.
The replacement in Eq. (14) leads to (suppressing the flavor index)
$`\left[f\overline{f}\right]`$ $``$ $`{\displaystyle \frac{d^2𝒃}{(2\pi )^2}e^{i𝒃𝒒_T}e^{S(b)}\stackrel{~}{f}(x_1,b)\stackrel{~}{\overline{f}}(x_2,b)}`$ (34)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}𝑑bbJ_0(bQ_T)e^{S(b)}\stackrel{~}{f}(x_1,b)\stackrel{~}{\overline{f}}(x_2,b).`$ (35)
The functions also have a renormalization and factorization scale dependence, which we will choose to be equal $`\mu _R=\mu _F=\mu `$. Hence, we have, for instance,
$$f(x;\mu )=d^2𝒑_Tf(x,p_T;\mu )\stackrel{~}{f}(x,b=0;\mu ),$$
(36)
where also the boundary of the integration gives a $`\mu `$ dependence. In Eq. (35) one usually takes $`\stackrel{~}{f}(x_1,b;\mu =b_0/b)`$ .
Of course, if one includes the effects of perturbative corrections, one should also include higher order corrections to the hard part not coming from soft gluons. But since the formalism, which means the factorized formula, is valid only for $`𝒒_T^2Q^2`$, we include only the effects of soft gluons, which should allow us to extend the range of applicability from the region of intrinsic transverse momentum to the region of moderate $`𝒒_T`$ values. An equivalent way of saying this is that one can perform a collinear expansion of the hard scattering part $`H^{\mu \nu }(x_1,x_2,𝒑_T,𝒌_T,𝒒_T,Q)H^{\mu \nu }(x_1,x_2,Q)`$, such that perturbative corrections to the hard part do not affect the transverse momentum structure of the tree level result and hence the transverse momentum weight in the asymmetry will be the same also beyond the tree level. Here we will consider $`H^{\mu \nu }(x_1,x_2,Q)`$ to lowest order in $`\alpha _s`$, therefore, only logarithmic $`Q^2`$ corrections to the results presented below are expected.
The numerator in Eq. (20) cannot be treated exactly like the denominator, so let us focus next on
$`\left[𝒑_T𝒌_Tf\overline{f}\right]`$ $``$ $`{\displaystyle \frac{d^2𝒃}{(2\pi )^2}e^{i𝒃𝒒_T}e^{S(b)}d^2𝒑_Td^2𝒌_T𝒑_T𝒌_Te^{i𝒃(𝒑_T+𝒌_T)}f(x_1,𝒑_T^2)\overline{f}(x_2,𝒌_T^2)}.`$ (37)
As mentioned before, we assume that the distribution functions are Gaussians (as a function of transverse momentum), all of equal width: $`f(x_1,𝒑_T^2)=f(x_1)𝒢(𝒑_T^2)`$ and $`\overline{f}(x_2,𝒌_T^2)=\overline{f}(x_2)𝒢(𝒌_T^2)`$. One can then change variables to $`𝒖=(𝒑_T+𝒌_T)/\sqrt{2}`$ and $`𝒗=(𝒑_T𝒌_T)/\sqrt{2}`$ and compute
$`{\displaystyle d^2𝒑_Td^2𝒌_T𝒑_T𝒌_Te^{i𝒃(𝒑_T+𝒌_T)}𝒢(𝒑_T^2)𝒢(𝒌_T^2)}`$ $`=`$ $`{\displaystyle \frac{b^2}{4R^4}}\mathrm{exp}\left({\displaystyle \frac{b^2}{2R^2}}\right),`$ (38)
which after application to Eq. (37) yields (see also Eq. (30))
$$\left[𝒑_T𝒌_Tf\overline{f}\right]=\frac{d^2𝒃}{(2\pi )^2}e^{i𝒃𝒒_T}\frac{b^2}{4R^4}e^{S(b)}\stackrel{~}{f}(x_1,b)\stackrel{~}{\overline{f}}(x_2,b),$$
(39)
which can be compared with Eq. (28). Both equations fulfill the property that the expression should vanish after $`𝒒_T`$ integration. One infers that the numerator of $`A_{TT}(Q_T)`$ and hence $`A_{TT}(Q_T)`$ itself oscillate, but in general $`𝑑Q_TA_{TT}(Q_T)0`$, because of the $`Q_T`$ dependence of the denominator.
For the asymmetry, Eq. (20), we then obtain
$`A_{TT}(Q_T)`$ $`=`$ $`{\displaystyle \frac{1}{8M^2R^4}}{\displaystyle \frac{_{a,\overline{a}}K_1^a(y)_0^{\mathrm{}}dbb^3J_0(bQ_T)e^{S(b)}\stackrel{~}{g}{}_{1T}{}^{a}(x_1,b)\stackrel{~}{\overline{g}}{}_{1T}{}^{a}(x_2,b)}{_{a,\overline{a}}K_1^a(y)_0^{\mathrm{}}dbbJ_0(bQ_T)e^{S(b)}\stackrel{~}{f}{}_{1}{}^{a}(x_1,b)\stackrel{~}{\overline{f}}{}_{1}{}^{a}(x_2,b)}}`$ (40)
$``$ $`{\displaystyle \frac{M^2}{2}}{\displaystyle \frac{_{a,\overline{a}}K_1^a(y)x_1g_1^a(x_1)x_2\overline{g}{}_{1}{}^{a}(x_2)}{_{a,\overline{a}}K_1^a(y)f_1^a(x_1)\overline{f}{}_{1}{}^{a}(x_2)}}{\displaystyle \frac{_0^{\mathrm{}}𝑑bb^3J_0(bQ_T)\mathrm{exp}\left[S(b)\frac{1}{2}b^2/R^2\right]}{_0^{\mathrm{}}𝑑bbJ_0(bQ_T)\mathrm{exp}\left[S(b)\frac{1}{2}b^2/R^2\right]}},`$ (41)
where the approximation arises from taking $`g_{1T}(x)=g_{1T}^{(1)}\mathrm{\hspace{0.17em}2}M^2R^2xg_1(x)\mathrm{\hspace{0.17em}2}M^2R^2`$.
In order to extend the above equation to the nonperturbative region of large values of $`b`$, one usually introduces $`b_{}=b/\sqrt{1+b^2/b_{\mathrm{max}}^2}`$ and an additional term $`\mathrm{exp}\left[S_{NP}(b)\right]`$, needed to describe the low $`𝒒_T`$ region properly. In part, $`S_{NP}(b)`$ is introduced to take care of the smearing due to the intrinsic transverse momentum, therefore, taking into account the term $`\mathrm{exp}\left(\frac{1}{2}b^2/R^2\right)`$ in addition will just produce a change in the coefficient of the $`b^2`$ term in $`S_{NP}(b)`$. To keep the unpolarized cross section unaffected, we will therefore introduce as nonperturbative term $`\mathrm{exp}\left[S_{NP}(b)+\frac{1}{2}b^2/R^2\right]`$ and study the following final expression for the asymmetry:
$$A_{TT}(Q_T)=\frac{1}{2}\frac{_{a,\overline{a}}K_1^a(y)x_1g_1^a(x_1)x_2\overline{g}{}_{1}{}^{a}(x_2)}{_{a,\overline{a}}K_1^a(y)f_1^a(x_1)\overline{f}{}_{1}{}^{a}(x_2)}𝒜(Q_T),$$
(42)
where we define
$$𝒜(Q_T)M^2\frac{_0^{\mathrm{}}𝑑bb^3J_0(bQ_T)\mathrm{exp}\left[S(b_{})S_{NP}(b)\right]}{_0^{\mathrm{}}𝑑bbJ_0(bQ_T)\mathrm{exp}\left[S(b_{})S_{NP}(b)\right]}.$$
(43)
The denominator is then the conventional unpolarized expression. Also, we note that $`𝒜(Q_T)`$ is dimensionless and, for simplicity, we will take $`M=1\text{GeV}`$.
The above approach of including Sudakov factors in the tree level azimuthal asymmetry expressions can also be applied to expressions derived in for electron-positron annihilation and in for lepton-hadron scattering.
### E Estimating the asymmetry
For the case of $`W`$ production the asymmetry becomes
$$A_{TT}^W(Q_T)=\frac{1}{2}\frac{_{a,\overline{a};b,\overline{b}}|V_{ab}|^2x_1g_1^a(x_1)x_2\overline{g}{}_{1}{}^{b}(x_2)}{_{a,\overline{a};b,\overline{b}}|V_{ab}|^2f_1^a(x_1)\overline{f}{}_{1}{}^{b}(x_2)}𝒜(Q_T).$$
(44)
If only the $`u`$ and $`d`$ quarks contribute, then also the CKM matrix elements drop out of the ratio.
For the nonperturbative Sudakov factor we use the parameterization of Ladinsky-Yuan, Ref. , which was fitted to relevant Fermilab data,
$$S_{NP}(b)=g_1b^2+g_1g_3b\mathrm{ln}(100x_1x_2)+g_2b^2\mathrm{ln}\left(\frac{Q}{2Q_0}\right),$$
(45)
with $`g_1=0.11\text{GeV}^2,g_2=0.58\text{GeV}^2,g_3=1.5\text{GeV}^1,Q_0=1.6\text{GeV}`$ and $`b_{\mathrm{max}}=0.5\text{GeV}^1`$. We will take $`x_1x_2=10^2`$, which is justified below. This leads to $`S_{NP}(b)=1.98b^2`$ at $`Q=80\text{GeV}`$ and $`S_{NP}(b)=0.77b^2`$ at $`Q=10\text{GeV}`$.
The result for the asymmetry factor $`𝒜(Q_T)`$ at $`Q=80\text{GeV}`$ is given in Fig. 1.
It is plotted up to $`Q_T=Q/10`$, since beyond that $`Q_T`$ range the magnitude only slowly decreases and also the approximation $`Q_T^2Q^2`$ is expected to become less valid. The asymmetry factor $`𝒜(Q_T)`$ at $`Q_T=0`$ is seen to be around $`0.32`$ and at $`Q_T`$ values of a few GeV –relevant for the majority of produced $`W`$ bosons– the asymmetry factor has a sign change and consequently a smaller magnitude. On top of that the asymmetry is proportional to $`|x_1g_1(x_1)x_2\overline{g}_1(x_2)|x_1f_1(x_1)x_2\overline{f}_1(x_2)`$. Therefore, the total asymmetry as a function of $`Q_T`$ is expected to be below the percent level, if one assumes that on average $`x_1=0.4`$ and $`x_2=0.07`$ for $`W`$ production at RHIC . Since we have implicitly included the gluons in the sum over flavors, the latter argument is not valid if $`\mathrm{\Delta }g`$ turns out to be extremely large at small $`x`$. Of course this will have even more serious implications for e.g. $`A_{LL}`$ in jet production, especially at low transverse momenta. We will just view this unlikely option as a proviso.
In Ref. it is demonstrated that the transversity double spin asymmetry, which is a $`𝒒_T`$ integrated asymmetry, is expected to be at most on the order of a few percent, which matches the level of sensitivity of RHIC. The present asymmetry is still a function of $`Q_T`$, requiring even more statistics. This will make the asymmetry $`A_{TT}^{W(Z)}(Q_T)`$ invisible at RHIC. Moreover, since it oscillates as a function of $`Q_T`$, one expects that the asymmetry partly integrated over $`Q_T`$, will not lead to any significant result either.
A few remarks about the dependence of the result on the nonperturbative parameters. The asymmetry factor is seen to decrease with increasing Gaussian smearing width. Taking a higher value of $`b_{\mathrm{max}}`$ and a lower value of $`x_1x_2`$ both increase this width. The above –optimistic– choices of $`b_{\mathrm{max}}=0.5\text{GeV}^1`$ and $`x_1x_2=10^2`$ are therefore expected to overestimate the asymmetry factor somewhat.
But at lower energies –where larger light cone momentum fractions can be achieved– this asymmetry for $`\gamma ^{}`$ production is still worth investigating.
In Fig. 2 we have given the function $`𝒜(Q_T)`$ at the scale $`Q=10\text{GeV}`$ and find that at low values of $`Q_T`$ it is around $`1`$. Measuring $`A_{TT}^\gamma ^{}(Q_T)`$ at larger values of $`x_1`$ and $`x_2`$ might make this asymmetry observable.
If one studies the values of $`𝒜(Q_T)`$ at $`Q_T=0`$ –where the asymmetry is largest– as a function of $`Q`$, one observes that inclusion of the Sudakov factor produces to good approximation a $`1/Q^\alpha `$ behavior, with $`\alpha 0.6`$. Even though this suppression is not very strong as a function of $`Q`$, one actually needs to compare the resulting asymmetry Eq. (44) with the tree level asymmetry Eq. (29). This leads to a comparison of $`𝒜(Q_T=0)`$ and $`2M^2R^2`$. In a tree level analysis $`R^2=1/𝒑_T^2`$ for a typical intrinsic transverse momentum squared value, i.e. $`R^2211\text{GeV}^2`$, corresponding to the range of $`𝒑_T^2(300700\text{MeV})^2`$. If we view $`R^2=1`$ as giving a lower bound for the tree level asymmetry factor, then $`𝒜(Q_T=0,Q=80)`$ is still an order of magnitude smaller.
We conclude that the Sudakov factors produce a strong suppression compared to the tree level result and the effect increases with energy. This will also have consequences for similar types of transverse momentum dependent azimuthal asymmetries appearing in for instance $`e^+e^{}`$ annihilation at the $`Z`$ mass scale , where the same strong suppression due to Sudakov factors is expected.
### F Physics beyond the standard model
It is now clear that the standard model (SM) mechanisms seem to produce negligible double transverse spin asymmetries in $`W`$ production ($`A_{TT}^W`$). In summary the reasons are the following: the transversity distribution $`h_1`$ does not contribute . At next-to-next-to-leading twist \[$`𝒪(M_1M_2/Q^2)`$\] the twist-3 distribution function $`g_T`$ (which is a helicity non-flip quark distribution) can contribute and its gluon analogue $`\mathrm{\Delta }_Tg`$ as well. At $`Q^2=M_W^2`$ these contributions will be negligible. Furthermore, one can also neglect contributions which are of higher order in the strong and/or weak coupling constants. In the case of perturbative QCD corrections double helicity flip will be accompanied by quark mass terms and therefore will also be suppressed by at least two factors of $`1/Q`$. In the case of weak corrections one also expects negligible contributions, at least at RHIC energies, say around $`Q^2=M_W^2`$. Any leftover asymmetry from incomplete averaging of the transverse momentum dependent asymmetry $`A_{TT}^W(Q_T)`$ we have estimated to be negligible as well. So at RHIC energies $`A_{TT}^W`$ is expected to be negligible within the SM.
It is now fair to address the question: if a significant asymmetry would nevertheless be found in the polarized proton-proton collisions at RHIC, can one really conclude something about physics beyond the SM? For instance, there could be scalar or tensor couplings of the $`W`$ to quarks that can generate an asymmetry $`A_{TT}^W`$. If the scale of such new physics is $`\mathrm{\Lambda }M_W`$, then one might need to compare effects of order $`M_1M_2/Q^2`$ with $`Q^2/\mathrm{\Lambda }^2`$. For instance, for $`Q^2=M_W^2`$ and $`\mathrm{\Lambda }=1\text{TeV}`$ the latter is a factor of 40 larger (although still quite small). But the problem of comparing to higher twist contributions disappears altogether if the new couplings violate symmetries.
This means that on top of the fact that the various SM mechanisms produce negligible asymmetries, one can also exploit the dependence on the orientation of the transverse spins compared to the lepton production plane to cancel out specific contributions exactly. We have seen that the transversity asymmetry $`A_{TT}^{\gamma /Z}`$ appears with an angular dependence $`\mathrm{cos}(\varphi _{S_1}^{\mathrm{}}+\varphi _{S_2}^{\mathrm{}})`$, whereas $`A_{TT}(Q_T)`$ and the $`1/Q^2`$ suppressed contribution from $`g_T\overline{g}_T`$ in Eq. (6) both appear with $`\mathrm{cos}(\varphi _{S_1}^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})`$. The latter does not depend on the lepton scattering plane, because the asymmetries are not double transverse spin asymmetries at the parton level. On the other hand, symmetry violation asymmetries can produce other angular dependences than any SM mechanism.
There might be $`T`$-odd asymmetries, for example the one of Ref. ,
$$A_{TT}^{}\mathrm{sin}(\varphi _{S_1}^{\mathrm{}}+\varphi _{S_2}^{\mathrm{}}),$$
(46)
which would arise due to $`CP`$-violating vector couplings of the quarks to the $`W`$, which is assumed not to be $`VA`$ anymore, but some complex linear combination of $`V`$ and $`A`$. It can clearly be distinguished from possible initial state interaction effects, which are $`P`$-even and only lead to asymmetries independent of the lepton scattering plane. To be more specific, if one assumes $`T`$-odd ($`P`$-even) distribution functions to be nonzero, then there will also be contributions proportional to \[cf. Eq. (A1) of \]
$`\mathrm{cos}(\varphi _{S_1}^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})[𝒑_T𝒌_Tf_{1T}^{}\overline{f}{}_{1T}{}^{}],`$ (47)
$`\mathrm{sin}(\varphi _{S_1}^{\mathrm{}}\varphi _{S_2}^{\mathrm{}})\left[𝒑_T𝒌_Tf_{1T}^{}\overline{g}_{1T}\right].`$ (48)
The function $`f_{1T}^{}`$ corresponds to the so-called Sivers effect . Here one usually assumes that such a function might arise due to initial state interactions and its contributions indeed do not depend on the lepton pair orientation as can be seen from the above two angular dependences.
Therefore, these structures are distinguishable from the $`T`$-odd asymmetry $`A_{TT}^{}`$. However, it is important to note that $`A_{TT}^{}`$ can also effectively arise due to SM $`CP`$ violation, hence this contribution must first be calculated before any conclusion about physics beyond the SM can be made. Also, this specific asymmetry will be accompanied by the product $`h_1^a(x_1)\overline{h}{}_{1}{}^{a}(x_2)`$, thus it will suffer from the same drawback as Eq. (1), namely that the transversity function for the antiquarks is presumably smaller than for the quarks, making this asymmetry hard (if not impossible) to detect at RHIC. But it illustrates how symmetry violation can be used in principle to disentangle SM asymmetries from new physics asymmetries.
## IV Conclusions
We have investigated a helicity non-flip double transverse spin asymmetry in vector boson production in hadron-hadron scattering, which was first considered by Ralston and Soper. It does not involve transversity functions and in principle also arises in $`W`$-boson production for which we have presented the expressions. The asymmetry requires observing the transverse momentum of the vector boson, but it is not suppressed by explicit inverse powers of the large energy scale $`Q`$. However, as we have shown, inclusion of Sudakov factors suppresses the asymmetry at least by an order of magnitude compared to the tree level result. This suppression increases with energy approximately as a fractional power, numerically found to be $`\alpha 0.6`$. Moreover, the asymmetry is shown to be approximately proportional to $`x_1g_1(x_1)x_2\overline{g}_1(x_2)`$, which gives rise to additional suppression at small values of the light cone momentum fractions. This implies that the asymmetry is negligible for $`Z`$ and $`W`$ production at RHIC and is mainly of interest at low energies (for $`\gamma ^{}`$ production). The strong suppression with respect to the tree level result will also have consequences for similar types of transverse momentum dependent azimuthal asymmetries in for instance $`e^+e^{}`$ annihilation at the $`Z`$ mass scale, where the same strong suppression due to Sudakov factors is expected.
We have also noted that unlike the transversity and $`CP`$-violating double transverse spin asymmetries, the helicity non-flip asymmetry $`A_{TT}(Q_T)`$ does not depend on the orientation of the transverse spin vectors compared to the lepton pair production plane orientation. This feature can be exploited to separate the different types of asymmetries.
## ACKNOWLEDGMENTS
I would like to thank Les Bland, Gerry Bunce, Bob Jaffe, Rainer Jakob, Piet Mulders, Naohito Saito and Werner Vogelsang for helpful comments and discussions. Furthermore, I thank RIKEN, Brookhaven National Laboratory and the U.S. Department of Energy (contract DE-AC02-98CH10886) for providing the facilities essential for the completion of this work.
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# BIPOLARON BINDING IN QUANTUM WIRES
## I INTRODUCTION
Landau’s idea of the auto-localized state of a charge carrier (polaron) in a homogeneous polar medium got a further development by Pekar who first studied a problem of a stable complex of two charge carriers of the same sign (bipolaron). The bipolaron binding energy was first calculated in Ref. . The bipolaron problem was widely discussed, see e. g. Refs. . A detailed outline of this subject is presented in the recent review .
Dimensionless constants of the Coulomb interaction $`U`$ and of the electron-phonon interaction $`\alpha `$ are related to each other by the equation :
$$U=\frac{\sqrt{2}\alpha }{1\eta },$$
(1)
where $`\eta =\epsilon _{\mathrm{}}/\epsilon _0`$ ($`\epsilon _0`$ and $`\epsilon _{\mathrm{}}`$ are static and optical dielectric constants, respectively). Due to the fact that $`\epsilon _0>\epsilon _{\mathrm{}}`$, the relation $`U\sqrt{2}\alpha `$ follows. When the distance $`l`$ between electrons is large or small compared with the characteristic polaron radius $`R_\mathrm{p}`$ (see Ref. ), the phonon-mediated attraction between electrons occurs to be weaker than the repulsion. At large distances $`lR_\mathrm{p}`$, both interaction potentials have similar spatial dependences but the Coulomb repulsion is stronger than the phonon-mediated attraction. In the opposite case, $`lR_\mathrm{p}`$, the Coulomb potential diverges at the zero distance, while the phonon-mediated attraction is always finite. Nevertheless, when two electrons move in such a way that the average distance between them is of the same order as the polaron radius, the bipolaron can be stable at $`\alpha 1`$ and $`\eta 1`$. When two electrons are together confined to a potential well, one can expect that the conditions of the bipolaron stability may be improved for relevant sizes of the well.
Two new circumstances have stimulated the bipolaron theory: the progress in the fabrication technology of mesoscopic nanostructures such as quasi-2D (quantum wells and superlattices), quasi-1D (quantum wires), quasi-0D (quantum dots), and the advancement of the hypothesis that bipolaron excitations might play a role in processes occurring in the high-temperature superconductors. The present research has been motivated also by the recent advances in creation of nanocrystals with a strong ionic coupling .
The basic bipolaron parameters are recalled in what follows. The bipolaron stability region is determined by the inequality $`W>0`$ for the bipolaron binding energy
$$W2E_\mathrm{p}E_{\mathrm{bip}}.$$
(2)
Here $`E_\mathrm{p}`$ and $`E_{\mathrm{bip}}`$ are the free polaron and bipolaron ground state energies, respectively. From the equation
$$W(\alpha ,\eta ,)=0,$$
(3)
where $``$ denotes the set of parameters determining the shape and the size of the confinement domain, the functions $`\alpha _c(\eta ,)`$ and $`\eta _c(\alpha ,)`$ describing the boundaries of the bipolaron stability region are found, for fixed $`\eta `$ and $`\alpha `$, respectively. According to different theoretical treatments the bipolaron binding energy is an increasing function of $`\alpha `$ and a decreasing function of $`\eta `$. It will be shown that the function $`\eta _c(\alpha ,)`$ starts from $`\eta _c=0`$ at $`\alpha =\alpha _{\mathrm{min}}\left(\right)0`$, grows with increasing $`\alpha `$ and tends to the upper limit $`\eta _{\mathrm{max}}`$ at $`\alpha \mathrm{}`$. The bipolaron stability region is then determined by the inequalities $`\alpha \alpha _{\mathrm{min}}\left(\right)`$ and $`0\eta <\eta _c(\alpha ,)`$.
Let us adduce typical values of the parameters $`\alpha _{\mathrm{min},3\mathrm{D}}`$ and $`\eta _{\mathrm{max},3\mathrm{D}}`$ of the bulk (3D) bipolaron: $`\alpha _{\mathrm{min},3\mathrm{D}}=6.8`$ and $`\eta _{\mathrm{max},3\mathrm{D}}=0.14`$ were found by Verbist, Peeters and Devreese and by Verbist, Smondyrev, Peeters and Devreese . Adamowski obtained $`\alpha _{\mathrm{min},3\mathrm{D}}=7.3`$ and $`\eta _{\mathrm{max},3\mathrm{D}}=0.14`$. The bipolaron theory developed for pure 2D and 1D models shows that the bipolaron stability region broadens when the dimensionality is reduced. For these systems, the following parameters were obtained: $`\alpha _{\mathrm{min},2\mathrm{D}}=2.9`$, $`\eta _{\mathrm{max},2\mathrm{D}}=0.158`$ (Ref. ); $`\alpha _{\mathrm{min},1\mathrm{D}}=0.9`$, $`\eta _{\mathrm{max},1\mathrm{D}}=0.764`$ (Ref. ). Bipolaron states were investigated in a quantum well and in a quantum wire as a function of the characteristic size of the system. The polaron theory for a quantum dot is developed in Refs. .
The goal of the present investigation is to determine the bipolaron stability region and to study the basic parameters characterizing the bipolaron ground state as a function of confinement. Two different types of confinement are considered and compared to each other: (i) a cylindrical quantum wire of the radius $`R`$, where continuous transitions from 3D to 1D are realized with decreasing $`R`$; (ii) a planar quantum wire of the width $`L`$, where a transition from 2D to 1D is realized with decreasing $`L`$. A unique approach, namely the Feynman variational method , is used throughout the paper for both systems under analysis.
The paper is organized as follows. In Section II, general formulae for parameters of a bipolaron in quantum wires are deduced. In Section III, particular cases of cylindrical and planar quantum wires are considered. The basic parameters of the bipolaron ground state are obtained. Limiting cases of strong and weak confinement are studied in detail. The obtained numerical and analytical results are discussed in Section IV. Section V contains conclusions about the influence of confinement on the bipolaron binding energy in quantum wires.
## II General theory
We analyze the bipolaron problem taking into account both the electron-phonon interaction and the Coulomb repulsion between two electrons confined to a quantum wire. The Lagrange function of the system is
$`L`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{D}}{}}}{\displaystyle \underset{n=1,2}{}}{\displaystyle \frac{m_i\dot{x}_{i,n}^2}{2}}{\displaystyle \underset{n=1,2}{}}𝒰\left(𝐫_n\right){\displaystyle \frac{e^2}{\epsilon _{\mathrm{}}\left|𝐫_1𝐫_2\right|}}`$ (5)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐤}{}}\left(\dot{w}_𝐤^2\omega _0^2w_𝐤^2\right){\displaystyle \underset{n=1,2}{}}{\displaystyle \underset{𝐤}{}}\gamma _𝐤\left(𝐫_n\right)w_𝐤,`$
where $`𝐫_n(x_{1n},x_{2n},x_{3n})`$ is the radius vector of the $`n`$-th electron ($`n=1,2`$); $`m_i`$ is the $`ii`$-th component ($`i=1,2,3`$) of the diagonal band mass tensor, $`𝒰\left(𝐫\right)`$ is the potential energy of an electron in the quantum wire, $`w_𝐤`$ are the normal coordinates of longitudinal optical (LO) phonon modes. Here, the parameter $`\mathrm{D}`$ determines the dimensionality of the electron subsystem: $`\mathrm{D}=3`$ and $`2`$ for cylindrical and planar quantum wires, respectively. Amplitudes of the electron-phonon interaction are taken in the Fröhlich form:
$$\gamma _𝐤\left(𝐫\right)=2\sqrt{\frac{2\pi \mathrm{}\omega _0\alpha }{V}}\frac{\omega _0}{k}\left(\frac{\mathrm{}}{2\overline{m}\omega _0}\right)^{1/4}\mathrm{exp}\left(i\mathrm{𝐤𝐫}\right),$$
(6)
where $`\overline{m}\left(m_1m_2m_3\right)^{1/3}`$, $`V`$ is the volume of the system, and the Fröhlich constant
$$\alpha =\frac{e^2}{2\mathrm{}\omega _0}\left(\frac{1}{\epsilon _{\mathrm{}}}\frac{1}{\epsilon _0}\right)\left(\frac{2\overline{m}\omega _0}{\mathrm{}}\right)^{1/2}$$
(7)
characterizes the strength of the coupling between an electron and bulk polar LO phonons with the long-wavelength frequency $`\omega _0`$. In this paper, the 3D phonon approximation is used, according to which the interaction of an electron with both bulk-like and interface phonons is replaced by that with 3D phonons. This often used approach is adequate because any integral polaron or bipolaron effect, resulting from a summation over all phonon modes, appears to be only weakly dependent on the details of the phonon spectrum. It should be also mentioned that the system under consideration simulates realistic structures with relatively smooth interface barriers, where interface-like phonon modes can appear, which are smoothly distributed in space rather than localized near a sharp boundary, as is the case for interface modes.
In order to study the bipolaron problem at arbitrary values of $`\alpha `$, the Feynman variational approach is the most appropriate method. The trial Lagrange function is written as
$`L_{tr}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{\mathrm{D}}{}}}{\displaystyle \underset{n=1,2}{}}\left[m_i\dot{x}_{i,n}^2+M_i\dot{X}_{i,n}^2k_i\left(x_{i,n}X_{i,n}\right)^2k_i^{}\left(x_{i,n}X_{i,\overline{n}}\right)^2\right]`$ (9)
$`+{\displaystyle \underset{i=1}{\overset{3}{}}}K_i\left(x_{i1}x_{i2}\right)^2{\displaystyle \underset{n=1,2}{}}𝒲\left(𝐫_n\right),`$
where $`X_{in}`$ are coordinates of the $`n`$-th “fictitious” particle ($`n=1,2`$). This model imitates the interaction of electrons with phonons and between each other by elastic bonds as shown in Fig. 1. The masses $`M_i`$ and the force constants $`k_i`$, $`k_i^{}`$, $`K_i`$ play the role of variational parameters. For $`n=1`$, $`\overline{n}`$ takes the value 2, and for $`n=2`$, $`\overline{n}`$ is equal to 1. The potential well $`𝒰\left(𝐫\right)`$ from Eq. (5) is simulated here by a parabolic function:
$$𝒲\left(𝐫\right)=\frac{1}{2}\underset{i=1}{\overset{q}{}}m_i\mathrm{\Omega }_i^2x_i^2.$$
(10)
The index $`q`$ characterizes the dimensionality of confinement and is determined as follows: for a planar quantum wire $`q=1`$ ($`\mathrm{\Omega }_10`$ and $`\mathrm{\Omega }_2=0`$) and for a cylindrical quantum wire $`q=2`$ ($`\mathrm{\Omega }_10`$, $`\mathrm{\Omega }_20`$, and $`\mathrm{\Omega }_3=0`$).
The basis of the Feynman variational method is the Jensen-Feynman inequality :
$$\mathrm{exp}(SS_{tr})_{S_{tr}}\mathrm{exp}SS_{tr}_{S_{tr}},$$
(11)
where the angular brackets denote averaging over electron paths:
$$G_{S_{tr}}=\frac{\mathrm{Tr}\mathrm{D}𝐫G\left[𝐫\right]\mathrm{exp}\left(S_{tr}\right)}{\mathrm{Tr}\mathrm{D}𝐫\mathrm{exp}\left(S_{tr}\right)}.$$
(12)
Here $`S`$ and $`S_{tr}`$ are the electron action functionals obtained after integration over phonon variables and over coordinates of “fictitious” particles, respectively. At low temperatures, the variational bipolaron energy is calculated using the expression
$$E_{bip}=E_{tr}\underset{\beta \mathrm{}}{lim}\frac{SS_{tr}_{S_{tr}}}{\beta },$$
(13)
where $`E_{tr}`$ is the ground state energy of the trial system with the Lagrangian (9), $`\beta =1/k_BT`$ is the inverse temperature.
The trial Lagrange function (9) consists of $`\mathrm{D}`$ independent parts: $`L_{tr}=\underset{i=1}{\overset{\mathrm{D}}{}}L_i`$. Each part $`L_i`$ is a function of four variables $`x_{i1}`$, $`x_{i2}`$, $`X_{i1}`$, $`X_{i2}`$. Let us introduce unified denotations for coordinates of electrons and of “fictitious” particles: $`\stackrel{~}{x}_{i1}=x_{i1}`$, $`\stackrel{~}{x}_{i2}=x_{i2}`$, $`\stackrel{~}{x}_{i3}=X_{i1}`$, $`\stackrel{~}{x}_{i4}=X_{i2}`$. It follows from the form of the trial Lagrangian (9) with Eq. (10) that the groups of variables $`\stackrel{~}{x}_{ij}`$ with different indices $`i`$ are dynamically independent from each other. They are related to normal variables $`\xi _{ij}`$ by the unitary transformation:
$$\stackrel{~}{x}_{ij}=\underset{j^{}=1}{\overset{4}{}}d_{i,jj^{}}\xi _{ij^{}},i=1,\mathrm{},\mathrm{D}$$
(14)
with $`4\times 4`$ matrices $`d_{i,jj^{}}`$ ($`j,j^{}=1,\mathrm{},4`$). From the equations of motion for the group of coordinates $`\stackrel{~}{x}_{ij}`$ ($`j=1,\mathrm{},4`$) with a fixed $`i`$, the following eigenfrequencies are obtained:
$`\omega _{ij}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{\left(1+{\displaystyle \frac{M_i}{m_i}}\right)v_i^2+\mathrm{\Omega }_i^2(1)^j\sqrt{\left[\left(1{\displaystyle \frac{M_i}{m_i}}\right)v_i^2\mathrm{\Omega }_i^2\right]^2+4{\displaystyle \frac{M_i}{m_i}}v_i^4}\right\},j=1,2,`$ (15)
$`\omega _{ij}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{\left(1+{\displaystyle \frac{M_i}{m_i}}\right)v_i^2+\mathrm{\Omega }_i^22{\displaystyle \frac{K_i}{m_i}}(1)^j\sqrt{\left[\left(1{\displaystyle \frac{M_i}{m_i}}\right)v_i^2\mathrm{\Omega }_i^2+2{\displaystyle \frac{K_i}{m_i}}\right]^2+4{\displaystyle \frac{(k_ik_i^{})^2}{m_iM_i}}}\right\},`$ (17)
$`j=3,4,`$
where $`v_i^2=(k_i+k_i^{})/M_i`$. Matrix elements of the unitary transformation (14) are
$`d_{i,11}^2={\displaystyle \frac{\omega _{i1}^2v_i^2}{2\left(\omega _{i1}^2\omega _{i2}^2\right)}},d_{i,12}^2={\displaystyle \frac{v_i^2\omega _{i2}^2}{2\left(\omega _{i1}^2\omega _{i2}^2\right)}},d_{i,13}^2={\displaystyle \frac{\omega _{i3}^2v_i^2}{2\left(\omega _{i3}^2\omega _{i4}^2\right)}},d_{i,14}^2={\displaystyle \frac{v_i^2\omega _{i4}^2}{2\left(\omega _{i3}^2\omega _{i4}^2\right)}},`$
$$d_{i,2j^{}}=s_j^{}d_{i,1j^{}};d_{i,3j^{}}=\frac{k_i+s_j^{}k_i^{}}{M_i\left(v_i^2\omega _{ij^{}}^2\right)}d_{i,1j^{}};d_{i,4j^{}}=s_j^{}d_{i,3j^{}};$$
(18)
$`s_j=1(j=1,2);s_j=1(j=3,4).`$
Note that the elastic repulsion imitating the Coulomb interaction gives a contribution to the eigenfrequencies with $`j=3`$ and $`4`$ through the force constants $`K_i`$. It is easy to see from Eq. (17) that under the conditions of a strong confinement along the $`i`$-th coordinate axis the eigenfrequencies with $`j=1`$ and $`3`$ corresponding to the motion of the bipolaron along this axis as a whole are determined mainly by the parameter $`\mathrm{\Omega }_i`$.
The action functionals $`S`$ and $`S_{tr}`$ in Eqs. (11) to (13) contain the potential energies $`𝒰`$ and $`𝒲`$, respectively. Though the shape of a real potential $`𝒰`$ may differ from that of the model quadratic potential (10), the averaged difference $`𝒰𝒲_{S_{tr}}`$ can be omitted as far as it is small when compared to the rest of $`SS_{tr}_{S_{tr}}/\beta `$.
The averaging procedure in Eq. (13) is carried out by the path integration and leads to the following form of the variational bipolaron energy:
$$E_{\mathrm{bip}}=\underset{i=1}{\overset{\mathrm{D}}{}}B_i+C+P.$$
(19)
Here the terms $`B_i`$ include the averaged kinetic energies of two electrons and of two “fictitious” particles as well as the averaged potential energy of the elastic interaction of these four particles:
$$B_i=\frac{1}{2}\underset{j=1}{\overset{4}{}}\omega _{ij}\left(1\frac{\omega _{ij}^2\mathrm{\Omega }_i^2}{\omega _{ij}^2}d_{i,1j}^2\right)v_i,i=1,\mathrm{},\mathrm{D}.$$
(20)
In Eqs. (19), (20) and further on, the Feynman units are used: $`\mathrm{}\omega _0`$ for energies; $`\omega _0`$ for frequencies; and $`\left(\mathrm{}/\overline{m}\omega _0\right)^{1/2}`$ for lengths.
The averaged potential energy of the Coulomb electron repulsion is
$$C=\frac{\alpha }{\left(1\eta \right)\pi ^2}𝒦_2\left(0\right),$$
(21)
and the averaged energy of the electron-phonon interaction is
$$P=\frac{\alpha }{\pi ^2}\underset{n=1,2}{}\underset{0}{\overset{\mathrm{}}{}}𝑑\tau \mathrm{e}^\tau 𝒦_n\left(\tau \right),$$
(22)
where
$$𝒦_n\left(\tau \right)=\underset{\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{k^2}\mathrm{exp}\left[\underset{i=1}{\overset{\mathrm{D}}{}}k_i^2A_{in}\left(\tau \right)\right]\underset{i=1}{\overset{\mathrm{D}}{}}dk_i.$$
(23)
The functions $`A_{in}\left(\tau \right)`$ are determined as follows:
$$A_{in}\left(\tau \right)=\frac{\overline{m}}{m_i}\left[\underset{j=1,2}{}\frac{d_{i,1j}^2}{\omega _{ij}}\left(1\mathrm{e}^{\omega _{ij}\tau }\right)+\underset{j=3,4}{}\frac{d_{i,1j}^2}{\omega _{ij}}\left(1+\left(1\right)^n\mathrm{e}^{\omega _{ij}\tau }\right)\right],n=1,2.$$
(24)
In order to find the bipolaron energy, it is necessary to minimize the function $`E_{\mathrm{bip}}`$ given by Eq. (19) over twelve independent variational parameters $`\omega _{ij},`$ $`i=1\mathrm{}3,`$ $`j=1\mathrm{}4`$, which are used instead of the masses $`M_i`$ and the force constants $`k_i`$, $`k_i^{}`$, $`K_i`$.
From Eq. (17) for the eigenfrequencies, the expression for components of the diagonal tensor of relative bipolaron effective mass is deduced straightforwardly:
$$\frac{\left(m_{\mathrm{bip}}\right)_i}{m_i}2\left(\frac{M_i}{m_i}+1\right)=\frac{2\left(\omega _{i1}^2+\omega _{i2}^2\mathrm{\Omega }_i^2\right)}{v_i^2},$$
(25)
where the values of parameters $`\omega _{i1}`$, $`\omega _{i2}`$, $`v_i`$ are taken which provide the bipolaron energy.
The number of phonons in the bipolaron cloud is determined by the general expression of Ref.
$$N_{\mathrm{ph}}=\frac{S}{\left(\mathrm{}\omega _0\right)}_{S_{tr}},$$
(26)
which gives in the case under consideration:
$$N_{\mathrm{ph}}=\frac{\alpha }{\pi ^2}\underset{n=1,2}{}\underset{0}{\overset{\mathrm{}}{}}𝒦_n\left(\tau \right)\mathrm{e}^\tau \tau 𝑑\tau .$$
(27)
Calculations of the average number of phonons according to Eq. (27) are performed using the results of minimization of the bipolaron energy.
## III Bipolaron in cylindrical and planar quantum wires
### A Variational problem
Here we write down the variational bipolaron energies in the cylindrical and planar quantum wires (see Fig. 2). Hereafter, the following denotation for the confinement parameter is used: $`\mathrm{\Omega }_i\mathrm{\Omega }_{}`$, $`i=1,q`$. The electron mass is taken to be isotropic, i. e. $`m_1=m_2=m_3=m`$. From Eqs. (19) to (22) we obtain the variational bipolaron energy
$$E_{\mathrm{bip}}=B_{}+B_{}+C+P,$$
(28)
where, in accordance with Eq. (20),
$$B_{}=\frac{1}{2}\underset{j=1}{\overset{3}{}}\omega _j\left(1d_{\mathrm{D},1j}^2\right)v_{},$$
(29)
$$B_{}=\frac{q}{2}\left[\underset{j=1}{\overset{4}{}}\omega _j\left(1\frac{\omega _j^2\mathrm{\Omega }_{}^2}{\omega _j^2}d_{1,1j}^2\right)2v_{}\right].$$
(30)
Here the frequencies of the motion along the $`z`$-axis (called below the longitudinal motion) are
$$\omega _{\mathrm{D1}}\omega _1\omega _{\mathrm{D2}}=0,\omega _{\mathrm{D}j}\omega _j,j=3,4,$$
(31)
$`v_\mathrm{D}v_{}`$and those of the motion in the $`xy`$-plane (the transverse motion) are
$$\omega _{ij}\omega _j,j=1,\mathrm{},4,v_iv_{}$$
(32)
with $`i=1,2`$ for $`q=2`$ and $`i=1`$ for $`q=1`$. In the case under consideration, the integrations containing the function $`𝒦_n\left(\tau \right)`$ are performed analytically. The calculation of the integrals in Eqs. (21) and (22) yields the averaged potential energy of the Coulomb repulsion between electrons
$$C=\frac{\sqrt{2}U}{\sqrt{\pi A_2(0)}}F_q\left(1\frac{A_2(0)}{A_2(0)}\right)$$
(33)
and the averaged energy of the electron-phonon interaction
$$P=\frac{2\alpha }{\sqrt{\pi }}\underset{n=1,2}{}\underset{0}{\overset{\mathrm{}}{}}𝑑\tau e^\tau \frac{1}{\sqrt{A_n(\tau )}}F_q\left(1\frac{A_n(\tau )}{A_n(\tau )}\right),$$
(34)
where
$`A_n\left(\tau \right)={\displaystyle \underset{j=3,4}{}}{\displaystyle \frac{d_{\mathrm{D},1j}^2}{\omega _j}}\left[1+(1)^n\mathrm{e}^{\omega _j\tau }\right]+{\displaystyle \frac{d_{\mathrm{D},11}^2}{\omega _1}}\left(1\mathrm{e}^{\omega _1\tau }\right)+d_{\mathrm{D},12}^2\tau ;`$
$`A_n\left(\tau \right)={\displaystyle \underset{j=1,2}{}}{\displaystyle \frac{d_{1,1j}^2}{\omega _j}}\left(1\mathrm{e}^{\omega _j\tau }\right)+{\displaystyle \underset{j=3,4}{}}{\displaystyle \frac{d_{1,1j}^2}{\omega _j}}\left[1+(1)^n\mathrm{e}^{\omega _j\tau }\right],n=1,2;`$
$$F_q(x)=\{\begin{array}{cc}\frac{\mathrm{tanh}^1\sqrt{x}}{\sqrt{x}},& q=2;\hfill \\ \underset{0}{\overset{\frac{\pi }{2}}{}}\frac{d\phi }{\left(1x\mathrm{sin}^2\phi \right)^{1/2}},& q=1.\hfill \end{array}$$
(35)
The minimization of the variational bipolaron energy $`E_{\mathrm{bip}}`$ determined by Eqs. (28) to (35) is carried out with respect to eight variational parameters $`v_{}`$, $`\omega _2`$ , $`\omega _j,\omega _j(j=1,3,4)`$. Setting $`k_i^{}`$ and $`K_i`$ equal to zero in these formulas, the twice value of the polaron energy is obtained from Eq. (28). The binding energy is then found according to Eq. (2). Results of the calculation of $`W`$ as a function of the quantum wire radius $`R=\mathrm{\Omega }_{}^{1/2}`$ are presented in Fig. 3 for different values of $`\alpha `$. Then the functions $`\alpha _{min}\left(R\right)\alpha _c\left(\eta =0,R\right)`$ (Fig. 4) and $`\eta _c(R,\alpha )`$ are obtained. The latter function is used for calculation of the critical value of the Coulomb repulsion constant $`U_c\left(\alpha \right)`$ in order to describe the bipolaron stability region shown in Fig. 5. From Eq. (25), taking into account Eq. (31), the relative bipolaron effective mass of the longitudinal motion is derived as
$$\frac{\left(m_{\mathrm{bip}}\right)_{}}{m}=\frac{2\omega _1^2}{v_{}^2}.$$
(36)
Plots of the relative bipolaron mass as a function of $`R`$ are shown in Fig. 6. A detailed discussion of the results will be given in Section IV.
### B Weak size quantization
For a weak size quantization, the eigenfrequencies of the transverse motion can be represented as expansion series in $`\mathrm{\Omega }_{}`$. In these series, we take into account only two first terms:
$`\omega _1^2=\omega _1^2+\mathrm{\Omega }_{}^2{\displaystyle \frac{\omega _1^2v_{}^2}{\omega _1^2}}+O(\mathrm{\Omega }_{}^4),\omega _2^2=\mathrm{\Omega }_{}^2{\displaystyle \frac{v_{}^2}{\omega _1^2}}+O(\mathrm{\Omega }_{}^4),`$
$$\omega _3^2=\omega _3^2+\mathrm{\Omega }_{}^2\frac{\omega _3^2v_{}^2}{\omega _3^2\omega _4^2}+O(\mathrm{\Omega }_{}^4),\omega _4^2=\omega _4^2+\mathrm{\Omega }_{}^2\frac{v_{}^2\omega _4^2}{\omega _3^2\omega _4^2}+O(\mathrm{\Omega }_{}^4).$$
(37)
As a consequence, Eq. (30) takes on the form
$$B_{}=q\left[B_{}+\mathrm{\Omega }_{}\frac{v_{}}{4\omega _1}\frac{3\omega _1^2v_{}^2}{\omega _1^2}\right]+O(\mathrm{\Omega }_{}^2),$$
(38)
and the parameters in the right-hand side of Eq. (34) satisfy the relations
$$A_n(\tau )=A_n(\tau )\frac{1}{4}\mathrm{\Omega }_{}\tau ^2\frac{v_{}^3}{\omega _1^3}+O(\mathrm{\Omega }_{}^2),n=1,2.$$
(39)
Substituting the expression (39) in Eqs. (33), (34), we obtain the energy of the Coulomb repulsion
$$C=\frac{\sqrt{2}U}{\sqrt{\pi }}\frac{f_{1q}}{\sqrt{A_{12}(0)}}+O(\mathrm{\Omega }_{}^2),$$
(40)
and the energy of the electron-phonon interaction
$$P=\frac{2\alpha }{\sqrt{\pi }}\underset{n=1,2}{}\underset{0}{\overset{\mathrm{}}{}}𝑑\tau \mathrm{e}^\tau \left[\frac{f_{1q}}{\sqrt{A_n(\tau )}}+\frac{\mathrm{\Omega }_{}}{4}\tau ^2\frac{v_{}^3}{\omega _1^3}\frac{f_{2q}}{\sqrt{A_n^3(\tau )}}\right]+O(\mathrm{\Omega }_{}^2),$$
(41)
where
$`f_{1q}=\{\begin{array}{c}1,q=2;\\ \frac{\pi }{2},q=1,\end{array}f_{2q}=\{\begin{array}{c}\frac{1}{3},q=2;\\ \frac{\pi }{8},q=1.\end{array}`$
The bipolaron energy in this limiting case can be represented as
$`E_{\mathrm{bip}}=E_{\mathrm{bip}}^0+\mathrm{\Delta }E_{\mathrm{bip}}+O(\mathrm{\Omega }_{}^2),`$where $`E_{\mathrm{bip}}^0`$ is the bipolaron energy in three or two dimensions for $`\mathrm{D}=3`$ or $`\mathrm{D}=2`$, respectively. The confinement-induced shift of the bipolaron energy $`\mathrm{\Delta }E_{\mathrm{bip}}`$ is given by
$$\mathrm{\Delta }E_{bip}=\mathrm{\Omega }_{}\left\{q\frac{v_{}}{4\omega _3}\frac{3\omega _1^2v_{}^2}{\omega _1^2}\frac{2\alpha }{\sqrt{\pi }}f_{2q}\frac{v_{}^3}{4\omega _1^2}\underset{0}{\overset{\mathrm{}}{}}𝑑\tau e^\tau \tau ^2\left(\frac{1}{\sqrt{A_{11}^3(\tau )}}+\frac{1}{\sqrt{A_{12}^3(\tau )}}\right)\right\}.$$
(42)
It is worth mentioning, that in the strong coupling regime, the integrals in Eq. (42) are calculated analytically, and the minimization of this variational function with respect to the frequencies is performed explicitly. For this purpose we use the results of Ref. , where the following analytical expressions for frequencies are obtained: $`\omega _i=\alpha ^2\stackrel{~}{\omega }_i`$ (for $`i=1,3`$), $`\omega _4=1,v_{}=1`$, where
$$\stackrel{~}{\omega }_1=\frac{128}{9\pi }\theta _\mathrm{D}\frac{\left[1\zeta ^2(U)\right]^4}{\zeta ^2(U)},\stackrel{~}{\omega }_3=\frac{128}{9\pi }\theta _\mathrm{D}\left[1\zeta ^2(U)\right]^3,$$
(43)
$`\zeta (U)={\displaystyle \frac{U}{16\alpha }}+{\displaystyle \frac{1}{2}}\sqrt{2+\left({\displaystyle \frac{U}{8\alpha }}\right)^2},\theta _\mathrm{D}=\{\begin{array}{cc}1,\hfill & \mathrm{D}=3;\hfill \\ \left(\frac{3\pi }{4}\right)^2,\hfill & \mathrm{D}=2.\hfill \end{array}`$Using these frequencies, we find the confinement-induced shift of the bipolaron energy to be
$$\mathrm{\Delta }E_{\mathrm{bip}}=q\frac{\mathrm{\Omega }_{}}{2\alpha ^2\stackrel{~}{\omega }_3}.$$
(44)
This result differs qualitatively from that deduced in Ref. for a bipolaron in a weak magnetic field, where the cyclotron frequency $`\omega _c`$ plays the role of $`\mathrm{\Omega }`$. Namely, as distinct from Eq. (44), in the equation from Ref. for the first-order correction to the bipolaron energy $`\alpha ^2`$ stands instead of $`\alpha ^4`$.
It is important to note, that this positive correction to the bipolaron energy due to the confinement is less than the twice value of the respective correction to the polaron energy . The confinement-induced variation of the bipolaron binding energy obeys the inequality
$$\mathrm{\Delta }W=\frac{q}{2}\frac{\mathrm{\Omega }}{}\alpha ^2\left[\frac{9\pi }{2\theta _D}\frac{1}{\stackrel{~}{\omega }_3}\right]>0.$$
(45)
Thus, the enhancement of the bipolaron binding takes place due to the confinement.
### C Strong size quantization
In the limiting case of a strong size quantization, the terms of the order of $`\mathrm{\Omega }_{}^2`$ play a determining role in Eq. (19). For the frequencies of the transverse motion, the expansion in inverse powers of $`\mathrm{\Omega }^2`$ gives:
$`\omega _1^2`$ $`=`$ $`\mathrm{\Omega }_{}^2+{\displaystyle \frac{M_1}{m}}v_{}^2+O(\mathrm{\Omega }_{}^2),\omega _2^2=v_{}^2+O(\mathrm{\Omega }_{}^2),`$ (46)
$`\omega _3^2`$ $`=`$ $`\mathrm{\Omega }_{}^2+{\displaystyle \frac{M_1}{m}}v_{}^22{\displaystyle \frac{K_{}}{m}}+O(\mathrm{\Omega }_{}^2),\omega _4^2=v_{}^2+O(\mathrm{\Omega }_{}^2).`$ (47)
Consequently, the bipolaron ground state energy is described by the expression
$`E_{\mathrm{bip}}`$ $`=`$ $`q\mathrm{\Omega }_{}+T_{}+{\displaystyle \frac{U}{\sqrt{2\pi A_2(0)}}}\mathrm{ln}\left({\displaystyle \frac{16}{q^2}}\mathrm{\Omega }_{}A_2(0)\right)`$ (49)
$`{\displaystyle \frac{\alpha }{\sqrt{\pi }}}{\displaystyle \underset{n=1,2}{}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\tau e^\tau {\displaystyle \frac{1}{\sqrt{A_n(\tau )}}}\mathrm{ln}\left({\displaystyle \frac{16}{q^2}}\mathrm{\Omega }_{}A_n(\tau )\right).`$
The first term in the right-hand side of Eq. (49) is the energy of two electrons in the parabolic potential. The last three terms in the right-hand side of Eq. (49) are due to the electron-phonon and Coulomb interactions. In the strong coupling regime, $`\omega _i=\alpha ^2\stackrel{~}{\omega }_i`$ ($`i=1,3`$); $`\omega _4`$ and $`v_{}`$ are proportional to $`\alpha ^0`$ with coefficients which are functions of $`\mathrm{\Omega }_{}`$. Then omitting the terms of the order of $`\alpha ^0`$ in the last three terms of the variational bipolaron energy (49) we obtain
$`E_{\mathrm{bip}}`$ $`=`$ $`q\mathrm{\Omega }_{}+\alpha ^2[{\displaystyle \frac{\stackrel{~}{\omega }_1+\stackrel{~}{\omega }_3}{4}}+{\displaystyle \frac{U\sqrt{\stackrel{~}{\omega }_1}}{\alpha \sqrt{2\pi }}}\mathrm{ln}\left({\displaystyle \frac{16}{q^2}}{\displaystyle \frac{\mathrm{\Omega }_{}}{\alpha ^2\stackrel{~}{\omega }_1}}\right)`$ (51)
$`{\displaystyle \frac{2\sqrt{2}}{\sqrt{\pi }}}\sqrt{{\displaystyle \frac{\stackrel{~}{\omega }_1\stackrel{~}{\omega }_3}{\stackrel{~}{\omega }_1+\stackrel{~}{\omega }_3}}}\mathrm{ln}\left({\displaystyle \frac{8}{q^2}}{\displaystyle \frac{\mathrm{\Omega }_{}}{\alpha ^2}}({\displaystyle \frac{1}{\stackrel{~}{\omega }_1}}+{\displaystyle \frac{1}{\stackrel{~}{\omega }_3}})\right)].`$
Note, that the second term in the right-hand side of Eq. (51) is proportional to $`\alpha ^2`$, as is expected in the strong coupling regime. When replacing $`\mathrm{\Omega }_{}\omega _c`$ at $`q=2`$, this polaronic term coincides with that of the bipolaron variational energy in a strong magnetic field from Ref. . The dependence of the bipolaron binding energy on the cylindrical confinement provides a possibility for a controllable enhancement of the bipolaron binding by decreasing the radius of a quantum wire.
Setting $`U=0`$ and $`\stackrel{~}{\omega }=\stackrel{~}{\omega }_1=\stackrel{~}{\omega }_3`$ in Eq. (51), one obtains the twice variational polaron energy $`2E_p`$ with the variational parameter $`\stackrel{~}{\omega }`$ (see Ref. ). The polaron energy $`E_p`$ depends on the confinement parameter similarly to $`E_{\mathrm{bip}}`$. The binding energy $`W`$ (which is not written explicitly to save space) increases logarithmically with increasing $`\mathrm{\Omega }_{}`$.
## IV DISCUSSION OF NUMERICAL RESULTS
Beyond the framework of the limiting cases which allow an analytical treatment as discussed above, the bipolaron stability is studied using the following computational procedure. First, we evaluate the bipolaron energy $`E_{\mathrm{bip}}`$ and the model bipolaron effective mass defined as $`m_{\mathrm{bip}}=2(M_{}+m)`$. Second, the functions $`\alpha _{\mathrm{min}}\left(R\right)`$ and $`\eta _c(\alpha ,R)`$ are found from Eq. (3). The region of the Fröhlich coupling constant ranging from 2 to 4 is chosen for the numerical work in order to include the values of $`\alpha `$ corresponding to some specific substances with small $`\eta `$ (for example, $`\mathrm{TiO}_2`$: $`\alpha =2.03`$, $`\eta =0.035`$ ; TlCl: $`\alpha =2.56`$, $`\eta =0.133`$ ; BaO: $`\alpha =3.23`$, $`\eta =0.118`$ ; LiBr: $`\alpha =4.15`$, $`\eta =0.24`$ ).
Fig. 3 illustrates the size dependence of the bipolaron binding energy in cylindrical and planar quantum wires, for $`\alpha `$ values corresponding to the above-mentioned substances and for $`\eta =0`$. As seen from Fig. 3, the bipolaron binding energy monotonously rises with increasing transverse confinement \[cf. Eqs. (28) to (34)\].
In Fig. 4, the minimal value $`\alpha _{\mathrm{min}}`$ is represented as a function of $`R`$ and $`L`$ for cylindrical and planar quantum wires, correspondingly. In the region of large $`R`$ and $`L`$, the minimal values $`\alpha _{\mathrm{min}}`$ for cylindrical and planar quantum wires tend to the three-dimensional and two-dimensional limits $`\alpha _{\mathrm{min},3\mathrm{D}}`$ and $`\alpha _{\mathrm{min},2\mathrm{D}}`$, respectively. When $`R`$ and $`L`$ decrease from 1.0 to 0.1, a rapid diminution of $`\alpha _{\mathrm{min}}`$ is seen. Note that at small values of $`R,L`$ (which are, however, still compatible with the continuum description), the bipolaron stability region extends to small values of $`\alpha `$. Note that the bipolaron parameters for quantum wires obtained in the formal limiting cases $`R,L0`$ differ substantially from those derived for the purely one-dimensional model (with 1D-electrons and 1D-phonons), which gives $`\alpha _{\mathrm{min},1\mathrm{D}}=0.9`$.
In Fig. 5, the ratio of the critical value of the Coulomb repulsion constant to the Fröhlich coupling constant $`\alpha `$,
$$\frac{U_c\left(\alpha \right)}{\alpha }=\frac{\sqrt{2}}{1\eta _c\left(\alpha \right)}$$
(52)
is plotted as a function of $`\alpha `$ for various radii (ranging from 0.01 to 20.0) of the cylindrical quantum wire. Since the parameter $`\eta _c`$ is non-negative, $`U_c\left(\alpha \right)/\alpha `$ cannot be less than the value $`\sqrt{2}`$ (shown by the line $`A`$). When increasing $`\alpha `$, the right-hand side of Eq. (52) tends to the three-dimensional limit $`\sqrt{2}/\left(1\eta _{c,3\mathrm{D}}\right)`$, marked by the line $`B`$. The physical sense of this trend consists in the following: when increasing the electron-phonon coupling, the electron confinement to the parabolic potential (10) is gradually replaced by the confinement to the polaronic potential well. The regions of bipolaron stability can exist only between the lines $`A`$ and $`B`$. The domain between a curve $`U_c\left(\alpha \right)/\alpha `$ and the line $`A`$ is the bipolaron stability region for a specific radius of the cylindrical quantum wire. This figure illustrates clearly an enlargement of the bipolaron stability region with decreasing the radius of the quantum wire. An analogous dependence of the bipolaron stability region on the width takes place for the planar quantum wire.
Bipolaron effective masses are represented in Fig. 6 as a function of the dimensionless sizes $`R`$ and $`L`$ of the cylindrical and planar quantum wires. The size dependence of the bipolaron effective mass is qualitatively similar to that of the ground state energy but appears to be substantially more pronounced. At small radii ($`0.1R0.2`$), the bipolaron mass strongly increases with decreasing $`R`$.
## V CONCLUSIONS
The conclusion of our analysis is that the confinement leads to an enlargement of the bipolaron stability region in cylindrical and planar quantum wires as compared to the corresponding regions of infinite three-dimensional and two-dimensional systems, respectively. For $`R1`$ or $`L1`$, the critical values $`\alpha _c`$ required for bipolaron stability are close to those for TiO<sub>2</sub>, TlCl, BaO and LiBr. For example, according to Fig. 4b, the bipolaron stability region sets in at the width $`L`$ of the planar quantum wire of about 8 nm for parameters of TlCl. In this view, manifestations of the bipolaron phenomena might be observed in the technically achievable planar quantum wire structures.
The performed analytical and numerical analysis of the influence of confinement on the bipolaron binding energy has shown that stable bipolaron states are possible even for intermediate values of $`\alpha `$ ($`\alpha 2`$) and for not too small values of $`\eta `$ ($`\eta 0.1`$) in nanostructures whose sizes are of the same order as the polaron radius $`R_\mathrm{p}`$. Among the considered systems, the most favorable conditions for the bipolaron stability take place in planar quantum wires, where the binding energy $`W`$ increases monotonously (logarithmically) with strengthening confinement. Nanostructures, whose sizes satisfy the conditions of the bipolaron stability, seem to be achievable for the modern technology.
###### Acknowledgements.
We thank V. N. Gladilin for valuable discussions. This work has been supported by the Interuniversitaire Attractiepolen — Belgische Staat, Diensten van de Eerste Minister — Wetenschappelijke, technische en culturele Aangelegenheden; Bijzonder Onderzoeksfonds (BOF) NOI of the Universiteit Antwerpen; PHANTOMS Research Network; F.W.O.-V. projects Nos. G.0287.95 and the W.O.G. WO.025.99N (Belgium). S.N.K. acknowledges a financial support from the UIA. E.P.P, S.N.K. and S.N.B. acknowledge with gratitude the kind hospitality during their visits to the UIA in the framework of the common research project supported by PHANTOMS.
FIGURE CAPTIONS
Fig. 1. A scheme of the trial system which contains two electrons connected with two “fictitious” particles through the elastic attraction and models the Coulomb interaction by the elastic repulsion.
Fig. 2. A scheme of cylindrical (a) and planar (b) quantum wires.
Fig. 3. The bipolaron binding energy $`W`$ in cylindrical (a) and planar (b) quantum wires plotted versus the dimensionless radius $`R`$ and width $`L`$, respectively.
Fig. 4. The minimal value (at $`\eta =0`$) of the critical electron-phonon coupling constant $`\alpha _c`$ plotted versus $`R`$ and $`L`$ in cylindrical (a) and planar (b) quantum wires, respectively.
Fig. 5. The ratio of the critical Coulomb repulsion constant $`U_c`$ and $`\alpha `$, as a function of $`\alpha `$, in cylindrical quantum wires for $`R=0.01`$ (1), 0.5 (2), 1.0 (3), and 20.0 (4).
Fig. 6. The bipolaron effective mass $`(m_{\mathrm{bip}})_{}/m`$ in cylindrical (a) and planar (b) quantum wires plotted versus $`R`$ and $`L`$, respectively. The curves for the effective mass are broken off as the bipolaron becomes unstable.
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# Generalized slow-roll inflation
## Acknowledgments
This work was supported by the EEC grant PSS\* 0992 and by OLAM, Fondation pour la Recherche Fondamentale, Brussels.
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# Renormalon Cancellation in Heavy Quarkonia and Determination of 𝑚_𝑏, 𝑚_𝑡Based on the invited talk “Top quark physics at future linear colliders” given at the Japan Physics Society Meeting, Osaka, Japan, March 30 - April 2, 2000.
## 1 Introduction
Recently there has been significant progress in our understanding of heavy quarkonia such as $`\mathrm{{\rm Y}}`$’s and remnant of toponium resonances. Developments in technologies of higher order calculations and the subsequent discovery of renormalon cancellation enabled extractions of $`m_b`$ and (in future experiments) of $`m_t`$ with high accuracy from the quarkonium spectra. In this article we review the notion of renormalon and its cancellation in the heavy quarkonium system. We demonstrate how it is useful in extracting the quark masses.<sup>*</sup><sup>*</sup>* See Ref. for a comprehensive review of renormalons.
We consider heavy quarkonia whose sizes (given by the Bohr radius $`(\alpha _Sm_q)^1`$) are much smaller than the hadronization scale $`\mathrm{\Lambda }_{\mathrm{QCD}}^1(0.3\mathrm{GeV})^1`$. In reality the candidates are $`\mathrm{{\rm Y}}(1S)`$ and (remnant of) toponium resonances, whose sizes are $`(1.5\mathrm{GeV})^1`$ and $`(20\mathrm{GeV})^1`$, respectively. In such a system, gluons participating in the binding of the $`q\overline{q}`$ boundstate have wavelengths much shorter than the hadronization scale, so theoretically nature of the boundstate can be described well using perturbative QCD. In particular the boundstate spectrum (the mass of boundstate) can be calculated as a function of the quark mass and $`\alpha _S`$. Consequently we can extract the quark masses, $`m_b`$, $`m_t`$, from the masses of the above quarkonia.
First let us state briefly the theoretical framework used in contemporary calculations of spectra of non-relativistic boundstates such as the heavy quarkonia. In old days people solved the celebrated Bethe-Salpeter equation to compute the boundstate spectrum. We no longer use this equation; instead we reduce the problem to a quantum mechanical one. Namely we solve the non-relativistic Schrödinger equation
$`\widehat{H}\psi _n(r)=E_n\psi _n(r)`$ (1)
to determine the boundstate wave functions and energy spectrum. The quantum mechanical Hamiltonian is determined from perturbative QCD order by order in expansion in $`1/c`$ (inverse of the speed of light):
$`\widehat{H}=\widehat{H}_0+{\displaystyle \frac{1}{c}}\widehat{H}_1+{\displaystyle \frac{1}{c^2}}\widehat{H}_2+\mathrm{}.`$ (2)
Since quark and antiquark inside the heavy quarkonium are non-relativistic, the expansion in $`1/c`$ leads to a reasonable systematic approximation. Presently the Hamiltonian is known up to $`𝒪(1/c^2)`$ :
$`\text{}\widehat{H}_0`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{p}^2}{m}}\underset{¯}{C_F{\displaystyle \frac{\alpha _S}{r}}},`$ (3)
$`\text{}\widehat{H}_1`$ $`=`$ $`\underset{¯}{C_F{\displaystyle \frac{\alpha _S}{r}}\left({\displaystyle \frac{\alpha _S}{4\pi }}\right)\{\beta _0\mathrm{log}(\mu ^2r^2)}+a_1\},`$ (4)
$`\widehat{H}_2`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{p}^4}{4m^3}}\underset{¯}{C_F{\displaystyle \frac{\alpha _S}{r}}\left({\displaystyle \frac{\alpha _S}{4\pi }}\right)^2\{\beta _0^2[\mathrm{log}^2(\mu ^2r^2)+{\displaystyle \frac{\pi ^2}{3}}]}+(\beta _1+2\beta _0a_1)\mathrm{log}(\mu ^2r^2)+a_2\}`$ (5)
$`+{\displaystyle \frac{\pi C_F\alpha _S}{m^2}}\delta ^3(\stackrel{}{r})+{\displaystyle \frac{3C_F\alpha _S}{2m^2r^3}}\stackrel{}{L}\stackrel{}{S}{\displaystyle \frac{C_F\alpha _S}{2m^2r}}\left(\stackrel{}{p}^2+{\displaystyle \frac{1}{r^2}}r_ir_jp_jp_i\right){\displaystyle \frac{C_AC_F\alpha _S^2}{2mr^2}}`$
$`{\displaystyle \frac{C_F\alpha _S}{2m^2}}\left\{{\displaystyle \frac{S^2}{r^3}}3{\displaystyle \frac{(\stackrel{}{S}\stackrel{}{r})^2}{r^5}}{\displaystyle \frac{4\pi }{3}}(2S^23)\delta ^3(\stackrel{}{r})\right\},`$
where $`m`$ denotes the pole mass of the quark; $`\alpha _S\alpha _S(\mu )`$; $`C_F=4/3`$, $`C_A=3`$ are color factors; $`\mu ^{}=\mu e^{\gamma _E}`$. The lowest-order Hamiltonian $`\widehat{H}_0`$ is nothing but that of two equal-mass particles interacting via the Coulomb potential.
In addition to the above Hamiltonians, some of the terms in higher order Hamiltonians $`\widehat{H}_n`$ (corresponding to the underlined terms) are known. From the analysis of these higher order terms, one finds that there exists a problem in extracting the quark mass from the boundstate spectrum. We will first address the problem, which is known as the “renormalon problem”, and then will see how it is solved.
## 2 The Renormalon Problem
The terms with underlines in Eqs. (3)–(5) stem from the running of the coupling constant dictated by
$`\mu ^2{\displaystyle \frac{d\alpha _S}{d\mu ^2}}={\displaystyle \frac{\beta _0}{4\pi }}\alpha _S^2+\mathrm{},`$ (6)
and all higher order terms can be determined using the renormalization-group equation. Thus, in the “large $`\beta _0`$ approximation”, the potential between quark and antiquark is given by It is the Coulomb potential with the “running charge”; cf. Eq. (3).
$`V_{\beta _0}(r)={\displaystyle }{\displaystyle \frac{d^3\stackrel{}{q}}{(2\pi )^3}}e^{i\stackrel{}{q}\stackrel{}{r}}C_F{\displaystyle \frac{4\pi \alpha _{1\mathrm{L}}(q)}{q^2}};q|\stackrel{}{q}|,`$ (7)
where the 1-loop running coupling is defined as a perturbation series in $`\alpha _S(\mu )`$:
$`\alpha _{1\mathrm{L}}(q)\alpha _S(\mu ){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{\beta _0\alpha _S(\mu )}{4\pi }}\mathrm{log}\left({\displaystyle \frac{q^2}{\mu ^2}}\right)\right\}^n.`$ (8)
We may not resum the geometrical series before the Fourier integration because then the integrand exhibits a pole at $`q=\mathrm{\Lambda }`$,
$`\alpha _{1\mathrm{L}}(q){\displaystyle \frac{\alpha _S(\mu )}{1+\frac{\beta _0\alpha _S(\mu )}{4\pi }\mathrm{log}\left(\frac{q^2}{\mu ^2}\right)}}={\displaystyle \frac{4\pi /\beta _0}{\mathrm{log}\left(\frac{q^2}{\mathrm{\Lambda }^2}\right)}}\text{ }q\mathrm{\Lambda }\text{ }{\displaystyle \frac{4\pi }{\beta _0}}{\displaystyle \frac{\mathrm{\Lambda }^2}{q^2\mathrm{\Lambda }^2}},`$ (9)
and the Fourier integral becomes ill-defined. Here,
$`\mathrm{\Lambda }\mu \mathrm{exp}\left[{\displaystyle \frac{2\pi }{\beta _0\alpha _S(\mu )}}\right]`$ (10)
is the $`\mu `$-independent integration constant of the 1-loop renormalization-group equation. Therefore, the above potential $`V_{\beta _0}(r)`$ can only be defined as a perturbation series in $`\alpha _S(\mu )`$. (Fourier integral of each term of the series is well-defined.)
When we examine the large-order behavior of this perturbation series, we find that it is an asymptotic series and has an intrinsic uncertainty
$`\delta V_{\beta _0}(r)\mathrm{\Lambda }=𝒪(300\mathrm{MeV}).`$ (11)
If we want to extract the quark mass from the spectrum of boundstates, this uncertainty in the potential is directly reflected to the uncertainty in the quark mass. It is because the quarkonium mass is determined as twice the quark pole mass minus the binding energy, and an uncertainty in the potential means an uncertainty in the binding energy. This implies that we cannot determine the quark mass to an accuracy better than $`𝒪(300\mathrm{MeV})`$.
Now we examine the series expansion of $`V_{\beta _0}(r)`$ and see how the uncertainty arises . If we perform the Fourier integration term by term,
$`\text{}V_{\beta _0}(r)`$ $`=`$ $`C_F\mathrm{\hspace{0.17em}4}\pi \alpha _S(\mu ){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d^3\stackrel{}{q}}{(2\pi )^3}\frac{e^{i\stackrel{}{q}\stackrel{}{r}}}{q^2}\left\{\frac{\beta _0\alpha _S(\mu )}{4\pi }\mathrm{log}\left(\frac{q^2}{\mu ^2}\right)\right\}^n}`$ (12)
$`=`$ $`C_F\mathrm{\hspace{0.17em}4}\pi \alpha _S(\mu ){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{\beta _0\alpha _S(\mu )}{4\pi }}\right\}^nf_n(r,\mu )\times n!,`$ (13)
the coefficients $`f_n(r,\mu )`$ can be determined from a generating function
$`\text{}F(r,\mu ;u)`$ $``$ $`{\displaystyle \frac{d^3\stackrel{}{q}}{(2\pi )^3}\frac{e^{i\stackrel{}{q}\stackrel{}{r}}}{q^2}\left(\frac{\mu ^2}{q^2}\right)^u}`$ (14)
$`=`$ $`{\displaystyle \frac{(\mu r)^{2u}}{4\pi r}}{\displaystyle \frac{1}{\mathrm{cos}(\pi u)\mathrm{\Gamma }(1+2u)}}`$ (15)
$`=`$ $`{\displaystyle \underset{n}{}}f_n(r,\mu )u^n.`$ (16)
Using this generating function, one easily obtains the asymptotic behavior of $`f_n(r,\mu )`$ for large $`n`$. The large-$`n`$ behavior of $`f_n(r,\mu )`$ determines the domain of convergence of the series expansion (16) at $`u=0`$.
Conversely from the structure of the pole of (15) nearest to $`u=0`$ (see Fig. 1), which limits the radius of convergence, one obtains The leading asymptotic behavior of $`f_n(r,\mu )`$ is same as that of the expansion coefficients of $`\mathrm{Res}[F;u=\frac{1}{2}]\times (u\frac{1}{2})^1`$. for $`n1`$
$`f_n(r,\mu ){\displaystyle \frac{1}{2\pi ^2}}\mu \times 2^n.`$ (17)
Note that this asymptotic behavior is independent of $`r`$. This means that, although each term of the potential $`V_{\beta _0}(r)`$ is a function of $`r`$, its dominant part for $`n1`$ is only a constant potential which mimics the role of the quark mass in the determination of the quarkonium spectrum.
Thus, asymptotically, the $`n`$-th term of $`V_{\beta _0}(r)`$ is given by
$`V_{\beta _0}^{(n)}C_F\mathrm{\hspace{0.17em}4}\pi \alpha _S(\mu )\times {\displaystyle \frac{\mu }{2\pi ^2}}\times \left\{{\displaystyle \frac{\beta _0\alpha _S(\mu )}{2\pi }}\right\}^n\times n!.`$ (18)
As we raise $`n`$, first the term decreases due to powers of the small $`\alpha _S`$; for large $`n`$ the term increases due to the factorial $`n!`$. Around $`n_0=2\pi /(\beta _0\alpha _S(\mu ))`$, $`V_{\beta _0}^{(n)}`$ becomes smallest. The size of the term scarcely changes within the range $`n(n_0\sqrt{n_0},n_0+\sqrt{n_0})`$; see Fig. 2.
We may consider the uncertainty of this asymptotic series as the sum of the terms within this range:
$`\delta V_{\beta _0}(r){\displaystyle \underset{n=n_0\sqrt{n_0}}{\overset{n_0+\sqrt{n_0}}{}}}\left|V_{\beta _0}^{(n)}\right|\mathrm{\Lambda }.`$ (19)
The $`\mu `$-dependence vanishes in this sum, and this leads to the claimed uncertainty.
In passing, we note that this asymptotic series is not Borel summable; a Borel summable series has terms alternating in sign, but the asymptotic series originating from the QCD infrared renormalon has terms with the same sign. We cannot circumvent the uncertainty by Borel summation of the series.
## 3 Renormalon Cancellation in the Total Energy of a $`q\overline{q}`$ system
Now we state how the problem can be circumvented. Consider the total energy of a color-singlet non-relativistic quark-antiquark pair:
$`E_{\mathrm{tot}}(r)2m_{\mathrm{pole}}+V_{\beta _0}(r).`$ (20)
It was found that the leading renormalon contained in the potential $`V_{\beta _0}(r)`$ gets cancelled in the total energy $`E_{\mathrm{tot}}(r)`$ if the pole mass $`m_{\mathrm{pole}}`$ is expressed in terms of the $`\overline{\mathrm{MS}}`$ mass. The potential and the pole mass are expressed in terms of the 1-loop running coupling $`\alpha _{1\mathrm{L}}(q)`$ as
$`V_{\beta _0}(r)={\displaystyle \frac{d^3\stackrel{}{q}}{(2\pi )^3}e^{i\stackrel{}{q}\stackrel{}{r}}C_F\frac{4\pi \alpha _{1\mathrm{L}}(q)}{q^2}},`$ (21)
$`m_{\mathrm{pole}}m_{\overline{\mathrm{MS}}}(\mu )+{\displaystyle \frac{1}{2}}\text{ }{\displaystyle }\text{ }q<\mu \text{ }{\displaystyle \frac{d^3\stackrel{}{q}}{(2\pi )^3}}C_F{\displaystyle \frac{4\pi \alpha _{1\mathrm{L}}(q)}{q^2}}.`$ (22)
The potential $`V_{\beta _0}(r)`$ is essentially the Fourier transform of the Coulomb gluon propagator exchanged between quark and antiquark; the difference of $`m_{\mathrm{pole}}`$ and $`m_{\overline{\mathrm{MS}}}`$ is essentially the infrared portion of the quark self-energy, see Fig. 3.
As we saw, the renormalon uncertainty is related to the “would-be pole” contained in $`\alpha _{1\mathrm{L}}(q)`$, cf. Eq. (9). The signs of the renormalon contributions are opposite between $`V_{\beta _0}(r)`$ and $`m_{\mathrm{pole}}`$ because the color charges are opposite between quark and antiquark while the self-enregy is proportional to the square of a same charge. Their magnitudes differ by a factor of two because both the quark and antiquark propagator poles contribute in the calculation of the potential whereas only one of the two contributes in the calculation of the self-energy. Expanding the Fourier factor in $`2m_{\mathrm{pole}}`$ in a Taylor series for small $`q`$,
$`e^{i\stackrel{}{q}\stackrel{}{r}}=1+i\stackrel{}{q}\stackrel{}{r}+{\displaystyle \frac{1}{2}}(i\stackrel{}{q}\stackrel{}{r})^2+\mathrm{},`$ (23)
the would-be pole contained in the leading term gets cancelled against $`V_{\beta _0}(r)`$<sup>§</sup><sup>§</sup>§ We are interested in the infrared region $`q\mathrm{\Lambda }`$. The expansion is justified since the typcal distance between quark and antiquark is much smaller than the hadronization scale $`\mathrm{\Lambda }r1`$. and consequently the renormalon contributions cancel.
As a result of this cancellation, the series expansion of the total energy in $`\alpha _S(\mu )`$ behaves better if we use the $`\overline{\mathrm{MS}}`$ mass instead of the pole mass. Residual uncertainty due to uncancelled pole can be estimated similarly as in the previous section and is suppressed as
$`\mathrm{\Lambda }\times (\stackrel{}{q}\stackrel{}{r})^2\mathrm{\Lambda }\times \left({\displaystyle \frac{\mathrm{\Lambda }}{\alpha _Sm_{\mathrm{pole}}}}\right)^2,`$ (24)
which is much smaller than the original uncertainty.
## 4 Extracting the $`\overline{\mathrm{MS}}`$ masses using the Full NNLO Result
To see how well the renormalon cancellation works, we examine extractions of the bottom and top quark masses using the full next-to-next-to-leading order (NNLO) result of the boundstate spectrum. The full NNLO formula for the lowest lying ($`1S`$) boundstate can be calculated from the Hamiltonian Eqs. (3)-(5) and is given by
$`M_{1S}=2m_{\mathrm{pole}}{\displaystyle \frac{4}{9}}\alpha _S(\mu )^2m_{\mathrm{pole}}[\mathrm{\hspace{0.17em}1}+\frac{\alpha _S(\mu )}{\pi }\{(11\frac{2}{3}n_l)L+(\frac{97}{6}\frac{11}{9}n_l)\}`$
$`+\left(\frac{\alpha _S(\mu )}{\pi }\right)^2\{(\frac{363}{4}11n_l+\frac{1}{3}n_l^2)L^2+(\frac{927}{4}\frac{193}{6}n_l+n_l^2)L`$
$`+(\frac{1793}{12}+\frac{2917\pi ^2}{216}\frac{9\pi ^4}{32}+\frac{275\zeta _3}{4})+(\frac{1693}{72}\frac{11\pi ^2}{18}\frac{19\zeta _3}{2})n_l+(\frac{77}{108}+\frac{\pi ^2}{54}+\frac{2\zeta _3}{9})n_l^2\}],`$ (25)
where $`L\mathrm{log}[\mu /(C_F\alpha _S(\mu )m_{\mathrm{pole}})]`$, and $`n_l`$ denotes the number of massless quarks. We may examine the size of each term of the above perturbation series. Alternatively we may rewrite the above expression in terms of the $`\overline{\mathrm{MS}}`$ mass and examine the series. Presently the relation between $`m_{\mathrm{pole}}`$ and $`\overline{m}m_{\overline{\mathrm{MS}}}(m_{\overline{\mathrm{MS}}})`$ is known up to three-loop order :
$`m_{\mathrm{pole}}=\overline{m}\times [\mathrm{\hspace{0.17em}1}+{\displaystyle \frac{4}{3}}\left(\frac{\alpha _S(\overline{m})}{\pi }\right)+\left(\frac{\alpha _S(\overline{m})}{\pi }\right)^2(1.0414n_l+13.4434)`$
$`+\left(\frac{\alpha _S(\overline{m})}{\pi }\right)^3(0.6527n_l^226.655n_l+190.595)].`$ (26)
First we apply the formula to the $`\mathrm{{\rm Y}}(1S)`$ state, for which $`n_l=4`$. Taking the input parameter as $`m_{\mathrm{pole}}=4.96`$ GeV/$`\overline{m}=4.22`$ GeV and setting $`\mu =\overline{m}`$ (i.e. expansion parameter is $`\alpha _S(\overline{m})=0.22`$),
$`M_{\mathrm{{\rm Y}}(1S)}`$ $`=`$ $`2\times (4.960.050.080.11)\mathrm{GeV}\text{(Pole-mass scheme)}`$ (27)
$`=`$ $`2\times (4.22+0.35+0.12+0.04)\mathrm{GeV}\text{(}\overline{\mathrm{MS}}\text{-scheme)}.`$ (28)
One sees that the series is not at all converging in the pole-mass scheme, whereas in the $`\overline{\mathrm{MS}}`$-scheme the series is converging quite nicely up to the calculated order. (See Sec. 6 for details of how we derived the series in the $`\overline{\mathrm{MS}}`$-scheme.) Comparing this with the experimental value $`M_{\mathrm{{\rm Y}}(1S)}=9.46037\pm 0.00021`$ GeV, one may extract the $`\overline{\mathrm{MS}}`$ bottom quark mass
$`\overline{m}_bm_{\overline{\mathrm{MS}}}(m_{\overline{\mathrm{MS}}})=4.22\pm 0.08\mathrm{GeV}.`$ (29)
One might think that, looking at the behavior of the above series, we may assign a smaller theoretical uncertainty. The present uncertainty is, however, dominated by non-perturbative uncertainties other than the renormalon contributions. Thus, presently $`\overline{m}_b`$ is determined to 2% accuracy . It seems to be fairly good in view of the fact that its major part is controlled by perturbative QCD.
Next we turn to the (remnant of) “toponium”, for which $`n_l=5`$. At future linear $`e^+e^{}`$ or $`\mu ^+\mu ^{}`$ colliders, the top quark mass will be determined to high accuracy from the shape of the $`t\overline{t}`$ total production cross section in the threshold region. The location of a sharp rise of the cross section is determined mainly from the mass of the lowest lying ($`1S`$) $`t\overline{t}`$ resonance, so we will be able to measure the resonance mass and extract the top quark mass. Similarly to the previous case, we set $`m_{\mathrm{pole}}=174.79`$ GeV/$`\overline{m}=165.00`$ GeV, $`\alpha _S(\overline{m})=0.1091`$ and obtain
$`M_{1S}`$ $`=`$ $`2\times (174.790.460.400.28)\mathrm{GeV}\text{(Pole-mass scheme)}`$ (30)
$`=`$ $`2\times (165.00+7.20+1.24+0.22)\mathrm{GeV}\text{(}\overline{\mathrm{MS}}\text{-scheme)}.`$ (31)
In Figs. 4 are shown the convergence properties of the above series together with the corresponding cross sections.
In the pole-mass scheme, the convergence is very slow. According to the renormalon argument, also uncalculated higher order terms would not become much smaller. On the other hand, in the $`\overline{\mathrm{MS}}`$-scheme the series shows a healthy convergence behavior. For top quark, non-perturbative uncertainties are much smaller than the present perturbative theoretical uncertainty. Thus, from the above series we estimate that $`\overline{m}_t`$ can be determined to around 100 MeV accuracy .
## 5 Physical Implications
Let us discuss some physical implications of renormalon cancellation in the heavy quarkonium system. Firstly, as already mentioned, we expect that gluons with wavelength much longer than the size of the quarkonium cannot couple to this system. Hence, we expect that infrared gluons with momentum transfer $`q\alpha _Sm`$ should decouple from the expression of $`E_{\mathrm{tot}}`$. This is a naive expectation based on classical dynamics. Such understanding should be valid when it is described by the bare QCD Lagrangian without large quantum corrections. The $`\overline{\mathrm{MS}}`$ mass is closely related to the bare mass of quark; only ultraviolet divergences are subtracted. On the other hand, the pole mass has much more intricate relation to the bare mass, because the relation includes in addition infrared dynamics of the quantum correction to the quark self-energy. In this sense it would be natural to expect the decoupling phenomenon to be realized when the $`\overline{\mathrm{MS}}`$ mass is used to express $`E_{\mathrm{tot}}`$.
Secondly, the pole mass of a quark is ill-defined beyond perturbation theory. It can be determined only when the quark can propagate an infinite distance. Generally accepted belief is that when quark and antiquark are separated beyond a distance $`\mathrm{\Lambda }^1`$ the color flux is spanned between the two charges due to non-perturbative effects and the free quark picture is no longer valid. On the other hand, it is natural to consider the total energy (or the mass) of a quarkonium which is a color-singlet state. It can propagate for a long time and the notion of its mass is not limited by the hadronization scale.
Thirdly, renormalon cancellation seems to be a universal feature which occurs process independently. The same phenomenon was known e.g. in the QCD corrections to the $`\rho `$-parameter and $`B`$ decays . Therefore, the $`\overline{\mathrm{MS}}`$ mass, which is determined accurately from the quarkonium spectum, would be more suited than the pole mass for an input parameter in describing other physical processes.
## 6 How to Cancel Renormalons in the Quarkonium Spectrum
There is one non-trivial point in realizing renormalon cancellation in the perturbation series of the quarkonium spectrum. When the pole mass and the binding energy are given as series in $`\alpha _S`$, renormalon cancellation takes place between the terms whose orders in $`\alpha _S`$ differ by one :
(32)
cf. Eqs. (25) and (26). Intuitively this can be seen from the diagrams shown in Fig. 5.
An additional power of $`\alpha _S`$ in the binding energy is provided by the inverse of the Bohr radius $`r^1\alpha _Sm`$,
$`\text{i.e.}C_F{\displaystyle \frac{\alpha _S^n}{r}}\alpha _S^{n+1}m.`$ (33)
Still, one might wonder how cancellation can ever take place between different orders in $`\alpha _S`$ for any value of $`\alpha _S`$. So we demonstrate the cancellations at large orders in a specific example. The mass of the $`1S`$ boundstate can be written in the form of an expectation value as
$`M_{1S}=1S\left|2m_{\mathrm{pole}}+\widehat{H}\right|\mathrm{\hspace{0.17em}1}S`$ (34)
using the Hamiltonian (2) and its 1S energy eigenstate $`|\mathrm{\hspace{0.17em}1}S`$. We are interested in the leading renormalon contributions, so we replace
$`2m_{\mathrm{pole}}+\widehat{H}2m_{\mathrm{pole}}+V_{\beta _0}(r).`$ (35)
The energy eigenstate can be expanded in $`1/c`$:
$`|\mathrm{\hspace{0.17em}1}S=|\mathrm{\hspace{0.17em}1}S^{(0)}+{\displaystyle \frac{1}{c}}|\mathrm{\hspace{0.17em}1}S^{(1)}+{\displaystyle \frac{1}{c^2}}|\mathrm{\hspace{0.17em}1}S^{(2)}+\mathrm{}`$ (36)
Renormalon cancellation takes place in arbitrary combination of $`1S^{(i)}|\mathrm{}|1S^{(j)}`$, but for simplicity we evaluate only the following part:
$`M_{1S}^{(0)}=1S^{(0)}|\mathrm{\hspace{0.17em}2}m_{\mathrm{pole}}+V_{\beta _0}(r)|1S^{(0)}=2m_{\mathrm{pole}}+1S^{(0)}|V_{\beta _0}(r)|1S^{(0)}.`$ (37)
The second term corresponds to the binding energy and we may evaluate it at each order of the perturbation series:
$`1S^{(0)}|V_{\beta _0}(r)|1S^{(0)}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑rr^2|R(r)|^2V_{\beta _0}(r)`$ (38)
$`=`$ $`{\displaystyle \frac{1}{2}}C_F^2\alpha _S(\mu )^2m_{\mathrm{pole}}\times {\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left\{{\displaystyle \frac{\beta _0\alpha _S(\mu )}{4\pi }}\right\}^ng_n(\mu a_0)\times n!,`$
where the zeroth-order $`1S`$ Coulomb wave function is given by
$`R(r)={\displaystyle \frac{2}{a_0^{3/2}}}e^{r/a_0},a_0=\left(\frac{1}{2}C_F\alpha _S(\mu )m_{\mathrm{pole}}\right)^1:\text{Bohr radius}.`$ (39)
$`g_n(\mu a_0)`$’s are polynomials of $`\mathrm{log}(\mu a_0)`$. Using the generating function method, one obtains the asymptotic form
$`g_n(\mu a_0){\displaystyle \frac{2}{\pi }}\mu a_0\times 2^n{\displaystyle \frac{1}{\alpha _S}}.`$ (40)
Thus, for $`n1`$, it becomes proportional to $`\alpha _S^1`$ and effectively shifts the order of $`\alpha _S`$. By setting $`\mu =\overline{m}`$, it is easy to check that in this example the leading renormalon cancels as in Eq. (32) and the residual piece behaves as
$`\left[n\text{-th term of}M_{1S}^{(0)}\right]\alpha _S(\overline{m})\overline{m}\times n!\times \left\{{\displaystyle \frac{\beta _0\alpha _S(\overline{m})}{6\pi }}\right\}^n.`$ (41)
It follows that
$`\delta M_{1S}^{(0)}\mathrm{\Lambda }\times \left({\displaystyle \frac{\mathrm{\Lambda }}{\alpha _S\overline{m}}}\right)^2\mathrm{\Lambda }.`$ (42)
From this example one learns that the cancellation of renormalon contribution between shifted orders $`A_n\alpha _S^n`$ and $`B_{n+1}\alpha _S^{n+1}`$ should be properly taken into account when expressing the boundstate mass as a perturbation series. There are many different prescriptions to accomplish this. We derived the series (28) and (31) in the following manner. We have rewritten Eq. (25) as
$`M_{1S}`$ $`=`$ $`2m_{\mathrm{pole}}\times \left[\mathrm{\hspace{0.17em}1}+{\displaystyle \underset{n=2}{\overset{4}{}}}P_n\alpha _S(\overline{m})^n\right]`$ (43)
$`=`$ $`2\overline{m}\times \left[\mathrm{\hspace{0.17em}1}+{\displaystyle \underset{n=1}{\overset{3}{}}}A_n\alpha _S(\overline{m})^n\right]\times \left[\mathrm{\hspace{0.17em}1}+{\displaystyle \underset{n=2}{\overset{4}{}}}P_n\alpha _S(\overline{m})^n\right],`$
where $`P_n`$’s are polynomials of $`\mathrm{log}\left(\alpha _S(\overline{m})\right)`$ and $`A_n`$’s are just constants independent of $`\alpha _S(\overline{m})`$. We identified $`P_n\alpha _S^n`$ as order $`\alpha _S^{n1}`$ and then reduced the last line to a single series in $`\alpha _S`$.
Parametric accuracy of the last terms of Eqs. (28) and (31) is $`\alpha _S^3\overline{m}`$. In order to improve the accuracy to $`\alpha _S^4\overline{m}`$, we need to know further (a) the exact 4-loop relation between $`m_{\mathrm{pole}}`$ and $`\overline{m}`$, and (b) the binding energy at order $`\alpha _S^5`$ in the large-$`\beta _0`$ approximation.
## 7 Some Questions
One can ask some interesting questions related to the renormalon problem and extraction of quark mass. In the case of top quark, its mass will also be measured from the invariant mass distribution of decay products of the top quark in future hadron collider experiments and in future $`e^+e^{}`$ collider experiments. What is the mass extracted from the peak position of the Breit-Wigner distribution? Naively it is the pole mass because the peak position is determined from the position of the pole of the top quark propagator in perturbative QCD. On the other hand, as we have seen, the pole mass suffers from a theoretical uncertainty of $`𝒪`$(300 MeV). Since future experiments will be able to determine the peak position of the invariant mass distribution to $`𝒪`$(100 MeV) accuracy, we will indeed face this serious conceptual problem. We believe that renormalon cancellation also takes place in this physical quantity. But the problem lies in the fact that we have no reliable theoretical method to calculate the invariant mass distribution of realistic color-singlet final states. Rather we calculate the invariant mass of final partons which do not combine to color-singlet state. Another question is whether the peak position of invariant mass distribution measured in hadron collider experiments will be the same as that measured in $`e^+e^{}`$ collider experiments. The answer is probably no, since at hadron colliders top and antitop are pair-created not necessarily in a color-singlet state.
These are very interesting questions which are worth further studies.
## Acknowledgements
The author is grateful to K. Melnikov, A. Hoang, M. Tanabashi and H. Ishikawa for fruitful discussions. Part of this work is based on a collaboration with T. Nagano.
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# A Luttinger’s theorem revisited
## I INTRODUCTION
Landau’s (1957) phenomenological theory for the low-lying excited states of many-particle systems is based on the assumption that these states are characterised by a distribution of a dilute gas of quasi-particles (QPs), implying that interaction amongst these QPs may, to a good approximation, be neglected (Nozières 1964, Pines and Nozières 1966, Abrikosov, Gorkov and Dzyaloshinski 1975) (for a comprehensive review see Platzman and Wolff (1973)). <sup>*</sup><sup>*</sup>* To be precise, Landau’s Fermi-liquid theory is applicable for temperatures $`T`$ in the range $`(T_c,T^{}`$), where $`T^{}`$ and $`T_c`$ stand, respectively, for the coherence and the transition temperature; for $`T`$ below $`T_c`$ the system is in a superconducting state (Kohn and Luttinger 1965). For details as well as a comprehensive discussion of instabilities of Fermi liquids see (Metzner, Castellani and Di Castro 1998). This ideal situation corresponds to the case in which QPs are the exact one-particle eigenstates of the many-body Hamiltonian, or, equivalently, the energies of the QPs are infinitely sharply defined. The success of Sommerfeld’s independent-particle model (for example Ashcroft and Mermin (1981, Ch. 2)) for conduction electrons in simple metals, despite the apparent non-negligible strength of the electron-electron Coulomb interaction, can be viewed as a most strong pillar of Landau’s phenomenological theory. Incorporation of the residual two-body interaction among the ideal QPs (at least beyond the Hartree-Fock or the static exchange approximation), which evidently should be weaker than the bare Coulomb interaction (insofar as low-energy scattering processes are concerned), Otherwise the independent-particle model would not be an appropriate zeroth-order approximation. renders the QP states non-stationary. In other words, the actual QPs, if such entities can at all be meaningfully defined (§ III), are merely approximate one-particle eigenstates of the many-body Hamiltonian; the non-stationary nature of a QP ‘eigenstate’ signifies that its energy, in contrast with that of an exact eigenstate, is un-sharp, or the spectral function of QPs consists of peaks with finite widths. Thus the true, that is interacting, QPs correspond to superpositions of nearly degenerate (as well as degenerate) many-body eigenstates, so that an initially well-defined QP ceases to be particle-like (in the sense of possessing a reasonably well-defined energy) with the passage of time.
The information with regard to QPs is contained in the single-particle Green function $`G`$. The formal Lehmann (1954) spectral representation of this function in the energy domain reveals that energies of the single-particle excitations, the QPs, can be identified with the singular points (often inappropriately designated as ‘poles’ — § V) along the real energy axis ($`\epsilon `$-axis) (Fetter and Walecka 1971). Such observation is of little practical relevance when systems in the thermodynamic limit are concerned. For macroscopic systems, the continuum of the energy levels From dimensional considerations it follows that for a system with linear dimension $`L`$, the separations between the energy levels scale like $`L^2`$ (Landau and Lifshitz 1980, p. 14). gives rise to branch points (which are not isolated singularities) and branch cuts along the energy axis. This is naturally consistent with the above-indicated observation with regard to the broad rather than delta-function-like spectral function of QPs. Such a dramatic change in the analytic structure of $`G(\epsilon )`$ gives rise to a number of effects the subtleties of which are often overlooked in the literature. In this work we discuss in some detail a number of these, insofar as they are relevant to the main objective of our work. One of these effects is associated with the fact that whereas a reasonably sharp peak in the spectral function may be described by a single pole with some small imaginary part, such description is in violation of the fact that $`G(\epsilon )`$ is a perfectly regular function on the entire complex energy plane, excluding the real $`\epsilon `$-axis. In § III we show that such complex poles do not simply correspond to the analytic continuation of $`G(\epsilon )`$ into the complex energy plane (i.e. the physical Riemann sheet (RS)), rather to the analytic continuation of the latter to a RS; we refer to this as a non-physical RS (see Appendix A).
The main complication, from the restricted viewpoint of the present paper, arises from the branch points of the Green function $`G(k;\epsilon )`$, in particular that at $`\epsilon =\epsilon _F`$ which can be shown to be also a branch point of the self-energy, $`\mathrm{\Sigma }(k;\epsilon )`$. From this it follows that, for instance, such expression as (Hugenholtz 1957, DuBois 1959b, Luttinger 1961) $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, cannot be considered as the leading term in the Taylor series expansion of the self-energy operator at $`\epsilon =\epsilon _F`$. Thus there exists no a priori reason for the dispersion of the QP energies <sup>§</sup><sup>§</sup>§ At places in the present text, such as here, we refer to ‘QP energies’ even though QPs may not be well-defined. This is justified by the fact that even in these cases the equation for the ‘QP energies’ (i.e. Eq. (18 below) does have a solution — the ill-defined nature of QPs in these cases is associated with the fact that the single-particle eigenstates corresponding to these energies (see Eq. (14) below) do not describe particle-like excitations (for details see in particular § IV.C.). to be a ‘smooth’ function of $`k`$ in a neighbourhood Throughout the present work by ‘neighbourhood’ we refer to an open non-vanishing interval (e.g. along the real axis) or domain (on the complex plane) which, however, may be arbitrarily small. of $`k_F`$ (for specification see further on) similar to that of non-interacting QPs, namely We have $`\epsilon _k^0:=\mathrm{}^2k^2/[2m_e]\epsilon _{k_F}^0+(\mathrm{}^2k_F/m_e)(kk_F)+(\mathrm{}^2/[2m_e])(kk_F)^2`$. The important feature of this energy dispersion is that it is of the general form $`\epsilon _k^0=\epsilon _{k_F}^0+\varsigma (kk_F)+o(kk_F)`$ for $`kk_F`$ (see § IV.C), where $`\varsigma `$ is a finite constant (explicitly, in the case at hand we have $`\varsigma \mathrm{}^2k_F/m_e`$) and $`o(kk_F)`$ stands for a function that in comparison with $`(kk_F)`$ is vanishingly small as $`kk_F`$. In this respect the property $`\epsilon _k^0k^2`$ is to us of relatively minor relevance. We point out that the coefficient of the linear term in the QP energy dispersion (here $`\varsigma `$), if such a coefficient at all exists, has the relevance that through it such attributes as the Fermi velocity $`v_F`$ and effective mass $`m_e^{}`$ can be assigned to the corresponding QP (see § IV.C, specifically Eqs. (24) and (25)). $`\epsilon _k^0k^2`$, with $`k:=𝐤`$. In fact, for the special case of ‘Luttinger liquids’ (Haldane 1980, 1981, Luttinger 1963, Mattis and Lieb 1965, Dover 1968, Anderson 1997), without hereby specifying the dimension of the spatial space to which the corresponding system is confined, or marginal Fermi liquids (Varma, et al., 1989, Littlewood and Varma 1991, Kotliar, et al., 1991), we explicitly demonstrate that $`\epsilon _k`$ is not differentiable at $`k=k_F`$ (see §§ IV.C and VI). As we shall emphasise later in this paper (§ IV.C), QPs in Fermi liquids are special in that their energy dispersion in the close vicinity of the Fermi surface is, apart from a scaling factor arising from the re-normalisation of the electron mass or the Fermi velocity, non-interacting-like. Our close inspection of the proof of the Luttinger (1961) theorem (§ IV), the contents of which we have spelled out in the Abstract, reveals that this ‘non-interacting-like’ assumption with regard to the dispersion of the QP energies is implicit in Luttinger’s (1961) proof, so that Luttinger’s theorem in essence amounts to a statement concerning the consistency of this assumption with the property $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, for $`\epsilon \epsilon _F`$. In view of our rigorous demonstration with regard to non-differentiability of $`\epsilon _k`$ at $`k=k_F`$ for Luttinger liquids, or marginal Fermi liquids, which invalidates Luttinger’s implicit assumption, we are left to conclude that in principle nothing precludes existence of a non-Fermi-liquid behaviour for metallic systems <sup>\**</sup><sup>\**</sup>\** Our emphasise throughout this work on the metallic nature of the systems under consideration corresponds to the fact that non-metallic systems are by definition non-Fermi liquids. in spatial dimensions larger than unity. In particular, contrary to the general belief, the break-down of the many-body perturbation theory (Luttinger 1961, Anderson 1988, 1989, Varma, et al. 1989, Anderson 1990a,b, 1991, Littlewood and Varma 1991, Kotliar, et al. 1991, Anderson 1992, 1993, 1997) for a metallic system and the nature of the low-lying single-particle excitations in this stand in no direct relationship. <sup>††</sup><sup>††</sup>†† We should like to emphasise that often from the contexts of arguments presented in the literature, in disapproval of the many-body perturbation theory, it is not evident what precisely is meant by the ‘perturbation theory’. In the present work we consider a series for the self-energy operator in terms of the interacting single-particle Green function, each term of which is diagrammatically represented by a skeleton diagram, as constituting a perturbation series. This we do entirely in the spirit of Luttinger’s (1961) work. We should point out that such a series is not a power series in the coupling constant of the electron-electron interaction and consequently does not conform to the orthodox definition for perturbation series.
For clarity of discussions, we briefly specify the scheme by which we classify a metallic system as a Fermi liquid. This scheme, which to our best knowledge has not been utilised earlier in a similar context, has the advantage of being general and independent of the details of the metallic system to which it is applied. To this end, we first point out that the behaviour $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, for $`\epsilon \epsilon _F`$, which is specific to Fermi liquids, leads to (see §§ IV.C, V and Appendix C) $`\mathrm{Re}\mathrm{\Sigma }(k,\epsilon )\mathrm{\Sigma }(k,\epsilon _F)\beta _k(\epsilon \epsilon _F)`$ for $`\epsilon \epsilon _F`$, with $`\beta _k`$ a constant to be specified in Appendix C. The distinguishing aspect of $`\mathrm{\Sigma }(k,\epsilon )`$ whose real and imaginary parts have such asymptotic forms is as follows.
Condition A: $`\mathrm{\Sigma }(k_F;\epsilon )`$ is a continuously differentiable function <sup>‡‡</sup><sup>‡‡</sup>‡‡ A function $`f(x)`$ is continuously differentiable in the closed interval $`[a,b]`$ provided firstly that it is continuous in $`[a,b]`$ and secondly that its derivative exists at all points of the open interval $`(a,b)`$ and coincides at all such points with a function which is continuous in $`[a,b]`$; $`f(x)`$ is continuously differentiable in a neighbourhood of $`x=x_0`$ if $`x_0`$ is interior to a finite interval, such as $`[a,b]`$, in which $`f(x)`$ is continuously differentiable. of $`\epsilon `$ in a neighbourhood of $`\epsilon _F`$.
This condition suffices to guarantee a finite discontinuity $`Z_{k_F}`$ in the momentum-distribution function $`𝗇(k)`$ at $`k=k_F`$ (§ V). A finite $`Z_{k_F}`$ is generally considered to signify a Fermi liquid. If this merely evidences a Fermi liquid, then condition (A) would constitute the necessary and sufficient condition for a metallic system to be one (see § V). However, condition (A) can be shown (see §§ IV.C and V) not to be sufficient for the QP energy to possess a ‘well-defined’ dispersion in a neighbourhood of $`k=k_F`$, that is one which is expressible in terms of a continuously differentiable function of $`k`$ in this neighbourhood. For the existence of such a dispersion, the following condition is also required to be satisfied.
Condition B: $`\mathrm{\Sigma }(k;\epsilon _F)`$ is a continuously differentiable function of $`k`$ in a neighbourhood of $`k=k_F`$.
We therefore choose to consider a metallic system as a Fermi liquid provided the corresponding $`\mathrm{\Sigma }(k;\epsilon )`$ satisfies both condition (A) and condition (B); this specifies the notion of ‘smooth’ introduced above. According to this, which is physically the most sound definition for Fermi liquids, <sup>\**</sup><sup>\**</sup>\** The notion of a QP without an associated energy dispersion is devoid of physical significance. a finite $`Z_{k_F}`$, although necessary, is not sufficient for a system to be a Fermi liquid. We explicitly demonstrate (§ VI) that whereas $`\mathrm{\Sigma }(k;\epsilon _F)`$ corresponding to a Fermi liquid is (by definition) a continuously differentiable function of $`k`$ in a neighbourhood of $`k=k_F`$, the same may not apply to $`\mathrm{\Sigma }(k;\epsilon )`$ when $`\epsilon \epsilon _F`$. We further demonstrate (§ VI) that $`\mathrm{\Sigma }(k;\epsilon )`$ pertaining to a metallic system (whether Fermi liquid or otherwise), which is not continuously differentiable with respect to $`k`$ in a neighbourhood of $`k=k_0`$ (where $`k_0`$ would possibly be $`k_F`$) for some $`\epsilon \epsilon _F`$, is not continuously differentiable with respect to $`k`$ in a neighbourhood of $`k=k_0`$ for any $`\epsilon `$, with the possible exception of $`\epsilon =\epsilon _F`$; this exceptional instance is (by definition) reserved for Fermi liquids (see above). To be explicit, consider $`\alpha _k`$ and $`\beta _k`$ which feature in the above expressions for the self-energy of a Fermi liquid (see also § V). Our above statements imply that these may not necessarily be continuously differentiable functions of $`k`$ in a neighbourhood of $`k=k_F`$; possible singularities associated with $`\alpha _k/k`$ and $`\beta _k/k`$, as $`kk_F`$, are suppressed by $`(\epsilon \epsilon _F)^2`$ and $`(\epsilon \epsilon _F)`$ respectively for $`\epsilon =\epsilon _F`$. In § VI we establish relationships between the behaviour of $`\mathrm{\Sigma }(k;\epsilon )`$ for $`k`$ and $`\epsilon `$ in neighbourhoods of $`k_F`$ and $`\epsilon _F`$ respectively and that of the momentum distribution function for $`k`$ approaching $`k_F`$.
Aside from the above-indicated ‘smoothness’ assumption with regard to the dispersion of the QP energies which is implicit in Luttinger’s (1961) proof, there is a second aspect specific to this proof that makes the theorem inapplicable to metallic systems of particles interacting through the long-range Coulomb interaction function. Luttinger’s approach is based on the perturbation expansion of the self-energy operator in terms of the exact single-particle Green function (the terms in this expansion are represented by means of skeleton self-energy Feynman diagrams (Luttinger and Ward 1960)) and the bare electron-electron interaction function. It can be shown (§ IV.D) that a certain class of terms in this expansion are unbounded as a consequence of non-integrable singularities corresponding to zero momentum-transfer scattering processes (infrared divergence). Therefore, for such systems a formal term-by-term analysis of the perturbation series for the self-energy operator can only be meaningful when the perturbation expansion is in terms of the dynamically-screened interaction function (Hubbard 1957) $`W(\epsilon )`$ (for a comprehensive discussion see Mattuck (1992, § 10.4)). In contrast to the static bare Coulomb interaction $`v_c`$, $`W(\epsilon )`$ is a function with singularities along the real energy (i.e. $`\epsilon `$) axis (branch points, branch cuts, etc. — similar to the Green function and the self-energy operator) and a Lehmann-type spectral representation of this function reveals that these singularities coincide with energies of the neutral elementary excitations of the system (i.e. energies of excited $`N`$-electron states as measured with respect to the energy of the $`N`$-electron ground state). It follows that an analysis of the self-energy operator along the lines of Luttinger (1961) requires additional knowledge with regard to the dispersion of these neutral excitations. In other words, in this case the dispersion of the QP energies and that of neutral excitations cannot be considered separately, but must be dealt with in conjunction. This is of relevance particularly to uniform two-dimensional systems of electrons whose spectrum of coherent neutral excitations (i.e. plasmons) is gap-less, implying some non-negligible amount of interference amongst the neutral and the single-particle, i.e. QP, excitations.
Our present work has been motivated by a long-standing discussion in the literature with regard to the unusual properties of layered high-$`T_c`$ compounds in their normal states. On the one hand these unusual properties have been ascribed to the fact that the normal metallic states of these systems not Fermi-liquid states (Anderson 1988, 1989, Varma, et al., 1989, Anderson 1990a,b, 1991, Littlewood and Varma 1991, Kotliar, et al., 1991, Anderson 1992, 1993, 1997) while, on the other hand, according to the Luttinger (1961) theorem dealt with in the present work, metallic interacting systems in spatial dimensions larger than unity must be Fermi liquids. As emphasised by Luttinger (1961), the validity of the Luttinger theorem is conditional to that of the many-body perturbation theory for the self-energy operator to all orders. In this light, it is asserted that a non-Fermi-liquid metallic behaviour of a system would necessarily imply break-down of the many-body perturbation theory for this system (Anderson 1997). These observations give rise to the following questions.
(1) Can interacting electron systems in spatial dimensions $`d`$ larger than unity, in particular for $`d=2`$, be non-Fermi liquids?
(2) does the many-body perturbation theory break down in cases where the metallic systems under consideration are non-Fermi liquids?
Concerning question (1), the analyses presented in this work lead us to the conclusion that in principle nothing stands in the way of realisation of non-Fermi-liquid states in spatial dimensions $`d`$ larger than unity. In fact the very strict conditions (A) and (B), introduced above, imposed on the self-energy of Fermi liquids, suggests that non-Fermi liquid metallic systems (in $`d>1`$) may not be as uncommon as generally perceived. We do not touch upon the above question (2) and refrain from statements that are likely to be speculative at this stage. In this connection we point out that, since the existence of non-Fermi liquids (in particular for $`d=2`$) and the validity of the many-body perturbation theory as applied to these are not a priori mutually exclusive, the question with regard to breakdown of this theory is not as urgent as it would be otherwise. This, however, does not diminish the relevance of question (2). We should like to emphasise that in calculating $`\mathrm{\Sigma }(k;\epsilon )`$, whether perturbatively or otherwise, account has to be taken of the fact that any approximation scheme that involves indiscriminate Taylor expansions of functions of energy and momentum (which expansions implicitly imply continuous differentiability of these functions in the pertinent regions of energy and momentum) may inhibit a non-differentiable $`\mathrm{\Sigma }(k;\epsilon )`$ as function of both $`k`$ and $`\epsilon `$, and therefore a non-Fermi liquid behaviour, from being obtained.
The above questions, (1) and (2), have been subject of intensive study in recent years, the main body of the results thus far obtained pointing towards a Fermi-liquid state in two spatial dimensions. In spite of these, it has as yet not proved possible unequivocally to rule out the existence of non-Fermi-liquid states in two-dimensional metallic systems. This can be ascribed to two main reasons.
(i) Owing to the complexity of the many-body problem at hand, application of reliable theoretical techniques must of necessity be accompanied by simplifying approximations, the integrity of which may be a matter of dispute.
(ii) Even relevance of certain calculated quantities to the problem at hand has been matter of debate.
Both these aspects are aptly represented in the following: Engelbrecht and Randeria (1990) have found no evidence for the breakdown of the many-body perturbation theory and the Fermi-liquid theory in a dilute two-dimensional system of fermions interacting through a short-range repulsive potential, contradicting the suggestion made by Anderson (1990b). It appears, however, that the phase shift as calculated by Engelbrecht and Randeria is not that which is encountered in the arguments by Anderson (see Engelbrecht and Randeria 1991, Anderson 1991). <sup>\*†</sup><sup>\*†</sup>\*† For a detailed discussion of the singular effective interaction amongst QPs, as encountered in the arguments by Anderson (1991), see Stamp (1993); see also Houghton, Kwon, and Marston (1994).
Work by Fujimoto (1990), Fukuyama, Narikiyo and Hasegawa (1991) and Fukuyama, Hasegawa and Narikiyo (1991) on the two-dimensional Hubbard model with repulsive on-site interaction $`U`$ within the $`t`$-matrix approximation (which is appropriate to the low-density limit) has indicated that, whereas for $`kk_F`$, $`\mathrm{Im}\mathrm{\Sigma }(k,\epsilon )`$ retains the conventional Fermi-liquid form in three spatial dimensions (presented above), <sup>\*‡</sup><sup>\*‡</sup>\*‡ Throughout, $`^{}`$ indicates that the corresponding asymptotic relations are correct up to multiplicative constants. $`\mathrm{Im}\mathrm{\Sigma }(k=k_F,\epsilon )^{}(\epsilon \epsilon _F)^2\mathrm{ln}|\epsilon \epsilon _F|`$ as $`\epsilon \epsilon _F`$. This singular contribution, obtained earlier by Hodges, Smith and Wilkins (1971) and Bloom (1975) concerning two-dimensional Fermi systems, is not sufficiently strong to render the self-energy non-Fermi-liquid like. It can be shown (see Appendix C) that the associated $`\mathrm{Re}\mathrm{\Sigma }(k_F;\epsilon )\mathrm{\Sigma }(k_F;\epsilon _F)^{}(\epsilon \epsilon _F)`$ so that $`\mathrm{\Sigma }(k_F;\epsilon )`$ is a continuously differentiable function of $`\epsilon `$ in a neighbourhood of $`\epsilon =\epsilon _F`$ and thus, on account of condition (A) discussed above, gives rise to a finite $`Z_{k_F}`$. Serene and Hess (1991), using the conserving ‘fluctuating-exchange (FLEX) approximation’ on a finite but large lattice, have also not found evidence for a non-Fermi-liquid behaviour in the two-dimensional Hubbard model. In a more recent work, Yokoyama and Fukuyama (1997) have arrived at the conclusion that in the two-dimensional Hubbard model with repulsive on-site interaction, the process of forward scattering gives rise to an anomalous contribution to the self-energy and consequently to a vanishing quasi-particle weight factor $`Z_{k_F}`$. This result, Yokoyama and Fukuyama (1997) asserted, would be “a microscopic demonstration of the claim by Anderson \[(1990a,b), (1993)\]”. We (Farid 1999a) have shown that the indicated anomalous contribution is entirely a consequence of the violation of a crucial symmetry in the momentum space (associated with the time-reversal symmetry of the ground state of the system under consideration) by the particle-particle correlation function employed by these workers.
Momentum-space perturbative renormalisation-group calculations by Shankar (1991, 1994) have also borne out a Fermi-liquid picture of fermions in two spatial dimensions. Shankar (1994) enumerated, however, a number of possibilities that in principle may render this finding, in its generality, invalid (see § XI in (Shankar 1994)). Work by Castellani, Di Castro (1994), Castellani, Di Castro and Metzner (1994), and Metzner, Castellani and Di Castro (1998) modelled on the treatment of one-dimensional interacting systems put forward by Dzyaloshinskiǐ and Larkin (1973), which exploits the conservation laws and the associated Ward identities, and treats $`d`$, the spatial dimension, as a real variable, establishes that metallic systems with strong forward scattering (i.e. that corresponding to small momentum transfers) are Fermi liquids in $`d>1`$, even though for $`d3`$ some properties of these systems differ from those pertaining to conventional Fermi liquids in $`d=3`$. Here the limited applicability of the formalism to systems with dominant forward scattering does not rule out non-Fermi-liquid systems in $`d>1`$. In a forthcoming work (Farid 1999b) we present a detailed study concerning some limiting aspects associated with the technique employed in these studies.
Systems of particles interacting through long-range repulsive interaction functions $`v(𝐫𝐫^{})`$ which for large $`𝐫𝐫^{}`$ behave like $`1/𝐫𝐫^{}^{2d}`$, when $`1<d<2`$, and like $`\mathrm{ln}𝐫𝐫^{}`$, when $`d=2`$, have been shown to resemble one-dimensional Luttinger liquids for $`1<d<2`$ and a “$`Z_{k_F}=0`$ Fermi liquid” for $`d=2`$ (Bares and Wen 1993, Kwon, Houghton and Marston 1995). We note that in metals the Coulomb interaction function, which behaves like $`1/𝐫𝐫^{}`$, is screened through the mediation of the charge polarisation fluctuations (Hubbard 1957); within the framework of the random-phase approximation (RPA), the screened interaction function in the static limit can be shown to behave like $`\mathrm{cos}(2k_F𝐫𝐫^{})/𝐫𝐫^{}^3`$ (Fetter and Walecka 1971, pp. 178 and 179; Ashcroft and Mermin 1981, p. 343). It follows that the above long-range interactions should have their origin in processes that are not related to the electrostatic electron-electron interaction. The only known interaction that remains long-ranged in metals corresponds to that of electronic currents mediated by the exchange of transverse photons (Holstein, Norton, and Pincus 1973, Reizer 1989, 1991). Lack of the static screening of this interaction leads to the asymptotic behaviour of the self-energy $`\mathrm{\Sigma }(k;\epsilon _k)`$ (on-the-mass-shell self-energy), For the equation satisfied by the QP energy dispersion $`\epsilon _k`$ see Eq. (22) below. as $`kk_F`$ (or $`\epsilon _k\epsilon _F`$), to be that characteristic of the marginal Fermi liquids (Varma, et al., 1989, Littlewood and Varma 1991, Kotliar, et al. 1991); on the other hand, for the fixed $`kk_F`$, $`\mathrm{\Sigma }(k;\epsilon )`$ behaves Fermi-liquid like as $`\epsilon \epsilon _F`$ (Holstein, Norton, and Pincus 1973). We note that the $`\mathrm{\Sigma }(k;\epsilon )`$ considered here does not involve the effects of the Coulomb interaction beyond a mean-field level; it merely describes the self-energy of the current-current interaction whose bare coupling constant is proportional to the square of the ratio of the Fermi velocity to the light velocity in vacuum ($`10^4`$) and therefore is not of dominating influence except in extremely pure metals and at very low temperatures. We shall not enter into further details concerning this subject here <sup>\*∥</sup><sup>\*∥</sup>\*∥ In particular we do not touch upon the subject of electrons interacting with a gauge field, such as is the case in two-dimensional electron systems exposed to external magnetic field (the gauge field here being the statistical Chern-Simons field), where in principle at fractional Landau-level filling factors $`\nu `$ with even denominators (specifically at and close to $`\nu =1/2`$) the states are metallic but may be non-Fermi liquids (Kalmeyer and Zhang 1992, Halperin, Lee and Read 1993). and refer the reader to the cited literature. <sup>\***</sup><sup>\***</sup>\*** For a concise review see (Tsvelik 1995, Ch. 12).
The organisation of this work is as follows. In § II we present the Lehmann spectral representation for the single-particle Green function and clarify certain elements that are of particular relevance to our analysis of the Luttinger (1961) theorem. In § III we rigorously define the notion of QP and derive equations from which QP energies and wavefunctions can be obtained. Here we indicate the necessary steps to be undertaken for obtaining complex-valued QP energies. We devote § IV to the main objective of our work, namely a detailed analysis of the Luttinger theorem. In § V we briefly review and discuss a theorem due to Migdal (1957) in light of our findings in § IV. In § VI we compare Fermi liquids with non-Fermi liquids in terms of the ‘smoothness’ properties of $`\mathrm{\Sigma }(k;\epsilon )`$, as a function of both $`k`$ and $`\epsilon `$. In § VII we conclude this work by a brief discussion and a review of our main results. In Appendix A we introduce and discuss a number of mathematical concepts that we repeatedly encounter in the present work. Here we further attempt to clarify the physical relevance of the introduced notions by means of some simple examples. In Appendix B we derive the leading asymptotic contribution to the self-energy operator $`\mathrm{\Sigma }(\epsilon )`$ at large $`|\epsilon |`$. For this, the next-to-the-leading-order asymptotic contribution to the single-particle Green function needs be calculated; in the same Appendix we present this contribution as well as some details underlying its derivation. In Appendices C and D we consider the next-to-the-leading-order asymptotic term to the self-energy pertaining to Fermi- and marginal Fermi-liquids (Appendix C) and the Luttinger liquid (Appendix D). Here we separately deal with the cases corresponding to $`k=k_F`$ and $`kk_F`$. In Appendix D we explicitly demonstrate that the self-energy of the one-dimensional Luttinger model is not continuously differentiable at the Fermi points, in full conformity with a general result established in Appendix C. In this work we identify electrons with spin-less fermions.
## II The single-particle Green function
Here we explicitly deal with the single-particle Green function. In addition to an exposition of the formal significance of the singularities of this function, we discuss three specific and distinct energies $`\mu `$, $`\mu _N`$ and $`\mu _{N+1}`$, that are invariably (but un-justifiably) identified in the literature concerning the free-electron system.
Consider the following Lehmann (1954) representation (Fetter and Walecka 1971) for the (‘physical’) single-particle Green function in the coordinate-free representation
$$G(\epsilon )=\mathrm{}\underset{s}{}\mathrm{\Lambda }_s\mathrm{\Lambda }_s^{}\left\{\frac{\mathrm{\Theta }(\mu \epsilon _s)}{\epsilon \epsilon _si\eta }+\frac{\mathrm{\Theta }(\epsilon _s\mu )}{\epsilon \epsilon _s+i\eta }\right\},\eta 0,$$
(1)
where
$`\mathrm{\Lambda }_s:=\{\begin{array}{cc}\mathrm{\Psi }_{N1,s}|\widehat{\psi }|\mathrm{\Psi }_{N,0},\hfill & \mathrm{when}\epsilon _s<\mu ,\hfill \\ \mathrm{\Psi }_{N,0}|\widehat{\psi }|\mathrm{\Psi }_{N+1,s},\hfill & \mathrm{when}\epsilon _s>\mu \hfill \end{array}`$ (4)
denotes a ‘Lehmann amplitude’, and
$`\epsilon _s:=\{\begin{array}{cc}E_{N,0}E_{N1,s},\hfill & \mathrm{when}\epsilon _s<\mu ,\hfill \\ E_{N+1,s}E_{N,0},\hfill & \mathrm{when}\epsilon _s>\mu .\hfill \end{array}`$ (7)
Above $`|\mathrm{\Psi }_{M,s}`$ stands for the exact $`M`$-particle eigenstate of the interacting system and $`E_{M,s}`$ for the corresponding eigenenergy; $`s`$ denotes the set of all quantum numbers that uniquely specify $`|\mathrm{\Psi }_{M,s}`$, (we choose $`s=0`$ to signify the ground state which we assume to be non-degenerate); $`\widehat{\psi }`$ is the annihilation field operator. <sup>\*††</sup><sup>\*††</sup>\*†† Throughout this work, $`N`$ denotes the actual number of the electrons in the system. Here the ‘chemical potential’ $`\mu `$ is any real value which satisfies $`\mu _N\mu \mu _{N+1}`$, where $`\mu _N:=E_{N,0}E_{N1,0}`$ and $`\mu _{N+1}:=E_{N+1,0}E_{N,0}`$. That such a $`\mu `$ should exist follows from the requirement of stability of the ground state, meaning that $`\epsilon _g:=(E_{N+1,0}E_{N,0})(E_{N,0}E_{N1,0})`$ $`\mu _{N+1}\mu _N`$ be non-negative. For a system with uniform electron distribution, which is thus a system in the thermodynamic limit, it holds $`\mu _{N+1}=\mu _N+𝒪(N^p)`$, with $`p>0`$. <sup>\*‡‡</sup><sup>\*‡‡</sup>\*‡‡ The value $`p=1`$ as given in, e.g., Fetter and Walecka (1971, p. 75), is incorrect. Therefore with regard to uniform systems, we are considering a case in which $`(\mu _N,\mu _{N+1})`$ is a finite infinitesimal open interval. In § IV.C we shall see that $`\mu _N\epsilon _F`$, the Fermi energy of the system of $`N`$ electrons.
It can be easily verified that (see Appendix A and in particular § A.3, for our notational conventions)
$$\stackrel{~}{G}(z):=\mathrm{}\underset{s}{}\frac{\mathrm{\Lambda }_s\mathrm{\Lambda }_s^{}}{z\epsilon _s}$$
(8)
is the analytic continuation (Whittaker and Watson 1927, pp. 96-98, Titchmarsh 1939, pp. 138-164) of the (‘physical’) single-particle Green function $`G(\epsilon )`$ into the complex $`z`$-plane, that it is $`\stackrel{~}{\stackrel{~}{G}}(z)`$ on the physical RS (see Appendix A); $`G(\epsilon )`$ is obtained from $`\stackrel{~}{G}(z)`$ according to the following prescription <sup>†\*</sup><sup>†\*</sup>†\* Below $`\stackrel{~}{f}(z)`$ ($`f(\epsilon )`$) is a representative for any function of $`z`$ ($`\epsilon `$) that we encounter in the present work. <sup>††</sup><sup>††</sup>†† It should be noted that the prescription given in Eq. (9) is connected to our convention with regard to the time-Fourier transform of the time-dependent functions, namely $`f(\epsilon ):=_{\mathrm{}}^{\mathrm{}}dtF(t)\mathrm{exp}(+i\epsilon t/\mathrm{})`$. Had we chosen $`\mathrm{exp}(i\epsilon t/\mathrm{})`$ rather than $`\mathrm{exp}(+i\epsilon t/\mathrm{})`$, ‘$`\pm `$’ would have to be ‘$``$’.
$$f(\epsilon ):=\underset{\eta 0}{lim}\stackrel{~}{f}(\epsilon \pm i\eta ),\text{for}\epsilon {}_{<}{}^{>}\mu .$$
(9)
For a finite system, the Lehmann representation indicates that $`\stackrel{~}{G}(\epsilon )G(\epsilon )`$ has poles <sup>†‡</sup><sup>†‡</sup>†‡ There is a subtlety involved here. For an open stable system, the spectrum consists of both a discrete and a continuous part. The former part corresponds to bound states, and the latter to scattering states; the continuous part of the spectrum gives rise to branch-cut discontinuities in the energy-dependent correlation functions of this system. Further, even though in such a system the number of particles is finite, the set of discrete one-particle excitation energies has an accumulation point which, as indicated in § A.1, is not an isolated singularity of $`\stackrel{~}{G}(z)`$. For closed and finite systems (that is those with impenetrable boundaries), the completeness of the one-particle eigenstates implies that the set of one-particle eigenenergies $`\{\epsilon _s\}`$ possesses at least one accumulation point (this follows from the Bolzano-Weierstrass theorem — see Whittaker and Watson (1927, pp. 12 and 13)). Therefore, the non-isolated nature of the singular points of, for example, $`\stackrel{~}{G}(z)`$ is not exclusively a peculiarity of systems in the thermodynamic limit. at $`\epsilon =\epsilon _s`$ for all $`s`$ (see Eq. (7)). In the thermodynamic limit, the singular points of $`\stackrel{~}{G}(z)`$ correspond to branch points (isolated singularities are not a priori excluded) and continua of branch cuts covering (parts of) the real energy axis (see § A.4). For finite systems as well as those in the thermodynamic limit, $`\stackrel{~}{G}(z)`$ is analytic for all $`z`$ with $`\mathrm{Im}(z)0`$.
In § III we shall encounter the inverse operator $`\stackrel{~}{G}^1(z)`$. In view of this, we point out that from the representation in Eq. (8) it can be easily demonstrated that $`\stackrel{~}{G}(z)`$ has no complex zeros (or, better, zero eigenvalues) on the $`z`$-plane (Luttinger 1961). Thus $`\stackrel{~}{G}^1(z)`$, similar to $`\stackrel{~}{G}(z)`$, is analytic for all $`z`$ with $`\mathrm{Im}(z)0`$.
## III Quasi-particles; their “energies” and wavefunctions
As we have mentioned in § I, to quasi-particles correspond one-particle wavefunctions which are eigenfunctions of some one-particle-like Schrödinger equation corresponding to an energy-dependent non-Hermitian one-body Hamiltonian, the potential-energy part of which consists of the self-energy operator <sup>†§</sup><sup>†§</sup>†§ In spite of the fact that according to our present convention $`\mathrm{\Sigma }(k;\epsilon )`$ has dimension $`\mathrm{s}^1`$, i.e. inverse second, we refer to it as self-energy. and of possibly a non-trivial contribution due to an external potential (such as an ionic potential in a solid). We devote this Section to derivation as well as to the discussion of the equations for QP energies and wavefunctions. Some specific properties of the equation for the QP energies are of considerable relevance to our discussion of the Luttinger (1961) theorem in § IV.
In the representation-free form and for the complex energy $`z`$, from the Dyson equation $`\stackrel{~}{G}(z)=\stackrel{~}{G}_0(z)+\stackrel{~}{G}_0(z)\stackrel{~}{\mathrm{\Sigma }}(z)\stackrel{~}{G}(z)`$, with $`\stackrel{~}{G}_0(z)`$ the single-particle Green function pertaining to the non-interacting system, one obtains
$$\stackrel{~}{G}^1(z)=\frac{1}{\mathrm{}}\{zI\stackrel{~}{}(z)\}$$
(10)
where $`I`$ denotes the unit operator in the single-particle Hilbert space and
$$\stackrel{~}{}(z):=H_0+\mathrm{}\stackrel{~}{\mathrm{\Sigma }}(z).$$
(11)
Any singular point of $`\stackrel{~}{G}(z)`$ is a solution of $`det\left(\stackrel{~}{G}^1(z)\right)=0`$, or, making use of Eq. (10), of
$$det\left(zI\stackrel{~}{}(z)\right)=0.$$
(12)
This is reminiscent of the equation for the energies of the ‘ideal’ (i.e. non-interacting) QPs, namely $`det\left(zIH_0\right)=0`$; the difference between the two arises from the difference between $`\stackrel{~}{}(z)`$ and $`H_0`$ which according to Eq. (11) is equal to $`\mathrm{}\stackrel{~}{\mathrm{\Sigma }}(z)`$. The latter, a function of the complex energy parameter $`z`$, is a non-Hermitian operator; even for $`z\epsilon \pm i\eta `$, $`\eta 0`$, for $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, leading to $`\stackrel{~}{\mathrm{\Sigma }}(z)\mathrm{\Sigma }(\epsilon )`$ (see Eq. (9) above), $`\stackrel{~}{}(z)(\epsilon )`$ is not in general Hermitian. This implies, among others, that, unless $`\stackrel{~}{}(z)`$ is Hermitian, the spectral decomposition of $`\stackrel{~}{}(z)`$ involves the distinct sets of left and right eigenvectors, $`\{\stackrel{~}{\varphi }_{\mathrm{}}(z)\}`$ and $`\{\stackrel{~}{\psi }_{\mathrm{}}(z)\}`$ respectively. Since the sets of eigenvalues pertaining to the latter two sets are up to ordering identical (Morse and Feshbach 1953, p. 885), by choosing the appropriate ordering, the spectral (or bi-orthonormal (Morse and Feshbach 1953, pp. 884-886, Layzer 1963) — see later) representation of $`\stackrel{~}{}(z)`$ is as follows
$$\stackrel{~}{}(z)=\underset{\mathrm{}}{}\stackrel{~}{E}_{\mathrm{}}(z)\stackrel{~}{\psi }_{\mathrm{}}(z)\stackrel{~}{\varphi }_{\mathrm{}}^{}(z),$$
(13)
where $`\stackrel{~}{E}_{\mathrm{}}(z)`$ denotes the common eigenvalue corresponding to $`\stackrel{~}{\varphi }_{\mathrm{}}(z)`$ and $`\stackrel{~}{\psi }_{\mathrm{}}(z)`$:
$$\stackrel{~}{\varphi }_{\mathrm{}}^{}(z)\stackrel{~}{}(z)=\stackrel{~}{E}_{\mathrm{}}(z)\stackrel{~}{\varphi }_{\mathrm{}}^{}(z);\stackrel{~}{}(z)\stackrel{~}{\psi }_{\mathrm{}}(z)=\stackrel{~}{E}_{\mathrm{}}(z)\stackrel{~}{\psi }_{\mathrm{}}(z).$$
(14)
For $`\stackrel{~}{E}_{\mathrm{}}(z)\stackrel{~}{E}_{\mathrm{}^{}}(z)`$ it holds $`\stackrel{~}{\varphi }_{\mathrm{}}(z),\stackrel{~}{\psi }_{\mathrm{}^{}}(z)=\delta _\mathrm{},\mathrm{}^{}`$ (Morse and Feshbach 1953, p. 885). In the case of degeneracy, that is $`\stackrel{~}{E}_{\mathrm{}}(z)=\stackrel{~}{E}_{\mathrm{}^{}}(z)`$ for $`\mathrm{}\mathrm{}^{}`$, the degenerate left and right eigenvectors can be made orthogonal through a Gram-Schmidt orthogonalisation procedure. When the left and right eigenfunctions are simultaneously bases of the unitary irreducible representations of the symmetry group of the QP Schrödinger equation, Eq. (14), the degenerate eigenvectors are automatically orthogonal as long as they belong to different unitary irreducible representations (Cornwell 1984, pp. 81-83). <sup>†¶</sup><sup>†¶</sup>†¶ Consider, for instance, the uniform-electron system with which we deal in detail in the present work. The symmetry group of this system is the continuous translation group (which is Abelian) and therefore the corresponding bases for the (one-dimensional) unitary irreducible representations are specified by wave-vectors $`𝐤`$. The time-reversal symmetry of the problem implies degeneracy of eigenfunctions at $`𝐤`$ and $`𝐤`$ (“Kramers’ degeneracy” (Landau and Lifshitz 1977, pp. 223-226)). Since for $`𝐤\mathrm{𝟎}`$, $`𝐤𝐤`$, the eigenfunctions corresponding to $`𝐤`$ and $`𝐤`$ are therefore automatically orthogonal. Note that, since our ‘electrons’ are spin-less, for the application of the Kramers theorem we cannot rely on the condition that the total spin of the system is half-integer. However, we can rely on the fact that all irreducible representations of the translation group corresponding to $`𝐤\mathrm{𝟎}`$ are essentially complex. Further, $`\{\stackrel{~}{\varphi }_{\mathrm{}}(z)\}`$ and $`\{\stackrel{~}{\psi }_{\mathrm{}}(z)\}`$ satisfy the completeness relation (Morse and Feshbach 1953, p. 886) <sup>†∥</sup><sup>†∥</sup>†∥ The left-hand side of this relation is the so-called idemfactor of the one-particle Hilbert space.
$$\underset{\mathrm{}}{}\stackrel{~}{\psi }_{\mathrm{}}(z)\stackrel{~}{\varphi }_{\mathrm{}}^{}(z)=I,$$
(15)
or $`_{\mathrm{}}\stackrel{~}{\psi }_{\mathrm{}}(𝐫;z)\stackrel{~}{\varphi }_{\mathrm{}}^{}(𝐫^{};z)=\delta (𝐫𝐫^{})`$.
Because $`\stackrel{~}{}^{}(z)=\stackrel{~}{}(z^{})`$, for $`\mathrm{Im}(z)0`$, a property that follows from $`\stackrel{~}{\mathrm{\Sigma }}^{}(z)=\stackrel{~}{\mathrm{\Sigma }}(z^{})`$ (Appendix B in DuBois (1959a), Luttinger 1961), it can be readily shown that, for $`\mathrm{Im}(z)0`$,
$$\stackrel{~}{\varphi }_{\mathrm{}}(z)\stackrel{~}{\psi }_{\mathrm{}}(z^{}),\stackrel{~}{\psi }_{\mathrm{}}(z)\stackrel{~}{\varphi }_{\mathrm{}}(z^{}),\stackrel{~}{E}_{\mathrm{}}(z^{})=\stackrel{~}{E}_{\mathrm{}}^{}(z).$$
(16)
As we shall see later, the last result <sup>†\**</sup><sup>†\**</sup>†\** This result is the expression of the Riemann-Schwarz reflection principle (Titchmarsh 1939, p. 155) which follows from the analyticity of $`\stackrel{~}{\mathrm{\Sigma }}(z)`$ over the entire complex $`z`$-plane, with the exception of the real $`\epsilon `$-axis, and the fact that over the finite (even though infinitesimal) open interval $`(\mu _N,\mu _{N+1})`$, $`\mathrm{\Sigma }(\epsilon )`$, or $`G(\epsilon )`$, is real valued. is of particular significance to our considerations. By defining $`E_{\mathrm{}}(\epsilon )`$ in accordance with Eq. (9), for $`\mathrm{Im}E_{\mathrm{}}(\epsilon )`$ the following must hold
$$\mathrm{Im}E_{\mathrm{}}(\epsilon ){}_{}{}^{}\mathrm{\hspace{0.33em}0},\mathrm{for}\epsilon {}_{<}{}^{>}\mu .$$
(17)
Violation of these inequalities signifies breakdown of causality, or instability of the ground state due to its collapse into a lower-energy state; to appreciate this, note that (combine Eqs. (10) and (15)) $`\stackrel{~}{G}(z)=\mathrm{}_{\mathrm{}}\stackrel{~}{\psi }_{\mathrm{}}(z)\stackrel{~}{\varphi }_{\mathrm{}}^{}(z)/\left(z\stackrel{~}{E}_{\mathrm{}}(z)\right)`$ and compare this with the Lehmann representation in Eq. (1) (see also Eq. (7)).
From the spectral representation for $`\stackrel{~}{}(z)`$ in Eq. (13) above, together with Eq. (15), the equivalence of Eq. (12) with the set of equations
$$\stackrel{~}{E}_{\mathrm{}}(z)=z,\mathrm{},$$
(18)
is readily established. As we have mentioned in § II, singular points of $`G(\epsilon )`$, according to the Lehmann representation, coincide with the real-valued quantities $`\epsilon _s`$ (see Eq. (7)). Since Eq. (18) is the equation for these singular points, it may be expected that by taking the appropriate limits $`z\epsilon \pm i\eta `$, with $`\eta 0`$ (see Eq. (9)), Eq. (18) should transform into an equation with real-valued solutions $`\epsilon _s`$; owing to the analyticity of $`\stackrel{~}{G}^1(z)`$ for $`\mathrm{Im}(z)0`$ (see § II), it follows that Eq. (18) cannot be satisfied for any $`z`$ with $`\mathrm{Im}(z)0`$ (for a physical interpretation of this statement see the following paragraph). As we shall demonstrate below, in the thermodynamic limit the number of real-valued solutions of Eq. (18), with $`z\epsilon \pm i\eta `$ (see Eq. (9)), is limited. Further, in the same limit, these real-valued solutions are in general not poles, in contrast with what the Lehmann representation in Eq. (1) (or Eq. (8)) may suggest. This is not a contradiction, since as we indicate in § A.4, analytic properties of correlation functions in the thermodynamic limit are fundamentally different from those of finite systems.
It is often mentioned that finite lifetimes of QPs is reflected in the energies of these being complex valued. This statement, without further qualification, is misleading. To clarify this, consider Eq. (18) and suppose that it were satisfied by $`z=z_0`$ with $`\mathrm{Im}(z_0)0`$. Owing to the reflection property of $`\stackrel{~}{E}_{\mathrm{}}(z)`$ presented in Eq. (16), it follows that $`z=z_0^{}`$ must also satisfy Eq. (18). This cannot be the case, since depending on whether $`\mathrm{Re}(z_0)`$ is less than or larger than $`\mu `$, one of the two solutions $`z_0`$ and $`z_0^{}`$ signals violation of the principle of causality. In this connection, recall that the single-particle Green function has been defined in terms of the time-ordered product of two field operators in the Heisenberg representation (Fetter and Walecka 1971). Consequently, no $`z_0`$ can satisfy Eq. (18) unless $`\mathrm{Im}(z_0)=0`$. Thus for a given $`\mathrm{}`$, Eq. (18) has either a real-valued solution or no solution at all; the possible real-valued solutions (see § A.1) are, in general, not isolated. To obtain complex-valued quasi-particle energies, the appropriate equation to be solved is the following
$$\stackrel{~}{\stackrel{~}{E}}_{\mathrm{}}(z)=z,\mathrm{},$$
(19)
where $`\stackrel{~}{\stackrel{~}{E}}_{\mathrm{}}(z)`$ represent the unique analytic function of which $`\stackrel{~}{E}_{\mathrm{}}(z)`$ is one specific branch (see § A.3). The analytic function $`\stackrel{~}{\stackrel{~}{E}}_{\mathrm{}}(z)`$ (i.e. any of its branches; there may be infinity of these, depending on the nature of the branch points of $`\stackrel{~}{E}_{\mathrm{}}(z)`$) can in principle be calculated from the knowledge of $`\stackrel{~}{E}_{\mathrm{}}(z)`$. To this end, one starts from $`\stackrel{~}{E}_{\mathrm{}}(z)`$ and analytically continues it across its branch cuts into a RS ‘above’ or ‘below’ the physical RS (see examples in §§ A.3 and A.4). This process can be repeated on any of the RSs and in this way one recovers all the possible branches of $`\stackrel{~}{\stackrel{~}{E}}(z)`$. From the specification presented in Eq. (9), it is readily observed that, to obtain the physically most relevant complex-valued solutions of Eq. (19), one has to analytically continue $`\stackrel{~}{E}_{\mathrm{}}(z)`$ from the first (third) quadrant of the physical RS into the fourth (second) quadrant of a non-physical RS when $`\mathrm{Re}(z)>\mu `$ ($`\mathrm{Re}(z)<\mu `$). We note that solution $`z_0`$ of $`\stackrel{~}{\stackrel{~}{E}}_{\mathrm{}}(z)=z`$ with $`\mathrm{Re}(z_0)>\mu `$ corresponds to a QP that propagates and attenuates forwards in time (i.e. it concerns a ‘particle-like’ QP) and that with $`\mathrm{Re}(z_0)<\mu `$ to one that propagates and attenuates backwards in time (a ‘hole-like’ QP).
## IV Discussion of the Luttinger theorem
We are now in a position to deal with the central objective of the present work. In doing so we distinguish three major steps in the proof by Luttinger (1961). We shall separately describe each of these steps and, when necessary, comment on the details. We conclude that Luttinger’s theorem does not exclude the existence of non-Fermi-liquid metallic states for systems in spatial dimensions larger than unity. Rather, it merely establishes that Landau’s Fermi-liquid theory for metals is consistent with the principles of the many-body theory of interacting systems.
### A Preliminaries
Before proceeding with the proof of the Luttinger (1961) theorem, several remarks are in order.
First, our considerations in § II have made evident that $`\mu _N`$ and $`\mu _{N+1}`$ are two branch points of $`\stackrel{~}{G}(z)`$ and $`\stackrel{~}{\mathrm{\Sigma }}(z)`$ and consequently of $`\stackrel{~}{E}_{\mathrm{}}(z)`$, for all $`\mathrm{}`$. For simplicity of notation, let $`\stackrel{~}{f}(z)`$ denote any of these functions of $`z`$. The fact that $`\mu _N`$ and $`\mu _{N+1}`$ are branch points of $`\stackrel{~}{f}(z)`$ implies that such expansion as $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, cannot correspond to the leading-order term in a Taylor expansion of $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ when $`z=\epsilon _F`$ is identified with either $`\mu _N`$ or $`\mu _{N+1}`$. Rather, it corresponds to the leading-order term in an asymptotic expansion (Whittaker and Watson 1927, Ch. VIII, Copson 1965, Dingle 1973, Lauwerier 1977) of $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ for $`\epsilon \epsilon _F`$ (see § A.2).
In view of the fact that for the uniform-electron system under consideration, the interval $`(\mu _N,\mu _{N+1})`$ is infinitesimally small, it is tempting to identify $`\epsilon _F`$, $`\mu _N`$ and $`\mu _{N+1}`$, as is commonly done in the literature. Such an identification is not justified. Our discussion of a Migdal’s (1957) theorem in § V further clarifies this statement. We note that, strictly speaking, $`\epsilon _F`$ is identical with $`\mu _N`$. Thus the precise statement of the Luttinger theorem under consideration is as follows. <sup>†††</sup><sup>†††</sup>††† In the work by Luttinger (1961), there is no explicit mention of the value or range of values of $`k`$ for which $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, is valid. This is, however, implicit in Eqs. (59)-(61) of Luttinger’s work: Eq. (61) implies $`u`$ to be small in magnitude ($`u:=(\mu x)`$; in our notation, $`u:=(\epsilon _F\epsilon )`$), so that by Eq. (60) one deduces $`t_1`$, $`t_2`$ and $`t_3`$ to be small, requiring by Eq. (59) that $`k_1,k_2,k_3k_F`$. From $`𝐤_3=𝐤_1+𝐤_2𝐤`$, using a geometrical construction, it is easily deduced that $`0k{}_{}{}^{<}3k_F`$. $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )0`$ for $`\epsilon (\mu _N,\mu _{N+1})`$; $`\mathrm{Im}\stackrel{~}{\mathrm{\Sigma }}(k;\epsilon +i\eta )`$ $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ $`\alpha _k(\epsilon \mu _{N+1})^2`$ for $`\epsilon \mu _{N+1}`$ and $`\mathrm{Im}\stackrel{~}{\mathrm{\Sigma }}(k;\epsilon i\eta )`$ $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ $`\alpha _k(\epsilon \mu _N)^2`$ for $`\epsilon \mu _N`$.
Second, the above asymptotic expansions are not uniform (see § A.2), as is evident from the sign of $`\mathrm{Im}\mathrm{\Sigma }(k;z)`$ which depends on the sign of $`\mathrm{Im}(z)`$ (or on $`\mathrm{arg}(z)`$). Thus $`z=\mu _N`$ and $`z=\mu _{N+1}`$ are indeed branch points.
Third, analysis of the next-to-leading order terms in the asymptotic expansions of $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ around $`z=\mu _N`$ and $`z=\mu _{N+1}`$ reveal that these terms involve the logarithm function, implying that both $`\mu _N`$ and $`\mu _{N+1}`$ are branch points of infinite order, that is to say $`\stackrel{~}{\stackrel{~}{f}}(z)`$ (see § A.3) consists of infinitely many branches. Already the fact that $`\mu _N`$ and $`\mu _{N+1}`$ are not regular singularities of $`\stackrel{~}{f}(z)`$ suggests that care should be exercised in defining the QP weights, in particular $`Z_{k_F}`$ which features in the Migdal (1957) theorem and which determines the value of the discontinuity in the momentum distribution functions at the Fermi momentum $`\mathrm{}k_F`$. We shall return to this point in § V.
In spite of our first remark above, for conformity with the existing texts in what follows we shall use the conventional notation (such as that in our Abstract) and proceed with the demonstration of the implicit assumption involved in the proof of the Luttinger (1961) theorem. To this end, we first summarise the main aspects of Luttinger’s (1961) proof. This proof is based on the assumption of validity of the many-body perturbation theory and further rests on the analyses of the terms in an implicit perturbation expansion of the self-energy operator; the implicit nature of this expansion is associated with the fact that contributions in this are in terms of the single-particle Green function pertaining to the interacting system which is, in turn, a functional of the self-energy operator itself. Diagrammatically, the terms in this perturbation series are represented by the proper skeleton self-energy diagrams (Luttinger and Ward 1960), that is those connected self-energy diagrams that do not accommodate any simpler self-energy parts in them.
In the proof by Luttinger (1961), the skeleton diagrams are in terms of the isotropic bare electron-electron interaction function $`v(𝐫𝐫^{})`$ and this, as we shall discuss in § IV.D, implies that Luttinger’s approach is not directly applicable to cases where $`v(𝐫𝐫^{})`$ is the long-range Coulomb interaction function.
Since the many-body perturbation theory underlies the proof by Luttinger (1961), the Luttinger theorem is valid insofar as this perturbation theory is valid. This aspect leads to the natural conclusion that metallic non-Fermi-liquid systems in spatial dimensions $`d`$ larger than unity (see footnote ‡‡IV B further on) fall outside the domain of applicability of the many-body perturbation theory, or that this theory should break down when applied to such systems. As has been emphasised earlier (Anderson 1993), validity or otherwise of a many-body perturbation series cannot be decided solely on the basis of its convergence or divergence, respectively. Despite convergence, the calculated limit may be unphysical. <sup>†‡‡</sup><sup>†‡‡</sup>†‡‡ The following example due to Simon (1970) should be clarifying: Consider the hydrogen-like Hamiltonian in three dimensions (here we use the Hartree atomic units): $`:=^2/2\lambda /r`$, $`r0`$. The discrete energies of this Hamiltonian for $`\lambda >0`$ are $`_n(\lambda )=\lambda ^2/[2n^2]`$, $`n=1,2,\mathrm{}`$. For $`\lambda 0`$ there are no discrete levels (no bound states). However, $`_n(\lambda )`$ being a finite-order polynomial (technically, an entire function) of $`\lambda `$, the Rayleigh-Schrödinger perturbation expansion for, say, the ground-state energy (i.e. $`_1(\lambda )`$), with $`\lambda /r`$ playing the role of the ‘perturbation’, yields exactly $`\lambda ^2/2`$, irrespective of whether $`\lambda >0`$ or $`\lambda 0`$ and independent of the magnitude of $`|\lambda |`$; the coefficients of $`\lambda ^m`$, for $`m=0,1`$ and $`m>2`$, in this expansion are identically vanishing. Evidently, the convergence of this perturbation series for $`\lambda 0`$ is a ‘bogus’ convergence.
### B Explicit analysis of the Luttinger theorem
After indicating that the perturbation contributions in the above-indicated expansion for $`\mathrm{\Sigma }(k;\epsilon )`$ in terms of the skeleton self-energy diagrams below second order in the electron-electron interaction $`v`$ are $`\epsilon `$-independent, Luttinger (1961) considered in detail the second-order self-energy diagram that we have depicted in Fig. 1 (for facilitating direct comparison with the results in the work by Luttinger, below we adopt the notation that he employed as closely as possible). In doing so, Luttinger first evaluated the imaginary part of the self-energy operator corresponding to this diagram, with the single-particle Green function pertaining to the interacting system, that is $`G(\epsilon )`$ (represented by full line in Fig. 1), replaced by that pertaining to the non-interacting system, i.e. $`G_0(\epsilon )`$. <sup>‡\*</sup><sup>‡\*</sup>‡\* Explicitly, we have $`G_0(k;\epsilon )\mathrm{}\left\{\mathrm{\Theta }(k_Fk)/(\epsilon \epsilon _k^0i\eta )+\mathrm{\Theta }(kk_F)/(\epsilon \epsilon _k^0+i\eta )\right\}`$, $`\eta 0`$; $`\epsilon _k^0`$ is defined in Eq. (20). From the well-known rules for transcribing Feynman diagrams into mathematical expressions in the momentum representation (Fetter and Walecka 1971, pp. 100-105), it is readily seen that the expression corresponding to the diagram in Fig. 1 involves three wave-vector integrations, over $`𝐤_1`$, $`𝐤_2`$ and $`𝐤_3`$. Following the fact that for the matrix elements of the electron-electron interaction, $`(𝐤,𝐤_3|v|𝐤_1,𝐤_2)`$ $`=\overline{v}(𝐤𝐤_1)\delta _{𝐤𝐤_1,𝐤_2𝐤_3}`$ holds, where $`𝐤`$ stands for the external wave-vector (i.e. the wave-vector associated with $`k:=𝐤`$ in the argument of $`\mathrm{\Sigma }(k;\epsilon )`$), it would be tempting to eliminate one of the wave-vector integrals. This elimination would make the subsequent algebra extremely cumbersome. Rather, Luttinger (1961) transforms the three wave-vector integrals into three one-dimensional integrals over energies $`\epsilon _{k_1}^0`$, $`\epsilon _{k_2}^0`$ and $`\epsilon _{k_3}^0`$ of the non-interacting electrons, where
$$\epsilon _k^0:=\frac{\mathrm{}^2}{2m_e}k^2,$$
(20)
with $`m_e`$ the bare-electron mass. <sup>‡†</sup><sup>‡†</sup>‡† The approach by Luttinger (1961) aims at obtaining the leading-order term in $`(\epsilon \epsilon _F)`$ (in Luttinger’s notation, $`u`$), so that other integrals which merely determine numerical pre-factors are not explicitly dealt with by him. For clarity, consider a three-dimensional non-interacting system with $`\epsilon _k^0`$ the one-particle energy dispersion presented in Eq. (20). In the spherical-polar coordinate system $`(k,\varphi ,\theta )`$ we have $`\mathrm{d}^3k=k^2\mathrm{sin}(\theta )\mathrm{d}k\mathrm{d}\varphi \mathrm{d}\theta `$ $`\frac{1}{2}(2m_e/\mathrm{}^2)^{3/2}\sqrt{\epsilon _k^0}\mathrm{sin}(\theta )\mathrm{d}\epsilon _k^0\mathrm{d}\varphi \mathrm{d}\theta `$. In the treatment by Luttinger (1961), $`\mathrm{d}^3k\mathrm{d}\epsilon _k^0`$ plays an all-important role.
A most crucial ingredient in Luttinger’s proof consists of the recognition that, although the momentum conservation, enforced by the above $`\delta _{𝐤𝐤_1,𝐤_2𝐤_3}`$, uniquely determines one of the three wave-vectors in terms of the other two, say $`𝐤_3`$ in terms of $`𝐤_1`$ and $`𝐤_2`$, nonetheless, barring one-dimensional systems (Luttinger 1961, footnote 5), $`\epsilon _{k_1}^0`$, $`\epsilon _{k_2}^0`$ and $`\epsilon _{k_3}^0`$ can be independently and continuously varied over some finite range of values. <sup>‡‡</sup><sup>‡‡</sup>‡‡ This is the instance where one-dimensional interacting systems are singled out as being in some fundamental way different from higher-dimensional interacting systems. The Luttinger model, which turns out to be akin to a large class of one-dimensional models (Haldane 1980, 1981), is however introduced in 1963 by Luttinger (1963) and was first correctly treated by Mattis and Lieb (1965). An earlier model to which the Luttinger model is similar, is due to Tomonaga (1950). For three different treatments of the Tomonaga-Luttinger models see Dzyaloshinskiǐ and Larkin (1973), Luther and Peschel (1974), and Everts and Schulz (1974). For reviews concerning one-dimensional systems of interacting fermions see Sólyom (1979), Mahan (1981, § 4.4), Voit (1994), Schönhammer (1997) and Schulz, Cuniberti and Pieri (1998). This can be seen as follows: using the above dispersion for $`\epsilon _k^0`$ as well as $`𝐤_3=𝐤_1+𝐤_2𝐤`$, it is readily seen that
$$\epsilon _{k_3}^0=\epsilon _k^0+\epsilon _{k_1}^0+\epsilon _{k_2}^0+\frac{\mathrm{}^2}{m_e}\left\{𝐤_1𝐤_2𝐤𝐤_1𝐤𝐤_2\right\}.$$
(21)
For fixed values of $`\epsilon _{k_1}^0`$ and $`\epsilon _{k_2}^0`$, and in spite of $`𝐤_3=𝐤_1+𝐤_2𝐤`$, for systems in spatial dimensions $`d>1`$, the orientational freedom of vectors $`𝐤_1`$ and $`𝐤_2`$ with respect to each other as well as $`𝐤`$ allows for continuous finite variations in $`\epsilon _{k_3}^0`$. A crucial element in our arguments that follow, is derived from the fact that freedom in the independent variations of $`\epsilon _{k_3}^0`$ is a consequence of the expression in Eq. (21), which, in turn, follows from the specific form of the energy dispersion pertaining to non-interacting electrons in Eq. (20). As a matter of course the latter form is not unique in that one can construct energy dispersions $`\epsilon _k^0`$ different from that in Eq. (20) that equally give rise to the possibility of independent and continuous variations of $`\epsilon _{k_j}^0`$, $`j=1,2,3`$, over some non-vanishing region. At the same time, it is not difficult to put forward energy dispersions that do not allow for such an independent continuous variation over a finite domain. We shall return to this point in § IV.C.
It is from the observation with regard to the possibility of independent variations of the three energies $`\epsilon _{k_j}^0`$, $`j=1,2,3`$, that Luttinger (1961) has established the leading-order contribution to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ of the second-order diagram in Fig. 1 to be proportional to $`(\epsilon \epsilon _F^0)^2`$ as $`\epsilon \epsilon _F^0`$. <sup>‡§</sup><sup>‡§</sup>‡§ We note in passing that ‘$`u^2`$’ in Eq. (60) of Luttinger (1960) is in error; the correct pre-factor is unity.
The second step and, from the viewpoint of our present work, the most crucial step in Luttinger’s (1961) proof consists in establishing that, upon employing $`G(\epsilon )`$ rather than the $`G_0(\epsilon )`$ of the first step, the leading-order contribution to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ pertaining to the second-order self-energy diagram in Fig. 1 is also proportional to $`(\epsilon \epsilon _F)^2`$, for $`\epsilon \epsilon _F`$. This conclusion is of fundamental importance to the third part in Luttinger’s proof, where, after having established the latter property, namely that the leading-order contribution to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$, as $`\epsilon \epsilon _F`$, is unaffected by evaluating the contributions of the skeleton self-energy diagrams in terms of $`G_0(\epsilon )`$ rather than $`G(\epsilon )`$, Luttinger demonstrated that to all orders in the perturbation theory $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )^{}(\epsilon \epsilon _F)^2`$ or explicitly, that any skeleton self-energy diagram (in terms of $`G_0(\epsilon )`$) of second and higher order in the bare electron-electron interaction $`v`$ has a contribution to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ proportional to $`(\epsilon \epsilon _F^0)^{2m}`$, where $`m1`$, for $`\epsilon \epsilon _F^0`$. <sup>‡¶</sup><sup>‡¶</sup>‡¶ It is known (Hodges, Smith and Wilkins 1971, Bloom 1975, Fujimoto 1990, Fukuyama, Narikiyo and Hasegawa 1991, Fukuyama, Hasegawa and Narikiyo 1991) that, in isotropic two-dimensional interacting systems, $`\mathrm{Im}\mathrm{\Sigma }(k=k_F;\epsilon )^{}(\epsilon \epsilon _F)^2\mathrm{ln}|\epsilon \epsilon _F|`$ as $`\epsilon \epsilon _F`$. Since the leading $`\epsilon `$-dependent contribution to $`\mathrm{Re}\mathrm{\Sigma }(k=k_F;\epsilon )`$ associated with this $`\mathrm{Im}\mathrm{\Sigma }(k=k_F;\epsilon )`$ is proportional to $`(\epsilon \epsilon _F)`$ (see Appendix C), we have that $`\mathrm{\Sigma }(k_F;\epsilon )`$ is a continuously differentiable function of $`\epsilon `$ in a neighbourhood of $`\epsilon =\epsilon _F`$. In view of our statements in § I (see also § IV.C), the above logarithmic contribution does not turn the system into a non-Fermi liquid.
In this second part of his proof, Luttinger (1961) proceeded as follows. On evaluating the contribution of the diagram in Fig. 1 to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ in terms of $`G(\epsilon )`$, he made a simplifying approximation, which, roughly speaking, corresponds to replacing a Lorentzian by a Dirac $`\delta `$-function. <sup>‡∥</sup><sup>‡∥</sup>‡∥ This approximation neglects the finite lifetimes of the QPs. From $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon _F)0`$ (note in passing that this together with $`\epsilon _k^0+\mathrm{}\mathrm{\Sigma }(k;\epsilon _F)=\epsilon _F`$ define the Fermi surface — see footnote §§IV C below) and the requirement of continuity of $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ for $`\epsilon \epsilon _F`$, we have $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )0`$ for $`\epsilon \epsilon _F`$. Consequently, for $`\epsilon \epsilon _F`$ life-time effects are not of relevance to the leading-order contribution to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ as $`\epsilon \epsilon _F`$. The mentioned ‘Lorentzian’ has its origin in the finite lifetimes of the QPs. For a summary of various ways in which $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ can vanish as $`\epsilon \epsilon _F`$, see the text following Eq. (35) in § V. Upon this, Luttinger arrived at an expression which has the same formal structure as the one discussed above. Similar to the second-order expression for $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ in terms of $`G_0`$, the simplified expression in terms of $`G`$ involves three wave-vector integrals over $`𝐤_1`$, $`𝐤_2`$ and $`𝐤_3`$. Now, at this place, Luttinger took a step which, as will become evident below, implicitly contains the property to be proven. Here Luttinger transforms the three wave-vector integrals into three energy integrals, the energy dispersions being $`\epsilon _{k_1}`$, $`\epsilon _{k_2}`$ and $`\epsilon _{k_3}`$, to be compared with $`\epsilon _{k_1}^0`$, $`\epsilon _{k_2}^0`$ and $`\epsilon _{k_3}^0`$ considered above (see § IV.C). Note that the dispersions $`\epsilon _{k_j}`$, $`j=1,2,3`$, correspond to the interacting (i.e. true) QPs and satisfy Eq. (18) or Eq. (22) below. <sup>‡\**</sup><sup>‡\**</sup>‡\** In Luttinger’s (1961) work there is no symbolic distinction between, for example, what we have denoted as $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ and $`\mathrm{\Sigma }(k;z)`$. Further, Luttinger does not discuss whether Eq. (22) can be exactly satisfied or not (for details see § IV.C). According to Luttinger, all the necessary mathematical steps that are to be taken in order to obtain, from the last-indicated integral, the leading-order contribution to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$, as $`\epsilon \epsilon _F`$, are identical with those taken in dealing with $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ evaluated in terms of $`G_0`$. This is an unjustified statement. The dependence on $`k`$ of $`\epsilon _k`$ being unknown, there is no a priori reason why $`\epsilon _{k_3}`$ can be varied continuously and independently of $`\epsilon _{k_1}`$ and $`\epsilon _{k_2}`$ over a finite range. <sup>‡††</sup><sup>‡††</sup>‡†† The contribution of the diagram in Fig. 1 to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ for one-dimensional metals, according to Luttinger (1961, p. 946, footnote 5) amounts to $`^{}(\epsilon \epsilon _F)`$, which we know to be characteristic of marginal-Fermi liquids (see Appendix C) rather than one-dimensional Luttinger liquids (see Appendix D) over a physically-relevant range of values for the anomalous dimension $`\alpha 2\gamma _0`$, that is $`0<\alpha <1`$. At the root of this incorrect inference by Luttinger (1961) lies the fact that in one-dimensional metals, $`\epsilon _k`$, contrary to Luttinger’s implicit assumption, is not a continuously-differentiable function of $`k`$ and therefore $`\mathrm{d}^1k\mathrm{d}\epsilon _k`$ does not apply, while $`\mathrm{d}^1k\mathrm{d}\epsilon _k^0`$ is valid. See § IV.C. If the effect of electron-electron interaction on the dispersion of the QP energies close to the Fermi surface were only to re-normalise the QP mass, then of course for $`k`$ in a neighbourhood of $`k=k_F`$ one invariably had $`\epsilon _k\mathrm{}^2k^2/(2m_e^{})`$, with $`m_e^{}`$ an electron’s renormalised mass (Pines and Nozières 1966, Abrikosov, Gorkov and Dzyaloshinski 1975). In such an event, for $`k`$ in a neighbourhood of $`k=k_F`$, Eq. (20) would hold for $`\epsilon _k`$ with $`m_e`$ replaced by $`m_e^{}`$, and thus the just-mentioned assertion by Luttinger (1961) were correct. However, as we shall demonstrate in § IV.C, $`\epsilon _k\mathrm{}^2k^2/(2m_e^{})`$ can only apply if $`\mathrm{\Sigma }(k;\epsilon )`$ is of the Fermi-liquid type, thus establishing that Luttinger’s proof amounts to a demonstration of consistency of the Fermi-liquid state with the many-body theory of interacting metals. This demonstration does not rule out existence of states different from the Fermi-liquid state.
### C On the dispersion of the QP energies
Here we analyse some general features of the QP energies pertinent to a uniform-electron system.
Consider Eq. (18) specialised to a uniform system. Any possible real-valued QP energy $`\epsilon _k`$ is a solution (which for an arbitrary $`k`$ may not exist; see § V) of the following equation (see also Eq. (30) below)
$$\epsilon _k=\epsilon _k^0+\mathrm{}\mathrm{\Sigma }(k;\epsilon _k).$$
(22)
The fact that $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ has a branch point at $`z=\mu _N\epsilon _F:=\epsilon _{k_F}`$ (§ IV.A), implies that the following equation is well satisfied by the real-valued energy $`\epsilon _F`$ (below $`\epsilon _F^0:=\epsilon _{k_F}^0`$): <sup>‡‡‡</sup><sup>‡‡‡</sup>‡‡‡ Isotropy of the system under consideration implies $`k_F`$ to be independent of the electron-electron interaction. We note in passing that this is a corollary to another celebrated theorem due to Luttinger (1960); see also Luttinger and Ward (1960).
$$\epsilon _F=\epsilon _F^0+\mathrm{}\mathrm{\Sigma }(k_F;\epsilon _F).$$
(23)
Assuming $`\mathrm{\Sigma }(k;\epsilon )`$ to be a continuously differentiable function of $`k`$ and $`\epsilon `$ in some neighbourhoods of $`k=k_F`$ and $`\epsilon =\epsilon _F`$, respectively, from Eq. (22) together with the manifest differentiability of $`\epsilon _k^0`$ in Eq. (20) with respect to $`k`$ it follows that $`\epsilon _k`$ is similarly a continuously differentiable function of $`k`$ in a neighbourhood of $`k=k_F`$. Consequently, <sup>§\*</sup><sup>§\*</sup>§\* As indicated earlier in this work, $`o`$ in $`f(x)=o\left(g(x)\right)`$ signifies $`f(x)/g(x)0`$ for $`x0`$.
$$\epsilon _k=\epsilon _F+\mathrm{}v_F(kk_F)+o\left(kk_F\right),$$
(24)
where
$$v_F:=\frac{1}{\mathrm{}}\frac{\epsilon _k}{k}|_{k=k_F}\frac{\mathrm{}k_F}{m_e^{}}$$
(25)
stands for the Fermi velocity. From Eq. (22) it can be shown that
$$\frac{\epsilon _k}{k}|_{k=k_F}=Z_{k_F}\frac{\{\epsilon _k^0+\mathrm{}\mathrm{\Sigma }(k;\epsilon _F)\}}{k}|_{k=k_F},$$
(26)
where
$$Z_{k_F}:=\left(1\mathrm{}\frac{\mathrm{\Sigma }(k_F;\epsilon )}{\epsilon }|_{\epsilon =\epsilon _F}\right)^1$$
(27)
is, according to a Migdal’s (1957) theorem (see § V), exactly the amount of discontinuity in the momentum distribution function $`𝗇(p)`$ at the Fermi momentum $`p_F\mathrm{}k_F`$. Eqs. (24) and (25) are readily seen to imply that for $`kk_F`$, $`\epsilon _k\mathrm{}^2k^2/(2m_e^{})`$. The above steps make explicit how the latter dispersion relation is rooted in the assumption of continuous differentiability of $`\mathrm{\Sigma }(k;\epsilon _k)`$ in a neighbourhood of $`k=k_F`$. For Fermi liquids $`0<Z_{k_F}1`$ holds, <sup>§†</sup><sup>§†</sup>§† Recall that $`0<Z_{k_F}1`$ although necessary, is not sufficient for rendering a metallic system a Fermi liquid (see § I, conditions (A) and (B)). where $`Z_{k_F}=1`$ corresponds to the case of non-interacting QPs. The condition $`Z_{k_F}=0`$, which is a signature of marginal- and Luttinger-liquid systems, corresponds, according to Eq. (27), to the case where $`\mathrm{\Sigma }(k_F;\epsilon )/\epsilon `$ is unbounded at $`\epsilon =\epsilon _F`$. In cases where $`Z_{k_F}`$ is vanishing, Eq. (26) formally implies that, for the Fermi velocity $`v_F`$ to be finite, $`\mathrm{\Sigma }(k;\epsilon _F)/k`$ is also to be unbounded at $`k=k_F`$. Although formally Eq. (24) can be maintained with $`v_F=\mathrm{}^1\mathrm{\Sigma }(k;\epsilon _F)/k|_{k=k_F}/\mathrm{\Sigma }(k_F;\epsilon )/\epsilon |_{\epsilon =\epsilon _F}`$, it should be evident, however, that in the case at hand the right-hand side of Eq. (24) neglects the fact that $`\mathrm{\Sigma }(k;\epsilon )`$ is singular at both $`k=k_F`$ and $`\epsilon =\epsilon _F`$ with the singularities showing up already in the first derivatives, thus invalidating this formal definition of the Fermi velocity. To appreciate the significance of this observation, note that the key element in obtaining Eq. (24) has been the assumption of continuous differentiability of $`\mathrm{\Sigma }(k;\epsilon )`$ in some neighbourhoods of $`k=k_F`$ and $`\epsilon =\epsilon _F`$, whereby it has been possible to write <sup>§‡</sup><sup>§‡</sup>§‡ Compare with Eq. (24). See also footnotes I and §\*IV C. $`\mathrm{\Sigma }(k;\epsilon _F)\mathrm{\Sigma }(k_F;\epsilon _F)+\beta (kk_F)+o(kk_F)`$ and $`\mathrm{\Sigma }(k_F;\epsilon )\mathrm{\Sigma }(k_F;\epsilon _F)+\gamma (\epsilon \epsilon _F)+o(\epsilon \epsilon _F)`$, with $`\beta `$ and $`\gamma `$ bounded, in some neighbourhoods of $`k=k_F`$ and $`\epsilon =\epsilon _F`$, respectively. These asymptotic relations signify the facts that $`\{\mathrm{\Sigma }(k;\epsilon _F)\mathrm{\Sigma }(k_F;\epsilon _F)\}/(kk_F)`$ and $`\{\mathrm{\Sigma }(k_F;\epsilon )\mathrm{\Sigma }(k_F;\epsilon _F)\}/(\epsilon \epsilon _F)`$ are vanishing for $`kk_F`$ and $`\epsilon \epsilon _F`$, respectively. With, for instance, $`\mathrm{\Sigma }(k;\epsilon _F)\mathrm{\Sigma }(k_F;\epsilon _F)+\beta |kk_F|^{\gamma _0}`$, $`0<\gamma _0<1`$ (in the present case $`\alpha :=2\gamma _0`$ is referred to as the anomalous dimension; see Appendix D, in particular paragraph following Eq. (D39)), for $`kk_F`$, or $`\mathrm{\Sigma }(k_F;\epsilon )\mathrm{\Sigma }(k_F;\epsilon _F)+\gamma (\epsilon \epsilon _F)\mathrm{ln}|\epsilon \epsilon _F|`$, for $`\epsilon \epsilon _F`$, which have not been accounted for by Eq. (24), $`\{\mathrm{\Sigma }(k;\epsilon _F)\mathrm{\Sigma }(k_F;\epsilon _F)\}/(kk_F)`$ and $`\{\mathrm{\Sigma }(k_F;\epsilon )\mathrm{\Sigma }(k_F;\epsilon _F)\}/(\epsilon \epsilon _F)`$ clearly diverge for $`kk_F`$ and $`\epsilon \epsilon _F`$, respectively, and consequently both Eq. (24) and the notion of Fermi velocity of QPs in the vicinity of the Fermi surface become altogether meaningless. <sup>§§</sup><sup>§§</sup>§§ We note in passing that for metallic systems the locus of the $`𝐤`$-points satisfying $`\epsilon _k^0+\mathrm{}\mathrm{\Sigma }(k;\epsilon _F)=\epsilon _F`$ (compare with Eq. (22)) is the Fermi surface, with $`\epsilon _k^0+\mathrm{}\mathrm{\Sigma }(k;\epsilon _F)<\epsilon _F`$ ($`>\epsilon _F`$) defining the interior (exterior) of the Fermi sea (Galitskii and Migdal 1958, Luttinger 1960, Eqs. (6) and (94)). Clearly, the existence of a Fermi surface is independent of the value of $`Z_{k_F}`$, as the latter is determined by the derivative with respect to $`\epsilon `$ of $`\mathrm{\Sigma }(k_F;\epsilon )`$; see Eq. (27). In particular $`Z_{k_F}=0`$ does not rule out a Fermi surface. We point out that $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon _F)0`$, for all $`k`$ (see footnote ‡∥IV B above), and that in an isotropic system the Fermi surface can consist of disconnected concentric surfaces; since $`\epsilon _k^0`$ is a monotonically increasing function of $`k`$, such Fermi surfaces can exist (in isotropic systems) only in consequence of the electron-electron interaction: for $`\epsilon _k^0+\mathrm{}\mathrm{\Sigma }(k;\epsilon _F)=\epsilon _F`$ to have more than one solution, it is necessary that $`\mathrm{}\mathrm{\Sigma }(k;\epsilon _F)`$ can counter $`\epsilon _k^0`$. In the present work we explicitly deal with a single Fermi surface of radius $`k_F`$.
The following two observations should be clarifying: first, any positive non-integer $`\gamma _0`$ in the above expression would signal a non-analytic singularity in the $`k`$-dependence of $`\mathrm{\Sigma }(k;\epsilon )`$, even when $`\gamma _0>1`$; second, even if one puts aside the possibility of non-differentiability with respect to $`k`$ of $`\mathrm{\Sigma }(k;\epsilon )`$ in a neighbourhood of $`k=k_F`$, it is certain that $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ possesses an essential singularity at $`z=\epsilon _F\mu _N`$. For instance, Eq. (24) with $`v_F`$ defined in accordance with Eqs. (25) and (26) takes no account of the possibility of, for example, $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )^{}(\epsilon \epsilon _F)`$ whose corresponding $`\mathrm{Re}\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)^{}(\epsilon \epsilon _F)\mathrm{ln}|\epsilon \epsilon _F|`$ as $`\epsilon \epsilon _F`$ (see text following Eq. (35) below as well as Appendix C) which is specific to ‘marginal’ Fermi liquids (Varma, et al. 1989, Littlewood and Varma 1991, Kotliar, et al. 1991).
Having demonstrated the difficulties that arise as a consequence of the fact that $`\mathrm{\Sigma }(k;\epsilon )`$ (and, owing to Eq. (22), $`\epsilon _k`$) is not a priori ‘smooth’ (in our above-indicated sense), it becomes evident that the implicit assumption in the proof of the Luttinger (1961) theorem, as though $`\mathrm{\Sigma }(k;\epsilon )`$ were invariably ‘smooth’, has no theoretical justification. In particular, we should like to emphasise that our above exposition makes evident that no theoretical treatments that aim to expose the behaviour of $`\mathrm{\Sigma }(k;\epsilon )`$, for $`\epsilon \epsilon _F`$, whether $`\mathrm{\Sigma }(k;\epsilon )`$ is evaluated by perturbative or non-perturbative techniques, should rely on the expansion in Eq. (24), since any consistent treatment that relies on this expansion is bound to arriving at the conclusion that $`\mathrm{\Sigma }(k;\epsilon )`$ is Fermi-liquid like (validity of the expansion in Eq. (24) is necessary and sufficient for the satisfaction of conditions (A) and (B) in § I); a different conclusion must of necessity signal some inconsistency (or inconsistencies) in the treatment, including plain algebraic errors.
One of the main conclusions that may be drawn from our above analyses is that a non-Fermi-liquid-type $`\mathrm{\Sigma }(k;\epsilon )`$ does not necessarily signal breakdown of the many-body perturbation theory. Consequently, we observe that non-Fermi-liquid behaviour can in principle be manifested in any spatial dimension higher than $`d=1`$.
### D A specific feature
Here we comment on one of the specific features in the treatment by Luttinger (1961) which brings out some additional aspects concerning the above-indicated implicit assumption. The diagram in Fig. 1 has the property that it involves a so-called ‘polarisation part’, that is a part which can be disconnected from the self-energy diagram through removing two interaction lines (indicated by broken lines). Momentum conservation implies that to both of these interaction lines the same momentum-transfer vector is associated. Consequently, any polarisation part in a self-energy diagram contributes at least a square of the electron-electron interaction to the respective integrand. As a consequence, the contribution of the diagram in Fig. 1 is divergent when the electron-electron interaction is the bare Coulomb interaction (for a comprehensive treatment see Mattuck (1992, § 10.4)) for which $`v(𝐫𝐫^{})1/𝐫𝐫^{}`$ holds and thus $`\overline{v}(k)1/k^2`$ in three spatial dimensions. Thus Luttinger’s analysis is not directly applicable to systems of particles interacting via the Coulomb interaction. <sup>§¶</sup><sup>§¶</sup>§¶ Here we are referring to metallic systems. For non-metallic ground states, this problem of divergent self-energy contributions does not arise. Since non-metallic systems are by definition non-Fermi liquids, they do not concern us here. However, as is well known, the contribution of the set of all divergent polarisation diagrams (the random-phase approximation, RPA, bubble-like diagrams) can be exactly calculated (as these diagrams give rise to a geometric series) and one observes that screening effects in the static limit remove the singularity of the bare Coulomb interaction function $`\overline{v}(k)`$ at $`k=0`$ (see, e.g., Fetter and Walecka 1971, pp. 178 and 179). In analysing the leading-order term in $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ as $`\epsilon \epsilon _F`$, the $`\epsilon `$-dependence of the dynamically-screened interaction function (Hubbard 1957) $`W(\epsilon )`$ necessitates knowledge not only of the energy dispersion of the QP excitations, i.e. $`\epsilon _k`$, but also of that of the bosonic neutral excitations, i.e. $`𝖾_k`$; here for the energies of the ‘neutral’ excitations we have $`𝖾_s:=E_{N,s}E_{N,0}0`$ (Fetter and Walecka 1971). For an interacting system, the exact dispersion of $`𝖾_k`$ is unknown, similar to that of $`\epsilon _k`$; $`𝖾_k`$ can be in principle determined from the dynamical density-density correlation function $`\chi (k;\epsilon )`$; for $`𝖾_k`$ this function plays a similar role as $`G(k;\epsilon )`$ does for $`\epsilon _k`$. It follows that for metallic Coulomb systems, the behaviour of $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ cannot be solely determined from the perturbation series expansion for the self-energy operator; one needs in addition the perturbation series expansion for $`\chi (k;\epsilon )`$.
From the above arguments we conclude that, even if the implicit assumption by Luttinger (1961), discussed in §§ IV.B and IV.C, were correct, Luttinger’s asymptotic expression for $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ as presented in the Abstract, would not be justifiably applicable to self-energies of metallic systems of particles interacting via the Coulomb interaction function (see § I for references to studies on systems of particles interacting through long-range repulsive interaction functions).
## V A brief discussion of a Migdal’s theorem
In § IV we have mentioned that the quantity $`Z_{k_F}`$, which coincides with weight of the Landau quasi-particles on the Fermi surface, is also equal to the amount of discontinuity of the momentum-distribution function $`𝗇(p)`$ at $`p_F=\mathrm{}k_F`$. This statement constitutes a theorem due to Migdal (1957) (Luttinger 1960). Here we briefly discuss this theorem in light of our above considerations. In what follows we deal with $`𝗇(k)`$ and use the commonly-employed designation ‘momentum-distribution function’, despite the fact that $`k=p/\mathrm{}`$ is not momentum. For $`𝗇(k)`$ we have
$$𝗇(k):=\mathrm{\Psi }_{N,0}|\widehat{a}_k^{}\widehat{a}_k|\mathrm{\Psi }_{N,0}\frac{1}{\mathrm{}}_𝒞\frac{\mathrm{d}z}{2\pi i}\stackrel{~}{G}(k;z),$$
(28)
where $`\widehat{a}_k^{}`$ and $`\widehat{a}_k`$ stand for creation and annihilation operators, respectively. The contour $`𝒞`$ of integration is depicted in Fig. 2. To prevent confusion, we point out that the right-hand side of Eq. (28) differs by a minus sign from the corresponding expression in Migdal’s (1957) work. This is because the single-particle Green function adopted here is defined according to the modern convention and is equal to minus the Green function employed by Migdal (1957). Further, our use of the symbol $`\stackrel{~}{G}(k;z)`$ (rather than $`G(k;\epsilon )`$ which is the common notation) is in accordance with our present conventions (see § A.3) and has the advantage that our following discussions will not suffer from mathematical ambiguities. One of these ambiguities can be found in the original work by Migdal (1957). Migdal stated namely that, for $`k`$ infinitesimally less than $`k_F`$ (denoted by $`k_F^{}`$), $`\stackrel{~}{G}(k;z)`$ would have a pole, with an infinitesimal imaginary part (apparently due to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ as indicated in the Abstract), enclosed by $`𝒞`$. Migdal (1957) further asserted that, through increasing $`k`$ from $`k_F^{}`$ to $`k_F^+`$ ($`k_F^+`$ denotes a value infinitesimally larger than $`k_F`$), the imaginary part of the mentioned pole would change sign, upon which this pole would leave the interior of $`𝒞`$. Since the regular or incoherent part of $`\stackrel{~}{G}(k;z)`$ does not contribute to the difference $`𝗇(k_F^{})𝗇(k_F^+)`$, from Eq. (28) it would follow that $`Z_{k_F}`$ would be the residue of an isolated pole of $`\stackrel{~}{G}(k_F^{};z)`$ at $`z=\epsilon _F`$. The Migdal theorem (Migdal 1957, Luttinger 1960) thus states
$$𝗇(k_F^{})𝗇(k_F^+)=Z_{k_F}.$$
(29)
From the Dyson equation we have
$$\stackrel{~}{G}(k;z)=\frac{\mathrm{}}{z\stackrel{~}{E}_k(z)},\text{where}\stackrel{~}{E}_k(z):=\epsilon _k^0+\mathrm{}\stackrel{~}{\mathrm{\Sigma }}(k;z).$$
(30)
Poles of $`\stackrel{~}{G}(k;z)`$ are thus seen to be the isolated (see § A.1) solutions of
$$\stackrel{~}{E}_k(z)=z,$$
(31)
which is equivalent to Eq. (18) specialised to the uniform isotropic system under consideration. According to our conventions, $`k_F^{}`$ has to be identified with $`k_F`$ and we need only to have $`k_F^+`$ for indicating a $`k`$ infinitesimally larger than $`k_F`$. As we have mentioned in § IV.C, indeed $`E_{k_F}(\epsilon _F)=\epsilon _F`$ must be satisfied. However, since $`\stackrel{~}{\mathrm{\Sigma }}(k_F;z)`$ has a branch point at $`z=\epsilon _F`$, it follows that $`\stackrel{~}{G}(k_F;z)`$ cannot have a pole at $`z=\epsilon _F`$, since a branch point is not an isolated singularity (see § A.1). Moreover, as we have discussed in § III, Eq. (18), or its equivalent, Eq. (31), cannot have a complex-valued solution, so that there cannot be any question of a ‘pole’ of $`\stackrel{~}{G}(k;z)`$ changing the sign of its imaginary part. The question, therefore, arises as to how, despite the fact that $`\stackrel{~}{G}(k_F;z)`$, or $`G(k_F;\epsilon )`$ in Migdal’s notation, has no pole at $`z=\epsilon _F`$, the above-introduced Migdal (1957) theorem could hold.
To answer the above question, we note that since, for $`k=k_F`$ and $`z\epsilon _F`$, $`z\stackrel{~}{E}_k(z)`$ approaches zero (by continuity and the fact that Eq. (23) applies), in principle it may be possible to write
$$\stackrel{~}{G}^1(k_F;z)\frac{1}{\mathrm{}}\left\{z[\epsilon _{k_F}^0+\mathrm{}\stackrel{~}{\mathrm{\Sigma }}(k_F;z)]\right\}\alpha (z\epsilon _F)+\stackrel{~}{L}(z),\text{as}z\epsilon _F,$$
(32)
with $`\stackrel{~}{L}(z)`$ satisfying $`lim_{z\epsilon _F}\stackrel{~}{L}(z)/(z\epsilon _F)=0`$, that is $`\stackrel{~}{L}(z)=o(z\epsilon _F)`$ for $`z\epsilon _F`$. <sup>§∥</sup><sup>§∥</sup>§∥ Because of $`\stackrel{~}{G}(z^{})=\stackrel{~}{G}^{}(z)`$, it follows that $`\stackrel{~}{L}(z^{})=\stackrel{~}{L}^{}(z)`$. Further, unless $`\epsilon (\mu _N,\mu _{N+1})`$, in general $`lim_{\eta 0}\{\stackrel{~}{L}(\epsilon +i\eta )\stackrel{~}{L}(\epsilon i\eta )\}0`$. This is a consequence of the fact that $`z=\mu _N\epsilon _F`$ and $`z=\mu _{N+1}`$ are branch points of $`\stackrel{~}{G}(k;z)`$, as discussed in § III. In Eq. (32), $`\alpha `$ stands for a constant to be specified below. Equation (32) implies $`\stackrel{~}{G}^1(k_F;z)`$, and thus $`\stackrel{~}{\mathrm{\Sigma }}(k_F;z)`$, to be at least once continuously differentiable with respect to $`z`$ at $`z=\epsilon _F`$ (see footnote ‡‡I); $`m`$th continuous differentiability of $`\stackrel{~}{G}^1(k_F;z)`$, and thus of $`\stackrel{~}{\mathrm{\Sigma }}(k_F;z)`$, with respect to $`z`$ at $`z=\epsilon _F`$ amounts to the condition $`\stackrel{~}{G}^1(k_F;z)_{j=0}^m\alpha _j(z\epsilon _F)^j+\stackrel{~}{L}(z)`$ with $`\{\alpha _j\}`$ finite constants and $`lim_{z\epsilon _F}\stackrel{~}{L}(z)/(z\epsilon _F)^m=0`$. <sup>§\**</sup><sup>§\**</sup>§\** We note that, when a function, say $`\stackrel{~}{f}(z)`$, is analytic inside and on the boundary of a simply-connected region $``$ of the complex $`z`$-plane, it is infinitely-many times differentiable in $``$, so that once-differentiability of $`\stackrel{~}{f}(z)`$ at a point, say $`z_0`$, interior to $``$ implies infinitely-many times differentiability of $`\stackrel{~}{f}(z)`$ at $`z=z_0`$. In the case under consideration, since $`z=\epsilon _F`$ is a branch point, it cannot be interior to any simply-connected region of analyticity of $`\stackrel{~}{G}(k_F;z)`$ (see § A.1). Consequently, $`m`$th differentiability of this function at $`z=\epsilon _F`$ does not imply its $`(m+1)`$th differentiability at this point. Consider for instance $`\stackrel{~}{f}(z):=(zz_0)^p`$ with $`p>0`$ a real non-integer. This function has two branch points, $`z=z_0`$ and $`1/z=0`$. With $`m`$ an integer satisfying $`m<p<m+1`$, $`\stackrel{~}{f}(z)`$ is seen to be $`m`$ times, but not $`m+1`$ times, differentiable at $`z=z_0`$. Consequently, $`\mathrm{\Sigma }(k_F;\epsilon )/\epsilon |_{\epsilon =\epsilon _F}`$ is well-defined and one has $`\alpha \mathrm{}^1\left(1\mathrm{}\mathrm{\Sigma }(k_F;\epsilon )/\epsilon |_{\epsilon =\epsilon _F}\right)1/(\mathrm{}Z_{k_F})`$ (see Eq. (27) above). Using the representation $`z\epsilon _F=\varrho \mathrm{exp}(i\phi )`$, with $`\varrho `$ a non-vanishing constant, to be let to approach zero, from Eq. (28) it can straightforwardly be shown that Eq. (29) remains intact, even though $`z=\epsilon _F`$ is not a pole (i.e. an isolated singularity) of $`\stackrel{~}{G}(k;z)`$.
Through the Kramers-Kronig relation (see Appendices B and C) for $`\mathrm{Re}\mathrm{\Sigma }(k;\epsilon )`$ in terms of $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$ and the asymptotic relation $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, for $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, it can readily be deduced that (see Appendix C)
$$\mathrm{Re}\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)+\beta _k(\epsilon \epsilon _F),\text{with}\beta _k0\text{for}kk_F.$$
(33)
Concerning the sign assigned to $`\beta _k`$ for $`k`$ close to $`k_F`$, as is seen from Eq. (35) below, a positive $`\beta _{k_F}`$ would imply $`Z_{k_F}>1`$ which is impossible on account of a combination of the following three reasons.
(i) By definition (see Eq. (28) above) $`0𝗇(k)1`$; this can also be traced back to the Pauli exclusion principle.
(ii) According to Eq. (29), $`Z_{k_F}=𝗇(k_F^{})𝗇(k_F^+)`$ so that in view of (i), $`|Z_{k_F}|1`$.
(iii) $`Z_{k_F}<0`$ is excluded owing to $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon ){}_{>}{}^{<}0`$, for $`\epsilon {}_{<}{}^{>}\mu `$.
In Appendix C (see text following Eq. (C9)) we demonstrate that $`\beta _k0`$, for $`k`$ in a neighbourhood of $`k_F`$, directly follows from a general sum rule involving $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$. We therefore have <sup>§††</sup><sup>§††</sup>§†† Consider the analytic function $`\stackrel{~}{f}(z)𝗎(x,y)+i𝗏(x,y)`$, where $`x+iyz`$ and $`𝗎(x,y)`$ and $`𝗏(x,y)`$ are real. Upon setting $`y`$ equal to zero, we obtain $`\stackrel{~}{f}(x)=𝗎(x,0)+i𝗏(x,0)`$. Replacing $`x`$ in the argument of $`\stackrel{~}{f}`$ by $`z`$, we obtain $`\stackrel{~}{f}(z)=𝗎(z,0)+i𝗏(z,0)`$ (in a similar manner one obtains $`\stackrel{~}{f}(z)=𝗎(0,iz)+i𝗏(0,iz)`$). This expression prescribes how to obtain $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ from the knowledge of $`\mathrm{\Sigma }(k;\epsilon )`$. In the present case, as $`z\epsilon _F`$, we have: $`\stackrel{~}{\mathrm{\Sigma }}(k;z)\mathrm{\Sigma }(k;\epsilon _F)+\beta _k(z\epsilon _F)i\alpha _k(z\epsilon _F)^2`$, for $`0\mathrm{arg}(z)<\pi `$, and $`\stackrel{~}{\mathrm{\Sigma }}(k;z)\mathrm{\Sigma }(k;\epsilon _F)+\beta _k(z\epsilon _F)+i\alpha _k(z\epsilon _F)^2`$, for $`\pi \mathrm{arg}(z)<0`$. The self-energy as specified in this way is seen to preserve the reflection property $`\stackrel{~}{\mathrm{\Sigma }}(k;z^{})=\stackrel{~}{\mathrm{\Sigma }}^{}(k;z)`$ for $`\mathrm{Im}(z)0`$ (see Eq. (16)).
$$\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)+\beta _k(\epsilon \epsilon _F)i\alpha _k(\epsilon \epsilon _F)^2,\epsilon {}_{<}{}^{>}\epsilon _{F}^{}.$$
(34)
A simple calculation reveals that for $`Z_{k_F}`$ in Eq. (29) (or Eq. (32))
$$Z_{k_F}(1\mathrm{}\beta _{k_F})^1$$
(35)
holds. A similar analysis based on the Kramers-Kronig relation for the self-energy (see Appendix C) yields that when <sup>§‡‡</sup><sup>§‡‡</sup>§‡‡ Recall that $`^{}`$ indicates that the corresponding asymptotic relation is correct up to a multiplicative constant. $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )^{}|\epsilon \epsilon _F|^\sigma `$ for $`\epsilon \epsilon _F`$, the following hold.
(i) $`\mathrm{Re}\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)^{}(\epsilon \epsilon _F)`$ when $`1<\sigma 2`$.
(ii) $`\mathrm{Re}\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)^{}(\epsilon \epsilon _F)\mathrm{ln}|\epsilon \epsilon _F|`$ when $`\sigma =1`$ (“marginal Fermi liquid” (Varma, et al. 1989, Littlewood and Varma 1991, Kotliar, et al. 1991) (see Appendix C).
(iii) $`\mathrm{Re}\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)^{}|\epsilon \epsilon _F|^\sigma `$ when $`0<\sigma <1`$ (see § IV.C, text following Eq. (27), and identify $`\sigma `$ with $`\gamma _0`$; see also Appendix D and identify $`\sigma `$ with $`12\gamma _0`$ in Eqs. (D28) and (D29) and with $`1\gamma _0`$ in Eqs. (D37) and (D38)).
It can straightforwardly be shown (using the same strategy as in the case of $`\sigma =2`$ employed above) that, for $`0<\sigma 1`$, $`Z_{k_F}`$ on the right-hand side of Eq. (29) is vanishing, whereas for $`1<\sigma 2`$ it takes a finite value. Vanishing of $`Z_{k_F}`$ corresponds to disappearance of the Landau QPs on the Fermi surface. Although for $`1<\sigma 2`$, $`Z_{k_F}`$ is non-vanishing, it is evident that the smaller the $`\sigma `$, the shorter are the life-times of the QPs close to the Fermi surface, if QPs can at all be meaningfully defined (recall condition (B) introduced in § I). We should emphasis that as $`\stackrel{~}{G}(k_F;z)`$ does not possess an isolated pole at $`z=\epsilon _F\mu _N`$, but a branch point, it follows that even in the case of the Landau Fermi liquids, the QPs on the Fermi surface are not truly infinitely long-lived but only so in an asymptotic sense.
We conclude that what Migdal (1957) describes as ‘leaving of a pole from the interior of contour $`𝒞`$ upon changing $`k`$ from $`k_F^{}`$ to $`k_F^+`$’ is to be understood as follows: for $`k=k_F^+`$, Eq. (22) (or Eq. (31)) cannot be satisfied when $`z=\epsilon _F\mu _N`$, that is $`E_{k_F^+}(\mu _N)\mu _N`$; rather we have $`E_{k_F^+}(\mu _{N+1})=\mu _{N+1}`$ and, as can be viewed from Fig. 2, $`z=\mu _{N+1}`$ indeed does not lie in the interior of $`𝒞`$.
In closing this Section, we indicate the following observation. In §§ I and IV.C we argued that Fermi liquids are distinguished by two specific aspects, namely that of condition (A), that is continuous differentiability of $`\mathrm{\Sigma }(k_F;\epsilon )`$ with respect to $`\epsilon `$ in a neighbourhood of $`\epsilon =\epsilon _F`$, and of condition (B), that is continuous differentiability of $`\mathrm{\Sigma }(k;\epsilon _F)`$ with respect to $`k`$ in a neighbourhood of $`k=k_F`$. It is owing to these that the Landau quasi-particle dispersion as presented in Eq. (24) can be rigorously obtained from the quasi-particle equation, Eq. (22), which is applicable to all isotropic metallic systems of spin-less fermions; conditions (A) and (B) are prerequisite to assigning a unique mass as well as velocity to an elementary excitation at and in a close vicinity of the Fermi surface (see Eq. (25) above). However, as we have observed in the present Section, in obtaining Eq. (29), only Eq. (32) combined with the condition $`lim_{z\epsilon _F}\stackrel{~}{L}(z)=0`$ (which together embody aspect (A)) played a role. Therefore, a non-vanishing $`Z_{k_F}`$ which is universally considered as the hallmark for Fermi liquids, is only a necessary condition for a metallic system to qualify as a Fermi liquid, but by no means a sufficient condition (see § I). In other words, metallic systems with $`Z_{k_F}0`$ can in principle be non-Fermi liquids in the above sense.
## VI (Non-)Fermi liquids and non-differentiability of the self-energy with respect to momentum close to and on the Fermi ‘surface’
From Eq. (28), making use of Eq. (30), one obtains
$$\frac{𝗇(k)}{k}=_𝒞\frac{\mathrm{d}z}{2\pi i}\frac{1}{\left(z\stackrel{~}{E}_k(z)\right)^2}\frac{\stackrel{~}{E}_k(z)}{k}.$$
(36)
In view of the singular behaviour of $`𝗇(k)`$ at $`k=k_F`$, in employing Eq. (36) we approach $`k_F`$ from either left or right, i.e. by $`𝗇(k)/k|_{k=k_F^{}}`$ we imply left and right derivatives, respectively.
For our further discussions it is convenient to introduce some notational conventions. We denote the circular contour circumscribing $`\epsilon _F`$ in Fig. 2 by $`𝒞_s`$ (subscript $`s`$ refers to ‘singular’, in view of the singularity of $`\stackrel{~}{G}(k_F;z)`$ at $`z=\epsilon _F`$) and its complement, $`𝒞\backslash 𝒞_s`$, by $`𝒞_r`$ (subscript $`r`$ refers to ‘regular’). Consequently, we write $`𝗇(k)=𝗇_s(k)+𝗇_r(k)`$ with $`𝗇_s(k)`$ and $`𝗇_r(k)`$ corresponding to the $`z`$-integrals along $`𝒞_s`$ and $`𝒞_r`$, respectively (here we assume the radius $`\varrho `$ characterising $`𝒞_s`$ to be infinitesimally small).
In arriving at the result in Eq. (29), we have made use of $`𝗇_r(k_F^+)=𝗇_r(k_F^{})`$ so that Eq. (29) in fact amounts to $`𝗇_s(k_F^{})𝗇_s(k_F^+)=Z_{k_F}`$. A similar property, namely $`𝗇_r(k)/k|_{k=k_F^{}}=𝗇_r(k)/k|_{k=k_F^+}`$, may not be applicable in general, in particular since $`𝗇_r(k)/k|_{k=k_F^{}}`$ may be unbounded. Below we demonstrate that: <sup>¶\*</sup><sup>¶\*</sup>¶\* Similar considerations can be given for the case $`k=k_F^+`$. i) for Fermi liquids $`𝗇(k)/k|_{k=k_F^{}}`$ is either finite, in which case it follows that $`\stackrel{~}{\mathrm{\Sigma }}(k;z)/k|_{k=k_F^{}}`$ is finite for all $`z`$, or infinite, in which case $`\beta _k/k|_{k=k_F^{}}`$ is also infinite and consequently $`\{\stackrel{~}{\mathrm{\Sigma }}(k;z)\mathrm{\Sigma }(k;\epsilon _F)\}/k|_{k=k_F^{}}`$ is infinite for all $`z`$ (note that, for Fermi liquids, $`\mathrm{\Sigma }(k;\epsilon _F)`$ is by definition a continuously differentiable function of $`k`$ in a neighbourhood of $`k_F`$); ii) for non-Fermi liquids characterised by $`Z_{k_F}=0`$, $`𝗇_s(k)/k|_{k=k_F^{}}=0`$. Further, although for these systems $`𝗇(k_F^{})𝗇(k_F^+)=0`$ holds, $`𝗇(k)/k`$ may or may not diverge as $`kk_F`$ ($`kk_F`$); in the one-dimensional Luttinger model (Luttinger 1963, Mattis and Lieb 1965), which is a prototype of one-dimensional systems of interacting spin-less fermions (Haldane 1980, 1981), with ‘weak’ interaction — as characterised by a small anomalous dimension, <sup>¶†</sup><sup>¶†</sup>¶† The non-interacting case corresponds to $`\alpha 2\gamma _0=0`$. See § IV.C and Appendix D. namely $`0<\alpha <1`$ —, $`𝗇(k)/k`$ does diverge for $`kk_F^{}`$, and consequently $`\stackrel{~}{\mathrm{\Sigma }}(k;z)/k|_{k=k_F^{}}`$ is infinite for all $`z`$; here, contrary to the case of Fermi liquids, $`\mathrm{\Sigma }(k;\epsilon _F)/k|_{k=k_F^{}}`$ is not necessarily finite. This aspect can be explicitly verified from the available explicit expression for the spectral function corresponding to the retarded single-particle Green function (see Appendix D).
Now we proceed with demonstrating the above statements. In cases where Eq. (32) and $`lim_{z\epsilon _F}\stackrel{~}{L}(z)=0`$ are satisfied, making use of the Cauchy residue theorem, one obtains
$$\frac{𝗇_s(k)}{k}|_{k=k_F^{}}=\mathrm{}Z_{k_F}^2\frac{^2\mathrm{\Sigma }(k;\epsilon )}{\epsilon k}|_{\genfrac{}{}{0pt}{}{\epsilon =\epsilon _F}{k=k_F^{}}},\frac{𝗇_s(k)}{k}|_{k=k_F^+}=0.$$
(37)
In view of the derivative with respect to $`\epsilon `$, $`\mathrm{\Sigma }(k;\epsilon )`$ in Eq. (37) may be replaced by $`\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)`$. From this it follows that a possible divergence of $`𝗇_s(k)/k`$ for $`kk_F`$ cannot be due to $`\mathrm{\Sigma }(k;\epsilon _F)`$ and therefore a divergent $`𝗇_s(k)/k`$ for $`kk_F`$ is not in conflict with the fundamental assumption of the Fermi-liquid theory, namely that of continuous differentiability of $`\mathrm{\Sigma }(k;\epsilon _F)`$ in a neighbourhood of $`k=k_F`$. We note in passing that, through employing Eqs. (34) and (35), we can also write $`𝗇_s(k)/k|_{k=k_F^{}}=Z_k/k|_{k=k_F}`$, where $`Z_k:=(1\mathrm{}\beta _k)^1`$. <sup>¶‡</sup><sup>¶‡</sup>¶‡ We should emphasise that here $`Z_k`$ is merely a generalisation of $`Z_{k_F}`$ and does not necessarily have the same physical significance as $`Z_{k_F}`$ does. To appreciate the relevance of this remark, one should realise that for $`kk_F`$, Eq. (22) may not have a real-valued solution. In Appendix C we demonstrate that divergence of $`\beta _k/k`$ (which in view of the latter result implies that of $`𝗇_s(k)/k`$) at any $`k=k_0`$, e.g. $`k=k_F`$, signals divergence of $`\stackrel{~}{\mathrm{\Sigma }}(k;z)/k|_{k=k_0}`$ for all $`z`$, $`\mathrm{Im}(z)0`$; in Appendix D we explicitly show this to be the case for the (one-dimensional) Luttinger model for spin-less fermions. With reference to Eq. (36), as well as Eq. (30) where $`\stackrel{~}{E}_k(z)`$ is defined in terms of $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$, divergence of $`𝗇_s(k)/k`$ at $`k=k_F^{}`$ is thus seen to imply divergence of $`𝗇(k)/k`$ at $`k=k_F^{}`$. This completes demonstration of our point (i) above. It is relevant to mention that for a three-dimensional isotropic system of fermions interacting through a hard-core potential, qualifying as a Fermi liquid, $`𝗇(k)/k`$ has been shown to be logarithmically divergent for $`kk_F^{}`$ (Belyakov 1961, Sartor and Mahaux 1980). <sup>¶§</sup><sup>¶§</sup>¶§ In our judgement, nothing precludes the possibility that this divergence is an artifact of the second-order perturbation theory employed in the treatment by Belyakov (1961). On the other hand, for a similar system of electrons, interacting through the long-range Coulomb interaction, no divergence is observed in $`𝗇(k)/k`$, $`kk_F^{}`$, within the framework of the random-phase approximation (RPA) (Daniel and Vosko 1960); calculations of $`𝗇(k)`$ beyond the RPA are also available (Geldart, Houghton and Vosko 1964), however the results do not provide information with regard to regularity or otherwise of $`𝗇(k)/k`$ at $`k=k_F^{}`$. <sup>¶¶</sup><sup>¶¶</sup>¶¶ For a calculation of $`𝗇(k)`$ at large $`k`$ as well as some general discussions concerning $`𝗇(k)`$ see also Yasuhara and Kawazoe (1976).
For non-Fermi liquids, such as the marginal Fermi liquid considered in Appendix C, the result $`𝗇_s(k)/k|_{k=k_F^{}}=0`$ follows exactly for the same reason that $`𝗇_s(k_F^{})=0`$ for these systems: through parametrising $`𝒞_s`$ in terms of the circular-polar angle $`\phi `$, one observes that despite singularity of the integrand at $`\epsilon _F`$, the $`\phi `$-integration over $`[\pi ,\pi )`$ yields an identically-vanishing contribution. As our considerations in § V have made explicit, at the root of this property lies the non-differentiability with respect to $`\epsilon `$ of $`\mathrm{\Sigma }(k_F;\epsilon )`$ at $`\epsilon _F`$. We emphasise that this is sufficient for disqualifying a system to be a Fermi liquid; our discussions in § IV.C have made evident that this non-differentiability condition stands in the way of obtaining the quasi-particle energy dispersion in Eq. (24) from Eq. (22).
For a general non-Fermi-liquid metallic system, in arbitrary $`d`$-dimensional spatial space (if indeed such a system exists in $`d>1`$), our knowledge with regard to behaviour of $`𝗇(k)/k`$ in some neighbourhoods of $`k=k_F^{}`$ cannot surpass that of $`𝗇(k)/k`$ corresponding to Fermi-liquid systems. <sup>¶∥</sup><sup>¶∥</sup>¶∥ Here we employ the notion of ‘neighbourhood’ in a loose sense, since, as we have mentioned earlier, here by $`𝗇(k)/k`$, $`kk_F`$, we express left or right derivative of $`𝗇(k)`$, thus excluding $`k=k_F`$ from the interior of the interval over which $`𝗇(k)`$ is being differentiated. Since, by definition, to a non-Fermi-liquid metallic system no quasi-particle energy dispersion similar to that in Eq. (24) can correspond (see above), and this can be effected through a non-continuous differentiable $`\mathrm{\Sigma }(k_F;\epsilon )`$ in a neighbourhood of $`\epsilon =\epsilon _F`$, the non-continuous-differentiability with respect to $`k`$ of $`\mathrm{\Sigma }(k;\epsilon _F)`$ pertaining to these systems cannot in general be ruled in but also not ruled out. In $`d=1`$ dimension, the exactly-solvable Luttinger (1963) model provides us with the opportunity to arrive at a rigorous statement, however. In this model, for $`k`$ close to $`k_F`$ one has (see, e.g., Voit 1993b) $`𝗇(k)1/2C_1\mathrm{sgn}(kk_F)|kk_F|^\alpha +C_2(kk_F)`$ where $`C_1`$ and $`C_2`$ are constants and $`\alpha 2\gamma _0`$, the ‘anomalous dimension’, is not universal but depends on the nature and strength of the particle-particle interaction in the system under investigation; for the one-dimensional Hubbard model corresponding to a finite repulsive on-site interaction, $`U`$, and away from the half-filling of the Hubbard band, $`0<\alpha <1/8`$, with the upper limit achieved for infinitely large $`U`$. One observes that for $`0<\alpha <1`$, which in view of the latter observation should amount to a wide range of values for $`\alpha `$, $`𝗇(k)`$, although continuous, is not continuously differentiable in a neighbourhood of $`k=k_F`$; $`𝗇(k)/k`$ diverges as $`kk_F^{}`$.
## VII Summary and concluding remarks
In this work we have presented a critical analysis of a celebrated theorem due to Luttinger (1961) concerning energies and life-times of low-energy QP excitations in interacting systems. This theorem has played the crucial role of classifying all metallic systems in spatial dimensions larger than unity as Landau Fermi liquids, with the explicit assumption with regard to applicability of the many-body perturbation theory to all orders to these systems. <sup>¶\**</sup><sup>¶\**</sup>¶\** This on account of the fact that the proof of this theorem has been based on an infinite-order perturbation expansion of the self-energy operator. In this work we have demonstrated that Luttinger’s (1961) proof involves an implicit assumption with regard to the dispersion of the QP energies, which assumption we have explicitly shown to be specific to Fermi-liquid systems. We have therefore shown that Luttinger’s proof amounts to a demonstration of consistency of the mentioned implicit assumption with the property $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )^{}(\epsilon \epsilon _F)^2`$, for $`\epsilon \epsilon _F`$, to all orders of the many-body perturbation theory. It follows that, contrary to the commonly-held view, the many-body perturbation theory does not necessarily break down when applied to systems whose self-energies are of non-Fermi-liquid type. In the absence of any theorem to replace the Luttinger theorem, one could reasonably conjecture that it is most likely that, in particular in two spatial dimensions, non-Fermi-liquid-like metallic systems if not abundant, should not be rare and that these may be correctly addressed within the framework of the many-body perturbation theory. In this connection it is important to point out the following two observations.
First, any static and local approximation to the self-energy operator implies a Fermi-liquid behaviour for the QPs, no matter how ingenious such approximation may be. The entire body of the energy-band methods based on the conventional density-functional theory, involving a local and energy-independent effective potential (for a comprehensive review see Dreizler and Gross 1990), pertains to this category of approximation frameworks. The (self-consistent) Hartree-Fock scheme which involves the non-local static exchange self-energy, may describe non-Fermi liquid metals, however only those whose $`Z_{k_F}=1`$. <sup>¶††</sup><sup>¶††</sup>¶†† For uniform systems of electrons interacting through the Coulomb potential, the derivative with respect to $`k`$ of the Hartree-Fock self- energy is logarithmically divergent at $`k=k_F`$ (Ashcroft and Mermin 1981, p. 334) — see Appendix B. Since $`\mathrm{\Sigma }^{HF}`$ is independent of $`\epsilon `$ (Appendix B), it is therefore analytic over the entire $`z`$-plane and thus the associated $`Z_{k_F}=1`$ (see § V). Therefore a metal within the Hartree-Fock scheme may be non-Fermi liquid solely on account of violation of condition (B) introduced in § I. Second, in evaluating the self-energy operator for a Coulomb system, within the framework of the many-body perturbation theory, it has to be realised that in metallic systems the long range of the bare Coulomb interaction gives rise to divergent self-energy contributions (from the second order in the interaction onwards). Thus, for such systems, perturbation expansion has to be in terms of the dynamically-screened electron-electron interaction function. In calculating this function, at least all polarisation diagrams of the random-phase approximation have to be taken into account, for all these diagrams have unbounded contributions and only in combination give rise to a finite result. Our considerations with regard to the problem of unbounded self-energy diagrams have led us to draw the additional conclusion that, irrespective of our above-indicated finding, Luttinger’s (1961) proof cannot have bearing on metallic systems of electrons interacting through the long-range Coulomb interaction function. This follows from the fact that Luttinger’s proof has been based on a perturbation expansion of the self-energy operator in terms of the bare electron-electron interaction. In this connection we point out that analytic properties of a series of functions of complex variable do not coincide with those of the constituent terms, unless the series be uniformly convergent (Whittaker and Watson 1927, pp. 91 and 92, Titchmarsh 1939, pp. 95-98). Further, the sum of a series which is not absolutely convergent depends on the order of summation of terms in the series (Whittaker and Watson 1927, pp. 18 and 25). It is evident that, in particular, perturbation series which involve singular terms cannot be absolutely or uniformly convergent.
We have further elaborated upon the analytic properties of the single-particle Green function and the self-energy operator as functions of a complex-valued energy parameter, $`z`$. In particular we have pointed out that the commonly-used equation for energies of the QPs in an interacting system can yield, if any, only real-valued solutions. To obtain complex-valued solutions (concerning systems in the thermodynamic limit), the self-energy operator has to be analytically continued into a non-physical RS of the complex energy plane.
We have paid careful attention to the precise nature of the singular points of the single-particle Green function and the self-energy operator. For instance we have pointed out that the Fermi energy is not an isolated singularity of the Green function and therefore cannot be a pole. In the light of this, we have considered a celebrated theorem due to Migdal (1957) and discussed how its proof is independent of the commonly-made assumption, that on the Fermi surface the Fermi energy were a pole of the single-particle Green function. We have explicitly shown that for the momentum distribution function $`𝗇(k)`$ pertaining to a metallic system to be discontinuous at $`k=k_F`$ by a finite amount — the magnitude of this discontinuity being, according to the Migdal theorem, equal to $`Z_{k_F}`$ —, it is only necessary that the self-energy $`\mathrm{\Sigma }(k_F;\epsilon )`$ be a continuously differentiable function of $`\epsilon `$ in a neighbourhood of $`\epsilon =\epsilon _F`$. Consequently metallic systems with $`Z_{k_F}0`$ are not necessarily Fermi liquids. Fermi liquids, we have shown, have in addition the property that their corresponding $`\mathrm{\Sigma }(k;\epsilon _F)`$ is a continuously differentiable function of $`k`$ in a neighbourhood of $`k=k_F`$. In spite of this, $`\mathrm{\Sigma }(k;\epsilon )`$ pertaining to Fermi liquids may not be continuously differentiable function of $`k`$ in a neighbourhood of $`k=k_F`$ when $`\epsilon \epsilon _F`$. These observations, which we have explicitly examined on the self-energies of Fermi-, marginal Fermi- and Luttinger-liquids, clearly demonstrate how uncritical Taylor expansions of various functions of $`k`$ and $`\epsilon `$ associated with the self-energy can lead to incorrect description of the physical properties of the systems under consideration at low energies.
## Acknowledgements
With appreciation I acknowledge support by the Max-Planck-Gesellschaft, Germany. It is a pleasure for me to thank Professor Peter B. Littlewood and members of the Theory of Condensed Matter Group for their kind hospitality at Cavendish Laboratory where this work was completed. I record my indebtedness to Girton College, Cambridge, for support in the course of my present association with Cavendish Laboratory. I take this opportunity and extend my heartfelt appreciations to Dr Martin W. Ennis for generously and patiently sharing his insight especially into the music of J.S. Bach and kind hospitality.
## A Mathematical preliminaries
Since in the present work we repeatedly encounter a number of specific mathematical notions, we devote this Appendix to a brief exposition of these.
### 1 On the types of singularity
A point at which a function $`g(z)`$ of complex variable $`z`$ is not analytic <sup>¶‡‡</sup><sup>¶‡‡</sup>¶‡‡ Analytic, regular and holomorphic are alternative but equivalent designations. is called a singular point (for details of what follows see, for example, Whittaker and Watson 1927, pp. 102 and 104, Titchmarsh 1939, pp. 89-95). Such a point is either isolated or non-isolated; $`z_0`$ is isolated if there exists a $`\delta >0`$ such that inside the region $`|zz_0|<\delta `$ there exists no other singular point of $`g(z)`$ than $`z_0`$. Otherwise $`z_0`$ is not isolated.
A singularity may be removable; such singularity corresponds to a point $`z_0`$ at which $`g(z)`$ is not defined, but $`lim_{zz_0}g(z)`$ exists. Thus $`z=0`$ is a removable singularity of $`g(z):=\mathrm{sin}(z)/z`$.
Limiting (or accumulation) point of an infinite sequence of poles, is not classified as a pole and thus is considered as an essential singularity, and here, a non-isolated essential singularity. For instance, the sequence of poles of $`g(z):=_{n=0}^{\mathrm{}}1/(n![1+a^{2n}z^2])`$, with $`a>1`$, have $`z=0`$ as their limiting point which is not isolated. This function has no Taylor or Laurent series expansion (see § A.2) over any domain of the $`z`$-plane which has $`z=0`$ as its interior.
Let $`g(z)`$ be a single-valued function throughout a domain $`𝒟`$ at whose interior point $`z_0`$, $`g(z)`$ is singular. Suppose that the principal part of the Laurent series expansion (see § A.2) of $`g(z)`$ around $`z=z_0`$ terminates with the term $`a_n/(zz_0)^n`$, with $`a_n`$ a non-vanishing constant. In such case, $`z_0`$ is called a pole of order $`n`$. Poles are thus by definition isolated singularities.
If the principal part of the Laurent expansion of $`g(z)`$ around $`z=z_0`$ does not terminate (i.e., if there exists no $`n_0`$ such that $`a_n=0`$ for all $`n>n_0`$), $`z_0`$ is an isolated essential singularity of $`g(z)`$. The function $`g(z):=\mathrm{exp}(z)`$ has one such singularity at the point of infinity, i.e. at $`z=1/\zeta `$ where $`\zeta =0`$.
Branch points belong to the class of singular points and concern multi-valued functions. Let $`g(z)`$ be one such function. By traversing a closed contour which circumscribes only one branch point of $`g(z)`$, one obtains, upon arriving at the starting point $`z_1`$, a different value for $`g(z_1)`$, indicating change of the initial branch of $`g(z)`$ into a different branch; for a branch point of order $`p`$, the original branch is recovered after completion of $`p`$ revolutions along the mentioned contour. Thus $`(zz_0)^{1/3}`$ has a third-order branch point at $`z=z_0`$. Functions can also possess branch points of infinite order; for $`g(z):=\mathrm{ln}(z)`$, $`z=0`$ and $`1/z=0`$ are such points.
### 2 Taylor, Laurent and asymptotic series
In the Abstract of this work we have referred to the expression $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, for $`\epsilon \epsilon _F`$, as an asymptotic relation. Here we specify what asymptotic relations and asymptotic series are and in what essential respects these series differ from the Taylor and the Laurent series (Whittaker and Watson 1927, pp. 93, 94 and 100). We also indicate the interrelation between branch points (see § A.1) and divergent asymptotic series.
When $`g(z)`$ is analytic at $`z=z_0`$, then, by definition, there exists an open domain $`𝒟`$ of which $`z_0`$ is an interior and over which $`g(z)`$ is analytic. Within a circle around $`z_0`$ embedded within $`𝒟`$, $`g(z)`$ can be represented in terms of a Taylor series: $`g(z)=_{n=0}^{\mathrm{}}a_n(zz_0)^n`$, with unique coefficients $`a_n`$. Thus $`\{(zz_0)^n|n=0,1,\mathrm{}\}`$ can be considered as a complete basis in a (finite) neighbourhood of $`z_0`$ for any function that is analytic at $`z=z_0`$ (note that this set does not contain such term as, for instance, $`(zz_0)^\sigma `$ with $`\sigma `$ non-integer). This basis can be extended to $`\{(zz_0)^n|n=m,m+1,\mathrm{},0,1,\mathrm{}\}`$ in order to form a complete basis in a (finite) neighbourhood of $`z_0`$ for representing any function which has a pole of order $`m`$ or lower at $`z_0`$ and is analytic in a neighbourhood of $`z=z_0`$; the resulting representation is the well-known Laurent series expansion (§ A.1). It cannot be extended, through increasing $`m`$ in $`\{(zz_0)^n|n=m,m+1,\mathrm{},0,1,\mathrm{}\}`$, for representing functions that possess a branch point at $`z=z_0`$ (see § A.3). It is important to point out that the Taylor and the Laurent series expansions are very special in that they provide uniform representations of an analytic function; when these series are convergent for any $`\varrho `$ and $`\phi `$ in $`zz_0\varrho \mathrm{exp}(i\phi )`$, they are uniformly convergent for all values of $`\phi `$, $`\phi [0,2\pi )`$. Branch points are peculiar in that they do not allow for any uniform representation (note the ‘$``$’ in the above expression for $`\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )`$). We shall return to this point further on in this Section.
Asymptotic series (in the sense of Poincaré) (Whittaker and Watson 1927, Ch. VIII, Copson 1965, Dingle 1973, Lauwerier 1977) are more general than the Taylor and the Laurent series in that they are in terms of “basis functions” that are not necessarily of the form $`(zz_0)^n`$, with $`n`$ possibly negative, zero and positive integer. To an asymptotic series corresponds a so-called asymptotic sequence: The set $`\{\varphi _n(z)|n=0,1,\mathrm{}\}`$ is an asymptotic sequence with respect to $`z_0`$ provided it possesses the property $`\varphi _{n+1}(z)/\varphi _n(z)0`$, denoted by $`\varphi _{n+1}(z)=o\left(\varphi _n(z)\right)`$, for $`zz_0`$. Thus $`\{(zz_0)^n|n=m,m+1,\mathrm{},0,1,\mathrm{}\}`$ is an asymptotic sequence with respect to $`z_0`$. We point out that some authors, for example Whittaker and Watson (1927, Ch. VIII), reserve the designation ‘asymptotic series’ for those series that are both asymptotic (in the above sense) and divergent. We do not follow this restrictive convention.
As we have mentioned above, an asymptotic series expansion of a function may not be convergent. It may not also be uniform; functions with non-uniform asymptotic series expansions for $`zz_0`$ are characterised by possessing different asymptotic series for different sectors of the $`z`$-plane around $`z=z_0`$. Thus, for instance, $`g(z):=\mathrm{exp}(z)+\mathrm{exp}(z)\mathrm{tanh}(1/z)`$ has the following asymptotic forms for $`z0`$: $`g(z)2\mathrm{cosh}(z)2+z^2+\mathrm{}`$, when $`\mathrm{Re}(z)>0`$, and $`g(z)2\mathrm{sinh}(z)2z+z^3/3+\mathrm{}`$, when $`\mathrm{Re}(z)<0`$ (Lauwerier 1977, p. 11). The fact that in different sectors around a point in the complex $`z`$-plane a function can have different asymptotic expansions with respect to the same asymptotic sequence, is referred to as the Stokes phenomenon (Watson 1952, Dingle 1973); ‘Stokes lines’ are thus the branch cuts in our discussions.
The Taylor series of analytic functions based on point $`z_0`$ are convergent uniform asymptotic series for $`zz_0`$ (similarly for the Laurent series). It can be shown (Lauwerier 1977, pp. 12-14) that an analytic function can be associated with even a divergent asymptotic series corresponding to the asymptotic sequence $`\{(zz_0)^n|n=0,1,\mathrm{}\}`$; the given (divergent) series is its asymptotic expansion. Thus, for instance, after Borel (or Euler) transformation (Whittaker and Watson 1927, pp. 154 and 155, Dingle 1973, pp. 405-408, Lauwerier 1977, pp. 45-50) of the divergent series $`f(z):=_{n=0}^{\mathrm{}}(1)^nn!z^n`$, one formally obtains $`f_B(z):=_{n=0}^{\mathrm{}}(1)^nz^n=(1+z)^1`$. Through the Borel back-transformation of $`f_B(z)`$, $`\overline{f}_B(z):=_0^{\mathrm{}}dx\mathrm{exp}(x)f_B(xz)`$, one obtains a function, i.e. $`\overline{f}_B(z)`$, which is analytic in the sector $`\pi <\mathrm{arg}(z)<\pi `$ of the $`z`$-plane. <sup>∥\*</sup><sup>∥\*</sup>∥\* It can be shown that $`\overline{f}_B(z)\mathrm{exp}(1/z)E_1(1/z)/z`$ where $`E_1(z)`$ stands for the exponential-integral function (see, e.g., Abramowitz and Stegun 1972, p. 228). We note that $`\mathrm{exp}(1/z)`$ has an isolated essential singularity at $`z=0`$ (see § A.1) and that $`E_1(1/z)`$ has branch points at $`z=0`$ and $`1/z=0`$. Through replacing $`f_B(xz)`$ by its formal geometric expansion $`1xz+(xz)^2\mathrm{}`$, the term-by-term evaluation of the latter integral yields the original divergent series. We observe that the divergence of the original asymptotic series is closely associated with the restricted sector of the $`z`$-plane around $`z=0`$ over which $`\overline{f}_B(z)`$ is analytic.
### 3 Many-valued functions: Physical and non-physical Riemann sheets
Many-valuedness (Whittaker and Watson 1927, pp. 96-98, Titchmarsh 1939, pp. 138-164) is a generic property of correlation functions of complex energy pertaining to systems in the thermodynamic limit (see § A.4). In particular, the single-particle Green function and the self-energy operator are many-valued (see § II).
An $`n`$-valued function of complex variable $`z`$ over domain $`𝒟`$ may be thought of as embodying $`n`$ branches of a single-valued function over the extended domain consisting of the union of $`n`$ replicas of $`𝒟`$; since these domains signify the same region of the complex plane, they are to be distinguished by considering them as being located on different sheets, namely the Riemann sheets (RSs), of the complex plane. In the present paper we single out one particular branch of the many-valued functions that we encounter, such as the single-particle Green function, because of its physical significance. To this particular branch we refer as the function in question on the physical RS or the ‘physical’ function (compare with $`G(\epsilon )`$ in Eq. (1) which we have designated as the ‘physical’ single-particle Green function).
Throughout this paper we employ the following notational conventions. Let $`f(\epsilon )`$ be a function of real energy variable $`\epsilon `$, which we consider to describe some dynamical physical property of the many-particle system under consideration. By $`\stackrel{~}{f}(z)`$ we denote the analytic function that uniformly approaches $`f(\epsilon )`$ when $`z\epsilon `$; the so-called ‘edge-of-the-wedge’ theorem (Streater and Wightman 1964) asserts uniqueness of $`\stackrel{~}{f}(z)`$. <sup>∥†</sup><sup>∥†</sup>∥† The condition of uniformity can be shown, a posteriori, to be redundant: the existence of the limit $`\stackrel{~}{f}(z)f(\epsilon )`$, as $`z\epsilon `$, with $`\epsilon `$ varying over an open interval where $`f(\epsilon )`$ is continuous, implies uniformity of $`\stackrel{~}{f}(z)f(\epsilon )`$; see p. 75 in Streater and Wightman (1964). We refer to $`\stackrel{~}{f}(z)`$ as $`\stackrel{~}{\stackrel{~}{f}}(z)`$ on the physical RS; $`\stackrel{~}{\stackrel{~}{f}}(z)`$ denotes the above-mentioned single-valued function over the union of all the RSs. We recall that in the main text we have referred to, for example, $`\stackrel{~}{G}(z)`$ in Eq. (8) as the single-particle Green function on the physical RS.
An example will clarify this. Let $`f(\epsilon ):=\mathrm{ln}(\epsilon )`$, for $`\epsilon >0`$. We can identify $`\mathrm{ln}(z)`$ by $`\stackrel{~}{f}(z)`$ since for $`z\epsilon `$, with $`\epsilon >0`$, indeed $`\mathrm{ln}(z)\mathrm{ln}(\epsilon )`$. One can easily verify that $`\mathrm{Ln}_n(z):=\mathrm{ln}|z|+i\{\mathrm{arg}(z)+2\pi n\}`$, with $`n`$ an integer (positive, zero and negative) and $`\pi \mathrm{arg}(z)<\pi `$, is the analytic function (i.e. $`\stackrel{~}{\stackrel{~}{f}}(z)`$) of which $`\stackrel{~}{f}(z):=\mathrm{ln}(z)`$ is the physical branch; thus this branch corresponds to $`n=0`$. The physical branch $`\stackrel{~}{f}(z)`$, like any other branch of $`\stackrel{~}{\stackrel{~}{f}}(z)`$, is a many-valued function (branch points are not removable \[see § A.1\] singularities). It has two branch points, one at $`z=0`$ and the other at $`1/z=0`$. Our above convention $`\pi \mathrm{arg}(z)<\pi `$ specifies the branch cut of $`\stackrel{~}{f}(z)`$, which connects the two branch points of $`\stackrel{~}{f}(z)`$, to be along the negative $`\epsilon `$-axis. One can determine all branches of $`\stackrel{~}{\stackrel{~}{f}}(z)`$ from the knowledge of any of its branches. Two of infinitely many branches of $`\stackrel{~}{\stackrel{~}{f}}(z)`$, $`\stackrel{~}{f}_1(z)`$ and $`\stackrel{~}{f}_2(z)`$ say, that can be obtained by means of a direct analytic continuation of $`\stackrel{~}{f}(z)`$, through its branch cut, into two different non-physical RSs are uniquely determined from the following requirements: $`lim_{\eta 0}\{\stackrel{~}{f}_1(\epsilon +i\eta )\stackrel{~}{f}(\epsilon i\eta )\}=0`$ and $`lim_{\eta 0}\{\stackrel{~}{f}_2(\epsilon i\eta )\stackrel{~}{f}(\epsilon +i\eta )\}=0`$, for $`\epsilon <0`$. It is not difficult to verify that $`\stackrel{~}{f}_1(z)\mathrm{Ln}_1(z)`$ and $`\stackrel{~}{f}_2(z)\mathrm{Ln}_{+1}(z)`$.
### 4 A physically-motivated example
Here we apply the above concepts to a simple model function which accommodates a number of salient features of the physical functions which we deal with in the main part of this paper. We shall particularly emphasise the role played by the process of evaluating the thermodynamic limit in modifying the nature of the singular points of dynamic correlation functions.
Consider $`f(\epsilon ;N,\mathrm{\Omega }):=\frac{1}{\mathrm{\Omega }}_{\mathrm{}}[\mathrm{\Theta }(\epsilon _{\mathrm{}}e_0)\mathrm{\Theta }(\epsilon _{\mathrm{}}e_1)]/(\epsilon _{\mathrm{}}\epsilon )`$, with $`\epsilon _{\mathrm{}+1}>\epsilon _{\mathrm{}}`$. Here $`e_0`$ and $`e_1`$ are finite constants for which we assume $`e_0<\epsilon _{\mathrm{}}<e_1`$ for some values of $`\mathrm{}`$; $`N`$ and $`\mathrm{\Omega }`$ indicate that $`f`$ is a function of the number $`N`$ of particles as well as the volume $`\mathrm{\Omega }`$ of the system. In the ‘thermodynamic limit’ ($`N\mathrm{}`$, $`\mathrm{\Omega }\mathrm{}`$ and finite concentration $`C:=N/\mathrm{\Omega }`$), $`f`$ is only a function of $`C`$. A brief glance at the contents of § II should clarify our present choice for $`f`$.
Suppose that in the thermodynamic limit $`(\epsilon _{\mathrm{}+M}\epsilon _{\mathrm{}})0`$ for any finite value of $`M`$. This condition implies that in taking the limit, the function $`f(\epsilon ;N,\mathrm{\Omega })`$ becomes ill-defined for any real $`\epsilon `$ in the interval $`[e_0,e_1]`$; an $`\epsilon `$ in this interval will be “pinched” (Itzykson and Zuber 1988, pp. 302 and 303) by poles of $`f(\epsilon ;N,\mathrm{\Omega })`$. To avoid this problem, the thermodynamic limit has to be effected after replacing the real energy variable $`\epsilon `$ by a complex energy variable $`z`$, since a complex $`z`$ cannot be “pinched” by the real poles of $`f(z;N,\mathrm{\Omega })`$ as $`N`$ and $`\mathrm{\Omega }`$ approach infinity.
To be specific, let us now assume that in the thermodynamic limit poles of $`f(z;N,\mathrm{\Omega })`$ populate the interval $`[e_0,e_1]`$ with a constant density equal to $`A\times \mathrm{\Omega }`$. We write
$`f(z;N,\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Omega }}}{\displaystyle \underset{\mathrm{}}{}}{\displaystyle \frac{\mathrm{\Theta }(\epsilon _{\mathrm{}}e_0)\mathrm{\Theta }(\epsilon _{\mathrm{}}e_1)}{\epsilon _{\mathrm{}}z}}`$ (A1)
$``$ $`A{\displaystyle _{e_0}^{e_1}}{\displaystyle \frac{\mathrm{d}\epsilon ^{}}{\epsilon ^{}z}}=:\stackrel{~}{f}(z;C),\text{for}N\mathrm{},\mathrm{\Omega }\mathrm{}.`$ (A2)
One trivially obtains
$$\stackrel{~}{f}(z;C)=A\{\mathrm{ln}(ze_1)\mathrm{ln}(ze_0)\}.$$
(A3)
Evidently, the process of evaluating the thermodynamic limit has led to a dramatic change in $`f(z;N,\mathrm{\Omega })`$, transforming it into $`\stackrel{~}{f}(z;C)`$ which analytically is distinct from $`f(z;N,\mathrm{\Omega })`$. For instance, $`\stackrel{~}{f}(z;C)`$ has no poles, rather it possesses two infinite-order branch points, at $`z=e_0`$ and $`z=e_1`$. Consequently, contrary to $`f(z;N,\mathrm{\Omega })`$, which is single-valued, $`\stackrel{~}{f}(z;C)`$ is a many-valued function of $`z`$.
We shall now explicitly demonstrate that $`\stackrel{~}{f}(z;C)`$ as presented in Eq. (A3) is the ‘physical’ branch of $`\stackrel{~}{\stackrel{~}{f}}(z)`$ (see § A.3). To this end, consider $`\stackrel{~}{f}(\epsilon \pm i\eta ;C)A_{e_0}^{e_1}d\epsilon ^{}/(\epsilon ^{}\epsilon i\eta )`$. For $`\eta 0`$ we have $`1/(\epsilon ^{}\epsilon i\eta )=𝒫\{1/(\epsilon ^{}\epsilon )\}\pm i\pi \delta (\epsilon ^{}\epsilon )`$, where $`𝒫`$ denotes the Cauchy ‘principal value’. It is then trivially seen that $`\stackrel{~}{f}(\epsilon \pm i\eta ;C)=A\{\mathrm{ln}|e_1\epsilon |\mathrm{ln}|e_0\epsilon |\}\pm i\pi A\mathrm{\Theta }(\epsilon e_0)\mathrm{\Theta }(e_1\epsilon )`$ for $`\eta 0`$. Hence, for $`\epsilon <e_0`$ and $`\epsilon >e_1`$ the ‘physical’ branch $`\stackrel{~}{f}(\epsilon \pm i\eta ;C)`$ must be real-valued, while for $`\epsilon (e_0,e_1)`$ it must satisfy $`lim_{\eta 0}\{\stackrel{~}{f}(\epsilon +i\eta ;C)\stackrel{~}{f}(\epsilon i\eta ;C)\}=2\pi iA`$; thus the interval $`[e_0,e_1]`$ constitutes the branch cut of $`\stackrel{~}{f}(z;C)`$. These conditions are exactly fulfilled by the expression in Eq. (A3); note that $`\mathrm{ln}(z)`$, as distinct from $`\mathrm{Ln}_n(z)`$ for $`n0`$, stands for the principal-value logarithm (see, e.g., Abramowitz and Stegun 1972, p. 67). Thus $`\stackrel{~}{f}(z;C)`$ in Eq. (A3) is indeed the analytic continuation of $`f(\epsilon ;C)`$ into the physical RS.
It is instructive to consider two examples concerning continuations of $`\stackrel{~}{f}(z;C)`$ into non-physical RSs. Let $`\stackrel{~}{g}(z):=A\{\mathrm{Ln}_1(ze_1)\mathrm{ln}(ze_0)\}`$. It is easily verified that for $`e_0<\epsilon <e_1`$, $`lim_{\eta 0}\{\stackrel{~}{f}(\epsilon +i\eta ;C)\stackrel{~}{g}(\epsilon i\eta )\}=0`$, which implies that $`\stackrel{~}{g}(z)`$ is the analytic continuation of $`\stackrel{~}{f}(z;C)`$ from the upper-half plane of the physical RS through the branch cut $`[e_0,e_1]`$ into the lower-half plane of a non-physical RS. That is $`\stackrel{~}{g}(z)`$ is a branch of $`\stackrel{~}{\stackrel{~}{f}}(z;C)`$ on a non-physical RS. As for the second example, consider $`\stackrel{~}{h}(z):=A\{\mathrm{Ln}_1(ze_1)\mathrm{ln}(ze_0)\}`$. For $`\epsilon (e_0,e_1)`$ we have $`lim_{\eta 0}\{\stackrel{~}{f}(\epsilon i\eta ;C)\stackrel{~}{h}(\epsilon +i\eta )\}=0`$ so that $`\stackrel{~}{h}(z)`$ is the analytic continuation of $`\stackrel{~}{f}(z;C)`$ from the lower-half plane through the branch cut $`[e_0,e_1]`$ into the upper-half plane of a non-physical RS.
## B Asymptotic behaviour of the self-energy at large energies
Here we demonstrate that, for $`|z|\mathrm{}`$, $`\stackrel{~}{\mathrm{\Sigma }}(z)\mathrm{\Sigma }^{HF}`$, the Hartree-Fock self-energy,
$$\mathrm{\Sigma }^{HF}(𝐫,𝐫^{})\mathrm{\Sigma }^H(𝐫,𝐫^{})+\mathrm{\Sigma }^F(𝐫,𝐫^{})\frac{1}{\mathrm{}}v_H(𝐫;[n])\delta (𝐫𝐫^{})\frac{1}{2\mathrm{}}v(𝐫𝐫^{})\varrho (𝐫,𝐫^{}),$$
(B1)
where $`v_H(𝐫;[n]):=\mathrm{d}^dr^{}v(𝐫𝐫^{})n(𝐫^{})`$ stands for the Hartree potential, which is a functional of the electronic number density $`n`$ in the ground state, and $`\varrho `$ for the reduced single-particle density matrix. We note that $`\varrho `$ in the definition for $`\mathrm{\Sigma }^F`$ pertains to the fully interacting system and is distinct from the Slater-Fock reduced density matrix $`\varrho _0`$; contrary to the former, the latter is idempotent, i.e. $`\varrho _0\varrho _0\varrho _0`$, whereas $`\varrho \varrho \varrho `$. This difference may have some far-reaching consequences. For instance, in a uniform isotropic system of electrons interacting through the long-range Coulomb interaction function, for $`\mathrm{\Sigma }^F(k)`$ evaluated in terms of $`\varrho _0`$, which we denote by $`\mathrm{\Sigma }_𝗌^F(k)`$, one has (Ashcroft and Mermin 1981, p. 334) <sup>∥‡</sup><sup>∥‡</sup>∥‡ For this system, $`\mathrm{\Sigma }^H(k)`$ is divergent but cancels an equally-divergent contribution due to the field of the positively-charged uniform background.
$$\mathrm{\Sigma }_𝗌^F(k)=\frac{2e^2}{\pi \mathrm{}}k_F(k/k_F),(x):=\frac{1}{2}+\frac{1x^2}{4x}\mathrm{ln}\left|\frac{1+x}{1x}\right|;$$
(B2)
the first derivative with respect to $`k`$ of $`\mathrm{\Sigma }_𝗌^F(k)`$ is seen to be logarithmically divergent at $`k=k_F`$, an aspect which may be an artifact of $`\mathrm{\Sigma }_𝗌^F(k)`$ and may not be shared by $`\mathrm{\Sigma }^F(k)`$.
From the Dyson equation we have $`\mathrm{\Sigma }(\epsilon )=G_0^1(\epsilon )G^1(\epsilon )`$. Therefore to deduce the asymptotic expansion of $`\mathrm{\Sigma }(\epsilon )`$ for $`|\epsilon |\mathrm{}`$, we need to determine those for $`G_0^1(\epsilon )`$ and $`G^1(\epsilon )`$. The asymptotic series for $`G_0^1(\epsilon )`$ and $`G^1(\epsilon )`$ can be deduced from those pertaining to $`G_0(\epsilon )`$ and $`G(\epsilon )`$, respectively, through reliance on the following result from the theory of asymptotic analysis (see Copson 1965, pp. 8 and 9): Let $`\stackrel{~}{f}(z)`$ have the following asymptotic expansion for $`|z|\mathrm{}`$, where $`z`$ is a complex variable: $`\stackrel{~}{f}(z)f_0+f_1/z+f_2/z^2+\mathrm{}`$, with $`f_0`$, $`f_1`$, $`\mathrm{}`$ constants, independent of $`z`$. Then, provided that $`f_00`$, $`1/\stackrel{~}{f}(z)1/f_0+\overline{f}_1/z+\overline{f}_2/z^2+\mathrm{}`$ holds for $`|z|\mathrm{}`$, where $`\overline{f}_1=f_1/f_0^2`$, $`\overline{f}_2=(f_1^2f_0f_2)/f_0^3`$, etc. <sup>∥§</sup><sup>∥§</sup>∥§ When $`f_0=0`$ and $`f_10`$, then one should apply this result to $`\stackrel{~}{g}(z):=z\stackrel{~}{f}(z)`$; from the asymptotic series for $`1/\stackrel{~}{g}(z)`$ that for $`1/\stackrel{~}{f}(z)`$ is obtained through a multiplication by $`z`$.
From the Lehmann representation for $`G(\epsilon )`$ in Eq. (1) it follows that
$$G(\epsilon )\frac{G_\mathrm{}_1}{\epsilon }+\frac{G_\mathrm{}_2}{\epsilon ^2}+\mathrm{},\text{for}|\epsilon |\mathrm{},$$
(B3)
where
$`G_\mathrm{}_1(𝐫,𝐫^{})`$ $`=`$ $`\mathrm{}{\displaystyle \underset{s}{}}\mathrm{\Lambda }_s(𝐫)\mathrm{\Lambda }_s^{}(𝐫^{})\mathrm{}\delta (𝐫𝐫^{}),`$ (B4)
$`G_\mathrm{}_2(𝐫,𝐫^{})`$ $`=`$ $`\mathrm{}{\displaystyle \underset{s}{}}\epsilon _s\mathrm{\Lambda }_s(𝐫)\mathrm{\Lambda }_s^{}(𝐫^{})\mathrm{}\left(\mathrm{\Xi }_<(𝐫,𝐫^{})+\mathrm{\Xi }_>(𝐫,𝐫^{})\right);`$ (B5)
for $`\mathrm{\Xi }_<`$ and $`\mathrm{\Xi }_>`$ see Eqs. (B9) and (B11) below. A similar expression to Eq. (B3) holds for $`G_0(\epsilon )`$. Therefore from our above considerations with regard to the asymptotic series for $`1/\stackrel{~}{f}(z)`$ in terms of the coefficients of that for $`\stackrel{~}{f}(z)`$, it follows that $`\mathrm{\Sigma }(\epsilon )`$ possesses the following asymptotic series expansion
$$\mathrm{\Sigma }(\epsilon )\mathrm{\Sigma }_\mathrm{}_0+\frac{\mathrm{\Sigma }_\mathrm{}_1}{\epsilon }+\frac{\mathrm{\Sigma }_\mathrm{}_2}{\epsilon ^2}+\mathrm{},\text{for}|\epsilon |\mathrm{}.$$
(B6)
Here we are interested in the leading asymptotic term $`\mathrm{\Sigma }_\mathrm{}_0`$ which our above considerations indicate to be determined from $`G_{0;\mathrm{}_2}`$ and $`G_\mathrm{}_2`$ as follows
$$\mathrm{\Sigma }_\mathrm{}_0=\frac{1}{\mathrm{}^2}\left\{G_{0;\mathrm{}_2}+G_\mathrm{}_2\right\}.$$
(B7)
For the many-body Hamiltonian of the form
$$\widehat{H}=\mathrm{d}^dr\widehat{\psi }^{}(𝐫)\left[\frac{\mathrm{}^2}{2m_e}^2+u(𝐫)\right]\widehat{\psi }(𝐫)+\frac{1}{2}\mathrm{d}^dr\mathrm{d}^dr^{}\widehat{\psi }^{}(𝐫)\widehat{\psi }^{}(𝐫^{})v(𝐫𝐫^{})\widehat{\psi }(𝐫^{})\widehat{\psi }(𝐫),$$
(B8)
with $`u(𝐫)`$ the local external potential, making use of the definition for the Lehmann amplitudes and energies as given in Eq. (7) and the fact that $`E_{M,s}|\mathrm{\Psi }_{M,s}=\widehat{H}|\mathrm{\Psi }_{M,s}`$, one readily obtains
$`\mathrm{\Xi }_<(𝐫,𝐫^{})`$ $`:=`$ $`{\displaystyle \underset{s}{}}\theta (\mu \epsilon _s)\epsilon _s\mathrm{\Lambda }_s(𝐫)\mathrm{\Lambda }_s^{}(𝐫^{})`$ (B9)
$`=`$ $`\mathrm{\Psi }_{N,0}|\widehat{\psi }^{}(𝐫^{})[\widehat{H},\widehat{\psi }(𝐫)]_{}|\mathrm{\Psi }_{N,0},`$ (B10)
$`\mathrm{\Xi }_>(𝐫,𝐫^{})`$ $`:=`$ $`{\displaystyle \underset{s}{}}\theta (\epsilon _s\mu )\epsilon _s\mathrm{\Lambda }_s(𝐫)\mathrm{\Lambda }_s^{}(𝐫^{})`$ (B11)
$`=`$ $`\mathrm{\Psi }_{N,0}|[\widehat{H},\widehat{\psi }(𝐫)]_{}\widehat{\psi }^{}(𝐫^{})|\mathrm{\Psi }_{N,0}.`$ (B12)
By applying the canonical anti-commutation relations for the field operators in the Schrödinger picture one arrives at the following result
$$G_\mathrm{}_2(𝐫,𝐫^{})=\mathrm{}\left\{\left[\frac{\mathrm{}^2}{2m_e}^2+u(𝐫)+v_H(𝐫;[n])\right]\delta (𝐫𝐫^{})\frac{1}{2}v(𝐫𝐫^{})\varrho (𝐫,𝐫^{})\right\}.$$
(B13)
The corresponding expression for the non-interacting Green function follows from this expression by setting the coupling constant of the electron-electron interaction equal to zero,
$$G_{0;\mathrm{}_2}(𝐫,𝐫^{})=\mathrm{}\left[\frac{\mathrm{}^2}{2m_e}^2+u(𝐫)\right]\delta (𝐫𝐫^{}).$$
(B14)
From Eqs. (B7), (B13), (B14) and (B1) we obtain the following result
$$\mathrm{\Sigma }_\mathrm{}_0\mathrm{\Sigma }^{HF},$$
(B15)
which completes our demonstration. We note that Eqs. (B3) and (B6) hold equally for $`\stackrel{~}{G}(z)`$ and $`\stackrel{~}{\mathrm{\Sigma }}(z)`$, respectively, with merely $`\epsilon `$ replaced by $`z`$.
## C The Kramers-Kronig relations for the self-energy; applications to isotropic Fermi and marginal-Fermi liquids
The analyticity of a function of the complex variable $`z=x+iy`$ in a domain $`𝒟`$ of the $`z`$-plane implies that in $`𝒟`$ its real (imaginary) part is up to a constant uniquely determined by its imaginary (real) part (Morse and Feshbach 1953, pp. 356-358). For functions, such as $`\stackrel{~}{f}(z)`$, whose singularities are located along the $`x`$-axis and whose real and imaginary parts, $`𝗎(x,y)`$ and $`𝗏(x,y)`$, approach zero not slower than $`1/\varrho ^p`$, with $`p>0`$, as $`\varrho :=(x^2+y^2)^{1/2}\mathrm{}`$, the Kramers-Kronig relations establish the mentioned unique relationships between $`𝗎(x,0^\pm )`$ and $`𝗏(x,0^\pm )`$.
The self-energy $`\stackrel{~}{\mathrm{\Sigma }}(z)`$ is analytic over the entire $`z`$-plane with the exception of the real axis. <sup>∥¶</sup><sup>∥¶</sup>∥¶ This follows from the Dyson equation $`\stackrel{~}{\mathrm{\Sigma }}(z)=\stackrel{~}{G}_0^1(z)\stackrel{~}{G}^1(z)`$ and the fact that neither $`\stackrel{~}{G}_0(z)`$ nor $`\stackrel{~}{G}(z)`$ possesses zero eigenvalues for $`\mathrm{Im}(z)0`$ (see last paragraph in § II). Two aspects should be addressed before a Kramers-Kronig pair of relations can be set up for the self-energy. First, as the considerations in Appendix B have demonstrated, $`\stackrel{~}{\mathrm{\Sigma }}(z)\mathrm{\Sigma }^{HF}`$ for $`|z|\mathrm{}`$, so that we need first to introduce
$$\stackrel{~}{\mathrm{\Sigma }}_{}(z):=\stackrel{~}{\mathrm{\Sigma }}(z)\mathrm{\Sigma }^{HF},$$
(C1)
which, following Eq. (B6), approaches zero according to $`\mathrm{\Sigma }_\mathrm{}_1/z`$. Second, as $`\stackrel{~}{\mathrm{\Sigma }}(z)`$ is an operator, it cannot be assigned real and imaginary parts; moreover, as $`\stackrel{~}{\mathrm{\Sigma }}(z)`$ is not gauge invariant, real and imaginary parts of the matrix elements $`\{\alpha |\stackrel{~}{\mathrm{\Sigma }}(z)|\beta \}`$, for a given representation, are dependent upon the choice of the gauge. We need therefore generalise for operators the notions ‘real part’ and ‘imaginary part’ before the corresponding Kramers-Kronig relations can be meaningfully defined. We introduce
$$\stackrel{~}{\mathrm{\Sigma }}^{}(z):=\frac{1}{2}\left\{\stackrel{~}{\mathrm{\Sigma }}(z)+\stackrel{~}{\mathrm{\Sigma }}^{}(z)\right\},\stackrel{~}{\mathrm{\Sigma }}^{\prime \prime }(z):=\frac{1}{2i}\left\{\stackrel{~}{\mathrm{\Sigma }}(z)\stackrel{~}{\mathrm{\Sigma }}^{}(z)\right\}.$$
(C2)
We similarly introduce $`\stackrel{~}{\mathrm{\Sigma }}_{}^{}(z)`$ and $`\stackrel{~}{\mathrm{\Sigma }}_{}^{\prime \prime }(z)`$; $`\stackrel{~}{\mathrm{\Sigma }}^{}(z)`$ and $`i\stackrel{~}{\mathrm{\Sigma }}^{\prime \prime }(z)`$ ($`\stackrel{~}{\mathrm{\Sigma }}_{}^{}(z)`$ and $`i\stackrel{~}{\mathrm{\Sigma }}_{}^{\prime \prime }(z)`$) are Hermitian and anti-Hermitian components of $`\stackrel{~}{\mathrm{\Sigma }}(z)`$ ($`\stackrel{~}{\mathrm{\Sigma }}_{}(z)`$) and are unique. This uniqueness enables us to deduce from
$$\stackrel{~}{\mathrm{\Sigma }}_{}(\epsilon \pm i\eta )=\pm \frac{1}{\pi i}𝒫_{\mathrm{}}^{\mathrm{}}d\epsilon ^{}\frac{\stackrel{~}{\mathrm{\Sigma }}_{}(\epsilon ^{}\pm i\eta )}{\epsilon ^{}\epsilon },$$
(C3)
which follows from the application of the Cauchy theorem together with the above-indicated analytic and asymptotic properties of $`\stackrel{~}{\mathrm{\Sigma }}_{}(z)`$, the following pair of representation-free Kramers-Kronig relations
$`\mathrm{\Sigma }_{}^{}(\epsilon )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}𝒫{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\epsilon ^{}{\displaystyle \frac{\text{sgn}(\mu \epsilon ^{})\mathrm{\Sigma }_{}^{\prime \prime }(\epsilon ^{})}{\epsilon ^{}\epsilon }},`$ (C4)
$`\mathrm{\Sigma }_{}^{\prime \prime }(\epsilon )`$ $`=`$ $`+{\displaystyle \frac{1}{\pi }}𝒫{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\epsilon ^{}{\displaystyle \frac{\text{sgn}(\mu \epsilon )\mathrm{\Sigma }_{}^{}(\epsilon ^{})}{\epsilon ^{}\epsilon }}.`$ (C5)
In obtaining these relations we have made use of our convention in Eq. (9). We note that $`\mathrm{\Sigma }^{}(\epsilon )\mathrm{\Sigma }^{HF}+\mathrm{\Sigma }_{}^{}(\epsilon )`$ and $`\mathrm{\Sigma }^{\prime \prime }(\epsilon )\mathrm{\Sigma }_{}^{\prime \prime }(\epsilon )`$.
### 1 Fermi liquids
Now we employ Eq. (C4) and obtain from <sup>∥∥</sup><sup>∥∥</sup>∥∥ Note that $`\epsilon _F\mu _N`$, $`\alpha _k0`$ and that this asymptotic relation is valid for $`k`$ in a neighbourhood of $`k_F`$, say for $`0k{}_{}{}^{<}3k_F`$ — see footnote †††IV A. $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon )\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)^2`$, $`\epsilon {}_{<}{}^{>}\epsilon _{F}^{}`$, the associated $`\mathrm{\Sigma }^{}(k;\epsilon )`$; in this way we obtain Eq. (33) as well as an explicit expression for $`\beta _k`$ in terms of $`\alpha _k`$ and $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon )`$. To this end, we subdivide the interval of the $`\epsilon ^{}`$-integration in Eq. (C4) into three subintervals, $`(\mathrm{},\mathrm{\Delta }]`$, $`(\mathrm{\Delta },+\mathrm{\Delta }]`$ and $`(+\mathrm{\Delta },+\mathrm{})`$, where we assume $`\mathrm{\Delta }>|\delta \epsilon |`$ with $`\delta \epsilon :=\epsilon \epsilon _F`$. Upon change of variables and transforming integrals over the negative $`\epsilon ^{}`$-axis into those over the positive $`\epsilon ^{}`$-axis we obtain
$$\mathrm{\Sigma }_{}^{}(k;\epsilon )\frac{2\alpha _k\delta \epsilon }{\pi }𝒫_0^\mathrm{\Delta }d\epsilon ^{}\frac{\epsilon _{}^{}{}_{}{}^{2}}{\epsilon _{}^{}{}_{}{}^{2}\delta \epsilon ^2}+\frac{1}{\pi }_\mathrm{\Delta }^{\mathrm{}}d\epsilon ^{}\left\{\frac{\mathrm{\Sigma }^{\prime \prime }(k;\epsilon ^{}+\epsilon _F)}{\epsilon ^{}\delta \epsilon }+\frac{\mathrm{\Sigma }^{\prime \prime }(k;\epsilon ^{}+\epsilon _F)}{\epsilon ^{}+\delta \epsilon }\right\}.$$
(C6)
The first integral on the right-hand side of Eq. (C6) is standard (Gradshteyn and Ryzhik 1965, p. 59) and one has $`𝒫_0^\mathrm{\Delta }d\epsilon ^{}\epsilon _{}^{}{}_{}{}^{2}/(\epsilon _{}^{}{}_{}{}^{2}\delta \epsilon ^2)=\mathrm{\Delta }(\delta \epsilon /2)\mathrm{ln}\left|(\mathrm{\Delta }+\delta \epsilon )/(\mathrm{\Delta }\delta \epsilon )\right|`$. Since, in the second integral on the right-hand side of Eq. (C6), $`\epsilon ^{}\mathrm{\Delta }`$ and since $`\mathrm{\Delta }>|\delta \epsilon |`$, we replace $`1/(\epsilon ^{}\delta \epsilon )`$ and $`1/(\epsilon ^{}+\delta \epsilon )`$ by their respective geometric series expansions, in powers of $`\delta \epsilon /\epsilon ^{}`$, and thus obtain a powers series in $`\delta \epsilon `$ for the second integral. Collecting terms of the zeroth and the first order in $`\delta \epsilon `$, while making use of Eq. (33), we obtain
$$\mathrm{\Sigma }(k;\epsilon _F)\mathrm{\Sigma }^{}(k;\epsilon _F)\mathrm{\Sigma }^{HF}(k)+S_0^{(\mathrm{\Delta })}(k),\beta _kS_1^{(\mathrm{\Delta })}(k)\frac{2\mathrm{\Delta }\alpha _k}{\pi },$$
(C7)
where
$$S_m^{(\mathrm{\Delta })}(k):=\frac{1}{\pi }_\mathrm{\Delta }^{\mathrm{}}\frac{\mathrm{d}\epsilon ^{}}{\epsilon _{}^{}{}_{}{}^{m+1}}\left\{\mathrm{\Sigma }^{\prime \prime }(k;\epsilon ^{}+\epsilon _F)+(1)^m\mathrm{\Sigma }^{\prime \prime }(k;\epsilon ^{}+\epsilon _F)\right\}.$$
(C8)
In Eq. (C7), ‘$``$’ signifies the approximate nature of the results in so far as $`\mathrm{\Delta }`$ is finite and, moreover, in evaluating the $`\epsilon ^{}`$-integral over $`(\mathrm{\Delta },+\mathrm{\Delta }]`$ we have employed a truncated asymptotic expansion for $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon )`$. As should be evident, our choice for a finite $`\mathrm{\Delta }`$ has it root in the requirement $`\mathrm{\Delta }>|\delta \epsilon |`$. Since Eqs. (C7) and (C8) are valid independently of the value of $`\delta \epsilon `$, it is possible to let $`\mathrm{\Delta }`$ in these equations approach zero; because $`\mathrm{\Sigma }^{\prime \prime }(k;\pm \epsilon ^{}+\epsilon _F)^{}\epsilon _{}^{}{}_{}{}^{2}`$ for $`\epsilon ^{}0`$, $`S_m^{(\mathrm{\Delta }0)}(k)`$, which involves $`1/\epsilon _{}^{}{}_{}{}^{m+1}`$, exists for $`m=0,1`$, and is equal to $`S_m^{(0)}(k)`$; on the other hand, for $`m2`$, $`S_m^{(\mathrm{\Delta })}(k)`$ does not have a limit for $`\mathrm{\Delta }0`$. Thus we can write
$$\mathrm{\Sigma }(k;\epsilon _F)\mathrm{\Sigma }^{}(k;\epsilon _F)=\mathrm{\Sigma }^{HF}(k)+S_0^{(0)}(k),\beta _k=S_1^{(0)}(k).$$
(C9)
We point out that, since, for $`\epsilon ^{}0`$, $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon ^{}+\epsilon _F)0`$ and $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon ^{}+\epsilon _F)0`$, <sup>∥\**</sup><sup>∥\**</sup>∥\** Violation of these inequalities implies breakdown of causality and instability of the ground state of the system (see Eq. (30) in comparison with Eq. (1)). it follows that $`S_m^{(\mathrm{\Delta })}(k)0`$ for all odd values of $`m`$. Specifically, following Eq. (C9), we have $`\beta _k0`$ (concerning the range of $`k`$ over which this result can be relied upon, see footnotes †††IV A and ∥∥C 1). As we have indicated in § V (see text following Eq. (33)), for $`kk_F`$ this inequality guarantees $`Z_{k_F}1`$.
It is interesting to note that $`S_m^{(\mathrm{\Delta })}(k)`$, $`m=0,1`$, can be written in the following appealing form
$`S_0^{(\mathrm{\Delta })}(k)`$ $`:=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}\left[{\displaystyle _{\mu +\mathrm{\Delta }}^i\mathrm{}}{\displaystyle \frac{\mathrm{d}z}{(z\mu )}}\left\{\stackrel{~}{\mathrm{\Sigma }}(k;z)+\stackrel{~}{\mathrm{\Sigma }}(k;2\mu z)2\mathrm{\Sigma }^{HF}(k)\right\}\right],`$ (C10)
$`S_1^{(\mathrm{\Delta })}(k)`$ $`:=`$ $`{\displaystyle \frac{1}{\pi }}\mathrm{Im}\left[{\displaystyle _{\mu +\mathrm{\Delta }}^i\mathrm{}}{\displaystyle \frac{\mathrm{d}z}{(z\mu )^2}}\left\{\stackrel{~}{\mathrm{\Sigma }}(k;z)\stackrel{~}{\mathrm{\Sigma }}(k;2\mu z)\right\}\right],`$ (C11)
where the term $`2\mathrm{\Sigma }^{HF}(k)`$, which is real-valued, has been introduced in order to render the corresponding integral existent (see Eq. (B15)). In Eqs. (C10) and (C11) we have further restored $`\mu `$ which originates from Eq. (C4) and which is not exactly equal to $`\epsilon _F`$ (see § IV.A). We point out that in these expressions, contrary to that in Eq. (C8), $`\mathrm{\Delta }`$ cannot be set equal to zero, since for $`\mathrm{\Delta }=0`$ the integrals enclosed by square brackets (logarithmically) diverge as a consequence of the integrands being proportional to $`1/(z\mu )`$ for $`z\mu `$; since, however, the singular contributions are real-valued, $`\mathrm{\Delta }`$ can be made as small as desired, provided that it is kept non-zero; thus $`S_m^{(0^+)}(k)`$, $`m=0,1`$, is well-defined.
The relevance of Eqs. (C10) and (C11) rests in the following: The fact that $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ is analytic over the entire $`z`$-plane (excluding the real axis outside $`(\mu _N,\mu _{N+1})`$), implies that if $`^{\mathrm{}}S_m^{(0^+)}(k)/k^{\mathrm{}}`$, $`m=0,1`$ (including $`\mathrm{}=0`$), is singular at some $`k`$, this singularity must be shared by the integrand of $`^{\mathrm{}}S_m^{(0^+)}(k)/k^{\mathrm{}}`$, $`m=0,1`$, for all $`z`$, $`\mathrm{Im}(z)0`$. In other words, Eqs. (C10) and (C11) make explicit that possible singularities in $`^{\mathrm{}}S_m^{(0^+)}(k)/k^{\mathrm{}}`$, $`m=0,1`$, are not specific to $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ for some $`z=z_0`$, but to $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ for all $`z`$, $`\mathrm{Im}(z)0`$. <sup>∥††</sup><sup>∥††</sup>∥†† Consider $`g(k):=_𝒞dz\stackrel{~}{f}(k;z)`$, where the contour of integration $`𝒞`$ lies inside an open domain $`𝒟`$ of the complex $`z`$-plane where $`\stackrel{~}{f}(k;z)`$ is analytic. Suppose that the $`\mathrm{}`$th-order derivative (including $`\mathrm{}=0`$) with respect to $`k`$ of $`g(k)`$ is divergent at $`k=k_0`$. This divergence cannot be due to some isolated singularity (singularities) of $`^{\mathrm{}}\stackrel{~}{f}(k;z)/k^{\mathrm{}}|_{k=k_0}`$ along $`𝒞`$, for analyticity of $`\stackrel{~}{f}(k;z)`$ enables one to deform $`𝒞`$ inside $`𝒟`$, thus avoiding the mentioned singularity (singularities), without changing the value of the integral (Cauchy’s theorem); exceptions to this concern the cases where $`^{\mathrm{}}\stackrel{~}{f}(k;z)/k^{\mathrm{}}|_{k=k_0}`$ is divergent either at the end-points of $`𝒞`$, which cannot be dislocated, or at two neighbouring points which ‘pinch’ $`𝒞`$. Note that in Eqs. (C10) and (C11) the contours of integration are arbitrary, as long as they connect $`\mu +\mathrm{\Delta }`$ with the point of infinity (the ‘end-points’) and are located on the upper-half of the $`z`$-plane (signifying $`𝒟`$). Therefore, excluding end-point and ‘pinch’ singularities, divergence of $`^{\mathrm{}}g(k)/k^{\mathrm{}}`$ at $`k=k_0`$ implies divergence of $`^{\mathrm{}}\stackrel{~}{f}(k;z)/k^{\mathrm{}}`$ at $`k=k_0`$ for all $`z`$ inside $`𝒟`$. In Eqs. (C10) and (C11) one of the end-points, namely the point of infinity, is harmless, as the integrands in both expressions are vanishing at this point; by choosing $`\mu =(\mu _N+\mu _{N+1})/2`$ and $`\mathrm{\Delta }=(\mu _{N+1}\mu _N)/2`$, owing to the assumption of continuous-differentiability with respect to $`k`$ of $`\mathrm{\Sigma }(k;\mu _N)`$ and $`\mathrm{\Sigma }(k;\mu _{N+1})`$ in a neighbourhood of $`k=k_F`$, a possible divergence of $`S_m^{(0^+)}(k)`$ or $`S_m^{(0^+)}(k)/k`$, $`m=0,1`$, at $`k=k_F`$ cannot be ascribed to an end-point singularity.
In order to demonstrate the utility of our above considerations, let us assume $`\mathrm{\Sigma }^{HF}(k)`$, similar to $`\mathrm{\Sigma }_𝗌^{HF}(k)`$ (see Eq. (B2)), has a divergent first-order derivative with respect to $`k`$ at $`k=k_F`$ (see Appendix B). For the system under consideration to be a Fermi liquid, $`\mathrm{\Sigma }(k;\epsilon _F)`$ must, by definition, be a continuously differentiable function of $`k`$ at $`k=k_F`$ (see §§ I and IV.C). From Eq. (C9) it follows that the assumed singularity due to $`\mathrm{\Sigma }^{HF}(k)/k`$ at $`k=k_F`$ must therefore be compensated by a counter contribution arising from $`S_0^{(0^+)}(k)`$. With reference to our above discussions, in such an event $`\{\stackrel{~}{\mathrm{\Sigma }}(k;z)+\stackrel{~}{\mathrm{\Sigma }}(k;2\mu z)2\mathrm{\Sigma }^{HF}(k)\}/k`$ must be singular at $`k=k_F`$ for all $`z`$, $`\mathrm{Im}(z)0`$ (see specifically footnote ∥††C 1). For the special choice of $`z=\mu +\mathrm{\Delta }`$, with $`\mu =(\mu _N+\mu _{N+1})/2`$ and $`\mathrm{\Delta }=(\mu _{N+1}\mu _N)/2`$, $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ and $`\stackrel{~}{\mathrm{\Sigma }}(k;2\mu z)`$ are continuously differentiable functions of $`k`$ in a neighbourhood of $`k=k_F`$ (by the Fermi-liquid assumption); nonetheless, in accordance with our expectation $`\{\stackrel{~}{\mathrm{\Sigma }}(k;z)+\stackrel{~}{\mathrm{\Sigma }}(k;2\mu z)2\mathrm{\Sigma }^{HF}(k)\}/k`$ remains divergent at $`k=k_F`$ owing to the assumed divergence of $`\mathrm{\Sigma }^{HF}(k)/k`$ at $`k=k_F`$. On the other hand, for an arbitrary $`z`$ the continuous differentiability with respect to $`k`$ of $`\stackrel{~}{\mathrm{\Sigma }}(k;z)`$ and $`\stackrel{~}{\mathrm{\Sigma }}(z;2\mu z)`$ is not guaranteed, so that divergence of $`\{\stackrel{~}{\mathrm{\Sigma }}(k;z)+\stackrel{~}{\mathrm{\Sigma }}(k;2\mu z)2\mathrm{\Sigma }^{HF}(k)\}/k`$ at $`k=k_F`$ can be in part due to that of $`\{\stackrel{~}{\mathrm{\Sigma }}(k;z)+\stackrel{~}{\mathrm{\Sigma }}(k;2\mu z)\}/k`$ at $`k=k_F`$. In view of the asymptotic relation in Eq. (34) (see footnote §††V), it follows that $`\beta _k`$ and $`\alpha _k`$ may not be necessarily continuously differentiable at $`k=k_F`$. The divergence of $`\beta _k/k`$ as $`kk_F^{}`$ gives rise to a divergent $`𝗇_s(k)/k`$, as well as a divergent $`𝗇(k)/k`$, at $`k=k_F^{}`$ (for a discussion of this case see § VI). Since for Fermi liquids (here as characterised by $`Z_{k_F}0`$), a divergent $`𝗇(k)/k`$, at $`k=k_F^{}`$, has proved possible (Belyakov 1961, Sartor and Mahaux 1980), we observe that in general even for Fermi liquids $`\mathrm{\Sigma }(k;\epsilon )`$ is not a continuously differentiable function of $`k`$ (in a neighbourhood of $`k=k_F`$) for $`\epsilon \epsilon _F^{}`$; at $`\epsilon =\epsilon _F^{}`$, $`\beta _k`$ and $`\alpha _k`$ are excluded from contributing to $`\mathrm{\Sigma }(k;\epsilon )`$ (see Eq. (34)).
### 2 Marginal Fermi liquids
Having considered the case of isotropic Fermi liquids in considerable detail, below we briefly deal with the case of the isotropic marginal Fermi liquids.
For marginal Fermi liquids we have $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon )\mathrm{Im}\mathrm{\Sigma }(k;\epsilon )\alpha _k(\epsilon \epsilon _F)`$, as $`\epsilon \epsilon _F`$, with $`\alpha _k0`$ (see § V). For the $`\mathrm{\Sigma }_{}^{}(k;\epsilon )`$ corresponding to this $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon )`$ one obtains an expression which except for the first term is identical with that presented in Eq. (C6). For this first term, which we denote by $`\mathrm{\Sigma }_{;1}^{}(k;\epsilon )`$, we have (as in the Fermi-liquid case, below $`\mathrm{\Delta }>|\delta \epsilon |`$)
$$\mathrm{\Sigma }_{;1}^{}(k;\epsilon ):=\frac{2\alpha _k\delta \epsilon }{\pi }𝒫_0^\mathrm{\Delta }d\epsilon ^{}\frac{\epsilon ^{}}{\epsilon _{}^{}{}_{}{}^{2}\delta \epsilon ^2}\frac{2\alpha _k\delta \epsilon }{\pi }\left\{\mathrm{ln}(\mathrm{\Delta })\mathrm{ln}|\delta \epsilon |\frac{1}{2\mathrm{\Delta }^2}\delta \epsilon ^2\right\};$$
(C12)
(see Gradshteyn and Ryzhik 1965, p. 59). One observes that here, contrary to the Fermi-liquid case, the limit $`\mathrm{\Delta }0`$ cannot be taken. This is related to the fact that here $`\mathrm{\Sigma }^{\prime \prime }(k;\pm \epsilon ^{}+\epsilon _F)^{}\epsilon ^{}`$ for $`\epsilon ^{}0`$, so that only $`S_0^{(\mathrm{\Delta })}(k)`$ has a limit for $`\mathrm{\Delta }0`$. Therefore, in the present case, only the first expression in Eq. (C9) is meaningful. Further, from Eq. (C12) we observe that the leading asymptotic contribution to $`\mathrm{\Sigma }_{}^{}(k;\epsilon )`$, as $`\epsilon \epsilon _F`$, involves $`\mathrm{ln}|\delta \epsilon |`$ which is singular for $`\delta \epsilon =0`$. Combining the above results, we have $`\mathrm{\Sigma }(k;\epsilon )\mathrm{\Sigma }(k;\epsilon _F)+\beta _k(\epsilon \epsilon _F)\mathrm{ln}|\epsilon \epsilon _F|+[\gamma _k+i\alpha _k](\epsilon \epsilon _F)`$, as $`\epsilon \epsilon _F`$, where
$`\beta _k`$ $``$ $`{\displaystyle \frac{2\alpha _k}{\pi }},`$ (C13)
$`\gamma _k`$ $``$ $`{\displaystyle \frac{2\alpha _k}{\pi }}\mathrm{ln}(\mathrm{\Delta })+S_1^{(\mathrm{\Delta })}(k).`$ (C14)
The first term on the right-hand side of Eq. (C14) cancels the logarithmically divergent contribution arising from the second term as $`\mathrm{\Delta }0`$. It can be readily verified that to the order in which $`\alpha _k(\epsilon \epsilon _F)`$ is an exact representation of $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon )`$ for $`\epsilon \epsilon _F`$, $`\gamma _k/\mathrm{\Delta }=0`$. This invariance property correctly reflects the fact that $`\gamma _k`$ must not depend on the cut-off energy $`\mathrm{\Delta }`$.
## D Single-particle Green function and Self-energy of the Luttinger model; some asymptotic expressions
In this Appendix we consider the single-particle Green function and the self-energy pertaining to the Luttinger model (Luttinger 1963, Mattis and Lieb 1965) for spin-less fermions in some detail. <sup>∥‡‡</sup><sup>∥‡‡</sup>∥‡‡ We note that an explicit expression for the single-particle Green function for the Luttinger model with $`\delta `$-function interaction in real space (this model is equivalent with the $`1+1`$ Thirring (1958) model, dealt with by Johnson (1961)) has been evaluated by Theumann (1967). The corresponding expression in Eq. (2.19)<sub>T</sub> (here subscript $`T`$ denotes an equation from Theumann’s work) turns out to be identically vanishing for the case of non-interacting Luttinger model. To verify this, note that the coupling constant $`\lambda `$ of interaction as presented in Eq. (1.4)<sub>T</sub>, determines $`\mathrm{\Omega }`$ in Eq. (2.3)<sub>T</sub> which in turn fixes $`\beta `$ in Eq. (2.14)<sub>T</sub>. It is seen that for non-interacting particles, corresponding to $`\lambda =0`$ and thus $`\beta =1`$, $`G_1(k,\omega )`$ is identically vanishing, owing to $`\mathrm{sin}\pi \beta `$ in Eq. (2.19)<sub>T</sub>. We have not attempted to identify the origin of this shortcoming of $`G_1(k,\omega )`$ in Eq. (2.19)<sub>T</sub>. From the expression for the former function we can explicitly test the validity of one of our findings in Appendix C, namely that when $`\stackrel{~}{\mathrm{\Sigma }}(k;z)/k`$ is divergent at $`k=k_0`$ and $`z=z_0`$, it is divergent at $`k=k_0`$ for all $`z`$. Further, we consider the asymptotic behaviour of both $`G(k;\epsilon )`$ and $`\mathrm{\Sigma }(k;\epsilon )`$ for $`\epsilon `$ in the close vicinity of $`\epsilon _F`$. In doing so we separately deal with the case where $`k=k_F`$ and where $`kk_F`$. In order to remain close to the available literature concerning the Luttinger model, in this Appendix we adopt the commonly-used notations and units and therefore in these deviate from the other parts of the present work. We specify the new notations as we proceed through this Appendix. In this Appendix $`\mathrm{}=1`$.
### 1 The single-particle Green function and its derivative with respect to momentum
We first construct the single-particle Green function $`G`$ from the spectral function $`\rho `$ corresponding to the retarded part $`G^R`$ of the Green function. For this spectral function we have (Voit 1993b, 1993a)
$`\rho _r(q,\omega )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{2v_0\mathrm{\Gamma }^2(\gamma _0)}}\mathrm{\Theta }(\omega +rv_0q)\mathrm{\Theta }(\omega rv_0q)\gamma (\gamma _0,{\displaystyle \frac{\mathrm{\Lambda }}{2v_0}}(\omega +v_0rq))`$ (D3)
$`\times \left({\displaystyle \frac{\mathrm{\Lambda }}{2v_0}}(\omega v_0rq)\right)^{\gamma _01}\mathrm{exp}\left({\displaystyle \frac{\mathrm{\Lambda }}{2v_0}}(\omega v_0rq)\right)`$
$`+\left(\omega \omega ,qq\right),\text{for}\gamma _00,`$
$`\rho _r(q,\omega )`$ $``$ $`\rho _{0;r}(q,\omega ):=\delta (\omega v_0rq),\text{for}\gamma _0=0.`$ (D4)
Here $`r=`$ specifies the left and right branches respectively of the single-particle spectrum in the Luttinger model, $`\omega \epsilon \mu `$ and $`kq+rk_F`$. It is important to note that contrary to the 2D and 3D cases, here $`k`$ and $`q`$ take on both positive and negative values. Moreover $`q`$ is measured with respect to <sup>\***</sup><sup>\***</sup>\*** Thus a better notation for $`q`$ would be $`q_r`$ in order to emphasise this $`r`$-dependence of the origin. $`rk_F`$; $`k`$, on the other hand, is measured with respect to the origin. Further, $`\mathrm{\Lambda }`$ stands for the (finite) cut-off on the range of the interaction in the momentum space — the precise value of $`\mathrm{\Lambda }`$ is not of relevance here (Sólyom 1979; see §§ 9 and 2 herein); $`\mathrm{\Gamma }(z)`$ stands for the Gamma function (Abramowitz and Stegun 1972, p. 255), $`\gamma _0\alpha /2`$, with $`\alpha `$ the ‘anomalous dimension’ (see § IV.C) and $`\gamma (a,z)`$ denotes the incomplete Gamma function (Abramowitz and Stegun 1972, p. 260). Unless we explicitly indicate otherwise, below $`0<\alpha <1`$ and thus $`0<\gamma _0<1/2`$. With
$$\rho _r(q,\omega ):=\frac{1}{\pi }\mathrm{Im}G_r^R(q,\omega ),$$
(D5)
where the $`G_r^R`$ denotes the retarded Green function, making use of the fact that $`\mathrm{Im}G_r(q,\omega )=\mathrm{sgn}(\omega )\mathrm{Im}G_r^R(q,\omega )`$, we have
$$G_r(q,\omega )=_{\mathrm{}}^{\mathrm{}}d\omega ^{}\frac{\rho _r(q,\omega ^{})}{\omega \omega ^{}+i\eta \mathrm{sgn}(\omega ^{})},\eta 0.$$
(D6)
From Eq. (30), introducing $`\omega _{r;q}^0:=rv_0q`$, we obtain (compare with Eq. (37)) <sup>\**†</sup><sup>\**†</sup>\**† It is readily verified that $`\stackrel{~}{G}_r(q,z)_{\mathrm{}}^{\mathrm{}}d\omega ^{}\rho _r(q,\omega ^{})/(z\omega ^{})`$, for $`\mathrm{Im}(z)0`$; in view of Eq. (9), $`G_r(q,\omega )=lim_{\eta 0}\stackrel{~}{G}_r(q,\omega \pm i\eta )`$, for $`\omega {}_{<}{}^{>}0`$, which can be easily shown to coincide with that presented in Eq. (D6). We point out that, with $`\rho _{0;r}(q,\omega )=\delta (\omega \omega _{r;q}^0)`$ (see Eq. (D4)), one obtains $`G_{0;r}(q,\omega )=1/[\omega \omega _{r;q}^0+i\eta \mathrm{sgn}(\omega _{r;q}^0)]`$ and consequently $`\stackrel{~}{G}_{0;r}(q,z)=1/[z\omega _{r;q}^0]`$, for $`\mathrm{Im}(z)0`$.
$$\frac{\stackrel{~}{G}_r(q,z)}{q}=\stackrel{~}{G}_r^2(q,z)\left\{\frac{\omega _{r;q}^0}{q}+\frac{\stackrel{~}{\mathrm{\Sigma }}_r(q,z)}{q}\right\}.$$
(D7)
Since $`\stackrel{~}{G}_r(q,z)`$ is an entire function of $`z`$ for all $`z`$, $`\mathrm{Im}(z)0`$ (see § II) and, since $`\omega _{r;q}^0/qrv_0`$ is finite, any possible divergence of $`\stackrel{~}{G}_r(q,z)/q`$, as function of $`q`$ and $`z`$ ($`\mathrm{Im}(z)0`$), must be due to a divergence of $`\stackrel{~}{\mathrm{\Sigma }}_r(q,z)/q`$ on the right-hand side of Eq. (D7). In order to be able to evaluate $`\stackrel{~}{G}_r(q,z)/q`$, which according to Eq. (D6) involves $`\rho _r(q,\omega )/q`$, some changes in the expression on the right-hand side of Eq. (D3) are necessary. To appreciate this, note that $`\mathrm{\Theta }(\omega rv_0q)/q=rv_0\delta (\omega rv_0q)`$, from which, on account of the term $`(\omega rv_0q)^{\gamma _01}`$, one immediately observes that independent of the value of $`q`$, $`\rho _r(q,\omega )/q=\mathrm{}`$. This would imply a peculiar situation where $`\rho _r(q,\omega )`$, which is continuous almost everywhere, would be nowhere on the $`q`$-axis differentiable. As we shall see, this is wholly attributable to the fact that the expressions on the right-hand side of Eq. (D3) which are multiplied by distributions $`\mathrm{\Theta }(\omega rv_0q)`$ and $`\mathrm{\Theta }(\omega \pm rv_0q)`$ do not qualify as test functions (Gelfand and Shilov 1964). For simplicity, in the following we shall explicitly deal with the case corresponding to $`r=+`$. We write
$`\rho _+(q,\omega )`$ $``$ $`\mathrm{\Theta }(\omega +v_0q)\mathrm{\Theta }(\omega v_0q)\varphi (q,\omega )+(\omega \omega ,qq),`$ (D8)
$`\varphi (q,\omega )`$ $`:=`$ $`𝒜_1\varphi _1(\omega +v_0q)\varphi _2(\omega v_0q),`$ (D9)
$`𝒜_\sigma `$ $`:=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }^2(\gamma _0)}}\left({\displaystyle \frac{\mathrm{\Lambda }}{2v_0}}\right)^\sigma ,`$ (D10)
$`\varphi _1(x)`$ $`:=`$ $`\gamma (\gamma _0,{\displaystyle \frac{\mathrm{\Lambda }}{2v_0}}x),`$ (D11)
$`\varphi _2(x)`$ $`:=`$ $`\left({\displaystyle \frac{\mathrm{\Lambda }}{2v_0}}x\right)^{\gamma _01}\mathrm{exp}\left({\displaystyle \frac{\mathrm{\Lambda }}{2v_0}}x\right).`$ (D12)
Since $`0<\gamma _0<1/2`$, it is evident that the problem at hand (i.e. $`\rho _+(q,\omega )/q=\mathrm{}`$) has its origin in the singularity of $`\varphi _2(x)`$ at $`x=0`$. Through application of the prescription as specified in Fig. 3 (see caption to this Figure), we define $`\stackrel{~}{G}_{+,\nu }(q,z)`$ and consider $`\stackrel{~}{G}_+(q,z)`$ as being obtained through $`lim_{\nu 0}\stackrel{~}{G}_{+,\nu }(q,z)`$. We have
$`\stackrel{~}{G}_{+,\nu }(q,z)`$ $``$ $`{\displaystyle _{\mathrm{}}^{v_0q}}d\omega ^{}{\displaystyle \frac{\varphi (q,\omega ^{})}{z\omega ^{}}}+{\displaystyle _{v_0q+\nu }^{\mathrm{}}}d\omega ^{}{\displaystyle \frac{\varphi (q,\omega ^{})}{z\omega ^{}}}`$ (D14)
$`+{\displaystyle _{v_0q\nu }^{v_0q}}d\omega ^{}{\displaystyle \frac{\varphi (q,v_0q+\nu )\left[1+(\omega ^{}v_0q)/\nu \right]}{z\omega ^{}}}+{\displaystyle _{v_0q}^{v_0q+\nu }}d\omega ^{}{\displaystyle \frac{\varphi (q,v_0q+\nu )}{z\omega ^{}}}.`$
It is from this expression that the derivative with respect to $`q`$ has to be taken; we then have $`\stackrel{~}{G}_+(q,\omega )/q:=lim_{\nu 0}\stackrel{~}{G}_{+,\nu }(q,\omega )/q`$. After some algebra, one obtains the following final result
$`{\displaystyle \frac{\stackrel{~}{G}_+(q,z)}{q}}=`$ $``$ $`v_0{\displaystyle \frac{\varphi (q,v_0q)}{z+v_0q}}+{\displaystyle _{\mathrm{}}^{v_0q}}d\omega ^{}{\displaystyle \frac{\varphi (q,\omega ^{})/q}{z\omega ^{}}}`$ (D15)
$`+`$ $`𝒜_1v_0{\displaystyle _{v_0q}^{\mathrm{}}}d\omega ^{}{\displaystyle \frac{\{\varphi _1(\omega ^{}+v_0q)/\omega ^{}\}\varphi _2(\omega ^{}v_0q)}{z\omega ^{}}}`$ (D16)
$`+`$ $`𝒜_1v_0{\displaystyle _{v_0q}^{\mathrm{}}}d\omega ^{}\varphi _2(\omega ^{}v_0q){\displaystyle \frac{}{\omega ^{}}}\left\{{\displaystyle \frac{\varphi _1(\omega ^{}+v_0q)}{z\omega ^{}}}\right\}.`$ (D17)
In arriving at this expression, we have made use of the properties $`\varphi _{1,2}(\omega ^{}\pm v_0q)/q=\pm v_0\varphi _{1,2}(\omega ^{}\pm v_0q)/\omega ^{}`$. It is evident that the integrands of the last two integrals on the right-hand side of Eq. (D15) are singular (owing to $`\varphi _2(\omega ^{}v_0q)`$) at the lower limits of the integration boundaries; however since $`0<\gamma _0<1/2`$, these singularities are integrable and therefore $`\stackrel{~}{G}_+(q,z)/q`$ as presented in Eq. (D15) is a well-defined function of $`q`$ and $`z`$.
We are now in a position to draw the important conclusion that since $`\varphi (q,v_0q)=𝒜_1\varphi _1(0)\varphi _2(2v_0q)q^{\gamma _01}`$, $`\stackrel{~}{G}_+(q,z)/q`$ diverges as $`q0`$, for all $`z`$. This is in full conformity with our statement in Appendix C that, when $`\stackrel{~}{G}_+(q,z)/q`$ is unbounded at a particular value of $`q`$, it is unbounded at that $`q`$ for all $`z`$.
### 2 Self-energy and its asymptotic forms close to the Fermi points and the Fermi energy
Now we proceed with evaluating the asymptotic behaviour of $`\mathrm{\Sigma }_r(q,\omega )`$ for $`\omega 0`$. We consider two cases, corresponding to $`q=0`$ and $`q0`$. We demonstrate that, as $`\gamma _0`$ approaches zero, the self-energy of the Luttinger (1963) model becomes vanishingly small in a manner that is specific to non-interacting Fermi systems.
Since we are interested in cases where $`0<\gamma _0<1/2`$, for the determination of the leading asymptotic contribution to $`G_r(q,\omega )`$, and thus to $`\mathrm{\Sigma }_r(q,\omega )`$, for $`q0`$ (i.e. $`kk_F`$) and $`\omega 0`$ (i.e. $`\epsilon \epsilon _F`$) we need to merely take account of the $`\omega ^{}`$-integration in Eq. (D6) restricted to the interval $`[\mathrm{\Delta },\mathrm{\Delta }]`$, where $`\mathrm{\Delta }`$ is some finite constant satisfying $`|\omega |<\mathrm{\Delta }`$ (see Appendix C). Thus we define (cf. Eq. (D6))
$$\stackrel{~}{G}_r^{(\mathrm{\Delta })}(q,z):=_\mathrm{\Delta }^\mathrm{\Delta }d\omega ^{}\frac{\rho _r(q,\omega ^{})}{z\omega ^{}}.$$
(D18)
Below we shall deal with $`G_r^{(\mathrm{\Delta })}(q,\omega )`$.
#### a The $`q=0`$ case
From the expression in Eq. (D3) it can readily be deduced that (see Eqs. (D8)-(D12))
$$\rho _r(q=0,\omega )𝒜_{\gamma _0+1}|\omega |^{2\gamma _01},\text{for}\omega 0.$$
(D19)
Making use of the general result ($`|\omega |<\mathrm{\Delta }`$)
$`𝒫{\displaystyle _0^\mathrm{\Delta }}d\omega ^{}{\displaystyle \frac{\omega _{}^{}{}_{}{}^{\sigma }}{\omega _{}^{}{}_{}{}^{2}\omega ^2}}`$ $`=`$ $`{\displaystyle \frac{\pi }{2}}\mathrm{tan}\left({\displaystyle \frac{\pi \sigma }{2}}\right)|\omega |^{\sigma 1}+{\displaystyle \frac{\mathrm{\Delta }^{\sigma 1}}{\sigma 1}}{}_{2}{}^{}F_{1}^{}(1,{\displaystyle \frac{1\sigma }{2}};{\displaystyle \frac{3\sigma }{2}};{\displaystyle \frac{\omega ^2}{\mathrm{\Delta }^2}})`$ (D20)
$`=`$ $`{\displaystyle \frac{\pi }{2}}\mathrm{tan}\left({\displaystyle \frac{\pi \sigma }{2}}\right)\omega ^{\sigma 1}+{\displaystyle \frac{\mathrm{\Delta }^{\sigma 1}}{\sigma 1}}+{\displaystyle \frac{\mathrm{\Delta }^{\sigma 3}}{\sigma 3}}\omega ^2+𝒪(\omega ^4),`$ (D22)
$`\sigma >1,\sigma 1,3,\mathrm{},`$
where $`{}_{2}{}^{}F_{1}^{}(a,b;c;z)`$ stands for the Gauss Hypergeometric function (Abramowitz and Stegun 1972, p. 556), for $`0<\gamma _0<1`$ and $`|\omega |<\mathrm{\Delta }`$ we obtain
$$G_r^{(\mathrm{\Delta })}(q=0,\omega )\pi 𝒜_{\gamma _0+1}\left\{\mathrm{tan}\left(\frac{\pi (2\gamma _01)}{2}\right)+i\right\}\mathrm{sgn}(\omega )|\omega |^{2\gamma _01}.$$
(D23)
Evidently, for $`0<\gamma _0<1/2`$, the right-hand side of this expression diverges as $`\omega 0`$. As a consequence, for $`0<\gamma _0<1/2`$, the possibly non-vanishing constant $`\{G_r(q=0,\omega )G_r^{(\mathrm{\Delta })}(q=0,\omega )\}`$ is asymptotically irrelevant for $`\omega 0`$ and we can write
$$G_r(q=0,\omega )\pi 𝒜_{\gamma _0+1}\left\{\mathrm{cot}(\pi \gamma _0)i\right\}\mathrm{sgn}(\omega )|\omega |^{2\gamma _01},\omega 0(0<\gamma _0<1/2).$$
(D24)
From the Dyson equation we have
$$\mathrm{\Sigma }_r(q,\omega )=G_{0;r}^1(q,\omega )G_r^1(q,\omega ).$$
(D25)
Since $`G_{0;r}(q,\omega )=1/[\omega v_0rq+i\mathrm{sgn}(v_0rq)]`$ (see footnote \**†D 1), it follows that
$$G_{0;r}^1(q,\omega )=\omega v_0rq,$$
(D26)
and thus $`G_{0;r}^1(q=0,\omega )=\omega `$. For $`0<\gamma _0<1`$, the leading-order contribution to $`\mathrm{\Sigma }_r(q=0,\omega )`$ is therefore entirely due to $`G_r^1(q=0,\omega )`$ on the right-hand side of Eq. (D25). From Eq. (D23) we thus obtain ($`0<\gamma _0<1/2`$)
$$G_r^1(q=0,\omega )\frac{1}{\pi 𝒜_{\gamma _0+1}\left(\mathrm{cot}^2(\pi \gamma _0)+1\right)}\left\{\mathrm{cot}(\pi \gamma _0)+i\right\}\mathrm{sgn}(\omega )|\omega |^{12\gamma _0}.$$
(D27)
As a consequence of our above arguments, for $`0<\gamma _0<1/2`$ we have (for $`\gamma _0`$ in this range, $`\mathrm{cot}(\pi \gamma _0)>0`$)
$`\mathrm{\Sigma }_r^{}(q=0,\omega )`$ $``$ $`{\displaystyle \frac{\mathrm{cot}(\pi \gamma _0)\mathrm{sgn}(\omega )}{\pi 𝒜_{\gamma _0+1}\left(\mathrm{cot}^2(\pi \gamma _0)+1\right)}}|\omega |^{12\gamma _0},`$ (D28)
$`\mathrm{\Sigma }_r^{\prime \prime }(q=0,\omega )`$ $``$ $`{\displaystyle \frac{\mathrm{sgn}(\omega )}{\pi 𝒜_{\gamma _0+1}\left(\mathrm{cot}^2(\pi \gamma _0)+1\right)}}|\omega |^{12\gamma _0}.`$ (D29)
We note that $`\mathrm{\Sigma }_r(q=0,\omega =0)=0`$ (see Eq. (D24) and text preceding it), which in view of our considerations in Appendix B amounts to the fact that in the Luttinger model $`\mathrm{\Sigma }_r^{HF}(q=0)`$ is entirely cancelled by correlation effects. It is further interesting to note that
$$\frac{\mathrm{\Sigma }_r^{\prime \prime }(q=0,\omega )}{\mathrm{\Sigma }_r^{}(q=0,\omega )\mathrm{\Sigma }_r(q=0,\omega =0)}\mathrm{tan}(\pi \gamma _0),\omega 0,$$
(D30)
which is seen to approach zero (like $`\pi \gamma _0`$) for $`\gamma _00`$. As we have mentioned in § IV.C, $`\gamma _0=0`$ corresponds to the non-interacting system, so that the expressions presented in Eqs. (D28) and (D29) reproduce the non-interacting limit in a continuous way.
#### b The $`q0`$ case
Here we confine our considerations to the case corresponding to $`r=+`$. From Eq. (D14) (taking the limit $`\nu 0`$) it readily follows that for $`0<\gamma _0<1`$ the leading asymptotic contribution to $`G_+(q,\omega )`$, as $`\omega 0`$, is due to
$$\stackrel{~}{𝒢}_+^{(\mathrm{\Delta })}(q,z):=_{v_0q}^{v_0q+\mathrm{\Delta }}d\omega ^{}\frac{\varphi (q,\omega ^{})}{z\omega ^{}}.$$
(D31)
For $`z=v_0q+\omega +i\eta `$, $`\eta 0`$, with $`0<\omega <\mathrm{\Delta }`$, making use of $`\gamma (\gamma _0,\mathrm{\Lambda }(\omega ^{}+2v_0q)/[2v_0])=\gamma (\gamma _0,\mathrm{\Lambda }q)+𝒪(\omega ^{})`$, $`\mathrm{exp}(\mathrm{\Lambda }\omega ^{}/[2v_0])=1+𝒪(\omega ^{})`$, while employing the standard result ($`0<\omega <\mathrm{\Delta }`$)
$`𝒫{\displaystyle _0^\mathrm{\Delta }}d\omega ^{}{\displaystyle \frac{\omega _{}^{}{}_{}{}^{\sigma 1}}{\omega ^{}\omega }}`$ $`=`$ $`\pi \mathrm{cot}(\pi \sigma )\omega ^{\sigma 1}+{\displaystyle \frac{\mathrm{\Delta }^{\sigma 1}}{\sigma 1}}{}_{2}{}^{}F_{1}^{}(1,1\sigma ;2\sigma ;{\displaystyle \frac{\omega }{\mathrm{\Delta }}})`$ (D32)
$`=`$ $`\pi \mathrm{cot}(\pi \sigma )\omega ^{\sigma 1}{\displaystyle \frac{\mathrm{\Delta }^{\sigma 1}}{1\sigma }}{\displaystyle \frac{\mathrm{\Delta }^{\sigma 2}}{2\sigma }}\omega +𝒪\left(\omega ^2\right),`$ (D34)
$`\sigma >0,\sigma 1,2,\mathrm{},`$
we obtain ($`0<\omega <\mathrm{\Delta }`$, $`0<\gamma _0<1`$ and $`q0`$)
$$𝒢_+^{(\mathrm{\Delta })}(q,\omega )G_+(q,\omega )\pi \gamma (\gamma _0,\mathrm{\Lambda }q)𝒜_{\gamma _0}\left\{\mathrm{cot}(\pi \gamma _0)i\right\}\omega ^{\gamma _01},\omega 0.$$
(D35)
Note that the exponent $`\gamma _01`$, as compared with $`2\gamma _01`$ for the case corresponding $`q=0`$, implies that even for $`\gamma _0`$ in the open interval $`(0,1)`$ (as opposed to $`(0,1/2)`$), $`G_+(q0,\omega )`$ is to leading order asymptotically identical with the corresponding auxiliary function $`𝒢_+^{(\mathrm{\Delta })}(q0,\omega )`$ as $`\omega 0`$; recall that the leading asymptotic terms of $`G_r^{(\mathrm{\Delta })}(q=0,\omega )`$ and $`G_r(q=0,\omega )`$ for $`\omega 0`$ (see Eqs. (D23) and (D24)) are identical only when $`0<\gamma _0<1/2`$. Further, similar to the case corresponding to $`q=0`$, the divergence of $`G_+(q,\omega )`$ for $`0<\gamma _0<1`$, as $`\omega 0`$, implies that the possibly non-vanishing constant $`\{G_+(q0,\omega )𝒢_+^{(\mathrm{\Delta })}(q0,\omega )\}`$ is asymptotically irrelevant.
Proceeding as in the case of $`q=0`$, for $`\omega 0`$ we obtain ($`0<\gamma _0<1`$)
$$G_+^1(q,\omega )\frac{1}{\pi \gamma (\gamma _0,\mathrm{\Lambda }q)𝒜_{\gamma _0}\left(\mathrm{cot}^2(\pi \gamma _0)+1\right)}\left\{\mathrm{cot}(\pi \gamma _0)+i\right\}\omega ^{1\gamma _0}.$$
(D36)
As a consequence of our above arguments, for $`0<\gamma _0<1`$ we have (for $`\gamma _0`$ in this range, $`\mathrm{cot}(\pi \gamma _0)`$ takes on both positive as well as negative values)
$`\mathrm{\Sigma }_+^{}(q,\omega )`$ $``$ $`{\displaystyle \frac{\mathrm{cot}(\pi \gamma _0)}{\pi \gamma (\gamma _0,\mathrm{\Lambda }q)𝒜_{\gamma _0}\left(\mathrm{cot}^2(\pi \gamma _0)+1\right)}}\omega ^{1\gamma _0},`$ (D37)
$`\mathrm{\Sigma }_+^{\prime \prime }(q,\omega )`$ $``$ $`{\displaystyle \frac{1}{\pi \gamma (\gamma _0,\mathrm{\Lambda }q)𝒜_{\gamma _0}\left(\mathrm{cot}^2(\pi \gamma _0)+1\right)}}\omega ^{1\gamma _0}.`$ (D38)
We point out that $`\mathrm{\Sigma }_+(q,\omega =0)=0`$ (see Eq. (D24) and text preceding it), which in view of our considerations in Appendix B leads to the conclusion that, in the Luttinger model $`\mathrm{\Sigma }_+^{HF}(q)`$ is entirely cancelled by correlation effects. As in the case corresponding to $`q=0`$, here we have (compare with Eq. (D30))
$$\frac{\mathrm{\Sigma }_+^{\prime \prime }(q,\omega )}{\mathrm{\Sigma }_+^{}(q,\omega )\mathrm{\Sigma }_+(q,\omega =0)}\mathrm{tan}(\pi \gamma _0),\omega 0,$$
(D39)
which approaches zero (like $`\pi \gamma _0`$) for $`\gamma _00`$; this result is consistent with the fact that $`\gamma _0=0`$ corresponds to the non-interacting system.
We further point out that since $`\gamma (a,z)/z=\mathrm{exp}(z)z^{a1}`$, from Eqs. (D37) and (D38) we obtain that $`\mathrm{\Sigma }_+(q,\omega )/q|q|^{\gamma _01}`$ which diverges as $`q0`$ for $`\gamma _0<1`$. This result is in full conformity with our finding above (see the last paragraph in § D.1, following Eq. (D15)).
In closing this Appendix, we point out that the apparent difference in the asymptotic behaviours of $`\mathrm{\Sigma }_r(q,\omega )`$ (specifically concerning $`r=+`$ explicitly considered here) for $`q0`$ and $`q=0`$, is not specific to the Luttinger model. It is well known (Hodges, Smith and Wilkins 1971, Bloom 1975, Fujimoto 1990, Fukuyama, Narikiyo and Hasegawa 1991, Fukuyama, Hasegawa and Narikiyo 1991) that for isotropic systems of interacting electrons in two spatial dimensions, under the assumption that they are Fermi liquids, while $`\mathrm{\Sigma }^{\prime \prime }(k;\epsilon )^{}(\epsilon \epsilon _F)^2`$ for $`kk_F`$, we have $`\mathrm{\Sigma }^{\prime \prime }(k_F;\epsilon )^{}(\epsilon \epsilon _F)^2\mathrm{ln}|\epsilon \epsilon _F|`$. $`\mathrm{}`$
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