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# SUSY DARK MATTER
## 1 INTRODUCTION
Over the recent past there has been considerable experimental activity in the direct detection of dark matter and further progress is expected in the ongoing experiments and new experiments that may come online in the future. At the same time there have been several theoretical developments which have shed light on the ambiguities and possible corrections that might be associated with the predictions on supersymmetric dark matter. These consist of the effects on the dark matter analyses of wimp velocity and of the rotation of the galaxy, the effects of the uncertainties of quark densities and the uncertainties of the SUSY parameters, effects of large CP violating phases, effects of scalar nonuniversalities, effects of nonuniversalities of gaugino masses and effects of coannihilation. In this paper we will discuss some of these briefly but mainly focus on the effects of nonuniversalities of the gaugino masses on dark matter. In the Minimal Supersymmetric Standard Model (MSSM) there are 32 supersymmetric particles and with R parity conservation the lowest mass supersymmetric particle (LSP) is absolutely stable. In many unified models, such as in the SUGRA models, one finds that the lightest neutralino is the LSP over most of the parameter space of the model. Thus the lightest neutralino is a candidate for cold dark matter. The quantity that constrains supersymmetric models is $`\mathrm{\Omega }_\chi h^2`$ where $`\mathrm{\Omega }_\chi =\rho _\chi /\rho _c`$ where $`\rho _\chi `$ is the density of relic neutralinos at the current temperatures, and $`\rho _c=3H_0^2/8\pi G_N`$ is the critical matter density, and h is the Hubble parameter $`H_0`$ in units of 100 km/sMpc. The most recent measurements of h from the Hubble Space Telescope give
$$h=0.71\pm 0.03\pm 0.07$$
(1)
Similarly the most recent analyses of $`\mathrm{\Omega }_m`$ give
$$\mathrm{\Omega }_m=0.3\pm 0.08$$
(2)
If we assume that the component of $`\mathrm{\Omega }_B`$ in $`\mathrm{\Omega }_m`$ is $`\mathrm{\Omega }_B0.05`$ which appears reasonable, then this leads to the result $`\mathrm{\Omega }_\chi h^2=0.126\pm 0.043`$. Perhaps a more cautious choice of the range would be a $`2\sigma `$ range which gives
$$0.02\mathrm{\Omega }_\chi h^20.3$$
(3)
The quantity of interest theoretically is
$`\mathrm{\Omega }_\chi h^22.48\times 10^{11}\left({\displaystyle \frac{T_\chi }{T_\gamma }}\right)^3\left({\displaystyle \frac{T_\gamma }{2.73}}\right)^3{\displaystyle \frac{N_f^{1/2}}{J(x_f)}}`$ (4)
Here $`T_f`$ is the freeze-out temperature, $`x_f=kT_f/m_{\stackrel{~}{\chi }}`$ where k is the Boltzman constant, $`N_f`$ is the number of degrees of freedom at the time of the freeze-out, $`(\frac{T_\chi }{T_\gamma })^3`$ is the reheating factor, and $`J(x_f)`$ is given by
$$J(x_f)=_0^{x_f}𝑑x\sigma \upsilon (x)GeV^2$$
(5)
where $`<\sigma v>`$ is the thermal average with $`\sigma `$ the neutralino annihilation cross-section and v the neutralino relative velocity.
## 2 Detection of Milky Way wimps
Both direct and indirect methods are desirable and complementary for the detection of Milky Way wimps. We shall focus here on the direct detection. In this case the fundamental detector is the quark and the relevant interactions are the supergravity neutralino-quark-squark interactions. The scattering of neutralinos from quarks contains squark poles in the s channel and the Z boson and the Higgs boson ($`h,H^0,A^0`$) poles in the t channel. Since the wimp scattering from quarks is occuring at rather low energies one may, to a good approximation, integrate on the intermediate squark, Z and Higgs poles to obtain a low energy effective Lagrangian which gives a four-Fermi interaction of the following form
$$_{eff}=\overline{\chi }\gamma _\mu \gamma _5\chi \overline{q}\gamma ^\mu (AP_L+BP_R)q+C\overline{\chi }\chi m_q\overline{q}q+D\overline{\chi }\gamma _5\chi m_q\overline{q}\gamma _5q$$
(6)
The contribution of D is generally small and thus the scattering is effectively governed by the terms A,B and C. Analysis of dark matter is affected by several factors. We discuss these briefly below.
### 2.1 Uncertainties in wimp density and velocity
Two of the quantities that control the detection of dark matter are the wimp mass density and the wimp velocity. Estimates of Milky Way wimp density lie in the range $`\rho _\chi =(0.20.7)GeVcm^3`$ and the event rates in the direct detection depend directly on this density. A second important factor regarding wimps that enters in the dark matter analyses is the wimp velocity. One typically assumes a Maxwellian velocity distribution for the wimps and the current estimates for the rms wimp velocity give $`v=270km/s`$ with, however, a significant uncertainty. Estimates for the uncertainty lie in the range of $`\pm 24`$ km/s to $`\pm 70`$ km/s . A reasonable estimate then is that the rms wimp velocity lies in the range
$$v=270\pm 50km/s$$
(7)
Analyses including the wimp velocity variations show that the detection rates can have a significant variation, i.e., a factor of 2-3 on either side of the central values.
### 2.2 Effects of uncertainties of quark densities
The scattering of neutralinos from quarks are dominated by the scalar interaction which is controlled by the term C in Eq.(6). The dominant part of the scattering thus arises from the scalar part of the $`\chi p`$ cross-section which is given by
$$\sigma _{\chi p}(scalar)=\frac{4\mu _r^2}{\pi }(\underset{i=u,d,s}{}f_i^pC_i+\frac{2}{27}(1\underset{i=u,d,s}{}f_i^p)\underset{a=c,b,t}{}C_a)^2$$
(8)
Here $`\mu _r`$ is the reduced mass in the $`\chi p`$ system and $`f_i^p`$ (i=u,d,s quarks) are quark densities inside the proton defined by
$$m_pf_i^p=<p|m_{qi}\overline{q}_iq_i|p>$$
(9)
There are significant uncertainties in the determination of $`f_i^p`$. To see the range of these uncertainties it is useful to parametrize the quark densities so that
$`f_u^p={\displaystyle \frac{m_u}{m_u+m_d}}(1+\xi ){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$
$`f_d^p={\displaystyle \frac{m_d}{m_u+m_d}}(1\xi ){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$
$`f_s^p={\displaystyle \frac{m_s}{m_u+m_d}}(1x){\displaystyle \frac{\sigma _{\pi N}}{m_p}}`$ (10)
where we have defined
$`x={\displaystyle \frac{\sigma _0}{\sigma _{\pi N}}}={\displaystyle \frac{<p|\overline{u}u+\overline{d}d2\overline{s}s|p>}{<p|\overline{u}u+\overline{d}d|p>}},\xi ={\displaystyle \frac{<p|\overline{u}u\overline{d}d|p>}{<p|\overline{u}u+\overline{d}d|p>}}`$
$`\sigma _{\pi N}=<p|2^1(m_u+m_d)(\overline{u}u+\overline{d}d|p>`$ (11)
The current range of determinations of $`\sigma _{\pi N}`$ give
$$\sigma _{\pi N}=48\pm 9MeV,x=0.74\pm 0.25,\xi =0.132\pm 0.035$$
(12)
With the above range of errors one finds that $`f_i^p`$ lie in the range
$$f_u^p=0.021\pm 0.004,f_d^p=0.029\pm 0.006,f_s^p=0.21\pm 0.12$$
(13)
Of these the errors in $`f_s^p`$ generates the largest variations. A detailed analysis shows that the scalar cross-section can vary by a factor of 5 in either direction due to errors in the quark densities.
### 2.3 CP violation effects on dark matter
The soft SUSY breaking parameters that arise in supersymmetric theories after spontaneous breaking are in general complex with phases O(1) and can lead to large electric dipole moment of the electron and of the neutron in conflict with current experiment. Recently a cancellation mechanism was proposed as a possible solution to this problem. With the cancellation mechanism the total EDM of the electron and of the neutron can be in conformity with data even with phases O(1) and sparticle masses which are relatively light. The presence of large CP phases can affect dark matter and other low energy phenomena. Effects of CP violating phases on dark matter have been investigated for some time and more recently such analyses have been extended to determine the effects of large CP phases under the cancellation mechanism. One finds that the EDM constraints play a crucial role in these analyses. Thus in the absence of CP violating phases one finds that the $`\chi p`$ cross-section can change by orders of magnitude when plotted as a function of $`\theta _\mu `$ (the phase of $`\mu `$). These effects are, however, significantly reduced when the constraints arising from the current experimental limits on the electron and on the neutron EDMs are imposed. After the imposition of the constraints the effects of CP violating phases are still quite significant in that the $`\chi p`$ cross-section can vary by a factor of $`2`$. Thus precision predictions of the $`\chi p`$ cross-section should take account of the CP phases if indeed such phases do exist in a given model. Indeed many string and D brane models do indeed possess such phases and thus inclusion of such phases is imperative in making predictions for direct detection in such models.
### 2.4 Effects of coannihilation
The effects of coannihilation may become important when the next to the lowest supersymmetric particle (NLSP) has a mass which lies close to the LSP mass. The size of the effects is exponentially damped by the factor $`e^{\mathrm{\Delta }_ix}`$ where $`\mathrm{\Delta }_i=(m_i/m_\chi 1)`$, $`x=m_\chi /kT`$ and where $`m_\chi `$ is the LSP mass. Because of this damping the coannihilation effects are typically important only for regions of the parameter space where the constraint $`\mathrm{\Delta }_i<0.1`$ is satisfied. Some of the possible candidates for NLSP are the light stau $`\stackrel{~}{\tau }_1`$, $`\stackrel{~}{e}_R`$, the next to the lightest neutralino $`\chi _2^0`$, and the light chargino $`\chi _1^+`$. An interesting result one finds is that in mSUGRA the upper limit on the neutralino mass consistent with the current experimental constraints on the relic density is extended from 200 GeV to 600 GeV when the effects of $`\chi \stackrel{~}{\tau }`$ coannihilation are included. In Secs.(2.6) and (2.7) we will show that the allowed range of the neutralino mass can also be extended by inclusion of nonuniversalities in the gaugino masses.
### 2.5 Nonuniversality of scalar masses
The minimal SUGRA model is based on the universality at the GUT scale. This includes the universality of the scalar masses, of the gaugino masses and of the trilinear couplings at the GUT scale. In supergravity unified models the universality of the soft SUSY breaking parameters arises from the assumption of a flat Kahler potential. However, the nature of physics at the GUT/Planck scale is not fully understood and a more general analysis of the soft SUSY breaking sector requires that one work with a curved Kahler potential. Such an analysis in general leads to nonuniversalities in the scalar sector of the theory. However, the nonuniversalities in the scalar sector cannot be completely arbitrary as there are very stringent constraints on the system from the limits on the flavor changing neutral currents (FCNC). A satisfaction of the constraints requires essentially a degeneracy in the scalar masses in the first two generations at the GUT scale. However, the constraints on the scalar masses in the Higgs sector and on masses in the third generation are far less severe and one could introduce significant amounts of nonuniversalities in these sectors without violating the FCNC constraints. It is found convenient to parametrize the nonuniversalities in the Higgs sector by $`\delta _1,\delta _2`$ so that at the GUT scale ($`M_G`$) one has $`m_{H_1}^2=m_0^2(1+\delta _1)`$, $`m_{H_2}^2=m_0^2(1+\delta _2)`$. Similarly one may parametrize the nonuniversalities in the third generation squark sector by $`\delta _3,\delta _4`$ so that at the scale $`M_G`$ one has $`m_{\stackrel{~}{Q}_L}^2=m_0^2(1+\delta _3)`$, $`m_{\stackrel{~}{U}_R}^2=m_0^2(1+\delta _4)`$. These nonuniversalities have a significant effect on the low energy physics. One of the main effects that occurs is through the effect on $`\mu `$ which is determined via the constraint of the radiative breaking of the electro-weak symmetry and is modified in the presence of the nonuniversalities in the Higgs sector and in the third generation sector. To one loop order it is given by
$$\mu ^2=\mu _0^2+\frac{m_0^2}{t^21}(\delta _1\delta _2t^2\frac{D_01}{2}(\delta _2+\delta _3+\delta _4)t^2)+\mathrm{\Delta }\mu ^2$$
(14)
Here $`\mu _0`$ is the value of $`\mu `$ in the absence of nonuniversalities, $`D_0`$ depends on the top Yukawa coupling and defines the position of the Landau pole, $`ttan\beta `$, and $`\mathrm{\Delta }\mu ^2`$ is the loop correction. We note that the entire effect of nonuniversalities is now explicity exhibited. One finds that the universalities can significantly affect the event rates. The effect on the event rates occurs specifically because of the effect on $`\mu `$. Thus one finds that for certain regions of the parameter space the nonuniversalities in the Higgs and in the third generation sector make a negative contribution to $`\mu ^2`$ which leads to larger higgsino components for the neutralino. Since in the direct detection the scattering is dominated by the scalar $`\chi p`$ cross-section which in turn depends on the product of the gaugino and the higgsino components one finds that a smaller $`\mu `$ leads to larger event rates in the direct detection. A detailed analysis of the effects of nonuniversalities of the scalar masses has been given in Refs.. We will discuss further this phenomena in the context of the nonuniversalities in the gaugino sector in the next section.
### 2.6 Gaugino nonuniversalities and dark matter in GUT models
Nonuniversality of gaugino masses arises in grand unified models via corrections to the gauge kinetic energy functions. Thus in grand unified models a non-trivial gauge kinetic energy function leads to a gaugino mass matrix which has the form
$$m_{\alpha \beta }=\frac{1}{4}\overline{e}^{G/2}G^a(G^1)_a^b(f_{\alpha \gamma }^{}/z^b)f_{\gamma \beta }^1$$
(15)
As an example if we consider the GUT group to be SU(5) then the gauge kinetic energy function $`f_{\alpha \beta }`$ transforms as follows
$$(\mathrm{𝟐𝟒}\times \mathrm{𝟐𝟒})_{symm}=\mathrm{𝟏}+\mathrm{𝟐𝟒}+\mathrm{𝟕𝟓}+\mathrm{𝟐𝟎𝟎}$$
(16)
where $`(\mathrm{𝟐𝟒}\times \mathrm{𝟐𝟒})_{symm}`$ stands for the symmetric product. The term that transforms like the singlet of SU(5) in the gauge kinetic energy function leads to universality of the gaugino masses, while the $`\mathrm{𝟐𝟒}`$ plet, the $`\mathrm{𝟕𝟓}`$ plet and the $`\mathrm{𝟏𝟐𝟎}`$ plet will generate corrections to universality. In general one could have an admixture of the various representations and this will lead to gaugino masses of the form
$$\stackrel{~}{m}_i(0)=m_{\frac{1}{2}}(1+\underset{r}{}c_rn_i^r)$$
(17)
where $`n_i^r`$ depend on r and for the representations $`\mathrm{𝟏},\mathrm{𝟐𝟒},\mathrm{𝟕𝟓},\mathrm{𝟐𝟏𝟎}`$ they are given in Table 1. The nonuniversality of the gaugino masses also leads to corrections of the gauge coupling constants at the GUT scale and in general one has $`g_i(M_G)=g_G(1+_rc_r^{}n_i^r)`$. We note, however, that the coefficients $`c_r^{}`$ that enter in $`g_i`$ are different than those that enter in $`m_i`$. This is so because the corrections to $`g_i`$ involve only the gauge kinetic energy function while the corrections to $`m_i`$ involve the gauge kinetic energy function as well as the nature of GUT physics.
| Table: nonuniversalities at $`M_X`$. | | | |
| --- | --- | --- | --- |
| SU(5) rep | $`n_1^r`$ | $`n_2^r`$ | $`n_3^r`$ |
| 1 | 1 | 1 | 1 |
| 24 | -1 | $`3`$ | $`2`$ |
| 75 | -5 | 3 | $`1`$ |
| 200 | 10 | 2 | 1 |
The gaugino sector nonuniversalities affect $`\mu `$. To exhibit this effect we can expand $`\mu `$ determined via the constraint of the radiative breaking of the electro-weak symmetry in terms of the parameter $`c_r`$. One finds the following expansion
$$\stackrel{~}{\mu }^2=\mu _0^2+\underset{r}{}\frac{\stackrel{~}{\mu }^2}{c_r}c_r+O(c_r^2)$$
(18)
and for $`c_{24}<0,c_{75}<0,c_{200}>0`$ one has
$$\frac{\mu _{24}^2}{c_{24}}>0,\frac{\mu _{75}^2}{c_{75}}>0,\frac{\mu _{200}^2}{c_{200}}<0$$
(19)
Thus in these cases the nonuniversalities lead to a smaller value of $`|\mu |`$. Now as already mentioned in the previous section the Higgsino components become more dominant as $`\mu `$ becomes smaller. We can exhibit this analytically for the case when $`\mu `$ is small but we are still in the scaling region where $`\mu ^2/M_Z^2>>1`$. In this case it is possible to analytically investigate the size of the gaugino-Higgsino components $`X_{n0}`$ of the LSP defined by
$$\chi =X_{10}\stackrel{~}{B}+X_{20}\stackrel{~}{W}_3+X_{30}\stackrel{~}{H}_1+X_{40}\stackrel{~}{H}_2$$
(20)
where $`\stackrel{~}{B}`$ is the Bino, $`\stackrel{~}{W}_3`$ is the Wino, and $`\stackrel{~}{H}_1`$, and $`\stackrel{~}{H}_2`$ are the two Higgsinos. In this case one finds that the gaugino components of the LSP are given by $`X_{11}1(\frac{M_Z^2}{2\mu ^2})sin^2\theta _W`$, and $`X_{12}\frac{M_Z^2}{2m_{\chi _1}^2\mu }sin2\theta _Wsin\beta `$ while the higgsino components are given by $`X_{13}\frac{M_Z}{\mu }sin2\theta _Wsin\beta `$, $`X_{14}\frac{M_Z}{\mu }sin2\theta _Wsin\beta `$. From the above one finds that the Higgsino components have a dependence on the inverse power of $`\mu `$ and thus a smaller $`\mu `$ will lead to a larger scalar $`\sigma _{\chi p}`$ cross-section.
The literature on the analyses of dark matter relic density and direct detection in MSSM and in SUGRA models is quite extensive. We discuss here the quantitative effects of the gaugino mass nonuniversality on dark matter. Some features of the effects of gaugino mass nonuniversalities have already been discussed in the literature and we review here the more recent developments. The techniques used in the analysis are as discussed in Ref. and in the analysis we impose the $`bs+\gamma `$ constraint. In Fig.1 we plot the scalar $`\chi p`$ cross-section as a function of the neutralino mass for the case of GUT scale nonuniversalities with values of $`c_{24}`$ in the range -0.1 to 0.08. One finds that the scalar cross-section is enhanced for negative values of $`c_{24}`$ just as one would expect from the general discussion above because it is for the case of $`c_{24}`$ negative that $`\mu `$ becomes small. One finds that in general the scalar cross-section increases systematically as $`|c_{24}|`$ increases for negative values of $`c_{24}`$ and an enhancement of the scalar cross-section by as much as a factor of 10 can be gotten relative to the universal case of $`c_{24}=0`$. One also finds an enhancement of the allowed range of the neutralino mass consistent with the constraints. In Fig.2 we plot the maximum and the minimum of the scalar cross-section as a function of the neutralino mass for the case of GUT scale nonuniversalities where the nonuniversalities arise from the 200 plet representation with $`c_{200}=0.1`$ when the other parameters are varied over their assumed naturalness range. The current experimental limits from DAMA and from CDMS are also plotted. Further, the currents limits would certainly be significantly improved in other dark matter detectors in the future and in Fig.2 we also plot the expected limits from future CDMS, and from GENIUS. One finds that the current experiment does constrain the theory in a small region of the parameter space. Further, the expected sensitivity in future experiment, i.e., in CDMS and in GENIUS will explore a major part of the parameter space of this model. We also note that the inclusion of nonuniversality significantly increases the allowed range of the neutralino parameter space.
### 2.7 Dark Matter on D Branes
Nonuniversality of gaugino masses is rather generic in string theory. However, the specific nature of the nonuniversality will depend on the details of the compactification. We discuss here the effects of gaugino mass nonuniversality on dark matter in the context of D brane models. The possibility of nonuniversal gaugino phases in brane models arises from the choice of embedding of the different gauge groups of the Standard Model on different branes. One may consider, from example, models that arise from Type IIB string compactified on a six-torus $`T^2\times T^2\times T^2`$ which contains 9 branes, $`7_i`$ and $`5_i`$ (i=1,2,3) branes and 3 branes. Not all the branes can be present simultaneuosly due to the constraint of N=1 supersymmetry which requires that one has either 9 branes and $`5_i`$ (i=1,2,3) branes or $`7_i`$ (i=1,2,3) branes and 3 branes. In the following we will make the choice of embedding on $`9`$ branes and 5 branes. One of the major problems in developing a sensible string phenomenology is that the mechanism of supersymmetry breaking in string theory is still lacking. However, some progress can be made by use of an efficient parametrization of supersymmetry breaking. Here we use the parametrization where the breaking of supersymmetry arises from the breaking generated by the dilaton and the moduli VEV’s of the following form $`F^S=`$$`\sqrt{3}m_{3/2}`$$`(S+S^{})sin\theta e^{i\gamma _S}`$, $`F^i=`$$`\sqrt{3}m_{3/2}`$$`(T_i+T_i^{})cos\theta _ie^{i\gamma _i}`$. We consider now a specific 9-5 brane model. Here one embeds the $`SU(3)_C\times U(1)_Y`$ gauge group on 9 branes and $`SU(2)_L`$ gauge group on a $`5_1`$ brane. The alternative possibility of embedding the Standard Model gauge group on five branes is discussed in the last two papers of Ref.. For the $`95_1`$ brane model the soft SUSY breaking sector of the theory is given by
$`\stackrel{~}{m}_1=\stackrel{~}{m}_3=\sqrt{3}m_{3/2}sin\theta e^{i\gamma _S}=A_0;\stackrel{~}{m}_2=\sqrt{3}m_{3/2}cos\theta e^{i\gamma _i}`$
$`\stackrel{~}{m}_9^2=m_{3/2}^2(13cos^2\theta \mathrm{\Theta }_1^2);\stackrel{~}{m}_{95_1}^2=m_{3/2}^2(1(3/2)cos^2\theta (1\mathrm{\Theta }_1^2))`$ (21)
Here $`\theta (\mathrm{\Theta }_i)`$ is the Goldstino direction in the dilaton S (moduli $`T_i`$) VEV space. We discuss now dark matter on D branes. In Fig.3 we give a plot of the scalar $`\chi p`$ cross section as a function of the neutralino mass for the $`95_1`$ D brane model. One of the interesting feature of the D brane model is that the scalar masses are in general not universal. However for $`\mathrm{\Theta }_1=1/\sqrt{3}`$ one has $`m_9=m_{95_1}`$ and the scalar masses are universal although the gaugino masses are still nonuniversal. Since we are mostly interested here in investigating the effects of nonuniversalities of the gaugino masses in this analysis, we impose universality of the scalar masses and set $`\mathrm{\Theta }_1=1/\sqrt{3}`$. In Fig.3 we give a plot of the minimum and the maximum of the scalar $`\chi p`$ cross section under this constraint. One finds that under the assumed constraints the allowed domain of the parameter space has the general features which are similar to the GUT scale nonuniversalities. One common feature is that the allowed domain of the parameter space is extended close to 500 GeV. One may note that if in addition to the constraint $`\mathrm{\Theta }_1=1/\sqrt{3}`$ one also sets $`\theta =\pi /6`$ one finds also universality of the gaugino masses. This situation is exhibited by the vertical dark line in the enclosed region on the left hand side in Fig.3.
## 3 Conclusion
In this paper we have given a brief review of the recent theoretical developments in the analyses of supersymmeteric dark matter. Our emphasis has been in exploring the effects of uncertainties of the input data and the effects of nonuniversalities of the gaugino masses on dark matter analyses. It is found that the uncertainties of the wimp velocities can change detection rates by up to factors of 2-3 while the uncertainties in quark masses and densities can change the $`\chi p`$ cross-section by up to factors of 5 in either direction. The effects of gaugino mass nonuniversalities on dark matter analyses is found to be quite dramatic. It is seen that gaugino mass nonuniversalities can increase the Higgsino components of the LSP and significantly increase the $`\chi p`$ cross-section from scalar interactions and also increase the allowed range of the LSP consistent with relic density constraints. Thus an increase in the scalar $`\chi p`$ cross-sections by up to a factor of 10 can occur while the allowed range of the neutralino masses can move up to 500 GeV consistent with the relic density constraints. Data from current dark matter experiments is beginning to put constraints on models with nonuniversalities. These constraints will become more severe as the sensitivity of dark matter experiments increase in the future.
Acknowledgements
This research was supported in part by NSF grant PHY-9901057.
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# Finding 𝜂_𝑐' and ℎ_𝑐(¹𝑃₁) at HERA-𝐵
## I Introduction
The dramatic discovery of Charmonia, the $`J/\psi `$ and its excited states, marked the beginning of a new era of particle physics. Till now the Charmonium physics remains to be one of the most exciting areas of high energy physics. As the ”hydrogen-like atoms” of strong interaction, the Charmonia could be investigated partly by virtue of the perturbative QCD (pQCD) in account of the large charm quark mass, which makes the study on a relatively solid ground, as well the study may give clues of the nature of non-perturbative QCD.
Although the first Charmonium state, $`J/\psi `$, was observed more than twenty years ago, the study of the Charmonium states is still far from satisfactory. Except for the $`J/\psi `$ itself, the knowledge of the other Charmonia is very limited. We do not even have a complete $`c\overline{c}`$ mass spectrum below the $`D\overline{D}`$ threshold , that is, the existence of the S-wave spin singlet $`\eta _c^{}`$ and the P-wave spin singlet $`h_c(^1P_1)`$ is still based on very weak experimental signals. To confirm the existing findings and give out more precise values of the mass, width, and other parameters of these two resonances are now a pressing task in experiment.
The $`\eta _c^{}`$ was first observed in the Crystal Ball experiment in the inclusive photon energy spectrum from $`\psi ^{}`$ decays at 3594 MeV , until now the signal was not observed by other experiments due to the low energy of the radiative photon and the relatively poor photon detection ability of other detectors comparing with that of the Crystal Ball’s . The $`h_c(^1P_1)`$ state was first observed at 3526.14 MeV in the proton antiproton annihilation experiment by E760 group performed at Fermilab , but with the statistical significance of the signal slightly more than three standard deviations, and also no other experiments definitely confirm the existence by now (the E705’s report was doubted by Barnes, Browder, and Tuan ).
Currently, the experiments suitable for Charmonium studies are: the BES detector running at the BEPC $`e^+e^{}`$ collider, the $`p\overline{p}`$ annihilation experiment represented by the E835 experiment at Fermilab, and the scarce studies of the two-photon process in high energy $`e^+e^{}`$ colliders like at LEP and CESR.
Due to the restriction of the quantum number at $`e^+e^{}`$ colliders, only the vector-like Charmonium states like $`J/\psi `$ and $`\psi ^{}`$ can be produced directly at lowest order, whereas the other Charmonium states, like $`\chi _{cJ}`$, $`\eta _c`$, and $`h_c(^1P_1)`$, can only be produced via either higher order processes or through the $`J/\psi `$($`\psi ^{}`$) electromagnetic and/or hadronic decays. For instance, the $`\eta _c^{}`$ may be produced via $`\psi ^{}\gamma +\eta _c^{}`$ and the $`h_c(^1P_1)`$ state via $`\psi ^{}\pi ^0+h_c(^1P_1)`$. Although BES detector has collected the largest $`\psi ^{}`$ data sample in the world, due to the limited energy resolution of the Electromagnetic Calorimeter and the rather small production rates of $`\eta _c^{}`$ and $`h_c(^1P_1)`$ in $`\psi ^{}`$ decays, the search of either $`\eta _c^{}`$ or $`h_c(^1P_1)`$ did not give significant results.
As for proton-antiproton annihilation experiments, although they can produce Charmonium states of various quantum numbers and can be used to determine the resonance parameters of the Charmonium produced, the study of Charmonium is limited by the detection of the electromagnetic final states and the low production rate. The E760 and its succeeding version E835 did a very good job in measuring the resonance parameters of the $`\chi _{c1}`$, $`\chi _{c2}`$ and some other Charmonium states, but the study on the $`h_c(^1P_1)`$ state is still insufficient and the existence of the $`\eta _c^{}`$ is not confirmed. Of course, the E835 will continue this work and look further with more data in the near future.
The HERA-$`B`$ , an experiment presently set up at DESY, which uses the HERA 920 GeV proton beam incident on various nuclear targets, is focused on the measurement of CP-violation in the $`B\overline{B}`$ system via mainly the final states containing $`J/\psi `$. The trigger system is designed to recognize events with $`J/\psi \mathrm{}^+\mathrm{}^{}`$ ($`\mathrm{}=e\text{or}\mu `$). Furthermore, the detector also is designed for precise measurement of photons with its Electromagnetic Calorimeter(ECAL), which makes the study of Charmonia very possible through detecting the final states of the Charmonium decays containing $`J/\psi `$ and neutral particles like $`\gamma `$ or $`\pi ^0`$.
The paper is organized as follows. In following section we present the formalism for $`\eta _c^{}`$ and $`h_c(^1P_1)`$ production in a general framework in fixed-target experiment. In section III the obtained formalism is applied to HERA-$`B`$ situation numerically; the direct and indirect production rates of these two states are evaluated. In section IV, we give a rough estimation of the signals and backgrounds in searching for these two states for experimentalists’ reference. In the last section some discussions and conclusions are made.
## II $`\eta _c^{}`$ and $`h_c(^1P_1)`$ Production
For $`\eta _c^{}`$ production, to leading order in $`\alpha _s`$ and $`v^2`$, the relative velocity of heavy quarks inside the bound state, it is a two to one process as shown in Figure 1. The parton level cross section can be easily calculated or just obtained from the corresponding $`\eta _c`$ producing process with the non-perturbative sector replaced. It is
$`\widehat{\sigma }_1={\displaystyle \frac{2\pi ^3\alpha _s^2}{9(2m_c)^5}}<0|𝒪_1^{\eta _c^{}}(^1S_0)|0>z\delta (1z).`$ (1)
Here, $`\alpha _s`$ is the strong coupling constant; $`<0|𝒪_1^{\eta _c^{}}(^1S_0)|0>`$ is the NRQCD Color-Singlet non-perturbative matrix element, which can be related to $`|R_{\eta _c^{}}(0)|`$, the radial wave function at the origin of the bound state, by $`<0|𝒪_1^{\eta _c^{}}(^1S_0)|0>=\frac{3}{2\pi }|R(0)|^2`$; and $`zM_{\eta _c^{}}^2/\widehat{s}`$, where $`\widehat{s}`$ denotes the c.m.s. energy in partonic system.
As for the $`h_c(^1P_1)`$ production, the situation is somewhat different from that of $`\eta _c^{}`$. Of the latter, at leading order in $`\alpha _s`$ and $`v^2`$ there is only one possible channel giving the contribution, but of the former, there are several to the same order of accuracy. To be more clearly, according to the BBL theory for Quarkonium production and decays , the Fock states of Quarkonium are ordered in $`v`$, i.e.,
$`|h_c(^1P_1)`$ $`=`$ $`𝒪(1)|c\underset{¯}{\text{c}}[^1P_1^{(1)}]+𝒪(v)|c\underset{¯}{\text{c}}[^1S_0^{(8)}]g+𝒪(v)|c\underset{¯}{\text{c}}[^1D_2^{(8)}]g+\mathrm{}.`$ (2)
Because for P-wave states the leading non-vanishing wave functions are the derivative of the wave functions at the origin, or in other words that the P-wave states are produced via the NRQCD dimension 8 operators or higher, the NRQCD scaling rules tell us that for $`h_c(^1P_1)`$ production the non-perturbative matrix elements stemming from the first two terms in Eq.(2) are of the same order in $`v^2`$. Based on this argument the leading order Color-Singlet and -Octet processes of the $`h_c(^1P_1)`$ production are shown in Figure 2.
As depicted in Figure 2(a), the Color-Octet process is also a two to one process. The cross section of the partonic scattering process can be straightforwardly obtained,
$`\widehat{\sigma }_2={\displaystyle \frac{5\pi ^3\alpha _s^2}{12(2m_c)^5}}<0|𝒪_8^{h_c}(^1S_0)|0>z\delta (1z),`$ (3)
where $`<0|𝒪_8^{h_c}(^1S_0)|0>`$ is the Color-Octet nonperturbative matrix element.
Of the Color-Singlet processes, Figure 2 (b)-(d), the two gluon fusion channel of (b) may survive only with at least an additional gluon in the final states from the Landau-Yang theorem, as shown in the Figure; the others are not restricted by this law, however ruled out by the properties of charge-conjugation of the processes. The reason for this is that heavy-quark-loop factor (including the projector for the quarkonium state) is odd under charge conjugation. That is, the C-odd $`h_c`$ state can not decay through two vector currents (C-even), and the direct calculation really shows they give no contributions. The cross section of Figure 2(b) reads as
$`{\displaystyle \frac{\widehat{\sigma }_3(g+gh_c[^1P_1^{(1)}])}{d\widehat{t}}}`$ $`=`$ $`{\displaystyle \frac{\pi ^2\alpha _s^3<0|𝒪_1^{h_c}(^1P_1)|0>}{108(2m_c)\widehat{s}^2}}\{24{\displaystyle \frac{4\widehat{t}^2\widehat{u}^2+\widehat{s}\widehat{t}\widehat{u}(\widehat{t}+\widehat{u})+2\widehat{s}^2(\widehat{t}^2+\widehat{t}\widehat{u}+\widehat{u}^2)}{(\widehat{s}+\widehat{t})^2(\widehat{s}+\widehat{u})^2(\widehat{t}+\widehat{u})^2}}`$ (4)
$`+`$ $`{\displaystyle \frac{40}{3(\widehat{s}+\widehat{t})^3(\widehat{s}+\widehat{u})^3(\widehat{t}+\widehat{u})^3}}(12\widehat{s}^6\widehat{t}+44\widehat{s}^5\widehat{t}^2+72\widehat{s}^4\widehat{t}^3+72\widehat{s}^3\widehat{t}^4`$ (5)
$`+`$ $`44\widehat{s}^2\widehat{t}^5+12\widehat{s}\widehat{t}^6+12\widehat{s}^6\widehat{u}+58\widehat{s}^5\widehat{t}\widehat{u}+149\widehat{s}^4\widehat{t}^2\widehat{u}+179\widehat{s}^3\widehat{t}^3\widehat{u}`$ (6)
$`+`$ $`140\widehat{s}^2\widehat{t}^4\widehat{u}+56\widehat{s}\widehat{t}^5\widehat{u}+12\widehat{t}^6\widehat{u}+46\widehat{s}^5\widehat{u}^2+157\widehat{s}^4\widehat{t}\widehat{u}^2+246\widehat{s}^3\widehat{t}^2\widehat{u}^2`$ (7)
$`+`$ $`231\widehat{s}^2\widehat{t}^3\widehat{u}^2+142\widehat{s}\widehat{t}^4\widehat{u}^2+44\widehat{t}^5\widehat{u}^2+78\widehat{s}^4\widehat{u}^3+198\widehat{s}^3\widehat{t}\widehat{u}^3+240\widehat{s}^2\widehat{t}^2\widehat{u}^3`$ (8)
$`+`$ $`178\widehat{s}\widehat{t}^3\widehat{u}^3+72\widehat{t}^4\widehat{u}^3+79\widehat{s}^3\widehat{u}^4+158\widehat{s}^2\widehat{t}\widehat{u}^4+149\widehat{s}\widehat{t}^2\widehat{u}^4+72\widehat{t}^3\widehat{u}^4`$ (9)
$`+`$ $`47\widehat{s}^2\widehat{u}^5+61\widehat{s}\widehat{t}\widehat{u}^5+44\widehat{t}^2\widehat{u}^5+12\widehat{s}\widehat{u}^6+12\widehat{t}\widehat{u}^6)\}.`$ (10)
Here in the above, the $`\widehat{s}(p_1+p_2)^2`$, $`\widehat{t}(p_1p_3)^2`$, and $`\widehat{u}(p_2p_3)^2`$ are ordinary Mandelstam variables; the universal non-perturbative matrix element $`<0|𝒪_1^{h_c}(^1P_1)|0>`$ related to the derivative of the radial wave function at original of $`h_c(^1P_1)`$ by $`<0|𝒪_1^{h_c}(^1P_1)|0>=\frac{27}{2\pi }|R_{h_c}^{}(0)|^2`$.
Except for the direct production of these two states given in above, another main source of their production is of the electromagnetic or hadronic decays of the $`\psi ^{}`$ in accompanying with one $`\gamma `$ or $`\pi ^0`$. The dominant partonic interaction processes of the $`\psi ^{}`$ production in $`pN`$ collision at HERA-$`B`$ energy are drawn as Figure 3.
The expression for gluon-gluon fusion processes, the Figure 3(a) and (b), can be written as
$`\widehat{\sigma }_4(g+g\psi ^{})=`$ (11)
$`{\displaystyle \frac{5\pi ^3\alpha _s^2}{12(2m_c)^5}}\{<0|𝒪_8^\psi ^{}(^1S_0)|0>+{\displaystyle \frac{3}{m_c^2}}<0|𝒪_8^\psi ^{}(^3P_0)|0>+{\displaystyle \frac{4}{5m_c^2}}<0|𝒪_8^\psi ^{}(^3P_2)|0>\}z\delta (1z)`$ (12)
$`+{\displaystyle \frac{20\pi ^2\alpha _s^3}{81(2m_c)^5}}(<0|𝒪_1^\psi ^{}(^3S_1)|0>z^2\{{\displaystyle \frac{1z^2+2z\mathrm{log}z}{(z1)^2}}+{\displaystyle \frac{1z^2+2z\mathrm{log}z}{(z+1)^3}}\}\theta (1z).`$ (13)
The expression for process of Figure 3(c) is quite simple, it is
$`\widehat{\sigma }_5(q+\overline{q}\psi ^{}[^3S_1^{(8)}])={\displaystyle \frac{16\pi ^3\alpha _s^2}{27(2m_c)^5}}<0|𝒪_8^\psi ^{}(^1S_0)|0>z\delta (1z).`$ (14)
Here, although the Octet processes are suppressed in $`v^2`$, they get compensation from the enhancement of $`1/\alpha _s`$ relative to the Color-Singlet process. So, it is proper to include them in the $`\psi ^{}`$ production rate estimation.
## III Numerical Estimation for $`\eta _c^{}`$ and $`h_c(^1P_1)`$ Production at HERA-$`B`$
In the above section we have calculated the necessary partonic cross sections at leading order in $`v^2`$ or/and $`\alpha _s`$ for $`\eta _c^{}`$ and $`h_c(^1P_1)`$ production in the proton-nucleon collision. According to the general factorization theorem the experimental cross sections can be obtained by convoluting the subprocess with the parton distribution functions in the nucleons. i.e.,
$`\sigma (A+BC+X)={\displaystyle G_a(x_a)G_b(x_b)\widehat{\sigma }(a+bC+Y)𝑑x_a𝑑x_b},`$ (15)
where the sum runs over all the possible initial interacting partons which involve in the interaction; the $`A`$ and $`B`$ represent nucleons; $`C`$ represents the Charmonium; $`X`$ and $`Y`$ are the remnants of the inclusive processes; $`G_a(x_a)`$ and $`G_b(x_b)`$ are the parton distribution functions of the colliding nucleons $`A`$ and $`B`$ with momentum fractions $`x_a`$ and $`x_b`$, respectively.
In doing the numerical estimation the following inputs are taken
$`\alpha _s(2m_c)=0.253,M_{\eta _c^{}}=3.6\mathrm{GeV},M_{h_c(^1P_1)}=3.5\mathrm{GeV},m_c=1.5\mathrm{GeV},`$ (16)
$`<0|𝒪_8^{h_c}(^1S_0)|0>=0.98\times 10^2\mathrm{GeV}^5\text{[14]},<0|𝒪_1^{h_c}(^1P_1)|0>=0.32\mathrm{GeV}^5\text{[15]},`$ (17)
$`<0|𝒪_8^\psi ^{}(^1S_0)|0>+{\displaystyle \frac{7}{m_c^2}}<0|𝒪_8^\psi ^{}(^3P_0)|0>=0.56\times 10^2\mathrm{GeV}^3\text{[16]},`$ (18)
$`<0|𝒪_1^\psi ^{}(^3S_1)|0>=0.44\mathrm{GeV}^3\text{[17]},<0|𝒪_8^\psi ^{}(^3S_1)|0>=6.2\times 10^3\mathrm{GeV}^3\text{[17]},`$ (19)
$`<0|𝒪_1^{\eta _c^{}}(^1S_0)|0>=0.20\mathrm{GeV}^3\text{[18]},`$ (20)
and the CTEQ 3M package for parton distributions is employed with the factorization scale chosen to be equal to the NRQCD scale $`\mu =2m_c`$. In making use of the present fitted matrix elements given in above, the spin symmetry relation $`<0|𝒪_8^\psi ^{}(^3P_J)|0>=(2J+1)<0|𝒪_8^\psi ^{}(^3P_0)|0>`$ has been applied.
With 920 GeV incident proton we find the magnitude of the cross sections given in the preceding section are
$`\sigma _1=1076.1\mathrm{nb}/\mathrm{n},\sigma _2=98.9\mathrm{nb}/\mathrm{n},\sigma _3=54.8\mathrm{nb}/\mathrm{n},\sigma _4=79.0\mathrm{nb}/\mathrm{n},\sigma _5=5.2\mathrm{nb}/\mathrm{n}.`$ (21)
Here, the nb/n means nb/nucleon for shorthand. The $`\psi ^{}`$ production cross section (84.2 nb/n) agrees well with the experimental measurement of $`(75\pm 5\pm 22)`$ nb/n by E789 , indicating the reliability of the other calculations in this paper. However, quarkonium production rates are often sensitive to the choice of $`m_c`$ and the parton distributions. To see the effect of the former, we assume the difference between calculated and measured $`\psi ^{}`$ production cross sections is a pure effect of $`m_c`$, to cover the error of the measured value, $`m_c`$ should vary from $`1.45`$ to $`1.65`$ GeV. By changing $`m_c`$ from $`1.5`$ to $`1.45`$ and $`1.65`$ GeV in all other cross section calculations, the relative uncertainties of the $`\sigma `$s are shown below. As for the latter, we simply take another parton distribution functions, the GRV , the deviations of the $`\sigma `$s are also listed below.
$`\mathrm{\Delta }\sigma _1=_{44.3}^{+40.1}+\mathrm{\hspace{0.17em}28.5}\%,\mathrm{\Delta }\sigma _2=_{46.6}^{+24.3}+\mathrm{\hspace{0.17em}28.5}\%,\mathrm{\Delta }\sigma _3=_{55.4}^{+32.6}+\mathrm{\hspace{0.17em}28.8}\%,`$ (22)
$`\mathrm{\Delta }\sigma _4=_{47.7}^{+25.3}+\mathrm{\hspace{0.17em}29.2}\%,\mathrm{\Delta }\sigma _5=_{41.4}^{+20.5}\mathrm{\hspace{0.17em}5.8}\%.`$ (23)
Here, the first deviations come from the the change of $`m_c`$ (”+” for $`m_c=1.45`$ GeV and ”-” for $`m_c=1.65`$ GeV); the second corresponds to the choice of a different parton distribution code (GRV results relative to the CTEQ ones); We can see from the above results that the deviations of the cross sections relative to different parton distributions agree within 30%, and the charm quark mass uncertainty changes cross sections around 50%.
Due to the projected high interaction rate, 40 MHz, the results in Eq. (21) means that in a running time of $`10^7`$s at HERA-$`B`$ using the $`Cu`$ target, for example, the directly produced $`\eta _c^{}`$ and $`h_c(^1P_1)`$ events number would be about $`3.3\times 10^{10}`$ and $`4.7\times 10^9`$. The $`\psi ^{}`$ events number would be about $`2.6\times 10^9`$, which is three orders higher than the present $`\psi ^{}`$ date sample collected at $`e^+e^{}`$ colliders.
Theoretical estimation of the branching fractions of the $`h_c(^1P_1)`$ production in $`\psi ^{}`$ decays are about $`10^{53}`$ from Refs. , and the $`\eta _c^{}`$ rate are about $`10^{43}`$ from the naive estimation of the M1 transition in non-relativistic limit . Therefore, the indirectly produced $`h_c(^1P_1)`$ and $`\eta _c^{}`$ would be of the order $`10^{46}`$ and $`10^{56}`$ correspondingly.
The indirect production of Charmonium in $`B`$ decays has been estimated in Ref. , the production rates of $`h_c(^1P_1)`$ and $`\eta _c^{}`$ (assuming the same as that of $`\eta _c`$) are of the order of $`10^3`$. Using the $`b\overline{b}`$ production cross section of 12 nb/n, the produced $`h_c(^1P_1)`$ and $`\eta _c^{}`$ events are of the order of $`10^{56}`$ in $`10^7`$s of the HERA-$`B`$ running time, which is the same order as via $`\psi ^{}`$ decays.
## IV Searching Strategy
As mentioned in the introduction part of this paper, the interested topologies of detecting these two states at HERA-B are $`\gamma J/\psi `$ and $`\pi ^0J/\psi `$ for $`\eta _c^{}`$ and $`h_c(^1P_1)`$ respectively, where $`J/\psi `$ decays into lepton pairs and $`\pi ^0`$ decays to two photons. Because of the charge-conjugation invariance, the decay modes $`\eta _c^{}\pi ^0J/\psi `$ and $`h_c\gamma J/\psi `$ are ruled out.
The $`h_c(^1P_1)`$ state was observed decaying to $`\pi ^0J/\psi `$ with branching ratio $`10^3`$ , which is of the same order of magnitude as the theoretical expectation , and the $`\eta _c^{}`$ decaying to $`\gamma J/\psi `$ is expected with a width of the order $`𝒪(1k\text{eV})`$ . Considering that the theoretical estimation of the decay width of the $`\eta _c^{}`$ is about 5 MeV, it has a branching ratio of $`𝒪(10^4)`$ in $`\gamma J/\psi `$ decay mode. Using the numbers listed above, TABLE I lists the estimation of produced events for $`h_c(^1P_1)`$ and $`\eta _c^{}`$ in all the production mechanisms, taken into account the branching ratios of $`J/\psi `$ leptonic decays and $`\pi ^0\gamma \gamma `$.
From the table, we can see the produced events of interested topologies from indirect productions are too low (of the order of 10 to 100) to produce meaningful signals for observing the two states. But instead, the direct productions of these two states are rather large, of the order of $`46\times 10^5`$. As we know the geometric acceptance of HERA-B detector is large and its trigger is optimized for $`J/\psi `$ events, we do expect high efficiency of detecting these two final states. Suppose the overall efficiency of detecting these two final states is around $`10`$%, one expects $`46\times 10^4`$ reconstructed events each channel, which are large numbers compared to those channels for observing CP violation (in the same running time, the reconstructed events of $`J/\psi K_s`$ is estimated to be around 1400!).
The main background channel for $`\eta _c^{}`$ observation is $`\chi _{c2}\gamma J/\psi `$, which has the same final states but much larger cross section and very near the expected $`\eta _c^{}`$ mass. Using the measured cross section of $`\chi _{c2}`$ by E771 , the number of reconstructed $`\chi _{c2}`$ events is estimated to be around $`10^8`$ (the combinational background at $`\chi _{c2}`$ mass region is about the same size as $`\chi _{c2}`$ events as shown in Ref. ). The significance of the observed $`\eta _c^{}`$ depends strongly on the mass resolution of $`\gamma J/\psi `$ system and the mass difference between $`\chi _{c2}`$ and $`\eta _c^{}`$. Theoratical estimations of the $`\eta _c^{}`$ mass ranges from 3589 to 3631 MeV , and only experimental hint is at mass of $`(3594\pm 5)`$ MeV. For a 3.6 GeV mass $`\eta _c^{}`$, if the mass resolution is around 10 MeV or less, $`\eta _c^{}`$ will produce a long tail at high mass side of $`\chi _{c2}`$, and at mass higher than 3.6 GeV, the events is almost free from $`\chi _{c2}`$ background. If the mass resolution reaches 15 MeV or even larger, it will be hard to distinguish $`\eta _c^{}`$ from $`\chi _{c2}`$. A larger $`m_{\eta _c^{}}`$ obviously will increase the possibility of resolving $`\eta _c^{}`$ from the $`\chi _{c2}`$ tail, while a low mass $`\eta _c^{}`$ will more depend on the mass resolution.
For $`h_c(^1P_1)`$, the main background is from the $`\pi ^0\pi ^0J/\psi `$ produced by $`\psi ^{}`$ decays. Compared with that in $`\eta _c^{}`$ case, here the $`h_c(^1P_1)`$ is at the phase space limit of $`\pi ^0J/\psi `$ system produced from $`\psi ^{}`$ decays and the cross section of the latter is smaller than $`\chi _{c2}`$ by at least a factor of 3%. Furthermore, there is no other nearby resonance decays to the same final states. All these make the observation of $`h_c(^1P_1)`$ easier than $`\eta _c^{}`$.
At the point of data analysis, for $`\eta _c^{}`$, instead of using the invariant mass of $`J/\psi `$ and the detected $`\gamma `$, using the mass difference between the $`\mathrm{}^+\mathrm{}^{}\gamma `$ system and the $`\mathrm{}^+\mathrm{}^{}`$ system would be better in finding the signal, since the latter can compensate some of the effects due to energy losses of radiation and bremsstrahlung of the lepton tracks.
In searching for the $`h_c(^1P_1)`$ state, the reconstruction of the $`\pi ^0`$ is also important for the event selection, and it is also a very good constraint to lower the background level greatly. As in the $`\gamma J/\psi `$ case, the mass difference method will be helpful to this channel as well.
Finally, to check the results, the sideband method maybe useful. In both cases the $`J/\psi `$ mass sidebands, and in $`h_c(^1P_1)`$ searching the $`\pi ^0`$ mass sidebands will tell us the shape of the background. The absence of the same peak in the mass spectrum of sidebands events will be a demonstration that the selection is reasonable.
It is important to note that all above discussions are based on a sample of $`10^7`$s running time. With more statistics, instead of reconstructing photon from ECAL, one can detect converted photon to reconstruct $`\eta _c^{}`$ and $`h_c(^1P_1)`$, as has been indicated by E771 . In this case, the mass resolution will be significantly improved ($`5.2\pm 2.0`$ MeV for $`\gamma J/\psi `$ system in E771 experiment), $`\eta _c^{}`$ will be resolved from $`\chi _{c2}`$ even it has a small mass.
## V Discussions and Conclusions
In this paper we have discussed the physics potential of HERA-$`B`$ in detecting the $`\eta _c^{}`$ and $`h_c(^1P_1)`$. Our numerical results reveal that there are about $`10^{10}`$ and $`10^9`$ of $`\eta _c^{}`$ and $`h_c(^1P_1)`$ events would be produced at HERA-$`B`$ in $`10^7`$s of running time. A rough estimation shows that $`h_c(^1P_1)`$ will be observed clearly in its $`\pi ^0J/\psi `$ decay mode, and $`\eta _c^{}`$ will be observed as a shoulder at high mass side of $`\chi _{c2}`$ in $`\gamma J/\psi `$ channel if the mass resolution is not too large.
The searching strategies of these two states at HERA-$`B`$ are given. The major backgrounds in the detection and the possible detecting measures are also discussed.
It should be mentioned that the theoretical basement of our calculation in this paper, the NRQCD factorization, may not work well in the inclusive quarkonium production at full phase space, that is at small $`p_T`$ ($`p_T`$ not much greater than $`\mathrm{\Lambda }_{\mathrm{QCD}}`$) region, which would cast some shadow on the validity of the results of the inclusive fix-target calculations. However, at least from our calculation on $`\psi ^{}`$ production, which agrees with the experiment value quite well, we are convinced to a certain degree of our other calculations in this paper.
Last, it should be noticed that although the study proceeded in the paper is just an order estimation because either input parameters, like the color-octet matrix elements, are more or less accurate just to an order, or the evaluation is based only on the first order calculation, or the factorization problem mentioned above, the results were well constrained by the known measurement of $`\psi ^{}`$ production, so the conclusion of the paper should hold. That is, the detection of $`\eta _c^{}`$ and $`h_c(^1P_1)`$ at HERA-$`B`$ is feasible and promising.
###### Acknowledgements.
C.-F. Q. thanks the Alexander von Humboldt Committee for financial support; C. Z. Y. thanks Prof. H. Kolanoski, Prof. C. H. Jiang and Prof. M. Davier for helpful discussions and comments.
FIGURE CAPTIONS
Figure 1. The leading order $`\eta _c^{}`$ production process at $`PN`$ collision in both $`\alpha _s`$ and $`v^2`$.
Figure 2. The generic diagrams of $`h_c(^1P_1)`$ production process at $`PN`$ collision at leading order in $`v^2`$; (a) the Color-Octet process, (b) the Color-Singlet process.
Figure 3. The generic diagrams of $`\psi ^{}`$ production process at $`PN`$ collision; (a) and (c) the Color-Octet processes at $`v^4`$ and leading order in $`\alpha _s`$, (b) the leading order Color-Singlet process in both $`\alpha _s`$ and $`v^2`$.
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# Automatic Loop Calculations with FeynArts, FormCalc, and LoopTools
## 1 Introduction
The precision of experimental data achieved at present colliders has in many cases reached or exceeded the per cent level. Obviously a comparable accuracy on the theoretical side is needed in order to draw significant conclusions from such precise measurements. For many observables this means that a one-loop calculation is the lowest acceptable approximation.
Doing one-loop calculations by hand is laborious and error-prone and in some cases simply impossible. So for some time already, software packages have been developed to automate these calculations (e.g. ). Incidentally, full automation is possible only up to one loop since no algorithms generic enough for the computation of arbitrary multi-loop Feynman diagrams are known at present. One remaining obstacle is that the existing packages generally tackle only part of the problem, and one still has to spend considerable effort adapting conventions etc. to make them work together.
In this paper the three Mathematica packages FeynArts, FormCalc, and LoopTools are presented which work hand in hand. The user has to supply only small driver programs whose main purpose is to specify the necessary input parameters. This makes the whole system very “open” in the sense that the results are returned as Mathematica expressions which can easily be manipulated, e.g. to select or modify terms.
Since one-loop calculations can range anywhere from a handful to several hundreds of diagrams (particularly so in models with many particles like the MSSM), speed is an issue, too. FormCalc, the program which does the algebraic simplification, therefore uses FORM for the time-consuming parts of the calculation. Owing to FORM’s speed, FormCalc can process, for example, the 1000-odd one-loop diagrams of W–W scattering in the Standard Model in about 5 minutes on an ordinary Pentium PC.
The following sections describe the main functions of each program. Furthermore, the FormCalc package contains two sample calculations in the electroweak Standard Model, $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}`$ and $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{e}$}^{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{t}$}`$ , which demonstrate how the programs are used together.
## 2 FeynArts
FeynArts is a Mathematica package for the generation and visualization of Feynman diagrams and amplitudes. The current version 2.2 is a much-expanded version of FeynArts 1 . The two most important new features are the generation of counter-term diagrams and the ability to deal with supersymmetric theories (cf. Sect. 5).
FeynArts works in three basic steps, sketched in Fig. 2. The first step is to create all different topologies for a given number of loops and external legs. For example, to create all one-loop topologies for a $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ process, the following call to CreateTopologies is used:
```
top = CreateTopologies[1, 1 -> 2]
```
In the second step, the actual particles in the model have to be distributed over the topologies in all allowed ways. E.g. the diagrams for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}\overline{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}}`$ are produced with
```
ins = InsertFields[ top,
V[2] -> {F[4, {3}], -F[4, {3}]} ]
```
where F\[4, {3}\] is the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$-quark, -F\[4, {3}\] its antiparticle, and `V[2]` the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{Z}$}`$ boson.
The fields, their propagators, and their couplings are defined in a special file, the model file, which the user can supply or modify. The following model files are included in FeynArts: the electroweak Standard Model (SM.mod) , the same including QCD (SMQCD.mod), and in the background-field formulation (SMbgf.mod). These model files all include the full set of one-loop counter terms. A model file for the MSSM is in preparation.
The diagrams can be drawn with `Paint[ins]`, depending on the options either on screen, or in a PostScript or file.
Finally, the analytic expressions for the diagrams are obtained by
```
amp = CreateFeynAmp[ins]
```
An important feature of FeynArts is that it distinguishes three levels of fields:
Generic level, e.g. the fermion F,
Classes level, e.g. the down-type quark F,
Particles level, e.g. the $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{b}$}`$-quark F\[4, {3}\].
This is useful for two reasons:
The kinematic structure of a coupling is fixed once the generic fields are specified. For example, all fermion–fermion–scalar couplings are of the form
$$\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{F}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{S}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\omega }$}_{\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{G}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\omega }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{+}$}$$
where $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\omega }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\pm }$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\gamma }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{5}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{/}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ are the chirality projectors. This means that most algebraic simplifications like the tensor reduction need to be carried out on the Generic-level amplitude only.
Furthermore, it is more economic to perform index summations (e.g. over the fermion-generation index) in a loop over Classes-level amplitudes instead of generating many Particles-level diagrams.
## 3 FormCalc
FeynArts produces very symbolic output which cannot straightforwardly be implemented in a numerical program. Its evaluation proceeds instead in two steps: first, algebraic simplification in Mathematica; then, translation into a Fortran program which computes the squared matrix element and from this the desired quantities (cross-sections, decay rates, asymmetries, etc.).
### 3.1 Algebraic simplification
The symbolic expressions for the diagrams are simplified algebraically with FormCalc which returns the results in a form well suited for numerical evaluation. Specifically, FormCalc performs the following simplifications:
– indices are contracted as far as possible,
– fermion traces are evaluated,
– open fermion chains are simplified using the
Dirac equation,
– colour structures are simplified using the
SU($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$) algebra,
– the tensor reduction is done,
– the results are partially factored,
– abbreviations are introduced.
The internal structure of FormCalc is simple: it prepares the symbolic expressions of the diagrams in an input file for FORM, runs FORM, and retrieves the results via the MathLink interface (see Fig. 3). This is done completely without intervention by the user, i.e. the user does not see the FORM code. FormCalc thus combines the speed of FORM with the powerful instruction set of Mathematica and the latter greatly facilitates further processing of the results.
The main function in FormCalc is OneLoop (the name is not strictly correct since it works also with tree graphs). It is used like this:
```
<< FormCalc‘
amps = << myamps.m
result = OneLoop[amps]
```
where it is assumed that the file myamps.m contains amplitudes generated by FeynArts. Note that OneLoop needs no declarations of the kinematics of the underlying process; it uses the information FeynArts hands down.
Even more comprehensive than OneLoop, the function ProcessFile can process entire files. It collects the diagrams into blocks such that index summations (e.g. over fermion generations) can later be carried out easily, i.e. only diagrams which are summed over the same indices are put in one block. This nicely complements the generation of Classes-level diagrams in FeynArts, which leaves the index summations to the numerical evaluation in order to reduce the number of diagrams. ProcessFile is invoked e.g. as
```
ProcessFile["vertex.amp", "vertex"]
```
which reads the FeynArts amplitudes from the input file vertex.amp and produces output files of the form vertexid.m, where id is some identifier for a particular block.
The output of OneLoop or ProcessFile is in general a linear combination of loop integrals with prefactors that contain model parameters, kinematic variables, and abbreviations introduced by FormCalc. Such abbreviations are introduced for spinor chains, scalar products of vectors, and epsilon tensors contracted with vectors. A term in the output could for instance look like
```
C0i[cc0, MW2, MW2, S, MW2, MZ2, MW2] *
( -4 Alfa2 CW2 MW2/SW2 S AbbSum16 +
32 Alfa2 CW2/SW2 S$`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\text{2}}$}`$ AbbSum28 +
4 Alfa2 CW2/SW2 S$`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\text{2}}$}`$ AbbSum30 -
8 Alfa2 CW2/SW2 S$`^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\text{2}}$}`$ AbbSum7 +
Alfa2 CW2/SW2 S (T - U) Abb1 +
8 Alfa2 CW2/SW2 S (T - U) AbbSum29 )
```
The first line represents the one-loop integral $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{C}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{0}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{(}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{W}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{W}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{s}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{W}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{Z}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{M}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{W}}$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{)}$}`$, which is multiplied with a linear combination of abbreviations like Abb1 or AbbSum29 with certain coefficients. These coefficients contain kinematical variables like the Mandelstam variables S, T, and U and model parameters, e.g. $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{𝙰𝚕𝚏𝚊𝟸}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{=}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\alpha }$}^\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$.
The automatic introduction of abbreviations is a very important feature which can drastically reduce the size of an amplitude, particularly so because the abbreviations are nested in three levels. Here is an example:
The definitions of the abbreviations can be retrieved by Abbreviations\[\] which returns a list of rules such that result //. Abbreviations\[\] gives the full, unabbreviated expression.
### 3.2 Translation to Fortran code
For numerical evaluation, the Mathematica expressions produced by FormCalc need to be translated into a Fortran program. (They could, in principle, be evaluated in Mathematica directly, but this becomes rather slow for large amplitudes.) The translation is done by the program NumPrep, which is part of the FormCalc package. The philosophy of NumPrep is that the user should not have to modify the generated code. This means that the code has to be encapsulated (i.e. no loose ends the user has to bother with), and that all necessary subsidiary files (include files, makefile) have to be produced, too.
From the point of view of the Fortran programmer who wants to use the generated code in his program, the output of NumPrep is a single subroutine called
squared\_me($`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{k}$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\epsilon }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\lambda }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{1}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mathrm{}}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{,}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\lambda }$}_\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{N}$}`$)
which takes as input the external momenta, polarization vectors, and helicities, and returns the squared matrix element. To obtain actual numerical results from the generated code, one needs in addition a driver program whose task is to initialize the model parameters, set up the kinematics, invoke the squared\_me subroutine, perform necessary phase-space integrations, and finally write the results to a file. A sample driver program for $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{}$}\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{2}$}`$ processes is included in FormCalc.
Finally, the generated code has to be linked with the LoopTools library which provides the one-loop functions.
## 4 LoopTools
LoopTools supplies the actual numerical implementations of the one-loop integrals needed for programs made from the FormCalc output. It is based on the reliable package FF and provides in addition to the scalar integrals of FF also the tensor coefficients in the conventions of . LoopTools offers three interfaces: Fortran, C++, and Mathematica, so most programming tastes should be served.
Using the LoopTools functions in Fortran and C++ is very similar. In Fortran it is necessary to include the file looptools.h in every function or subroutine (for the common blocks). In C++, clooptools.h must be included once. Before using any LoopTools function, ffini must be called and at the end of the calculation ffexi may be called to obtain a summary of errors. It is of course possible to change parameters like the scale $`\colorbox[rgb]{1,1,1}{$\textcolor[rgb]{0,0,0}{\mu }$}`$ from dimensional regularization; this is described in detail in the manual .
A very simple Fortran program would for instance be
```
program simple
#include "looptools.h"
call ffini
print *, B0(1000D0,50D0,80D0)
call ffexi
end
```
The C++-version of this program is
```
#include "clooptools.h"
```
```
main()
{
ffini();
cout << B0(1000.,50.,80.) << "\n";
ffexi();
}
```
The Mathematica interface is even simpler to use:
```
In[1]:= Install["LoopTools"]
```
```
In[2]:= B0[1000, 50, 80]
```
```
Out[2]= -4.40593 + 2.70414 I
```
## 5 Calculations in Supersymmetric Models
Special emphasis has been placed on the possibility to do calculations in supersymmetric models with FeynArts and FormCalc. In particular the following two fundamental problems become relevant in supersymmetric theories:
Problem 1: SUSY theories in general contain Majorana fermions and hence fermion-number-violating couplings (e.g. quark–squark–gluino). The textbook prescription of ordering the Dirac matrices opposite to their occurrence along the arrows on fermionic lines obviously breaks down in this case since one cannot define a fermion-number flow. (Put differently, Majorana-fermion lines have no arrow.)
Solution: FeynArts uses the “flipping-rule” algorithm : instead of traversing the fermion lines along the fermion-number flow imposed from the outside, it chooses a particular direction for each fermion chain. If it later turns out that, for a Dirac fermion, the chosen direction is opposite to the actual fermion flow, it applies a so-called flipping rule.
Problem 2: Dimensional regularization, the default regularization scheme employed by FormCalc, is known to break supersymmetry .
Solution: FormCalc has two regularization schemes built in which are chosen with the variable $Dimension. The default is $Dimension = D which corresponds to dimensional regularization. Putting $Dimension = 4 switches to constrained differential renormalization . The latter technique is equivalent at the one-loop level to regularization by dimensional reduction and is hence suited for calculations in SUSY models.
## 6 Requirements and Availability
All three packages require Mathematica 2.2 or above; FormCalc needs in addition FORM, preferably version 2 or above; LoopTools needs a Fortran compiler and gcc/g++.
The packages should compile and run without change on any Unix platform. They are specifically known to work under DEC Unix, HP-UX, Linux, Solaris, and AIX. A comprehensive manual which explains installation and usage is included in each package. All three packages are open-source programs and stand under the GNU library general public license. They are available from
http://www.feynarts.de
http://www.feynarts.de/formcalc
http://www.feynarts.de/looptools
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# Detecting the First Objects in the Mid-IR with NGST
## 1 Introduction
Detecting the first luminous objects in the universe will be the primary goal of several future space- and ground-based telescopes. The importance of these experiments consists in the fact that they could be able to test current cosmological scenarios, study the properties of these (supposedly very small) galaxies, and their effects on the surrounding environment (as for example reionization, heating and metal enrichment of the IGM). Due to the predicted low luminosity of such objects, this task will be at the capability edge of even the most advanced and powerful instruments. Apart from the indirect probes of their effects, as for example the secondary anisotropies in the CMB left by reionization (Knox, Scoccimarro & Dodelson 1998; Gruzinov & Hu 1998; Bruscoli et al. 2000; Benson et al. 2000), some search strategies have already been suggested in the literature. Marri & Ferrara (1998) and Marri, Ferrara & Pozzetti (2000) have suggested that Type II supernovae occurring in the first objects might outshine their parent galaxy by more than a hundred times and become visible by instruments like the Next Generation Space Telescope (NGST). Schneider et al. (2000) have investigated the possibility of detecting gravitational wave emission from high-redshift very massive objects with LISA. Oh (1999) proposes direct imaging of the ionized halos around primordial objects either via their free-free emission (possibly detectable with the Square Kilometer Array) or Balmer line emission again with NGST.
One point that is particularly important when dealing with these low mass systems is that they are strongly affected by feedback mechanisms both of radiative and stellar type; these have been extensively investigated by Ciardi et al. (2000). Molecular hydrogen, being the only available coolant in a plasma of nearly primordial composition, is a key species in the feedback network as it regulates the collapse and star formation in these objects. Ferrara (1998, F98) pointed out that H<sub>2</sub> is efficiently formed in cooling gas behind shocks produced during the blowaway (i.e. the complete ejection of the galactic gas due to SN explosions) process thought to occur in the first objects, with typical H<sub>2</sub> fractions $`f_{H_2}6\times 10^3`$. We shall see that the conditions in these cooling blastwaves are such that a noticeable amount of the explosion energy is carried away by infrared (redshifted into the Mid-IR \[MIR\] spectral region) H<sub>2</sub> molecular lines, which therefore might provide us with a superb tool to detect and trace these very distant primordial galactic blocks.
The use of molecular lines as diagnostics of moderate redshift ($`z\stackrel{<}{}3`$) galaxies has already been proposed by some authors (Frayer & Brown 1997; Blain et al. 2000); at higher redshifts pioneering calculations were carried on by Shchekinov & Éntel’ (1985); more recently Silk & Spaans (1997) concentrated on the CO and dust emission from HII regions inside larger galaxies. All these studies have emphasized the power of the molecular line emission as a probe of distant sources. Based on previous calculations (F98), in this paper we calculate MIR fluxes and number counts for these sources in various bands and assess if forthcoming observing facilities will be able to eventually unveil the beginning of the cosmic star formation era and study the feedback processes in the young universe.
## 2 Emission model
Due to its symmetry, the H<sub>2</sub> molecule has no electric dipole moment in the ground state. Therefore, the first detectable H<sub>2</sub> emission lines are produced by quadrupole radiation and they are purely rotational. The line properties of interest to the present work, both for rotational and roto-vibrational lines, are listed in Table 1. There we see that the excitation temperatures, $`T_{ex}`$, of these lines fall in the range 500-7000 K. The temperature range spanned during the cooling of the post-shock IGM gas produced during the blowaway of low-mass primordial galaxies is (F98):
$$300\mathrm{K}\stackrel{<}{}T\stackrel{<}{}2.3M_6^{2/5}(1+z)^{18/5}\mathrm{K}$$
(1)
where $`M=10^6M_6M_{}`$ is the total mass of the galaxy.<sup>1</sup><sup>1</sup>1 We adopt a $`\mathrm{\Lambda }`$CDM (cluster normalized) cosmology with $`\mathrm{\Omega }_M=0.35,\mathrm{\Omega }_\mathrm{\Lambda }=0.65,\mathrm{\Omega }_b=0.04,h=0.65`$. The upper limit is set by the condition that the cooling time of the post-shock gas is shorter than the Hubble time. Strictly speaking the above upper limit holds until the IGM cooling is dominated by the inverse Compton radiative losses ($`z\stackrel{>}{}6`$), but we use it down to $`z=4`$, where it still represents a very good approximation for our purposes. Given the range in eq. 1, it is then conceivable that the above molecular lines are excited during the process in which a cold shell forms, thus producing potentially detectable radiation. The flux observed on the ground in a given line is derived as follows.
Let $`\nu _{ik}`$ be the restframe frequency of photons emitted by the molecule during the transition between the energy levels $`k`$ and $`i`$; then the line emissivity (erg s<sup>-1</sup> cm<sup>-3</sup> sr<sup>-1</sup>) is (Spitzer 1978):
$$j_{\nu _{ik}}=\frac{h}{4\pi }\nu _{ik}n_kA_{ki},$$
(2)
where $`n_k`$ is the number density of molecules in the $`k`$ level and $`A_{ki}`$ is the Einstein coefficient for spontaneous emission. The typical H<sub>2</sub> densities found are much lower than the critical one ($`10^4`$ cm<sup>-3</sup>), hence we neglect collisional de-excitations.
As the conditions for thermodynamic equilibrium are not satisfied, the population of the various levels must be obtained by solving the detailed balance equations (Spitzer 1978). Using this approach, the number density of molecules in the vibrational level $`v`$, $`n_J`$, of even and odd rotational levels decouple, and they are obtained by iteration using the following formula:
$$n_{J+2}(v)=n_J(v)g_J\frac{\gamma _J\mathrm{e}^{\mathrm{\Delta }E/kT}}{1+\gamma _J},$$
(3)
where $`g_J`$ is the statistical weight of level $`J`$, $`\gamma _J`$ is the ratio between the collisional excitation rate and the Einstein $`A`$ coefficient, $`\mathrm{\Delta }E`$ is the energy difference between the two levels. Obviously, in the limit of large $`\gamma _J`$ eq. 3 approaches the Boltzmann distribution.
Given the post-shock temperatures found (eq. 1), it is only necessary to consider the two vibrational levels $`v=0,1`$. The relative number of molecules in these two levels depends on the total molecular hydrogen density, $`n_{H_2}`$. For a given redshift, the value of $`n_{H_2}`$ (F98) is:
$$n_{H_2}(z)=pf_{H_2}(M,z)n_spf_{H_2}(M,z)\delta n_H(z),$$
(4)
where $`f_{H_2}`$ is the molecular fraction in the shell, and $`n_H(z)`$ is the IGM hydrogen density. We allow for a density enhancement in the shell with respect to $`n_H`$ equal to $`\delta `$; this is produced both by the shock compression and the possible occurrence of the explosion inside an overdense region of the universe (e.g. a cosmological filament). Because of the first effect, and as the expanding blastwave becomes rapidly radiative, $`\delta `$ is given by the square of the shock Mach number, $``$, with respect to the ambient IGM. Following F98, it is easy to show that the density enhancement is given by:
$$\delta =4.8M_6^{2/5}(1+z)^{13/5}.$$
(5)
For simplicity, we have neglected the additional density increase produced by an explosion in a overdense region; thus, the results presented above should be interpreted as a conservative lower limit to the detectability of the sources under investigation. Finally, we assume a branching ratio $`p=0.75`$ (0.25) for ortho (para) transitions which is valid under LTE conditions; however, non-equilibrium conditions might lead to slightly lower values of the ortho-to-para ratio (Chrysostomou et al. 1993; Rodríguez-Fernández et al. 2000).
With the emissivity given by eq. 2, the observed flux (erg s<sup>-1</sup> cm<sup>-2</sup>) is:
$$F(\nu _o)=\frac{L_{\nu _{ik}}(1+z)}{d_L^2},$$
(6)
where $`\nu _o=\nu _{ik}/(1+z)`$ is the observed frequency, $`L_{\nu _{ik}}`$ the luminosity and $`d_L`$ is the cosmological luminosity distance. For a flat universe ($`\mathrm{\Omega }_0=\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }=1`$), $`d_L`$ is given by
$$d_L=(1+z)_0^z(1+z^{})\left|\frac{dt}{dz^{}}\right|𝑑z^{},$$
(7)
$$\left|\frac{dt}{dz}\right|^1=H_0(1+z)\sqrt{(1+\mathrm{\Omega }_Mz)(1+z)^2\mathrm{\Omega }_\mathrm{\Lambda }z(2+z)},$$
(8)
where $`H_0=100h`$ km s<sup>-1</sup> Mpc<sup>-1</sup> is the current Hubble constant. The corresponding luminosity is:
$$L_{\nu _{ik}}=4\pi Vj_{\nu _{ik}},$$
(9)
where $`V=(4/3)\pi R_s^2dR_s`$ is the physical volume occupied by the H<sub>2</sub> forming shell. The values of $`R_s`$ and $`dR_s`$ are obtained from the formulae in F98; the only difference consists in the assumption of a single burst of star formation, as opposed to the quiescent star formation prescription of F98. By assuming a star formation efficiency of 10%, and one SN every 100 $`M_{}`$ of stars formed, it is easy to show that the total energy of the explosion is equal to $`E=1.3\times 10^{53}M_6`$ erg. Then, from F98 it follows:
$$R_s0.88M_6^{1/5}(1+z)^{11/5}\mathrm{Mpc},$$
(10)
$$dR_s=\frac{N_{H_2}}{n_{H_2}}\frac{f_{H_2}M_s}{m_HR_s^2n_{H_2}},$$
(11)
where we have assumed that all the swept IGM mass is in the cool shell of total mass $`M_s`$, together with the galaxy interstellar medium. This result is somewhat dependent on the assumption made to derive the energy of the explosion, while it is not strongly dependent on the cosmology.
## 3 Number Counts
The number of objects whose observed flux is larger than the threshold value $`F_{min}`$ is:
$`N(>F_{min})`$ $`=`$ $`{\displaystyle _{z_{min}}^{z_{max}}}𝑑z{\displaystyle \frac{dV}{dz}}{\displaystyle _{L_{min}}^{L_{max}}}𝑑L{\displaystyle \frac{dn}{dM}}{\displaystyle \frac{dM}{dL}}`$ (12)
$`=`$ $`{\displaystyle _{z_{min}}^{z_{max}}}𝑑z{\displaystyle \frac{dV}{dz}}{\displaystyle _{M_{min}}^{M_{max}}}𝑑M{\displaystyle \frac{dn}{dM}},`$
where $`dn/dM`$ is the comoving number density of halos with masses between $`M`$ and $`M+dM`$ (Press & Schechter 1974) and $`dV/dz`$ is the comoving volume element per unit redshift:
$$\frac{dV}{dz}=\frac{4\pi cd_L^2}{(1+z)}\left|\frac{dt}{dz}\right|,$$
(13)
where $`d_L`$ and $`dt/dz`$ are given in eqs. 7 and 8, respectively. The limits $`z_{min}`$ and $`z_{max}`$ depend on the observational wavelength band; $`L_{max}`$ is the luminosity corresponding to the maximum mass value, $`M_{max}`$, experiencing blowaway at redshift $`z`$ (Ciardi et al. 2000); $`M_{min}=\mathrm{min}[M_H,M(L_{min})]`$. Here $`M_H`$ is the minimum halo mass in which stars can form at a given redshift, which is related to the minimum virial temperature of $`10^4`$ K below which atomic hydrogen cooling of the gas is suppressed and fragmentation into stars is inhibited (Haiman, Rees & Loeb 1997; Ciardi et al. 2000); $`M(L_{min})`$ is the mass of the halo producing the luminosity $`L_{min}`$ corresponding to the observed flux $`F_{min}`$. Finally, the luminosity, $`L`$, is taken as the maximum one reached by the object (usually corresponding to temperatures equal to the excitation ones) that experiences a blowaway at redshift $`z`$, providing that it is emitted at a redshift $`>z_{min}`$. These calculations give lower limits, as they do not include objects blowing away at $`z>z_{max}`$, but having a peak luminosity at $`z<z_{max}`$.
## 4 Results
The above model allows us to determine the observational properties of the MIR emission from the first objects. In particular, we present results concerning the expected flux in the various lines considered (see Table 1) and the corresponding number counts. These are then compared with the sensitivities of various planned instruments.
### 4.1 Expected MIR Fluxes
Fig. 1 shows the expected flux as a function of the observed wavelength, $`\lambda _o`$, in the five different lines 2.12 $`\mu `$m, 6.9 $`\mu `$m, 9.7 $`\mu `$m, 17 $`\mu `$m and 28 $`\mu `$m. As an example, we plot the fluxes calculated for halos of mass $`10^810^9M_{}`$, corresponding to the mean one among those suffering blowaway, as given by the Press-Schechter formalism, at the explosion redshift. The curves refer to emission occurring when the shocked gas temperature has decreased down to values close to the corresponding line excitation temperature. Thus, the actual emission redshift is somewhat lower than the explosion redshift.
From Fig. 1 we see that the typical observed fluxes are in the range $`10^{22}10^{17}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, depending on the wavelength band (and therefore emission redshift). For the 2.12 $`\mu `$m, 6.9 $`\mu `$m and 9.7 $`\mu `$m emission lines the flux is increasing with wavelength, whereas the flux of the 17 $`\mu `$m and 28 $`\mu `$m excitation lines is decreasing towards higher $`\lambda _o`$. From a general point of view this can be understood as follows. As the IGM density increases with redshift, this causes a corresponding increase in the emissivity. In addition, the postshock gas tends to be warmer at high $`z`$, hence inducing preferentially the excitation of the higher excitation temperature lines. These effects are dominant for the 2.12 $`\mu `$m, 6.9 $`\mu `$m and 9.7 $`\mu `$m lines, while they are overwhelmed by a larger value of $`d_L`$ and a smaller shell emission volume for the redder ones, producing a decreasing intensity.
As detection chances are highest for the 2.12 $`\mu `$m line, we have derived its luminosity evolution and the corresponding fraction of the SN energy radiated in this line by objects at different redshifts. As for Fig. 1, the emission is derived for an average halo mass. In Fig. 2a we show the luminosity evolution for selected redshifts in the range $`4z10`$. As the explosion reshift is increased, the peak of the emission is shifted at earlier times and it becomes narrower. Both effects are due to the faster evolution of the temperature driven by the enhanced (i.e. $`(1+z)^4`$) Compton cooling rates. The more pronounced peak occurring at higher $`z`$ is due to the $`(1+z)^3`$ density increase; this trend is partially counterbalanced by the decreasing average halo masses with redshift. Fig. 2b shows the time dependence of the cumulative fraction of the SN mechanical energy $`f_e(t)=E_e(t)/E`$ emitted in the $`\lambda _e=2.12`$ $`\mu `$m line. This quantity increases over a time scale inversely related to the explosion redshift and reaches a plateau, determined by the position of the $`L(t)`$ peak seen in Fig. 2a. At high redshift up to 10% of the SN mechanical energy is carried away by the considered roto-vibrational line; at lower redshifts this value rapidly decreases below 1%.
### 4.2 Number Counts and Detectability
By using the expressions given in eq. 12 we calculate the expected number of objects per unit sky area as a function of the limiting flux, $`F_{min}`$, of a given experiment. The corresponding curves are shown in Fig. 3 for the 2.12 $`\mu `$m H<sub>2</sub> line. The number counts have been estimated in two different MIR bands in the sensitivity range of planned imaging detectors on board of the Next Generation Space Telescope<sup>2</sup><sup>2</sup>2http://augusta.stsci.edu (NGST). For our study, we consider the $`R=3`$ imaging mode of NGST, in the two bands centered at $`\lambda _o=15`$ and 20 $`\mu `$m. In the $`20\mu `$m band, the largest fluxes are at the $`6\times 10^{18}`$ erg s<sup>-1</sup> cm<sup>-2</sup> level, where about $`10^4`$ objects/deg<sup>2</sup> should be seen; the surface density of sources increases rapidly as $`F_{min}`$ decreases. In the $`15\mu `$m band objects are fainter and a rapid drop is seen in the number counts above $`10^{18}`$ erg s<sup>-1</sup> cm<sup>-2</sup>.
In the above bands, NGST will reach a sensitivity of $`4\times 10^{18}`$ erg s<sup>-1</sup> cm<sup>-2</sup> (shown in the Figure) for a $`3\sigma `$ detection in one hour. At the above sensitivity level, and with its field of view of 4’$`\times `$4’, NGST should be able to detect about 200 sources in the $`20\mu `$m band in one hour integration at the $`3\sigma `$ level. These objects are all located in the redshift range $`\mathrm{\Delta }z710`$, a very intriguing time during cosmic evolution, perhaps bracketing the reionization epoch. Assuming a diffraction limited telescope (angular resolution $`0.25\mathrm{"}`$) at this wavelength, field crowding should not be a problem, the number of sources being well below the confusion limit.
Longer exposure times might also allow detections in the 15 $`\mu `$m band. Sensitivities comparable to the NGST ones are also expected for the SPICA HII/L2 future mission<sup>3</sup><sup>3</sup>3http://www.ir.isas.ac.jp.. Other MIR planned facilities, such as SIRTF or FIRST will only reach sensitivities $`10^{14}10^{15}`$ erg s<sup>-1</sup> cm<sup>-2</sup>, depending on observational bands.
In conclusion, NGST appears to constitute the best tool to reveal the emission from the first objects; while other above mentioned instruments will not reach the required sensitivities.
## 5 Conclusions
We have investigated the possibility to detect the first objects in the universe through the MIR line emission associated with the stellar feedback (i.e. blowaway of the gas) induced by their supernovae. A noticeable fraction of the explosion energy (up to 10% in the 2.2$`\mu `$m line alone) is eventually converted and carried away by molecular hydrogen roto-vibrational lines. By using a detailed treatment of the non-equilibrium formation and emission of H<sub>2</sub> molecules behind these cosmological blastwaves we have calculated the observed fluxes and number counts of primordial objects and compared them with the expected sensitivities of future instruments. At the limiting flux of NGST, we do predict that about 200 primordial objects can be detected 20 $`\mu `$m in barely one hour of observation time.
This detection would allow to directly image the first star formation in the universe and our results show that MIR observations, when compared with estimates of previous studies in other bands and/or exploiting different strategies, represent a superb tool for this study. In addition, the proposed excitation mechanism will also allow to test and calibrate the stellar feedback process in the dark ages and finally assess the degree at which the first galaxies (and the IGM) have been influenced by supernova energy deposition. This obviously holds the key for the understanding of the subsequent evolution and formation of larger galactic blocks in most cosmological models. MIR line emission allows the high redshift universe ($`z\stackrel{>}{}5`$) to be much more easily explored than in the near IR bands, where the intergalactic absorption might be found to strongly affect the detectability of objects located close or beyond the reionization epoch: for example, the predicted dust opacity to sources located at redshift $`5`$ is as high as $`0.13`$ at the observed wavelength $`\lambda _01\mu `$m, and could considerably affect observations of the distant universe in that band (Ferrara et al. 1999; Loeb & Haiman 1997). It is worth noting that our predictions constitute lower limits to the number of observable objects as we have not considered possible density enhancements descending from explosions occurring inside overdense regions, as the filaments of the cosmic web (and also because of some details of the calculations, see Sec. 3).
Molecular hydrogen lines might also be excited during the collapse of the so-called PopIII objects, which rely on this molecule to collapse and form stars. However, there are several reasons to suspect that the contribution of this process will be negligible with respect to the emission due to the stellar feedback. First, if a UV background is present, molecular hydrogen can be destroyed by photodissociation. Even if there is no UV background, dissociation of molecular hydrogen by internal UV radiation emitted from massive stars, formed in the high density regions of the objects, is very efficient and the evolution of the objects and star formation are strongly affected (Omukai & Nishi 1999). This regulation effect is efficient as long as the line emission of molecular hydrogen is the main cooling process, i.e. if $`Z\stackrel{<}{}10^2Z_{}`$ and $`T\stackrel{<}{}8000`$ K (Nishi & Tashiro 2000). Second, the mass of these objects is very small and therefore the amount of gas in the required thermodynamic state for the emission is very limited. Interestingly, the relative little importance of PopIII objects in terms of their radiation power in ionizing photons with respect to objects of mass above $`M_H`$ has already been established by Ciardi et al. (2000). The latter objects emit a very large fraction of their binding energy in the hydrogen Ly$`\alpha `$ line which is therefore a much better tracer of their formation (Haiman, Spaans & Quataert 2000).
## Acknowledgments
We are indebted to the referee, R. Schneider, for insightful comments. It is also a pleasure to acknowledge several useful discussions with A. Loeb, R. Maiolino, R. Nishi, F. Palla, N. Scoville, Y. Shchekinov and T. Takeuchi. AF acknowledges support from Ecole Normale Superieure, Paris.
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# Classical relativistic systems of charged particles in the front form of dynamics and the Liouville equation
## 1 Introduction
All fundamental interactions in physics have got gauge nature and demand the use of singular Lagrangians. An adequate method for dealing with such systems was developed by Dirac (see ). His approach, i.e., constrained Hamiltonian mechanics, has been elaborated in many directions and usefully applied to the problems of quantum field theory (see, e.g., ). But an application of this technique to the relativistic statistical mechanics seems rather obscure, although, for example, the thermo field formulation of the condensed matter theory admits a variational approach with the same gauge structure as in quantum electrodynamics and in related theories. We only note the paper , where the use of Faddeev’ measure was suggested for the construction of the relativistic partition function of directly interacting particle system described by means of the constrained Hamiltonian formalism.
The present paper is concerned with the classical relativistic system of point charges coupled with an electromagnetic field. The first attempt to analyse the constraint contents of such a system was made by Dirac . The corresponding action functional in the four-dimensional representation has got two kinds of the gauge freedom. The first is connected with arbitrariness in the parametrisation of particle world lines (chronometrical invariance); the second is generated by the proper gauge transformations of electromagnetic potentials. We reduce the gauge freedom of the first kind by means of the farther great Dirac’s invention, namely, the concept of the forms of relativistic dynamics . Here we use the notion of the form of dynamics as denoting the description of relativistic system, which corresponds to a given global simultaneity relation defined by means of some foliation of the Minkowski space by space-time or isotropic hypersurfaces. Moreover, we consider only the case when the corresponding simultaneity relation is independent of particle or field configurations (cf. ). Using Dirac’s constraint formalism, we try to isolate the gauge degrees of freedom and to formulate the statistical description of the system in a gauge-invariant manner.
The general features of the formalism were presented in within the instant form of dynamics. It will be our purpose in the present paper to explore the further possibilities connected with other forms of relativistic dynamics; in particular, we shall consider the front form description of the charged particle system.
The paper is organised as follows. In Sec. 2 we will first establish the structure of Hamiltonian description of the system of particles plus field in an arbitrary form of relativistic dynamics. Our results are then applied to the front form of dynamics. In Sec. 3 we perform the analysis of the corresponding constraints. Section 4 contains the elimination of the gauge degrees of freedom by suitable canonical transformation. Sections 5 and 6 are devoted to the application in relativistic statistical mechanics. We formulate a Liouville equation for distribution function in the front form of dynamics and write down its equilibrium solution corresponding to classical Gibbs ensemble accounting particle and field degrees of freedom. Peculiarity of the front form of dynamics allow us to perform an integration over field variables in a relativistic partition function. Some calculations of a purely technical nature are collected in appendices.
## 2 Charged particle system in an arbitrary form of dynamics
We shall consider a system of $`N`$ charged particles, which is described by their (time-like) world lines in the Minkowski space-time<sup>1</sup><sup>1</sup>1The Minkowski space-time $`𝕄_4`$ is endowed with a metric $`\eta _{\mu \nu }=\mathrm{diag}(1,1,1,1)`$. The Greek indices $`\mu ,\nu ,\mathrm{}`$ run from 0 to 3; the Latin indices from the middle of alphabet, $`i,j,k,\mathrm{}`$, run from 1 to 3 and both types of indices are subject of the summation convention. The Latin indices from the beginning of alphabet, $`a,b`$, label the particles and run from 1 to $`N`$. The sum over such indices is pointed explicitly. The velocity of light and the Planck constant $`\mathrm{}`$ are equal to unity. addtoresetequationsection
$$\gamma _a:𝕄_4,\tau x_a^\mu (\tau ).$$
(2.1)
An interaction between charges is assumed to be mediated by an electromagnetic field $`F=dA`$ with the electromagnetic potential 1-form
$$A=\stackrel{~}{A}_\mu (x)dx^\mu ;$$
(2.2)
$$F=\frac{1}{2}\stackrel{~}{F}_{\mu \nu }(x)dx^\mu dx^\nu ,\stackrel{~}{F}_{\mu \nu }(x)=_\mu \stackrel{~}{A}_\nu _\nu \stackrel{~}{A}_\mu ,$$
(2.3)
$`_\nu /x^\nu `$. The dynamical properties of such a system are completely determined by the following action functional
$$S=\underset{a=1}{\overset{N}{}}𝑑\tau _a\left\{m_a\sqrt{u_a^2(\tau _a)}+e_au_a^\nu (\tau _a)\stackrel{~}{A}_\nu [x_a(\tau _a)]\right\}\frac{1}{16\pi }d^4x\stackrel{~}{F}_{\mu \nu }(x)\stackrel{~}{F}^{\mu \nu }(x),$$
(2.4)
where $`m_a`$ and $`e_a`$ denote the mass and the charge of particle $`a`$, respectively, $`u_a^\mu (\tau _a)=dx_a^\mu (\tau _a)/d\tau _a`$. We are interested in constructing the Hamiltonian description of the system in a given form of relativistic dynamics. The particle and field degrees of freedom will be treated on equal level.
The form of relativistic dynamics in its geometrical definition is specified by the foliation $`\mathrm{\Sigma }=\{\mathrm{\Sigma }_t|t\}`$ of the Minkowski space-time with space-like or isotropic hypersurfaces
$$\mathrm{\Sigma }_t=\{x𝕄_4|\sigma (x)=t\};$$
(2.5)
$$\eta ^{\mu \nu }(_\mu \sigma )(_\nu \sigma )0.$$
(2.6)
(see ). As it follows from condition (2.6), $`_0\sigma 0`$ and therefore the hypersurface equation (2.5) can be solved with respect to $`x^0`$ in the form:
$$x^0=\phi (t,𝐱),𝐱=(x^1,x^2,x^3).$$
(2.7)
For definiteness we put
$$_0\sigma >0\mathrm{and}\frac{\phi (t,𝐱)}{t}\phi _t(t,𝐱)>0.$$
(2.8)
Making use of the identities
$$_0\sigma |_{\mathrm{\Sigma }_t}=\phi _t^1,_i\sigma |_{\mathrm{\Sigma }_t}=\phi _i\phi _t^1,\phi _i\phi /x^i,$$
(2.9)
it is easy to see that condition (2.6) implies the following inequality for the function $`\phi `$:
$$𝝋^21,𝝋=(\frac{\phi }{x^1},\frac{\phi }{x^2},\frac{\phi }{x^3}).$$
(2.10)
Then one can consider the particle world lines $`\gamma _a`$ as being determined by a set of points $`x_a(t)=\gamma _a\mathrm{\Sigma }_t`$ of intersections with the elements of the foliation $`\mathrm{\Sigma }`$. The parametrical equation of the world line in a given form of dynamics is
$$x^0=x_a^0(t)=\phi (t,𝐱_a(t))\phi _a,x^i=x_a^i(t).$$
(2.11)
The variable $`t`$ serves as a common evolution parameter of the particle system. Therefore, the choice of the form of dynamics is equivalent to the fixing of the parameters $`\tau _a=t`$ of the particle world lines in the reparametrization-invariant action (2.4).
On the other hand, the foliation $`\mathrm{\Sigma }`$ specifies a certain $`3+1`$ splitting of the Minkowski space-time $`𝕄_4\times \mathrm{\Sigma }_0`$, because by the definition of a foliation all hypersurfaces $`\mathrm{\Sigma }_t`$ for different $`t`$ are diffeomorphic to $`\mathrm{\Sigma }_0`$. That splitting, $`f:𝕄_4\times \mathrm{\Sigma }_0`$, is really determined by (2.7):
$$f:(x^0,𝐱)(\phi (t,𝐱),𝐱).$$
(2.12)
Accounting that the determinant of the transformation (2.12) is
$$\frac{(x^0,𝐱)}{(t,𝐱)}=\phi _t(t,𝐱),$$
(2.13)
one can immediately rewrite the action functional (2.4) into the form
$$S=𝑑tL$$
(2.14)
with the Lagrangian function
$$L=\underset{a=1}{\overset{N}{}}\left\{m_a\sqrt{(D\phi _a)^2𝐯_a^2}+e_a\left[A_0(t,𝐱_a)D\phi _a+v_a^iA_i(t,𝐱_a)\right]\right\}+d^3x;$$
(2.15)
$$=\frac{1}{16\pi }\phi _t(t,𝐱)F_{\mu \nu }F^{\mu \nu }.$$
(2.16)
Here $`𝐯_a=d𝐱_a/dt`$, $`D=d/dt=/t+_av_a^i/x_a^i`$, and $`A_\mu =\stackrel{~}{A}_\mu f^1`$, $`F_{\mu \nu }=\stackrel{~}{F}_{\mu \nu }f^1`$. The dynamical variables of our variational problem will be the functions $`𝐱_a(t)`$, $`A^\mu (t,𝐱)`$ and their first order derivatives with respect to the evolution parameter, $`𝐯_a(t)`$ and $`A^\mu {}_{,t}{}^{}(t,𝐱)`$. Using the obvious relations
$$_0\stackrel{~}{A}_\mu =A_{\mu ,t}\phi _t^1,_i\stackrel{~}{A}_\mu =A_{\mu ,i}A_{\mu ,t}\phi _i\phi _t^1,$$
(2.17)
we find the following expression for the field Lagrangian density:
$$=\frac{\phi _t}{8\pi }\left(e_ie^ih_ih^i\right),$$
(2.18)
where
$$e_i=A_{0,i}+\left(A_{i,t}+A_{0,t}\phi _i\right)\phi _t^1,h_i=\epsilon _{ijk}\left(A^{k,j}A^k{}_{t}{}^{}\phi _{}^{j}\phi _t^1\right).$$
(2.19)
The canonical momenta of our problem are given by
$$p_{ai}=\frac{L}{v_a^i}=\frac{m_a(v_{ai}\phi _{ai}D\phi _a)}{\sqrt{(D\phi _a)^2𝐯_a^2}}+e_a\left[A_i(t,𝐱_a)+A_0(t,𝐱_a)\phi _{ai}\right],$$
(2.20)
$$\mathrm{\Pi }_i(t,𝐱)=\frac{}{A^i_{,t}}=\frac{1}{4\pi }\left(e_i\epsilon _{ijk}\phi ^jh^k\right),$$
(2.21)
$$\mathrm{\Pi }(t,𝐱)=\frac{}{A_{0,t}}=\frac{1}{4\pi }e_i\phi ^i.$$
(2.22)
The basic Poisson brackets are
$`\{x_a^i,p_{bj}\}=\delta _{ab}\delta _j^i,`$
$`\{A_i(t,𝐱),\mathrm{\Pi }_j(t,𝐲)\}=\delta _{ij}\delta ^3(𝐱𝐲),`$ (2.23)
$`\{A_0(t,𝐱),\mathrm{\Pi }(t,𝐲)\}=\delta ^3(𝐱𝐲);`$
all other brackets vanish. Equations (2.21) and (2.22) imply the primary constraint
$$\xi \mathrm{\Pi }+\mathrm{\Pi }^i\phi _i=0.$$
(2.24)
In the isotropic forms of dynamics, however, some additional primary constraints may occur. To see that, we rewrite equation (2.21) into the form
$$\mathrm{\Pi }_i(t,𝐱)=\frac{1}{4\pi }\left(𝔖_i+\beta _{ij}A^j{}_{,t}{}^{}\phi _{t}^{1}+A_{0,t}\phi _i\phi _t^1\right),$$
(2.25)
where $`𝔖_i`$ does not contain the derivatives with respect to the evolution parameter $`t`$,
$$𝔖_i=A_{0,i}\phi ^jh_{ij};h_{ij}A_{j,i}A_{i,j},$$
(2.26)
and
$$\beta _{ij}=(1𝝋^2)\delta _{ij}+\phi _i\phi _j.$$
(2.27)
Upon using the identity $`\beta _{ij}\phi ^j=\phi _i`$, one obtains equation
$$\beta _{ij}\left(A^j{}_{,t}{}^{}+A_{0,t}\phi ^j\right)=\phi _t(4\pi \mathrm{\Pi }_i𝔖_i).$$
(2.28)
The matrix $`\beta =\beta _{ij}`$ has determinant $`det\beta =(1𝝋^2)^2`$.
In the case of space-like forms of dynamics ($`𝝋^2<1`$) the matrix is nonsingular and possesses an inverse
$$\gamma _{ij}=\frac{\delta _{ij}\phi _i\phi _j}{1𝝋^2},\beta ^{ij}\gamma _{jk}=\delta _k^i,$$
(2.29)
so that equation (2.28) can be solved as
$$A_{i,t}+A_{0,t}\phi _i=\gamma _{ij}\phi _t\left(4\pi \mathrm{\Pi }^j𝔖^j\right).$$
(2.30)
For isotropic forms of dynamics we have $`\beta _{ij}=\phi _i\phi _j`$ and (2.28) implies constraints
$$\xi _i\epsilon _{ijk}\phi ^j\left(4\pi \mathrm{\Pi }^k𝔖^k\right)=0.$$
(2.31)
Because it holds identically $`\phi ^i\xi _i=0`$, really only two additional primary constraints occur.
Next, we consider the canonical Hamiltonian of our system, which is defined as
$$H_\mathrm{c}=\underset{a=1}{\overset{N}{}}p_{ai}v_a^i+d^3x\left(\mathrm{\Pi }^iA_{i,t}\mathrm{\Pi }A_{0,t}\right)L.$$
(2.32)
The immediate calculations give
$$H_\mathrm{c}=\underset{a=1}{\overset{N}{}}H_a(t,𝐱_a,𝐩_a)+d^3x_\mathrm{c}$$
(2.33)
with
$$H_a=\frac{m_a\phi _{at}D\phi _a}{\sqrt{(D\phi _a)^2𝐯_a^2}}e_a\phi _{at}A_0(t,𝐱_a)$$
(2.34)
and
$$_\mathrm{c}=\frac{1}{4\pi }\left[\frac{1}{2}\phi _t\left(e_ie^i+h_ih^i\right)+e^iA_{0,i}\epsilon _{ijk}A^i{}_{,t}{}^{}\phi _{}^{j}h^k\right].$$
(2.35)
The problem of solution of (2.20) with respect to the particle velocities is analogous to the case of free particles . For completeness we collect the corresponding expressions in Appendix A. Replacing into (A.9) $`p_i`$ by
$$k_{ai}p_{ai}e_a\left[A_i(t,𝐱_a)+A_0(t,𝐱_a)\phi _{ai}\right]$$
(2.36)
yields the one-particle Hamiltonian for the space-like forms of dynamics:
$$H_a=\frac{\phi _{at}}{1𝝋_a^2}\left[\sqrt{(1𝝋_a^2)(m_a^2+𝐤_a^2)+(k_a^i\phi _{ai})^2}+k_a^i\phi _{ai}\right]e_a\phi _{at}A_0(t,𝐱_a).$$
(2.37)
Similarly, for the isotropic forms of dynamics we obtain from (A.11)
$$H_a=\frac{1}{2}\phi _{at}\frac{m_a^2+𝐤_a^2}{k_a^i\phi _{ai}}e_a\phi _{at}A_0(t,𝐱_a).$$
(2.38)
To express the field Hamiltonian (2.35) in the terms of canonical variables we rewrite it into the form
$$_\mathrm{c}=\frac{1}{8\pi }\{\phi _t[\frac{1}{2}h_{ij}h^{ij}+A_{0,i}A^{0,i}]+\phi _t^1\beta _{ij}\left(A^i{}_{,t}{}^{}+A_{0,t}\phi ^i\right)\left(A^j{}_{,t}{}^{}+A_{0,t}\phi ^j\right)\}.$$
(2.39)
Then for the space-like forms of dynamics, taking into account (2.29) and (2.30), we obtain
$$_\mathrm{c}=\frac{1}{8\pi }\phi _t\left\{\gamma ^{ij}\left(4\pi \mathrm{\Pi }_i𝔖_i\right)\left(4\pi \mathrm{\Pi }_j𝔖_j\right)+\frac{1}{2}h_{ij}h^{ij}+A_{0,i}A^{0,i}\right\}.$$
(2.40)
In the case of isotropic forms of dynamics ($`𝝋^2=1`$), when $`\beta _{ij}=\phi _i\phi _j`$, (2.28) gives
$$\phi _j\left(A^j{}_{,t}{}^{}+A_{0,t}\phi ^j\right)=\phi _t\phi _j(4\pi \mathrm{\Pi }^j𝔖^j),$$
(2.41)
so that (2.39) takes the form
$$_\mathrm{c}=\frac{1}{8\pi }\phi _t\left\{\left[\phi ^i\left(4\pi \mathrm{\Pi }_i𝔖_i\right)\right]^2+\frac{1}{2}h_{ij}h^{ij}+A_{0,i}A^{0,i}\right\}.$$
(2.42)
Detailed analysis of the Hamiltonian description for the charged particle system in the instant form of dynamics was carried out in with the application to the classical relativistic statistical mechanics. In the following sections we shall be concerned with the front form of dynamics, analyzing additional primary constraints and ensuring the corresponding Liouville equation.
## 3 Constraint analysis in the front form of dynamics
Let us consider the family of the forms of dynamics, which is given by
$$x^0=t+\mathrm{𝐧𝐱},𝐧^2=1.$$
(3.1)
According to (2.24), (2.31), the set of primary constraints is
$`\xi \mathrm{\Pi }+\mathrm{\Pi }^in_i=0`$ (3.2)
$`\xi _i\epsilon _{ijk}n^j(4\pi \mathrm{\Pi }^k+A_0{}_{}{}^{,k}n_lh^{lk})=0.`$ (3.3)
It is easy to see that (3.3) is equivalent to
$$(4\pi \mathrm{\Pi }^k+A_0{}_{}{}^{,k})_{}n_jh^{jk}=0,$$
(3.4)
where we define orthogonal and longitudinal projections of an arbitrary 3-vector $`f^k`$ with respect to the vector $`n^k`$ as
$$f^k=f_{}^k+f_{}^k,f_{}^k=n^kn_lf^l,f_{}^k=f^kn^kn_lf^l.$$
(3.5)
The canonical Hamiltonian of the system is determined by
$`H_\mathrm{c}={\displaystyle \underset{a=1}{\overset{N}{}}}\left[{\displaystyle \frac{m_a^2+𝐤_a^2}{2𝐤_a𝐧}}+e_aA_0(t,𝐱_a)\right]`$
$`+{\displaystyle \frac{1}{8\pi }}{\displaystyle \left\{\left[n^i\left(4\pi \mathrm{\Pi }_i+A_{0,i}\right)\right]^2+\frac{1}{2}h_{ij}h^{ij}+A_{0,i}A^{0,i}\right\}d^3x}.`$ (3.6)
Then we get Dirac Hamiltonian which takes into account primary constraints (3.2) and (3.4):
$$H_\mathrm{D}=H_\mathrm{c}+[\lambda (\mathrm{\Pi }+\mathrm{\Pi }^in_i)+\lambda _k((4\pi \mathrm{\Pi }^k+A_0{}_{}{}^{,k})_{}n_jh^{jk})]d^3x,$$
(3.7)
where $`\lambda `$, $`\lambda _k`$ are the Dirac multipliers.
The preservation of the constraint (3.2) in time produces the secondary constraint:
$`0=\{\mathrm{\Pi }+\mathrm{\Pi }^in_i,H_\mathrm{c}\}=\rho _i\mathrm{\Pi }_{}^i+{\displaystyle \frac{_i}{4\pi }}[(A_0{}_{}{}^{,i})_{}n_jh^{ji}]`$
$`\mathrm{\Pi }^i{}_{,i}{}^{}\rho ,`$ (3.8)
where $``$ means ”weak equality” in the sense of Dirac and
$$\rho (t,𝐱)=\underset{a=1}{\overset{N}{}}e_a\delta ^3(𝐱𝐱_a(t))$$
(3.9)
is a charge density.
Next, we consider commutation relations between constraints (3.2), (3.4), and (3.8):
$`\{\mathrm{\Pi }+\mathrm{\Pi }^in_i,(4\pi \mathrm{\Pi }^k+A_0{}_{}{}^{,k})_{}n_jh^{jk}\}=0,`$ (3.10)
$`\{\mathrm{\Pi }^i{}_{,i}{}^{}+\rho ,(4\pi \mathrm{\Pi }^k+A_0{}_{}{}^{,k})_{}n_jh^{jk}\}=0,`$ (3.11)
$`\{(4\pi \mathrm{\Pi }^k(t,𝐱)+A_0{}_{}{}^{,k}(t,𝐱))_{}n_jh^{jk}(t,𝐱),`$
$`(4\pi \mathrm{\Pi }^i(t,𝐲)+A_0{}_{}{}^{,i}(t,𝐲))_{}n_jh^{ji}(t,𝐲)\}`$
$`=8\pi \left(\delta ^{ki}n^kn^i\right)\left(n^j{\displaystyle \frac{}{x^j}}\right)\delta ^3(𝐱𝐲)\mathrm{\Omega }^{ki}(𝐱𝐲).`$ (3.12)
We can check directly, that
$$\{\mathrm{\Pi }^i{}_{,i}{}^{}+\rho ,H_\mathrm{c}\}=0.$$
(3.13)
Therefore, the two constraints
$$\xi \mathrm{\Pi }+\mathrm{\Pi }^in_i=0,\mathrm{\Gamma }\mathrm{\Pi }^i{}_{,i}{}^{}\rho 0,$$
(3.14)
belong to the first class.
Taking into account (3.12), we come to the conclusion that the constraints (3.4) are of the second class. However, there are two independent second class constraints only. Then we can reduce them by means of Dirac bracket:
$`\{F,G\}_\mathrm{D}=\{F,G\}{\displaystyle }d^3xd^3y\{F,(4\pi \mathrm{\Pi }^\alpha (t,𝐱)+A_0{}_{}{}^{,\alpha }(t,𝐱))_{}n_ih^{i\alpha }(t,𝐱)\}`$
$`\times C_{\alpha \beta }(𝐱𝐲)\{(4\pi \mathrm{\Pi }^\beta (t,𝐲)+A_0{}_{}{}^{,\beta }(t,𝐲))_{}n_jh^{j\beta }(t,𝐲),G\},\alpha ,\beta =1,2,`$ (3.15)
where $`C_{\alpha \beta }(𝐱𝐲)`$ is an inverse matrix to $`\mathrm{\Omega }^{\alpha \beta }(𝐱𝐲)`$:
$$C_{\alpha \gamma }(𝐱𝐳)\mathrm{\Omega }^{\gamma \beta }(𝐳𝐲)d^3z=\delta _\alpha ^\beta \delta ^3(𝐱𝐲).$$
(3.16)
In the next chapter we shall consider the elimination of the first class constraints and the formulation of the Hamiltonian description in the terms of independent physical variables.
Now let us canonically transform the field variables:
$`(A_0,\mathrm{\Pi },A_i,\mathrm{\Pi }_i)(A_0,,𝒜_i,_i),`$
$`=\mathrm{\Pi }+\mathrm{\Pi }^in_i,_i=\mathrm{\Pi }_i,𝒜_i=A_i+A_0n_i.`$ (3.17)
After the transformation the set of constraints of our system becomes
$`=0,\mathrm{\Gamma }^i{}_{,i}{}^{}\rho 0,`$ (3.18)
$`4\pi _{}^kn_j(𝒜^{k,j}𝒜^{j,k})=0.`$ (3.19)
Using (3.19), we can rewrite the canonical Hamiltonian as
$`H_\mathrm{c}={\displaystyle \underset{a=1}{\overset{N}{}}}\left[{\displaystyle \frac{m_a^2+𝐤_a^2}{2𝐤_a𝐧}}+e_aA_0(t,𝐱_a)\right]`$
$`+{\displaystyle }\{2\pi (n^i_i)^2{\displaystyle \frac{1}{8\pi }}𝒜_{i,j}(𝒜^{j,i}𝒜^{i,j})A_0^i{}_{,i}{}^{}\}d^3x,`$ (3.20)
where now $`k_{ai}=p_{ai}e_a𝒜_i(t,𝐱_a)`$.
## 4 Elimination of the gauge degrees of freedom
Let us consider the first class constraints. We see immediately that $`A_0,`$ are a pair of conjugated gauge canonical variables. A second such pair is formed by $`Q=\mathrm{\Delta }^1𝒜^i_{,i}`$ and $`\mathrm{\Gamma }`$, where $`\mathrm{\Delta }=_i_i,\mathrm{\Delta }\mathrm{\Delta }^1=1,\mathrm{\Delta }_𝐱^1\delta ^3(𝐱)=\mathrm{\Delta }^1(𝐱)1/(4\pi |𝐱|)`$.
Therefore, we can separate the gauge degrees of freedom and gauge-invariant ones by means of the Hodge decomposition (see, e.g., ):
$$𝒜_i=\stackrel{}{𝒜_i}+_iQ,_i=\stackrel{}{_i}+_i\mathrm{\Delta }^1(\mathrm{\Gamma }+\rho ),$$
(4.1)
where
$$\stackrel{}{𝒜_i}=(\delta _i^j_i\mathrm{\Delta }^1^j)𝒜_j,\stackrel{}{_i}=(\delta _i^j_i\mathrm{\Delta }^1^j)_j.$$
(4.2)
Then we have
$$\{Q(t,𝐱),\mathrm{\Gamma }(t,𝐲)\}=\delta ^3(𝐱𝐲),\{\stackrel{}{𝒜_i}(t,𝐱),\stackrel{}{_j}(t,𝐲)\}=(\delta _{ij}^i^j\mathrm{\Delta }^1)\delta ^3(𝐱𝐲).$$
(4.3)
Since $`\stackrel{}{^i_{,i}}=0`$, we can define $`\stackrel{}{^i}`$ as follows
$$\stackrel{}{^i}=(\delta _\alpha ^i\delta _3^i\frac{_\alpha }{_3})\frac{\stackrel{~}{e}^\alpha }{\sqrt{4\pi }},\stackrel{~}{e}^\alpha =\sqrt{4\pi }^\alpha ,\alpha =1,2.$$
(4.4)
Now we have to do a canonical transformation to the new variables $`((x_a^i,\stackrel{~}{\pi }_{ai})`$, $`(a_\alpha ,\stackrel{~}{e}_\alpha )`$, $`(Q,\mathrm{\Gamma })`$, $`(A_0,))`$, which is generated by the functional
$$F=\underset{a=1}{\overset{N}{}}x_a^ip_{ai}A_i\left[\left(\delta _\alpha ^i\delta _3^i\frac{_\alpha }{_3}\right)\frac{\stackrel{~}{e}^\alpha }{\sqrt{4\pi }}+^i\mathrm{\Delta }^1(\mathrm{\Gamma }+\rho )\right]d^3x.$$
(4.5)
We obtain
$`a_\alpha ={\displaystyle \frac{\delta F}{\delta e^\alpha }}=\left(\delta _\alpha ^i\delta _3^i{\displaystyle \frac{_\alpha }{_3}}\right){\displaystyle \frac{𝒜_i}{\sqrt{4\pi }}},`$ (4.6)
$`\stackrel{~}{\pi }_{ai}={\displaystyle \frac{F}{x_a^i}}=p_{ai}e_a_iQ(𝐱_a).`$ (4.7)
The transverse part of $`𝒜_i`$, is connected with the new canonical variables as
$$\stackrel{}{𝒜_i}=\sqrt{4\pi }(\delta _i^\alpha _i\mathrm{\Delta }^1^\alpha )a_\alpha ,\alpha =1,2.$$
(4.8)
Now it is convenient to take $`𝐧=(0,0,1)`$ and perform the following canonical transformation of the field and particle momenta:
$$\stackrel{~}{e}_\alpha =e_\alpha +\sqrt{4\pi }_\alpha \mathrm{\Delta }^1\rho ,\stackrel{~}{\pi }_{ai}=\pi _{ai}\sqrt{4\pi }e_a_i\mathrm{\Delta }^1^\alpha a_\alpha (𝐱_a).$$
(4.9)
The set of constraints for our system in the terms of the new variables can be written as follows:
$$(,\mathrm{\Gamma },e_\alpha a_{\alpha ,3})0.$$
(4.10)
We reduce the second class constraints by means of the Dirac bracket:
$`\{F,G\}_\mathrm{D}=\{F,G\}{\displaystyle d^3xd^3y\{F,e_\alpha (t,𝐱)a_{\alpha ,3}(t,𝐱)\}}`$
$`\times C^{\alpha \beta }(𝐱𝐲)\{e_\beta (t,𝐲)a_{\beta ,3}(t,𝐲),G\}.`$ (4.11)
Here $`C^{\alpha \beta }(𝐱𝐲)`$ is an inverse matrix to
$`𝕎=\{e_\alpha (t,𝐱)a_{\alpha ,3}(t,𝐱),e_\beta (t,𝐲)a_{\beta ,3}(t,𝐲)\}`$
$`=2\delta _{\alpha \beta }{\displaystyle \frac{}{x^3}}\delta ^3(𝐱𝐲),`$ (4.12)
$`{\displaystyle C^{\alpha \gamma }(𝐱𝐳)W_{\gamma \beta }(𝐳𝐲)d^3z}=\delta _\beta ^\alpha \delta ^3(𝐱𝐲).`$ (4.13)
It is given by
$$C^{\alpha \beta }(𝐱𝐲)=\delta ^{\alpha \beta }\delta (x^1y^1)\delta (x^2y^2)\mathrm{sgn}(x^3y^3),$$
(4.14)
because $`C^{\alpha \beta }(𝐱𝐲)`$ must be antisymmetric.
Let us introduce the following denotation:
$$\frac{1}{2}\delta (x^1)\delta (x^2)\mathrm{sgn}(x^3)\frac{1}{_3}\delta ^3(𝐱).$$
(4.15)
The final form of the Dirac bracket is
$`\{F,G\}_\mathrm{D}=\{F,G\}_{(x_a,\pi _a)}+\{F,G\}_{(A_0,)}+\{F,G\}_{(Q,\mathrm{\Gamma })}+{\displaystyle \frac{1}{2}}\{F,G\}_{(a,e)}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle d^3x\left[\frac{\delta F}{\delta e_\alpha (𝐱)}_3\frac{\delta G}{\delta e_\alpha (𝐱)}\frac{\delta F}{\delta a_\alpha (𝐱)}\frac{1}{_3}\frac{\delta F}{\delta a_\alpha (𝐱)}\right]},`$ (4.16)
where $`\{F,G\}_{(x,\pi )}`$ denote the standard Poisson bracket in the terms of $`x`$ and $`\pi `$.
Elimination of constraints (4.10) into $`H_\mathrm{c}`$ leads to the physical Hamiltonian:
$`H_{\mathrm{ph}}={\displaystyle \frac{1}{2}}{\displaystyle \underset{a=1}{\overset{N}{}}}\left[\pi _{a3}+{\displaystyle \frac{(\pi _{a\alpha }\sqrt{4\pi }e_aa_\alpha (𝐱_a))^2+m_a^2}{\pi _{a3}}}\right]+{\displaystyle \frac{1}{2}}{\displaystyle \underset{a,b=1}{\overset{N}{}}}{\displaystyle \frac{e_ae_b}{|𝐱_a𝐱_b|}}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \left(a_\alpha \sqrt{4\pi }\frac{_\alpha }{_3}\mathrm{\Delta }^1\rho \right)\mathrm{\Delta }\left(a_\alpha \sqrt{4\pi }\frac{_\alpha }{_3}\mathrm{\Delta }^1\rho \right)d^3x},`$ (4.17)
which generates evolution of an arbitrary function $`f`$ depending on the gauge-invariant variables $`x_a^i`$, $`\pi _{ai}`$ and $`a_\alpha `$, $`e_\alpha `$ in the terms of the Dirac bracket (4.16):
$$\frac{\mathrm{d}f}{\mathrm{d}t}=\frac{f}{t}+\{f,H_{\mathrm{ph}}\}_\mathrm{D}.$$
(4.18)
The gauge-invariant volume element of the constrained field phase space $`(a_\alpha ,e^\alpha )`$ is written as
$$\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}^\mathrm{f}(t)=\gamma \sqrt{\mathrm{Det}𝕎}\underset{\alpha =1,2}{}\underset{𝐱}{}\delta [e_\alpha (t,𝐱)a_{\alpha ,3}(t,𝐱)]\mathrm{d}a_\alpha (t,𝐱)\mathrm{d}e_\alpha (t,𝐱),$$
(4.19)
where $`\gamma `$ is defined as a normalisation constant of the Gauss integral :
$`\gamma {\displaystyle \mathrm{exp}\left(\frac{1}{2}a_\alpha (t,𝐱)L^{\alpha \beta }(t,𝐱,𝐲)a_\beta (t,𝐲)d^3xd^3y\right)\underset{\alpha =1,2}{}\underset{𝐱}{}\mathrm{d}a_\alpha (t,𝐱)}`$
$`=\mathrm{Det}^{1/2}L^{\alpha \beta }(t,𝐱,𝐲).`$ (4.20)
Taking into account (4.19), we can write the volume element of the physical phase space of the described system:
$$\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}(t)=\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}^\mathrm{p}(t)\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}^\mathrm{f}(t).$$
(4.21)
Here
$$\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}^\mathrm{p}(t)=\underset{a=1}{\overset{N}{}}\underset{i=1}{\overset{3}{}}\mathrm{d}x_a^i(t)\mathrm{d}\pi _{ia}(t).$$
(4.22)
is the volume element of the particle phase space.
Now we need check the Liouville theorem: $`\mathrm{\Gamma }_{\mathrm{ph}}(t)=\mathrm{\Gamma }_{\mathrm{ph}}(t_0)`$.
It is well known from the classical mechanics, that system evolution in the phase space can be described by means of a canonical transformation in the terms of the Poisson bracket, which immediately leads to the conservation of the phase space volume. This proves the conservation of $`\mathrm{\Gamma }_{\mathrm{ph}}^\mathrm{p}`$. In our case the field evolution is generated by the Dirac bracket. Nevertheless, the conservation of the volume of the constrained field phase space can be proved and we shall demonstrate its possible proof for the considered Dirac bracket (4.16) in Appendix B.
The volume element of the full phase space is
$$\mathrm{d}\mathrm{\Gamma }_{\mathrm{full}}(t)=\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}(t)\mathrm{d}\mathrm{\Gamma }_\mathrm{g}(t),$$
(4.23)
where $`\mathrm{d}\mathrm{\Gamma }_\mathrm{g}(t)`$ is the gauge volume element of the phase space:
$$\underset{𝐱}{}\delta [(t,𝐱)]\delta [\mathrm{\Gamma }(t,𝐱)]\mathrm{d}A_0(t,𝐱)\mathrm{d}(t,𝐱)\mathrm{d}Q(t,𝐱)\mathrm{d}\mathrm{\Gamma }(t,𝐱).$$
(4.24)
However, we turn off $`\mathrm{d}\mathrm{\Gamma }_\mathrm{g}(t)`$ from the further description, because all thermodynamical characteristics of the system do not depend on dynamics of the gauge degrees of freedom (see ).
## 5 Statistical mechanics
Now let us imagine that the physical initial data $`x_a^i(t_0)`$, $`\pi _{ai}(t_0)`$, $`a_\alpha (t_0,𝐱)`$, $`e_\alpha (t_0,𝐱)`$ are not precisely known. Hence, we can introduce a probability density $`\varrho (t_0)\varrho (t_0,x_a^i(t_0),\pi _{ai}(t_0),a_\alpha (t_0,𝐱),e_\alpha (t_0,𝐱))`$ for having various initial states. This function satisfies the condition
$$\varrho (t_0)d\mathrm{\Gamma }_{\mathrm{ph}}(t_0)=1.$$
(5.1)
Then average value of a general dynamical variable $`f`$ is defined by
$$\overline{f}(t)=f(t,x_a^i(t),\pi _{ai}(t),a_\alpha (t,𝐱),e_\alpha (t,𝐱))\varrho (t)d\mathrm{\Gamma }_{\mathrm{ph}}(t).$$
(5.2)
Since $`t_0`$ has been randomly selected, we come to the relation
$$1=\varrho (t_0)d\mathrm{\Gamma }_{\mathrm{ph}}(t_0)=\varrho (t)d\mathrm{\Gamma }_{\mathrm{ph}}(t).$$
(5.3)
We have already seen that $`\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}(t_0)=\mathrm{d}\mathrm{\Gamma }_{\mathrm{ph}}(t)`$, so one immediately obtains the Liouville equation
$$0=\frac{\mathrm{d}\varrho (t)}{\mathrm{d}t}=\frac{\varrho (t)}{t}+\{\varrho (t),H_{\mathrm{ph}}(t)\}_\mathrm{D}.$$
(5.4)
Taking into account (5.4), we interpret $`\varrho `$ as an integral of motion.
Let us consider the canonical Gibbs ensemble in equilibrium. In this case we have
$$\varrho (t)=C\mathrm{e}^{\beta H_{\mathrm{ph}}(t)},$$
(5.5)
where $`\beta =1/kT`$ and $`C`$ is a normalisation constant.
Partition function can be found as
$$Z=\frac{e^{\beta H_{\mathrm{ph}}(t)}}{(2\pi )^{3N}N!}d\mathrm{\Gamma }_{\mathrm{ph}}(t).$$
(5.6)
We shall find below the value of $`Z`$.
Let us first rewrite the physical Hamiltonian as follows
$`H_{\mathrm{ph}}={\displaystyle \frac{1}{2}}{\displaystyle \underset{a=1}{\overset{N}{}}}\left(\pi _{a3}+{\displaystyle \frac{\pi _{a\alpha }^2+m_a^2}{\pi _{a3}}}\right)+{\displaystyle \frac{1}{2}}{\displaystyle \underset{a,b=1}{\overset{N}{}}}{\displaystyle \frac{e_ae_b}{|𝐱_a𝐱_b|}}`$
$`2\pi {\displaystyle \frac{_\alpha }{_3}\rho (t,𝐱)\mathrm{\Delta }^1(𝐱𝐲)\frac{_\alpha }{_3}\rho (t,𝐲)d^3xd^3y}`$
$`+2\pi {\displaystyle \left(\frac{_\alpha }{_3}\rho (t,𝐱)+j_\alpha ^0(t,𝐱)\right)G(t,𝐱𝐲)\left(\frac{_\alpha }{_3}\rho (t,𝐲)+j_\alpha ^0(t,𝐲)\right)d^3xd^3y}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \stackrel{~}{a}_\alpha (t,𝐱)\left(\mathrm{\Delta }+4\pi \underset{a=1}{\overset{N}{}}\frac{e_a^2}{\pi _{a3}(t)}\delta ^3(𝐱𝐱_a(t))\right)\stackrel{~}{a}_\alpha (t,𝐱)d^3x},`$ (5.7)
where
$`\stackrel{~}{a}_\alpha (t,𝐱)=a_\alpha (t,𝐱)\sqrt{4\pi }{\displaystyle G(t,𝐱𝐲)\left(\frac{_\alpha }{_3}\rho (t,𝐲)+j_\alpha ^0(t,𝐲)\right)d^3y},`$ (5.8)
$`j_\alpha ^0(t,𝐱)={\displaystyle \underset{a=1}{\overset{N}{}}}e_a{\displaystyle \frac{\pi _{a\alpha }(t)}{\pi _{a3}(t)}}\delta ^3(𝐱𝐱_a(t)),`$ (5.9)
$`\left[\mathrm{\Delta }+4\pi {\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{e_a^2}{\pi _{a3}(t)}}\delta ^3(𝐱𝐱_a(t))\right]G(t,𝐱)=\delta ^3(𝐱).`$ (5.10)
Since $`H_{\mathrm{ph}}`$ does not depend on field momenta $`e^\alpha `$, after an integration over $`e^\alpha `$ in (5.6) we obtain the expression for the partition function:
$`Z=\gamma \sqrt{\mathrm{Det}𝕎}{\displaystyle \frac{1}{N!}}{\displaystyle }\mathrm{exp}(\beta )[{\displaystyle \frac{1}{2}}{\displaystyle \underset{a=1}{\overset{N}{}}}(\pi _{a3}+{\displaystyle \frac{\pi _{a\alpha }^2+m_a^2}{\pi _{a3}}})+{\displaystyle \frac{1}{2}}{\displaystyle \underset{a,b=1}{\overset{N}{}}}{\displaystyle \frac{e_ae_b}{|𝐱_a𝐱_b|}}`$
$`2\pi {\displaystyle \frac{_\alpha }{_3}\rho (t,𝐱)\mathrm{\Delta }^1(𝐱𝐲)\frac{_\alpha }{_3}\rho (t,𝐲)d^3xd^3y}`$
$`+2\pi {\displaystyle \left(\frac{_\alpha }{_3}\rho (t,𝐱)+j_\alpha ^0(t,𝐱)\right)G(t,𝐱𝐲)\left(\frac{_\alpha }{_3}\rho (t,𝐲)+j_\alpha ^0(t,𝐲)\right)d^3xd^3y}`$
$`{\displaystyle \frac{1}{2}}{\displaystyle }\stackrel{~}{a}_\alpha (t,𝐱)(\mathrm{\Delta }+4\pi {\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{e_a^2}{\pi _{a3}(t)}}\delta ^3(𝐱𝐱_a(t)))\stackrel{~}{a}_\alpha (t,𝐱)d^3x]`$
$`\times {\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{d}x_a^i\mathrm{d}\pi _{ai}}{2\pi }}{\displaystyle \underset{\alpha =1,2}{}}{\displaystyle \underset{𝐱}{}}\mathrm{d}a_\alpha (t,𝐱)`$ (5.11)
Here we can replace $`da_\alpha `$ by $`d\stackrel{~}{a}_\alpha `$. Taking into account (4.8), the integration over $`\stackrel{~}{a}_\alpha `$ yields
$`Z=Z^\mathrm{f}{\displaystyle \frac{1}{N!}}{\displaystyle }{\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{3}{}}}{\displaystyle \frac{\mathrm{d}x_a^i\mathrm{d}\pi _{ai}}{2\pi }}\mathrm{exp}(\beta )[{\displaystyle \frac{1}{2}}{\displaystyle \underset{a=1}{\overset{N}{}}}(\pi _{a3}+{\displaystyle \frac{\pi _{a\alpha }^2+m_a^2}{\pi _{a3}}})`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{a,b=1}{\overset{N}{}}}{\displaystyle \frac{e_ae_b}{|𝐱_a𝐱_b|}}2\pi {\displaystyle \frac{_\alpha }{_3}\rho (t,𝐱)\mathrm{\Delta }^1(𝐱𝐲)\frac{_\alpha }{_3}\rho (t,𝐲)d^3xd^3y}`$
$`+2\pi {\displaystyle }({\displaystyle \frac{_\alpha }{_3}}\rho (t,𝐱)+j_\alpha ^0(t,𝐱))G(t,𝐱𝐲)({\displaystyle \frac{_\alpha }{_3}}\rho (t,𝐲)+j_\alpha ^0(t,𝐲))d^3xd^3y]`$
$`\times \mathrm{Det}^{1/2}\delta ^{\alpha \beta }\left(\delta ^3(𝐱𝐲)+4\pi \mathrm{\Delta }^1(𝐱𝐲){\displaystyle \underset{a=1}{\overset{N}{}}}{\displaystyle \frac{e_a^2}{\pi _{a3}}}\delta ^3(𝐲𝐱_a)\right),`$ (5.12)
where
$$Z^\mathrm{f}=\frac{\sqrt{\mathrm{Det}𝕎}}{\sqrt{\mathrm{Det}\beta \delta ^{\alpha \beta }\mathrm{\Delta }\delta ^3(𝐱𝐲)}}$$
(5.13)
represents the free field partition function.
## 6 Conclusions
In this paper we have once more demonstrated the usefulness of the various forms of relativistic dynamics in treating the relativistic particle systems. Constraint analysis of the considered problem shows significantly different structures of Hamiltonian description of charged particles, interacting by means of an electromagnetic field, in the space-like and isotropic forms of dynamics. Specifically, the use of front form of relativistic dynamics allows us to exclude the electromagnetic field variables from classical partition function of the system of charged particles. It is obvious that such an exclusion must depend on the particular boundary conditions for the field variables, but at present the proper sense of these conditions remains unclear. The obtained representation for the partition function contains highly nonlocal expressions, and their further analysis consist of a complicated task. The consideration of the various approximation schemes at this step seems to be inevitable.
It should be noted that exclusion of the field variables transform the problem into the domain of relativistic direct interaction theory (see, e.g., ). An application of such a theory to the consistent formulation of the relativistic statistical mechanics is just at the beginning.
On the other hand, the established Liouville equation may be used in nonequilibrium situation as well. It serves an useful starting point for deriving various new forms of kinetic equations for charged particle system.
## Acknowledgment
We are greatly indebted to Yu. Yaremko and V. Shpytko for many stimulating discussions. Numerous helpful conversation with A. Duviryak and his support in the performance of this work are especially acknowledged.
## Appendix A <br>Relativistic free particle in an arbitrary form of dynamics
The Lagrangian of relativistic free particle in a given form of dynamics (2.5) is
$$L=m\sqrt{(D\phi )^2𝐯^2}m\mathrm{\Gamma }^1.$$
(A.1)
The canonical momentum and Hamiltonian are given by
$$p_i=m\mathrm{\Gamma }(g_{ij}v^j\phi _i\phi _t),$$
(A.2)
$$H=m\mathrm{\Gamma }\phi _tD\phi ,$$
(A.3)
where the matrix $`g_{ij}=\delta _{ij}\phi _i\phi _j`$ has been introduced. The determinant of the matrix is
$$gg_{ij}=1𝝋^2,$$
(A.4)
and $`g0`$ as result of the condition (2.10).
Consider firstly the case $`g>0`$. Then the matrix $`g_{ij}`$ has an inverse, $`\stackrel{~}{g}_{ij}=\delta _{ij}+g^1\phi _i\phi _j`$ and from (A.2) it follows
$$\stackrel{~}{g}_{ij}p_ip_j=m^2\left(\mathrm{\Gamma }^2g^1\phi _t^21\right).$$
(A.5)
Using (2.8), we find
$$m\mathrm{\Gamma }=\phi _t^1\sqrt{g(m^2+𝐩^2)+(p_i\phi _i)}\phi _t^1B,$$
(A.6)
$$v_i=B^1\phi _t\left[p_i+g^1\phi _i(B+p_j\phi ^j)\right],$$
(A.7)
$$D\phi =B^1\phi _tg^1\left(B+p_i\phi ^i\right).$$
(A.8)
Combining these results with (A.3), one gets
$$H=\phi _tg^1\left[\sqrt{g(m^2+𝐩^2)+(p_i\phi ^i)^2}+p_i\phi ^i\right].$$
(A.9)
The case of isotropic forms of mechanics can be treated by taking the limit $`g+0`$ and using that, in view of (A.2),
$$p_i\phi ^i=m\mathrm{\Gamma }\left[gv_i\phi ^i+\phi _t(g1)\right],$$
(A.10)
so that $`p_i\phi ^i0`$ as $`g+0`$. It gives
$$H=\frac{1}{2}\phi _t\frac{m^2+𝐩^2}{p_i\phi ^i}.$$
(A.11)
## Appendix B <br>Proof of the Liouville theorem for a given Dirac bracket
Let us show that evolution of the field variables as the generalized canonical transformation in the terms of the Dirac bracket conserves the volume of the constrained field phase space.
We first consider transformation generated by some functional $`G`$ with an arbitrary parameter $`\xi `$:
$`\alpha _\alpha =a_\alpha +\xi \{a_\alpha ,G\}_\mathrm{D}=a_\alpha +{\displaystyle \frac{\xi }{2}}{\displaystyle \frac{\delta G}{\delta e^\alpha }}{\displaystyle \frac{\xi }{2}}{\displaystyle \frac{1}{_3}}{\displaystyle \frac{\delta G}{\delta a^\alpha }},`$
$`ϵ_\alpha =e_\alpha +\xi \{e_\alpha ,G\}_\mathrm{D}=e_\alpha {\displaystyle \frac{\xi }{2}}{\displaystyle \frac{\delta G}{\delta a^\alpha }}+{\displaystyle \frac{\xi }{2}}_3{\displaystyle \frac{\delta G}{\delta e^\alpha }}.`$ (B.1)
We immediately see that $`e_\alpha a_{\alpha ,3}=ϵ_\alpha \alpha _{\alpha ,3}`$. So, the constancy condition of $`\mathrm{\Gamma }_{\mathrm{ph}}^\mathrm{f}`$ leads to the relation:
$`{\displaystyle \underset{\alpha =1,2}{}\underset{𝐱}{}\delta [ϵ_\alpha (t,𝐱)\alpha _{\alpha ,3}(t,𝐱)]\mathrm{d}\alpha _\alpha (t,𝐱)\mathrm{d}ϵ_\alpha (t,𝐱)}`$
$`={\displaystyle J(\xi ,0)\underset{\alpha =1,2}{}\underset{𝐱}{}\delta [e_\alpha (t,𝐱)a_{\alpha ,3}(t,𝐱)]\mathrm{d}a_\alpha (t,𝐱)\mathrm{d}e_\alpha (t,𝐱)}.`$ (B.2)
Here Jacobian $`J(\xi ,\eta )`$ is defined by
$$J(\xi ,\eta )\mathrm{Det}\frac{\delta (\alpha _\alpha (t,𝐱;\xi )ϵ_\alpha (t,𝐱;\xi )}{\delta (\alpha _\beta (t,𝐲;\eta )ϵ_\beta (t,𝐲;\eta )}.$$
(B.3)
If $`\eta =0`$, then we have
$$J(\xi ,0)=\mathrm{Det}\frac{\delta (\alpha _\alpha (t,𝐱;\xi )ϵ_\alpha (t,𝐱;\xi )}{\delta (a_\beta (t,𝐲)e_\beta (t,𝐲))}$$
(B.4)
It is evident that $`J(0,0)=1`$. Let us compute $`(\mathrm{d}J(\xi ,0)/\mathrm{d}\xi )|_{\xi =0}`$. We get
$`{\displaystyle \frac{\mathrm{d}J(\xi ,0)}{\mathrm{d}\xi }}|_{\xi =0}=\mathrm{Tr}{\displaystyle \frac{\delta _\xi \alpha _\alpha (t,𝐱;\xi )}{\delta a_\beta (t,𝐲)}}|_{\xi =0}+\mathrm{Tr}{\displaystyle \frac{\delta _\xi ϵ_\alpha (t,𝐱;\xi )}{\delta e_\beta (t,𝐲)}}|_{\xi =0}`$
$`=\mathrm{Tr}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta ^2G}{\delta a_\beta (t,𝐲)\delta e_\alpha (t,𝐱)}}+\mathrm{Tr}{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{_3}}{\displaystyle \frac{\delta ^2G}{\delta a_\beta (t,𝐲)\delta a_\alpha (t,𝐱)}}`$
$`+\mathrm{Tr}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta ^2G}{\delta e_\beta (t,𝐲)\delta a_\alpha (t,𝐱)}}+\mathrm{Tr}{\displaystyle \frac{1}{2}}_3{\displaystyle \frac{\delta ^2G}{\delta e_\beta (t,𝐲)\delta e_\alpha (t,𝐱)}}.`$ (B.5)
It is obvious that the first and the third terms vanish. One of possible ways demonstrating a cancellation of the second and fourth terms is based on commutation of the Dirac bracket action and replacement of $`a_{\alpha ,3}`$ by $`e_\alpha `$ (or $`e_\alpha `$ by $`a_{\alpha ,3}`$ in view of $`e_\alpha a_{\alpha ,3}=0`$).
Therefore, we have
$$\frac{\mathrm{d}J(\xi ,0)}{\mathrm{d}\xi }|_{\xi =0}=0.$$
(B.6)
Since $`J(\xi ,0)=J(\xi ,\xi _1)J(\xi _1,0)`$, then
$$\frac{\mathrm{d}J(\xi ,0)}{\mathrm{d}\xi }|_{\xi =\xi _1}=\frac{\mathrm{d}J(\xi ,\xi _1)}{\mathrm{d}\xi }|_{\xi =\xi _1}J(\xi _1,0)=0.$$
(B.7)
Thus, $`J(\xi ,0)=1`$ for all $`\xi `$.
Now if we take $`\xi =t`$ and $`G=H_{\mathrm{ph}}`$, we come to conclusion
$$\mathrm{\Gamma }_{\mathrm{ph}}(t)=\mathrm{\Gamma }_{\mathrm{ph}}(t_0),$$
(B.8)
namely, the time evolution preserves the phase space volume $`\mathrm{\Gamma }_{\mathrm{ph}}`$.
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# The 2dF QSO Redshift Survey - I. The Optical QSO Luminosity Function
## 1 Introduction
The QSO optical luminosity function (OLF) and its evolution with redshift provides fundamental information on the overall demographics of the AGN population. It provides constraints on the physical models for QSOs (Haehnelt & Rees 1993; Terlevich & Boyle 1993), information on models of structure formation in the early Universe (Efstathiou & Rees 1988) and a picture of the ionizing UV/optical luminosity density from QSOs as a function of redshift (Meiksin & Madau 1993; Boyle & Terlevich 1998).
Based on the major ultra-violet excess (UVX) QSO surveys of the 1980s (Bracessi: Marshall et al. 1983, Palomar–Green: Green, Schmidt & Liebert 1986, Durham–AAT: Boyle et al. 1990), a picture emerged in which the low-intermediate redshift ($`z<2.2`$) QSO OLF was modelled by a two-power-law function with a steep bright end ($`\mathrm{\Phi }(L_\mathrm{B})L_\mathrm{B}^{3.6\pm 0.1}`$) and a much flatter faint end ($`\mathrm{\Phi }(L_\mathrm{B})L_\mathrm{B}^{1.2\pm 0.1}`$) whose redshift dependence was best fit by pure luminosity evolution, i.e., a uniform increase in luminosity toward high redshift (see e.g., Marshall 1985; Boyle et al. 1988; Hartwick & Schade 1990). This evolution was modelled as a power-law in redshift of the functional form $`L_\mathrm{B}^{}(1+z)^{k_\mathrm{L}}`$, where $`3<k_\mathrm{L}<3.4`$ for an Einstein-de Sitter universe.
Later QSO surveys, particularly those that focussed on bright magnitudes (LBQS: Hewett, Foltz & Chaffee 1993, EQS: Goldschmidt & Miller 1998, HBQS: La Franca & Cristiani 1997) have reported evidence for a more complex form of evolution in which the bright end of the OLF showed a significant steepening with increasing redshift. Indeed the steepening was sufficiently dramatic that, at the very lowest redshifts studied, the OLF showed a single featureless power law (Koehler et al. 1997) with little evidence for any cosmic evolution amongst the brightest QSOs ($`M_\mathrm{B}<27`$). Coupled with this, further claims were made (Hawkins & Veron 1995) that the observed break in the luminosity function at higher redshifts was much less dramatic than had previously been reported.
With the increasing number of high redshift ($`z>2`$) QSOs discovered in surveys with well-defined selection criteria, evidence was also found that the strong power-law evolution does not continue on beyond $`z2`$. The nature of this change was modelled in a variety of ways; from a simple halt to constant comoving density (see Boyle et al. 1991) between $`z2`$ to $`z3`$ to a slackening off of the evolution rate over the broad redshift range $`1.6<z<3.3`$ (Hewett et al. 1993; Warren, Hewett & Osmer 1994). From surveys of higher redshift QSOs ($`3.5<z<5.0`$), strong evidence emerged for the onset of a dramatic decline in the QSO space density at $`z>3.5`$ (see e.g Warren et al. 1994).
However, compared to our knowledge of the galaxy luminosity function, QSO OLF determinations were based on relatively few objects. This was particularly true for the extremes of the redshift distribution $`z<0.3`$ and $`z>4`$ or the high luminosity end of the OLF ($`M_\mathrm{B}<27`$), where departures from the pure luminosity evolution model were first noted. In these regimes, the numbers of objects in suitable surveys could be counted in the tens. Even in the most well-sampled regions of the QSO ($`M_\mathrm{B},z`$) plane, surveys comprised a few hundred or up to one thousand QSOs in total.
The long-heralded arrival of the new generation of large QSO surveys compiled with facilities such as the 2dF (Lewis et al. 1998) or Sloan Digital Sky Survey (Loveday et al. 1998) will shortly generate many tens of thousands of QSOs with which to carry out the definitive study of the QSO OLF and its evolution with redshift. These surveys also provide the opportunity to study the QSO OLF as a function of various physical parameters (e.g., emission-line strength, continuum slope) or types (e.g., BAL QSOs).
In this paper, we present a new determination of the QSO OLF based on the first $`6000`$ QSOs identified in the 2dF QSO redshift survey (2QZ) (see Smith et al. 1998). The 2QZ is currently more than a factor of 10 larger than previous QSO surveys to a similar magnitude limit ($`b_\mathrm{J}<20.85`$ mag). When complete, it is planned that the 2QZ will comprise 25000–30000 QSOs. In section 2 we present a brief overview of the survey and in section 3 we describe the analysis of the survey. We present our conclusions in section 4.
## 2 data
### 2.1 The 2dF QSO redshift survey
#### 2.1.1 Input catalogue
For the purposes of the analysis we have used the current version (as of September 1999) of the 2QZ catalogue containing 6684 QSOs in total. The identification of the QSO candidates for the 2QZ was based on broadband $`ub_\mathrm{J}r`$ colour selection from APM measurements of UK Schmidt (UKST) photographic plates. The survey area comprises 30 UKST fields, arranged in two $`75^{}\times 5^{}`$ declination strips centred on $`\delta =30^{}`$ and $`\delta =0^{}`$. The $`\delta =30^{}`$ strip extends from $`\alpha `$ = 21<sup>h</sup>40 to $`\alpha `$ = 3<sup>h</sup>15 and the equatorial strip from $`\alpha `$ = 9<sup>h</sup>50 to $`\alpha `$ = 14<sup>h</sup>50. The total survey area is 740 deg<sup>2</sup>, when allowance is made for regions of sky excised around bright stars. The 2QZ area forms an exact subset of the 2dF Galaxy Redshift Survey (GRS: see Colless 1998) area, with identical ‘holes’ used for both surveys. In a typical 2dF field, approximately 225 fibres are devoted to galaxies, 125 to QSOs and 25–30 fibres are devoted to sky. The data are reduced using the standard 2dF pipeline reduction system (Bailey & Glazebrook 1999).
In each UKST field, APM measurements of one $`b_\mathrm{J}`$ plate, one $`r`$ plate and up to four $`u`$ plates/films were used to generate a catalogue of stellar objects with $`b_\mathrm{J}<20.85`$. A sophisticated procedure was devised to ensure catalogue homogeneity (Smith 1998). Corrections were made for vignetting and field effects due to variable desensitization in the UKST plates, these effects being particularly noticeable at the edges of plates. The criteria for inclusion in the catalogue were ($`ub_\mathrm{J})0.36;(ub_\mathrm{J})<0.120.8(b_\mathrm{J}r);(b_\mathrm{J}r)<0.05`$. Based on the colours of QSOs which have previously been identified in the survey region, we estimate the catalogue is $``$ 90% complete for $`z<2`$ QSOs and comprises $``$55% QSOs (see Croom 1997). Subsequent spectroscopic observations with 2dF have confirmed this QSO fraction, with the principal contamination arising from galactic subdwarfs and compact blue galaxies. Full details of the construction of the input catalogue may be found in Croom (1997) and Smith (1998).
#### 2.1.2 Spectroscopic observations
2dF spectroscopic observations for the 2QZ began in January 1997, although the bulk of the redshifts have been obtained in the observing runs after October 1998 as the 2dF system has gained functionality (increased number of fibres, faster field re-configuration times). Over 6000 QSO redshifts have been obtained with the 2dF, and the 2QZ is now the largest single homogeneous QSO catalogue in existence. Each 2dF field in the survey is observed for typically one hour. The spectral resolution is 4Å pixel<sup>-1</sup> and the spectra cover the wavelength range 3700Å–8000Å. This set-up gives a typical signal-to-noise ratio of approximately seven or greater between 4000Å–6500Å in the continua of the faintest objects ($`b_\mathrm{J}=20.85`$) in the 2QZ. This was sufficient to identify up to 85 per cent of the objects in each 2dF field down to this magnitude limit. Where poor conditions prevented us from achieving this identification rate in any single 2dF observation, we included in the final catalogue only those QSOs that were brighter than the magnitude at which 80 per cent of the QSO candidates in the field had a positive identification.
#### 2.1.3 Spectroscopic Incompleteness
In the OLF analysis below, we correct for incompleteness in the 2QZ by using an effective area which is a function of both apparent $`b_\mathrm{J}`$ magnitude and redshift. Hereinafter, all $`b_\mathrm{J}`$ magnitudes are corrected for galactic extinction using the Schlegel, Finkbeiner & Davis (1998) values.
The 2QZ is subject to three forms of spectroscopic incompleteness. First, a small fraction of survey fields were observed in poor conditions, and do not reach the target spectroscopic completeness of 80 per cent. As described above, the magnitude limits for these fields were made brighter until the desired completeness level was reached. Secondly, we corrected the inevitable trend to increasing spectroscopic incompleteness at fainter magnitudes by correcting the actual area observed by the fraction of objects with a reliable spectroscopic identification. This was done in 0.025mag bins to track accurately the magnitude-dependent incompleteness.
Finally, we applied a uniform correction factor of 0.71 to the effective area to account for the 29 per cent objects in the input catalogue that have not yet had been observed by 2dF, despite being located in areas already observed in the 2dF survey. This is because the numbers of candidates in the combined QSO/galaxy catalogues in a typical 2dF field is $``$ 50 per cent higher than the number of fibres available. Near-complete ($`>`$ 95 per cent) spectroscopic coverage of the survey area by 2dF is achieved by a complex tiling algorithm, with each 2dF field being visited 1.5 times on average. Since the observations are still at a relatively early stage, the coverage of the survey area with 2dF is relatively patchy and the completeness fraction is still relatively low (71 per cent). The resulting effective area of the survey is plotted in Figure 1.
#### 2.1.4 Photometric Incompleteness
Photometric incompleteness, arising from the errors in the photographic magnitudes ($`\pm `$0.1 mag in each band) and variability (due to the noncontemparaneous nature of the $`ub_\mathrm{J}r`$ plates on each field) will cause QSOs to exhibit $`ub_\mathrm{J}/b_\mathrm{J}r`$ colours outside our selection criteria. The incompleteness is a complex function of both magnitude and redshift. An estimate of this incompleteness has been made by Croom et al. (in preparation) using mean colour-redshift relations (zero-pointed to the 2QZ system) from the non-colour-selected QSO survey of Hawkins & Veron (1995). At $`z<2.3`$ the Hawkins & Veron colours accurately trace the mean colours of the 2QZ QSOs. The dispersion in the colours as a function of magnitude was derived from the observed dispersion of colours for the 2QZ QSOs in the redshift interval $`1<z<2`$, the regime of highest completeness. A Monte Carlo simulation using the measured dispersion in colours and the Hawkins & Veron mean colours, with $`10^6`$ QSOs in each $`\mathrm{\Delta }b_\mathrm{J}`$ and $`\mathrm{\Delta }z`$ bin, was then used to predict the photometric completeness as a function of both magnitude and redshift. The completeness contours are shown in Figure 2. The photometric completeness is largely independent of magnitude and is at least 85 per cent or greater over the redshift range $`0.4<z<2.1`$. At higher redshifts the completeness rapidly drops, falling to below 50 per cent at $`z>2.3`$. We have therefore chosen this redshift as the upper limit to our analysis below. Although the catalogue contains many hundreds of QSOs with $`z>2.3`$ (which will be used for the clustering analysis that is not so dependent on photometric completeness), any small errors in the completeness estimates at these redshifts can lead to large variations in the computed effective area.
At the lowest redshifts, the derived photometric completeness is still relatively high ($``$ 60 per cent at $`z=0.2`$). However, low redshift, low luminosity QSOs ($`M_\mathrm{B}>23`$) dominated by their host galaxy light will be lost from the 2QZ as a result of the stellar selection criterion applied to the input catalogue. Although it may be possible to correct for such selection effects, for the present analysis we have simply chosen to exclude all low luminosity QSOs ($`M_\mathrm{B}>23`$) from the OLF fitting procedure below. Given the bright apparent magnitude limit of the 2QZ ($`b_\mathrm{J}=18.25`$ mag), this effectively limits the minimum redshift in the complete 2QZ sample to $`z>0.35`$. We have therefore adopted this as our low redshift limit for all QSO surveys used in the OLF analysis below. Thus, out of a total of 6684 QSOs identified in the 2QZ (13552 objects observed in 219 2dF fields), 5067 fulfil the criteria for inclusion in the complete sample ($`q_0=0.5`$, $`H_0=50`$km s$`^1`$Mpc<sup>-1</sup>).
### 2.2 Other Samples
We are currently in the process of extending the bright limit of the 2QZ from $`b_\mathrm{J}=18.25`$ mag to $`b_\mathrm{J}=17`$ mag using observations made with the FLAIR spectrograph on the UK Schmidt Telescope. These observations are not yet complete and we have incorporated a number of independent QSO surveys at brighter magnitudes to extend the present analysis of the OLF to higher luminosities. Table 1 lists the areas, magnitude limits, and the number of QSOs within the completeness limits ($`M_\mathrm{B}<23`$, $`0.35<z<2.3`$) for each QSO survey used in this paper. The absolute magnitudes and redshifts for the QSOs in these surveys are plotted in Figure 3.
The Large Bright QSO Survey (LBQS: Hewett et al. 1995) provides a complementary sample to the 2QZ. The details and completeness of the survey are well established and it provides a large number of QSOs in the $`1.52`$mag interval brighter than the $`b_\mathrm{J}>18.25`$ 2QZ survey limit.
At the brightest magnitudes $`B<16.5`$mag, the largest currently published survey is the Palomar–Green survey (Green et al. 1986). However, the completeness of this survey has been called into question by a number of recent surveys including the Edinburgh Quasar Survey (EQS: Goldschmidt et al. 1998) and the Hamburg/ESO Quasar Survey (HEQS: Koehler et al. 1997). Unfortunately, full details of both these catalogues have yet to be published, and although we have access to the unpublished EQS, details of its completeness have yet to be accurately determined.
Given the remaining uncertainty over their details, and apparent discrepancy between QSO surface densities derived from these brighter surveys, we chose to carry out our analysis of the OLF for three separate survey combinations: a) 2QZ + LBQS; b) 2QZ + LBQS + PG; c) 2QZ + LBQS + EQS + HEQS. Given the available information, the 2QZ + LBQS data sample constitutes our primary data sample in this paper. When more extensive published catalogues are available from brighter surveys (in particular the HEQS) they will clearly provide a powerful test of the OLF models at the brightest absolute magnitudes (see Figure 3).
A further survey at intermediate magnitudes, the Homogeneous Bright QSO Survey (HBQS) is currently under construction (see Cristiani et al. 1995). There is considerable overlap between the survey areas for the HBQS and EQS. The currently published HBQS samples contains fewer QSOs than the EQS sample to which we have access. We have therefore chosen not to include the HBQS in the analysis below.
With one exception, the surveys are all largely independent of one another. The the current coverage of the 2QZ results in a very small overlap with the LBQS; there is only one QSO in common between the 2QZ and LBQS. Given the complexity of the current 2dF coverage function in the spatial domain (due to incomplete tiling), we therefore chose to treat the LBQS and 2QZ as independant surveys. At brighter magnitudes, the PG, HEQS and EQS are all independant of one another.
The only exception is the overlap between the LBQS and EQS. There are 32 QSOs in the EQS sample which are also in the LBQS sample, or approximately 30 per cent of the EQS sample used in this analysis. The QSOs are distrubuted throughout the four UK Schmidt fields in common between the two surveys. In the analysis below we have therefore excluded these four fields from the EQS during the model fitting procedure, but we test the acceptibilty of the best-fit model against the full EQS sample. The number of QSOs for the EQS quoted in Table 1 refers to the full sample.
## 3 ANALYSIS
### 3.1 $`1/V`$ estimator
As the first step in our analysis, we obtained a graphical representation of the OLF and its evolution with redshift. We used the new $`1/V`$ estimator devised by Page & Carrera (1999) to derive a binned estimate of the OLF. Although it addresses one of the potential biases of the traditional $`1/V_a`$ statistic (Schmidt 1968), like most binned LF estimators the Page & Carrera (1999) method assumes that absolute magnitude and redshift bins can be chosen sufficently small so that the effects of evolution and a steeply rising OLF are negligable across each bin. For previous samples this has been difficult to achieve whilst still retaining a sufficient number of QSOs in each bin.
With the large numbers of QSOs now available for this analysis, we were able to minimise these effects by choosing much smaller bins in redshift and absolute magnitude. We used 0.25 mag bins and 10 redshift bins equally spaced in $`\mathrm{log}(1+z)`$ over the interval $`0.35<z<2.3`$ to compute the $`1/V`$ estimate of the OLF.
For an OLF that evolves significantly across a $`M_B`$,$`z`$ bin, $`1/V`$ methods will still yield biased estimates of the OLF in bins cut by the magnitude limit of the sample. In such bins, the estimated OLF will only represent the mean space density of objects sampled in bright $`M_B`$, low $`z`$ region of the bin. In the absence of any magnitude limit most objects would tend to lie in the faint $`M_B`$, high $`z`$ part of the bin (at least for a power law LF undergoing strong redshift evolution) and so the estimate of the OLF in bins which are not fully sampled will be biased low. The extent of the bias will, of course, depend on the extent to which the OLF evolves across the bin. Our choice of small bins alleviates, but does not remove this bias and so we have also excluded all $`M_B,z`$ bins which contain the $`b_\mathrm{J}20.8`$ mag survey limit in our binned estimate of the OLF. A similar bias also affects (but in the opposite sense) bins cut by the bright mangitude limit of the sample. However, we also imposed the constraint that we would not compute the OLF in bins where there were five or fewer objects. This latter conditions prevents any bins cut by the bright survey limit from appearing in estimate of the OLF below.
The resulting OLF, calculated for a flat universe with $`q_0=0.5`$ and a $`H_0=50\mathrm{k}\mathrm{m}\mathrm{s}^1\mathrm{Mpc}^1`$, in plotted Figure 4. Absolute $`M_\mathrm{B}`$ magnitudes were derived for the QSOs using the k-corrections derived by Cristiani & Vio (1990). The median number of QSOs in each plotted bin is 40, although some bins contain up to 170 QSOs. The shape and redshift dependence of the OLF in Figure 4 are strongly suggestive of a luminosity evolution model similar to those derived previously from the Durham/AAT sample (Boyle et al. 1988).
### 3.2 Maximum Likelihood Analysis
#### 3.2.1 Method
To obtain a more quantitative descriptions of suitable models, we carried out a maximum likelihood fitting procedure for a number of models to the data (see Boyle et al. 1988). This technique relies on minimizing the likelihood function $`S`$ corresponding to the Poisson probability distribution function for both model and data (Marshall et al. 1983). We tested the goodness-of-fit of the model to the data using the 2D Kolmogorov-Smirnoff (KS) statistic. The KS is notoriously insensitive to discrepancies between the data and the model predictions in the wings of the distributions. This problem is particularly severe in the present analysis, where the weakness of the KS test at the brightest absolute magnitudes is compounded by the vanishingly small fraction of objects in this region of the $`M_B,z`$ plane in the combined sample (see Figure 3).
To alleviate this problem, we derived 2D KS statistics from each survey separately; combining them into a final KS probability using the $`Z`$ statisitic described by Peacock (1983). This approach enables each of the model predictions to be tested against the data in specific regions of the ($`M_\mathrm{B},z`$) plane sampled by different surveys. A significant rejection of the model, even by a relatively small sample, will thus have major impact on the overall acceptability of the fit.
In the fitting procedure, we used no more free parameters in any model than was required to obtain an acceptable fit. We defined $`apriori`$ an acceptable fit as one which could not be rejected at the 99 per cent confidence level or greater i.e., a KS probability ($`P_{\mathrm{KS}}`$) of 1 per cent or greater. Errors on the fit parameters correspond to the $`\mathrm{\Delta }S=1`$ contours around each parameter, or, equivalently, the 68 per cent confidence contour for one interesting parameter.
#### 3.2.2 Models
Guided by the appearance of the OLF in Fig 4 we chose to model the OLF $`\mathrm{\Phi }(L,z)`$ as a two-power-law in luminosity,<sup>1</sup><sup>1</sup>1All simpler forms of the OLF e.g., single power law, Schechter function were strongly ruled out ($`P_{\mathrm{KS}}<10^{10}`$) in preliminary fits.
$$\mathrm{\Phi }(L_\mathrm{B},z)=\frac{\mathrm{\Phi }(L_\mathrm{B}^{})}{[(L_\mathrm{B}/L_\mathrm{B}^{})^\alpha +(L_\mathrm{B}/L_\mathrm{B}^{})^\beta ]}.$$
Expressed in magnitudes this becomes
$$\mathrm{\Phi }(M_\mathrm{B},z)=\frac{\mathrm{\Phi }(M_\mathrm{B}^{})}{10^{0.4[(\alpha +1)(M_\mathrm{B}M_\mathrm{B}^{}(z))]}+10^{0.4[(\beta +1)(M_\mathrm{B}M_\mathrm{B}^{}(z))]}},$$
where the evolution is given by the redshift dependence of the break luminosity $`L_\mathrm{B}^{}`$ or magnitude, $`M_\mathrm{B}^{}(z)`$.
We first attempted to fit the evolution by using ‘standard’ models, i.e., an exponential luminosity evolution with fractional look-back time ($`\tau `$), such that:
$$L_\mathrm{B}^{}(z)=L_\mathrm{B}^{}(0)\mathrm{exp}(k_1\tau ),$$
and a power-law luminosity evolution model with a redshift cut-off:
$$L_\mathrm{B}^{}(z)=L_\mathrm{B}^{}(0)(1+z)^{k_1}z<z_{\mathrm{max}}$$
$$L_\mathrm{B}^{}(z)=L_\mathrm{B}^{}(0)(1+z_{\mathrm{max}})^{k_1}zz_{\mathrm{max}}$$
In contrast to previous analyses, the power-law evolution model provided a poor fit to all datasets ($`P_{\mathrm{KS}}0.01`$) for all cosmological models. This implied that any decline in the evolution at high redshifts $`(z2`$) was not well represented by an abrupt cut-off at $`z_{\mathrm{max}}`$. We therefore adopted a more general exponential evolution law incorporating a second-order polynomial function (hereinafter referred to as polynomial evolution) that gives a smoother transition in the power law behaviour of the evolution at $`z>2`$ as allows for the possibility of negative evolution at high redshift <sup>2</sup><sup>2</sup>2Polynomial evolution models were first fit to the radio galaxy/QSO LFs by Dunlop & Peacock (1990):
$$L_\mathrm{B}^{}(z)=L_\mathrm{B}^{}(0)10^{k_1z+k_2z^2},$$
or equivalently,
$$M_\mathrm{B}^{}(z)=M_\mathrm{B}^{}(0)2.5(k_1z+k_2z^2).$$
#### 3.2.3 Results
The best-fit parameters, and KS probabilities for the polynomial and exponential evolution fits are given in Table 2. The statistical errors on individual parameters are $`\mathrm{\Delta }\alpha \pm 0.05`$, $`\mathrm{\Delta }\beta \pm 0.1`$, $`\mathrm{\Delta }M_\mathrm{B}^{}\pm 0.2`$, $`\mathrm{\Delta }k_1(\mathrm{exponential})\pm 0.1`$, $`\mathrm{\Delta }k_1(\mathrm{polynomial})\pm 0.05`$, $`\mathrm{\Delta }k_2\pm 0.02`$.
Acceptable fits ($`P_{\mathrm{KS}}>0.01`$) to the primary 2QZ + LBQS sample were found for both evolution models. The exponential model favoured low $`q_0`$ and the polynomial model favoured high $`q_0`$.
The predicted differential number-magnitude, n(m), and number-redshift, n(z), relations for the 2QZ survey based on the best-fit $`q_0=0.5`$ polynomial evolution model to the 2QZ + LBQS samples are shown in Figure 5. The model predictions provide a good fit to the derived n(m) and observed n(z) relations for QSOs with $`0.35<z<2.3`$ and $`M_B<23`$ from the 2QZ survey also plotted in this figure.
The extrapolated maximum $`L_\mathrm{B}^{}`$ in the polynomial evolution model occurs in the range $`2.46<z<2.50`$ for $`0.05<q_0<0.5`$. The behaviour of $`L_\mathrm{B}^{}`$ as a function of fractional look-back time for the $`q_0=0.5`$ universe is shown in Figure 6.
In Figure 7 we have also plotted the Poisson significance of the residual difference between the best-fit $`q_0=0.5`$ polynomial evolution model and 2QZ + LBQS data-set for the $`M_B,z`$ bins used in the $`1/V`$ analysis of the OLF. It can be seen that, within the range of luminosity-redshift parameter space covered by the data, there are no significant ($`>3\sigma `$) differences between the model and data. Although there may be a possible weak trend at low redshifts for the model to over-predict the numbers of QSOs at low luminosities and to under-predict numbers at high luminosities, any such discrepancies are not yet statistically significant with this current dataset.
We found no significant difference in the values of the best-fit parameters when the fits were restricted to $`z>0.5`$ or $`z<2.0`$, i.e., excluding the sample regions where we apply the largest correction for incompleteness. The extrapolated peak in $`L_\mathrm{B}^{}`$ occurs at the same redshift ($`z=2.5`$) for model fits to both the $`z<2`$ and $`z<2.3`$ samples see (Figure 6).
The addition of the bright samples yielded much poorer overall fits. Inclusion of the PG sample gave, at best, marginally acceptable fits for the exponential evolution model ($`q_0=0.05`$) and polynomial evolution model ($`q_0=0.5`$). No acceptable fit was found for a $`q_0=0.5`$ universe with the incorportion of HEQS and EQS, and only a barely acceptable fit ($`P_{\mathrm{KS}}=0.01`$) was obtained for the $`q_0=0.05`$ exponential evolution model.
Previous attempts to characterize departures from pure luminosity evolution (PLE) have led to models with a redshift-dependent bright-end slope (La Franca & Cristiani 1997). We attempted to fit such a model for the polynomial evolution models to the data-set including the HEQS and EQS samples using a simple redshift dependence of the form
$$\alpha (z)=\alpha (0)+\kappa _3z.$$
For $`q_0=0.5`$ the best fit value of $`\kappa _3=0.36`$, while for $`q_0=0.05`$, $`\kappa _3=0.0`$. However, this model did not provide an acceptable fit to the data-set ($`P_{\mathrm{KS}}<0.01`$) for either $`q_0=0.5`$ or $`q_0=0.05`$.
By virtue of its size, rejection of the fits by the brighter data-sets was dominated by the EQS. Nevertheless the EQS contains relatively few QSOs, and a full understanding of the evolution of the very brightest QSOs clearly still requires both larger samples and a more detailed knowledge of their properties. This should be available with the completed HEQS.
There could be a variety of reasons why luminous QSOs may exhibit deparatures from pure luminosity evolution models. For example, bright QSO samples are known to contain a greater fraction of radio-loud QSOs (Peacock, Miller & Longair 1986) and also contain a greater fraction of gravitationally lensed QSOs (Kochanek 1991). Both effects could give rise to systematic departures from simple luminosity evolution models fit predominantly to less luminous QSOs. Equally, these results may imply that luminous QSOs simply evolve differently from the bulk of ’normal’ QSOs.
We also explored the effect of different cosmological models on the OLF for our primary 2QZ + LBQS sample. The results are reported in Table 3. For an $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ universe we fit the polynomial evolution model to the data for a variety of different values of $`q_0`$. All values of $`q_0`$ in the range $`0<q_0<0.5`$ yielded acceptable fits, with probabilty reaching a broad maximum in the the region $`q_0=0.30.4`$. However, this value is highly dependent on the evolution model chosen to fit the data. For example the exponential evolution model favours lower values of $`q_0`$ ($`q_0<0.1`$, see Table 2). It is unlikely that the QSO OLF can be used to define useful contraints on the value of $`q_0`$ until a meaningful physical model for QSO evolution is available.
The best-fit OLF parameters for a non-zero cosmological constant in a flat universe (i.e., $`\mathrm{\Omega }_\mathrm{m}+\mathrm{\Omega }_\mathrm{\Lambda }=1`$) were also obtained using the expressions for comoving distance $`r(z)`$ and comoving volume element $`dV/dz`$ given by:
$$r(z)=\frac{c}{H_0a_0}_0^z\frac{dz}{\sqrt{\mathrm{\Omega }_\mathrm{m}(1+z)^3+\mathrm{\Omega }_\mathrm{\Lambda }}}$$
$$\frac{dV}{dz}(z)=r(z)^2\sqrt{\mathrm{\Omega }_\mathrm{m}(1+z)^3+\mathrm{\Omega }_\mathrm{\Lambda }}$$
A variety of fits for different values of $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are presented in Table 3. Luminosity evolution remains a good fit to the 2QZ + LBQS dataset in all cases. As above, it is apparent the OLF provides little potential for discrimination between cosmological models with a non-zero $`\mathrm{\Omega }_\mathrm{\Lambda }`$.
## 4 CONCLUSIONS
For absolute magnitudes $`26<M_B<23`$ ($`q_0=0.5`$) and redshifts $`0.35<z<2.3`$, pure luminosity evolution (PLE) has been shown to produce an acceptable fit to the QSO distribution. This region of parameter space contains the QSOs reponsible for the vast majority of the AGN luminosity density in the redshift range between $`0.35<z<2.3`$. The best-fitting PLE models exhibit an exponential increase in luminosity with either look-back time ($`q_0=0.05`$) or as a second-order polynomial function of redshift ($`q_0=0.5`$). In the near future, the space density of $`z>4`$ QSOs predicted from extrapolations of these models may be usefully compared with the significant numbers of such QSOs now being identified at these redshifts by the Sloan Digital Sky Survey (Fan et al. 1999).
Inclusion of bright surveys such as the EQS and HEQS do result in departures from luminosity evolution similar in form to those seen by other authors (e.g. Hewett et al. 1993, La Franca & Cristaini 1997, Goldschmidt & Miller 1998). The principal region of parameter space in which such departures are seen lies at $`z<0.5`$ where those authors claim a significant flattening of the luminosity function. The 2QZ contains relatively few QSOs with $`0.35<z<0.5`$ and does not yet probe the redshift range $`z<0.35`$. Hence there is at present no inconsistency between the various studies. Our analysis of the 2dF survey does however show that, over the redshift range $`0.35<z<2.3`$ and for absolute magnitude $`M_B<23`$, pure luminosity evolution does provide an accurate phenomenological model of QSO evolution.
## 5 ACKNOWLEDGMENTS
The 2QZ was based on observations made with the Anglo-Australian Telescope and the UK Schmidt Telescope. We are indebted to Mike Hawkins for providing us with colour information on his QSO sample prior to publication. NL was supported by a PPARC studentship during the course of this work.
This paper has been produced using the Blackwell Scientific Publications macros.
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# Effects of Fermi energy, dot size and leads width on weak localization in chaotic quantum dots
## I Introduction
Experimental studies of magnetoconductance in quantum dots show that, at low magnetic fields (typically below one flux quantum), the conductance increases with the field . The effect has been investigated theoretically and related to a similar behavior observed in disordered metallic conductors in the diffusive regime, that is referred to as Weak Localization (WL) .
There is also fairly conclusive experimental evidence which indicates that the average magnetoconductance $`<G(B)>`$ behaves in a qualitatively different way in regular and chaotic cavities, namely, whereas in the former it increases linearly with $`B`$, in chaotic cavities the WL peak has a Lorentzian shape . Semiclassical analyses ascribe this difference to the distributions of the areas $`A`$ enclosed by the trajectories of the carriers . While in regular systems the probability distribution of enclosed areas larger than $`A`$ is $`1/A`$ , in fully chaotic systems it is exponential . As a consequence, in chaotic cavities the increment in the magnetoconductance as a function of magnetic flux $`\mathrm{\Phi }`$ is given by:
$$\delta G=G(\mathrm{\Phi })G(0)=\frac{a\mathrm{\Phi }^2}{1+b\mathrm{\Phi }^2},$$
(1)
where the conductance and the magnetic flux are given in units of their respective quanta, $`G_0=e^2/h`$ and $`\mathrm{\Phi }_0=h/e`$. The constant $`b`$ gives the critical flux at which the time–reversal symmetry is effectively destroyed, $`\mathrm{\Phi }_c=1/\sqrt{b}`$, whereas the ratio $`a/b`$ gives the weak localization term, i.e., $`G(\mathrm{})G(0)=a/b`$. The supersymmetric $`\sigma `$–model predicts that $`a`$ and $`b`$ should be inversely proportional to the number of channels $`N_{ch}`$ that contribute to the current ,
$$b=c\frac{\nu D_0}{N_{ch}}\frac{L}{N_{ch}}$$
(2)
where $`c`$ is a constant (for the present geometry $`c=2\pi /3`$), $`\nu `$ the density of states, $`D_0`$ the diffusion coefficient and $`L`$ the linear size of the cavity. The size dependence arises from the standard expression for the diffusion coefficient $`D_0=v_Fl/2`$, where $`v_F`$ is the Fermi velocity, and $`l`$ the elastic mean free path, and from the fact that in a two–dimensional ballistic system, $`lL`$ . The qualitative behavior of Eq. (2) is similar to the Random Matrix Theory (RMT) result for the critical flux at which the time–reversal symmetry is broken (GOE–GUE transition) reported in . The two constants $`a`$ and $`b`$ are proportional to each other. In particular RMT gives ,
$$\frac{a}{b}=\frac{N_{ch}}{4N_{ch}+2},$$
(3)
where $`N_{ch}`$ is related to the zero field conductance through,
$$G_{\mathrm{RMT}}(0)=\frac{N_{\mathrm{ch}}}{2}\frac{N_{\mathrm{ch}}}{4N_{\mathrm{ch}}+2},$$
(4)
On the other hand, a fitting of the numerical results obtained from a random matrix model Hamiltonian gave
$$b=2k\frac{2N_{ch}1}{N_{ch}^2},$$
(5)
$`k`$ being a constant which, as in Eq. (2), depends on the Fermi energy. Although Eq. (5) gives the same dependence on the number of channels than Eq. (2) in the large $`N_{ch}`$ limit, it does not explicitly reproduce neither its size nor its energy dependence. Moreover, as remarked in , Eq. (5) is only valid for few channel ballistic cavities. It should also be mentioned that the size dependence of the constant $`b`$ has also been obtained within RMT (see ).
At present there is no published numerical study of the effects of the size of the cavity, the leads width, and the Fermi energy on weak localization in reasonably realistic models of quantum chaotic cavities. The purpose of this work is to discuss the results of such an investigation. Quantum dots are described by means of a tight-binding Hamiltonian on $`L\times L`$ clusters of the square lattice. Non–regular (chaotic) behavior is induced by introducing a number of bulk vacancies proportional to the linear size of the system . This model has been shown to behave similarly to dots in which chaoticity is induced by introducing disorder at the surface . The effects of leads width $`W`$, system size, and number of channels that contribute to the current are discussed in detail. Our results show that the the critical flux is not simply proportional to the square root of the number of open channels as concluded in ; it turns out that this relationship is obscured by the strong energy dependence of the proportionality constant already implicit in Eq. (2). Significant deviations from RMT are observed for large leads width ($`W`$ of the order of the system size $`L`$). In particular the weak localization term decreases as $`W`$ approaches $`L`$. However, our numerical data for $`W=L`$ indicate that this term does not vanish as $`L`$ increases.
The paper is organized as follows. Section II includes a description of our model of chaotic quantum dot and of the method we used to compute the current. The results are discussed in Section III. We first briefly consider the case of zero field, comparing our results with those derived from RMT. The results concerning the effects of Fermi energy, leads width and dot size are presented and discussed thereafter. Again, comparison with RMT is highlighted. Section IV is devoted to summarize the conclusions of our work.
## II Model and Procedures
### A Model of quantum chaotic dot
Our model of a quantum chaotic dot is described by means of a tight–binding Hamiltonian with a single atomic orbital per lattice site,
$`\widehat{H}={\displaystyle \underset{<m,n;m^{},n^{}>}{}}t_{m,n;m^{},n^{}}|m,n><m^{},n^{}|,`$ (6)
where $`|m,n>`$ represents an atomic orbital on site $`(m,n)`$. Indexes run from 1 to $`L`$, and the symbol $`<>`$ denotes that the sum is restricted to the existing nearest-neighbors of site $`(m,n)`$. Using Landau’s gauge the hopping integral is $`t_{m,n;m^{},n^{}}=\mathrm{exp}\left(2\pi i\frac{m}{(L1)^2}\frac{\mathrm{\Phi }}{\mathrm{\Phi }_0}\right)`$, for $`m=m^{}`$, and 1 otherwise. Therefore, the difference between our Hamiltonian $`H`$ and the one corresponding to an ideal $`L\times L`$ cluster on the square lattice is the absence of hopping to and from $`L`$ sites chosen at random among the $`L^2`$ sites defining the lattice. A full discussion of the properties of this model for the case of a closed system and zero field can be found in Ref. .
### B Conductance
The conductance (measured in units of the quantum of conductance $`G_0=e^2/h`$) was computed by using the implementation of Kubo formula described in Ref. (applications to mesoscopic systems can be found in Refs. and ). For a current propagating in the $`x`$–direction, the static electrical conductivity is given by:
$$G=2\left(\frac{e^2}{h}\right)\mathrm{Tr}\left[(\mathrm{}\widehat{v}_x)\mathrm{Im}\widehat{}G(E)(\mathrm{}\widehat{v}_x)\mathrm{Im}\widehat{}G(E)\right],$$
(7)
where $`\mathrm{Im}\widehat{}G(E)`$ is obtained from the advanced and retarded Green functions:
$$\mathrm{Im}\widehat{}G(E)=\frac{1}{2i}\left[\widehat{}G^R(E)\widehat{}G^A(E)\right],$$
(8)
and the velocity (current) operator $`\widehat{v}_x`$ is related to the position operator $`\widehat{x}`$ through the equation of motion $`\mathrm{}\widehat{v}_x=[\widehat{H},\widehat{x}]`$, $`\widehat{H}`$ being the Hamiltonian.
Numerical calculations were carried out connecting quantum dots to semiinfinite leads of width W in the range 1–$`L`$. The hopping integral inside the leads and between leads and dot at the contact sites is taken equal to that in the quantum dot (ballistic case). Assuming the validity of both the one-electron approximation and linear response, the exact form of the electric field does not change the value of $`G`$. An abrupt potential drop at one of the two junctions provides the simplest numerical implementation of the Kubo formula since, in this case, the velocity operator has finite matrix elements on only two adjacent layers and Green functions are just needed for this restricted subset of sites. Assuming this potential drop to occur at the left contact ($`lc`$) side, the velocity operator can be explicitly written as,
$$i\mathrm{}v_x=\underset{j=1}{\overset{W}{}}\left(|lc,j><1,j||1,j><lc,j|\right)$$
(9)
where $`(|lc,j>`$ are the atomic orbitals at the left contact sites nearest neighbors to the dot.
Green functions are given by:
$$[E\widehat{I}\widehat{H}\widehat{\mathrm{\Sigma }}_1(E)\widehat{\mathrm{\Sigma }}_2(E)]\widehat{}G(E)=\widehat{I},$$
(10)
where $`\widehat{\mathrm{\Sigma }}_{1,2}(E)`$ are the selfenergies introduced by the two semiinfinite leads . The explicit form of the retarded selfenergy due to the mode of wavevector $`k_y`$ is:
$$\mathrm{\Sigma }(E)=\frac{1}{2}\left(Eϵ(k_y)i\sqrt{4(Eϵ(k_y))^2}\right),$$
(11)
for energies within its band $`|Eϵ(k_y)|<2`$, where $`ϵ(k_y)=2\mathrm{c}\mathrm{o}\mathrm{s}(k_y)`$ is the eigenenergy of the mode $`k_y`$ which is quantized as $`k_y=(n_{k_y}\pi )/(W+1)`$, $`n_{k_y}`$ being an integer from 1 to $`W`$. The transformation from the normal modes to the local tight–binding basis is obtained from the amplitudes of the normal modes, $`<n|k_y>=\sqrt{2/(W+1)}\mathrm{sin}(nk_y)`$. Note that in writing Eq. (11) we assumed that the magnetic field was zero outside the dot .
### C Numerical Procedures
Input/output leads were attached at opposite corners of the dot as follows: input lead connected from site $`(1,1)`$ to site $`(1,1+W)`$, and output lead from site $`(L,1)`$ to site $`(L,1+W)`$. We have checked that changing the sites at which leads are attached does not qualitatively modify the results discussed here. The conductance was averaged over disorder realizations (local distribution of vacancies) and within selected energy ranges. The latter were chosen to fit the number of channels in the leads. More specifically, for leads with $`N_{\mathrm{ch}}`$ channels energy averages were carried in the range,
$$E[E_{N_{\mathrm{ch}}},E_{N_{\mathrm{ch}}+1}],$$
(12)
for $`N_{\mathrm{ch}}`$ channels in the leads, where,
$$E_n=2\left(1+\mathrm{cos}\frac{\pi n}{W+1}\right).$$
(13)
Some calculations were also carried out at a fixed Fermi energy. In all cases averages were done over at least 1200 values of the conductance.
## III Results
### A Zero field conductance
Fig. 1 shows relative deviations of the conductance with respect to the RMT result (see Eq. (4)) for narrow and rather wide leads as a function of the dot size $`L`$. It is noted that for small $`W`$ deviations are always smaller than 5%, and typically below 2%. The results fluctuate more appreciably for the narrower lead ($`W`$=1) as expected . Relative deviations from the RMT result are significantly larger for $`W=`$ 9 and 18. The results of Fig. 1 suggest that the difference with respect RMT is not a size effect. The larger deviation, and stronger variation in the explored range of $`L`$, observed for $`N_{\mathrm{ch}}=W=9`$ is likely a consequence of the important contribution that the center of the band ($`E=0`$) has in that case. Both the center of the band and its bottom ($`E=4`$) show rather odd behaviors. In particular at $`E=0`$ no weak localization effect was observed (see below).
The change in the zero field conductance as the number of channels is varied, for fixed dot size, is illustrated in Fig. 2. The results for $`N_{\mathrm{ch}}=W/2`$ can be accurately fitted by means of a straight line, as expected, although the slope is smaller than the RMT prediction (see caption of Fig. 2 and Eq. 4). The slope of the straight line varies with the ratio $`N_{\mathrm{ch}}/W`$, or alternatively the average Fermi energy; for instance at $`E=0`$ it is actually larger than 0.5. For fixed leads width $`W=22`$ and a variable number of channels a large deviation with respect to a straight line is instead observed. This deviation, which is a consequence of the concomitant change in the Fermi energy as the number of channels is varied, increases with the number of channels, likely due to the increasing importance of the contribution of the band center. Numerical results indicate that at the band center the conductance shows a much stronger dependence (increase) on the dot size that at any other energy within the band, probably due to the building up of the singularity in the density of sates characteristic of the square lattice at that energy. These results suggest that if the $`N_{\mathrm{ch}}`$ dependence has to be investigated it is more reliable to work at a fixed $`N_{\mathrm{ch}}/W`$ ratio and vary the leads width.
### B Magnetoconductance: Weak Localization
We first discuss the energy dependence of the critical flux and of the weak localization term. This was done by investigating rather large $`W`$ and varying the number of channels in each lead. This is equivalent to vary the energy range over which the energy was calculated (see above). Fig. 3 depicts numerical results for cavities of linear size $`L`$=78 and leads of width $`W`$= 22 (a) and 44 (b). The conductance was obtained by averaging over 60 disorder realizations and 21 energies in the ranges corresponding to the number of channels in the leads (see subsection IIC). The weak localization peak shows the expected behavior. The numerical results were fitted by means of Eq. (1). At this stage it is worth noting that the conductance remains constant in a wide range of fluxes only for small $`W`$. For large $`W`$ the conductance follows Eq. (1) in a rather narrow range of $`\mathrm{\Phi }`$ and then increases steadily. This deviation from the Lorentzian–like law hinders the fitting of the numerical results. The fitted parameters are reported in Table I. The parameters derived from RMT (Eqs. (35)) are also given in the Table. Eq. (5) was used with $`k=1`$ as its explicit dependence on energy was not given in Ref. ; this will suffice, however, to illustrate our point concerning the strong energy dependence of that constant. We first note that the weak localization term is smaller than the values predicted by RMT likely due to the large values of $`W`$ (see below). However, the dependence of $`a/b`$ on the number of channels is the correct one (it increases with $`N_{\mathrm{ch}}`$) but for $`W=44`$ and $`N_{\mathrm{ch}}=28`$. The latter deviation is a consequence of the increasing importance of the band center. At that energy the results indicate that $`G(\mathrm{\Phi })`$ decreases as a function of the flux, i.e., there is no weak localization effect. Due to computing limitations we have not been able to check whether this result is a size effect. The numerical results for parameter $`b`$ indicate that it increases with the number of open channels, a behavior opposite to that given by RMT with $`k`$= constant. This suggests that it cannot be safely concluded that the critical flux is proportional to the square root of the number of open channels, as, increasing the Fermi energy, not only increases $`N_{ch}`$ but it also dramatically changes constant $`k`$ in Eq. (5). The energy dependence of $`k`$ already appears in the supersymmetric $`\sigma `$–model result . We have checked that if the energy dependent factor in Eq. (2), namely, $`\nu D_0`$, is included (the mean free path was calculated following the procedure of , see also ) the dependence of the RMT result on the number of channels is reversed, in agreement with our numerical results.
The effects of the dot size on the weak localization peak were investigated for $`L`$ in the range of $`L=27`$–137 and three combinations of $`(W,N_{\mathrm{ch}})`$. Averages were identical to those mentioned in the preceding paragraph. The results are depicted in Figs. 4 and (5). The weak localization term ($`a/b`$) shows a slight size dependence at small $`L`$, saturating for $`L`$ approximately larger than 50 (see Fig. 4). This indicates that the smaller values of $`a/b`$ obtained in our calculations, with respect to RMT, is probably not a size effect. The results for $`(W,N_{\mathrm{ch}})`$=(1,1) are slightly smaller than those for (2,1) surely due to the contribution of the band center in the first case. The results for (10,5) are larger than the other two, in agreement with RMT. On the other hand our results for constant $`b`$ increase with $`L`$ as expected (see Eq. (2)). The numerical results can be reasonably fitted by means of straight lines as shown in Fig. 5. The differences in the slopes is a consequence of the energy dependence discussed above.
In order to get rid as much as possible of the strong energy dependence of the shape of the weak localization peak, we have carried out the study of the effects of the leads width at a fixed energy. We have chosen $`E=2.001`$ (away from the band center and bottom) which approximately correspond to $`N_{\mathrm{ch}}=W/2`$. We fixed the dot size at $`L=78`$ and varied the leads width in the range $`W`$=4–78. The results are shown in Figs. 6-8. The conductance versus the magnetic field for small and large $`W`$ is depicted in Fig.(6). It is readily noted that both the weak localization term and constant $`b`$ (or the inverse of the square root of the critical flux) sharply decreases with $`W`$. Although the results are nicely fitted by means of Eq. 1) the deviation of the numerical results with respect to that equation which occurs at large $`W`$ (see above) is already observed for $`W=78`$ (note that the fitting closely follows the numerical results only up to $`\mathrm{\Phi }1.5`$). The decrease of $`b`$ with $`W`$, or, alternatively, with $`N_{\mathrm{ch}}`$ is illustrated in Fig. 7. The results can be satisfactorily fitted by means of the RMT result (see caption of Fig. 7). On the other hand the weak localization term shows a size dependence that has not been previously anticipated. At small $`W`$ (or number of channels) it increases as predicted by Eq. (3). However, beyond $`W0.2L`$ it begins to decrease sharply reaching a value slightly larger than 0.05 for $`W=L`$. To explore the possibility that the weak localization term vanishes in the large $`L`$ limit we have calculated the magnetoconductance for $`W=L`$, $`E=2.001`$ and $`L`$ in the range 30–126. The numerical results were fitted by means of Eq. (1) with the parameters reported in Table II. The results clearly indicate that $`a/b`$ does not vanish as $`L`$ increases. The fact that $`b`$ is almost independent of $`L`$ is a consequence of the dependence of $`b`$ on the ratio $`L/N_{\mathrm{ch}}`$ (note that by taking $`W=L`$ and a fixed energy the number of channels is proportional to $`L`$).
## IV Concluding Remarks
Summarizing, we have presented what we believe to be the first detailed numerical study of the effects of Fermi energy, leads width and dot size on the shape of the weak localization peak in quantum chaotic cavities. The study was carried out on a model that was recently proposed by us which shows all the expected features of closed chaotic quantum billiards. Although the conclusions of our investigation qualitatively agree with most predictions of random matrix theory, some significant issues have to be highlighted. We first note that our results show that albeit the critical flux is proportional to the square root of the number of open channels, as predicted by RMT, the proportionality constant strongly depends on the Fermi energy in agreement with Efetov’s analysis . This introduces a model (system) dependence which makes theoretical (experimental) comparisons with RMT rather delicate. Our results clearly illustrate the size dependence of the critical flux, in particular $`\mathrm{\Phi }_c1/\sqrt{L}`$, in agreement with Efetov results and the RMT results reported in Refs. (note that this size dependence was not found in a previously published RMT study, see Ref.). Finally, we have investigated the effects of the leads width concluding that the weak localization term sharply decreases with the ratio $`W/L`$ (a result that has not been previously reported), although it is likely that it does not vanish in the infinite $`L`$ limit. This suggest that RMT is probably not valid for sufficiently open systems.
###### Acknowledgements.
This work was supported in part by the Spanish CICYT (grants PB96-0085 and 1FD97–1358). Useful discussions with E. Cuevas and M. Ortuño are gratefully acknowledged.
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# Grothendieck groups and tilting objects
## Introduction
Let $`k`$ be an algebraically closed field and $`𝒞`$ a hereditary abelian $`\mathrm{Ext}`$-finite $`k`$-category. That $`𝒞`$ is hereditary means that $`\mathrm{Ext}^i(,)`$ vanishes for $`i2`$, and we say that $`𝒞`$ is $`\mathrm{Ext}`$-finite if $`\mathrm{Ext}^i(A,B)`$ is finite dimensional over $`k`$ for all $`A,B`$ in $`𝒞`$ and all $`i`$. A central problem in the representation theory of artin algebras is to describe such $`𝒞`$ which have a tilting object. This is important in connection with the investigation of quasitilted algebras, as introduced in .
When a hereditary abelian $`\mathrm{Ext}`$-finite $`k`$-category $`𝒞`$ has a tilting object, it is a consequence that the Grothendieck group $`K_0(𝒞)`$ is free abelian of finite rank . This suggests the problem of describing the $`𝒞`$ for which $`K_0(𝒞)`$ is finitely generated (or free abelian of finite rank), and to decide to which extent having a finitely generated Grothendieck group implies the existence of a tilting object. Relating the existence of a tilting object to properties of the more widely known notion of Grothendieck group provides a better insight into the meaning of the condition of the existence of a tilting object. In particular, it is interesting to understand in terms of Grothendieck groups the special role the category $`\text{coh}𝒳`$ of coherent sheaves on a weighted projective line plays within the larger class of quotient categories of finitely generated graded modules over commutative noetherian isolated singularities of Krull dimension two.
We will use the general classification results from to solve the above problems under the additional hypotheses that $`𝒞`$ is noetherian and has a Serre functor (see section 1). The latter hypothesis is natural since it is a consequence of the existence of a tilting object. Both additional hypotheses are satisfied for the quotient categories mentioned above. Since the properties of $`𝒞`$ having a tilting object, a finitely generated Grothendieck group or a Serre functor are preserved under derived equivalence of hereditary categories our results apply more generally to the $`𝒞`$ with the additional hypothesis of having a Serre functor and being derived equivalent to a noetherian hereditary category.
Some of the results on Grothendieck groups of quotient categories proved in this paper are inspired by similar results for two-dimensional complete noetherian rings, used as a tool for classifying maximal orders of finite representation type in .
Section 1 is devoted to discussing background material from various sources, collected together for the benefit of the reader. In section 2 we describe the $`𝒞`$ with finitely generated Grothendieck group. A new criterion for an object to be a tilting object is given in section 3, which it is interesting also in its own right. In section 4 we construct exceptional collections of modules over a hereditary order over a discrete valuation ring. This is used in section 5, along with the criterion from section 3, to construct a tilting object in the category $`\text{coh}𝒪`$ of coherent modules over a sheaf $`𝒪`$ of hereditary orders over $`^1`$. In section 6 we give our main result on the connection between the existence of a tilting object and the Grothendieck group being finitely generated. Under our assumptions the conditions turn out to be equivalent if $`𝒞`$ is connected and has some object of infinite length. In section 7 we give some examples and comments.
We give an appendix proving directly the relationship between the category of coherent sheaves on a weighted projective line and the above mentioned category $`\text{coh}𝒪`$.
Hereditary abelian $`k`$-categories which are $`\mathrm{Ext}`$-finite, noetherian and have a Serre functor were classified in . It is also of interest to investigate hereditary abelian categories which do not satisfy the additional assumptions, for example with respect to when the Grothendieck group is finitely generated. In Appendix B we give some sources of examples of hereditary abelian categories.
## 1. Background
In this section we provide some background material from various sources, to provide a better understanding of how our work fits in.
### 1.1. Tilting objects
Let $`𝒜`$ be a triangulated category with the property that between two objects only a finite number of $`\mathrm{Ext}`$ are non-zero. If $`T𝒜`$, then add(T) is by definition the smallest additive category containing $`T`$ which is closed under finite direct sums and summands. We say that $`T`$ is a tilting object if $`\mathrm{Ext}_𝒜^i(T,T)=0`$ for $`i0`$ and $`\mathrm{add}T`$ generates $`𝒜`$ (in the sense that $`𝒜`$ is the smallest subcategory of $`𝒜`$ containing $`\mathrm{add}T`$ which is closed under shifts and cones).
If $`𝒞`$ is an abelian category of finite homological dimension then $`T𝒞`$ is a tilting object if it is a tilting object in $`D^b(𝒞)`$. This definition of a tilting object in $`𝒞`$ is equivalent to the usual notion of tilting module (of finite projective dimension) when $`𝒞`$ is the category $`\mathrm{mod}\mathrm{\Lambda }`$ of finitely generated modules for an artin algebra of finite global dimension, as is seen directly or by using . For an $`\mathrm{Ext}`$-finite hereditary abelian $`k`$-category $`𝒞`$ it is equivalent to the following definition used in (reformulating conditions from ): An object $`T`$ in $`𝒞`$ is a tilting object if $`\mathrm{Ext}^1(T,T)=0`$ and if $`\mathrm{Hom}(T,X)=0=\mathrm{Ext}^1(T,X)`$ implies $`X=0`$. In general the definition in is modelled on the definition of a tilting module of projective dimension at most one, and is hence different from ours when $`𝒞`$ is not hereditary.
When $`T`$ is a tilting object in $`\mathrm{mod}\mathrm{\Lambda }`$ for an artin algebra $`\mathrm{\Lambda }`$ of finite global dimension, there is an induced equivalence $`D^b(\mathrm{mod}\mathrm{\Lambda })D^b(\mathrm{mod}\mathrm{End}(T)^{\mathrm{opp}})`$ between bounded derived categories, and similarly if $`T`$ is a tilting object in the category $`\text{coh}X`$ of coherent sheaves on a smooth projective variety $`X`$ . This can easily be extended to the case of $`\text{coh}𝒪`$ where $`𝒪`$ is a coherent $`𝒪_X`$-algebra locally of finite global dimension. The analogous result for tilting objects in $`\mathrm{Ext}`$-finite hereditary abelian $`k`$-categories is given in .
In general if $`T`$ is a tilting object in a triangulated category $`𝒜`$ one may expect an equivalence $`𝒜D^b(\mathrm{End}(T)^{\mathrm{opp}})`$. There is recent work in this direction by Keller, using $`A_{\mathrm{}}`$-categories, and building on . In particular it follows from his results that if $`T`$ is a tilting object in an $`\mathrm{Ext}`$-finite abelian $`k`$-category $`𝒞`$ of finite homological dimension, then there is an equivalence between $`D^b(𝒞)`$ and $`D^b(\mathrm{End}(T)^{\mathrm{opp}})`$.
It follows from this derived equivalence that when $`𝒞`$ is an $`\mathrm{Ext}`$-finite abelian $`k`$-category of finite homological dimension having a tilting object $`T`$ , the Grothendieck group $`K_o(𝒞)`$ is isomorphic to $`^n`$, where $`n`$ is the number of nonisomorphic summands of $`T`$ (see for the hereditary case). Given $`T`$ in $`𝒞`$ with $`\mathrm{Ext}^i(T,T)=0`$ for $`i>0`$, it is an important problem, open even for artin algebras, whether $`T`$ having $`n`$ nonisomorphic summands is sufficient for $`T`$ to be a tilting object. It is known in the case of artin algebras when the projective dimension of $`T`$ is at most one , and for $`𝒞`$ $`\mathrm{Ext}`$-finite hereditary if $`𝒞`$ has some tilting object .
### 1.2. Hereditary categories and weighted projective lines
Hereditary abelian $`\mathrm{Ext}`$-finite categories are of special interest in connection with quasitilted algebras, as introduced in . The quasitilted algebras are by definition the endomorphism algebras End$`(T)^{\mathrm{opp}}`$ when $`T`$ is a tilting object in a hereditary category. Equivalently, an algebra $`\mathrm{\Lambda }`$ is quasitilted if and only if gl.dim.$`\mathrm{\Lambda }2`$ and each indecomposable $`X`$ in $`\mathrm{mod}\mathrm{\Lambda }`$ has projective or injective dimension at most one . Main examples of categories $`𝒞`$ are $`\mathrm{mod}\mathrm{\Lambda }`$ where $`\mathrm{\Lambda }`$ is a hereditary artin algebra and $`\text{coh}X`$ when $`X`$ is a smooth projective curve. More generally there are the coherent sheaves on weighted projective lines , which we discuss next.
Let $`\lambda =(\lambda _1,\lambda _2,\mathrm{},\lambda _n)`$, $`n2`$, be a finite number of points in $`^1`$, with $`\lambda _1=0`$, $`\lambda _2=\mathrm{}`$, $`\lambda _3=1`$. For a sequence $`e=(e_1,\mathrm{},e_n)`$ of positive integers consider the associated ring
(1.1)
$$R=k[x_1,\mathrm{},x_n]/[(x_i^{e_i}x_2^{e_2}+\lambda _ix_1^{e_1})]_{i3}.$$
Let $`H`$ be the abelian group generated by $`h_1,\mathrm{},h_n`$ and with relations $`e_1h_1=\mathrm{}=e_nh_n`$. Then $`H`$ is isomorphic to $`G`$, where $`G`$ is a finite group . $`R`$ is then a $`H`$-graded ring, and Geigle and Lenzing write $`\text{coh}𝒳`$ for the hereditary category gr$`{}_{H}{}^{}R/\text{finite length}`$, where $`\mathrm{gr}_HR`$ denotes the category of finitely generated $`H`$-graded $`R`$-modules with degree zero homomorphisms. Geometrically $`\text{coh}𝒳`$ can be viewed as the coherent sheaves over a (hypothetical) space $`𝒳`$ which is a generalization of $`^1`$. Hence Geigle and Lenzing call $`𝒳`$ a “weighted projective line”.
The category $`\text{coh}𝒳`$ is a noetherian hereditary abelian category with finite dimensional homomorphism and extension spaces, which has a tilting object $`T`$ such that End$`{}_{}{}^{}(T)_{}^{\mathrm{opp}}`$ is a canonical algebra in the sense of Ringel , and actually all canonical algebras are obtained this way. Like for tilting for finite dimensional algebras , there is also induced an equivalence of derived categories $`D^b(\text{coh}𝒳)D^b(\text{ End}{}_{}{}^{}(T)_{}^{\mathrm{opp}})`$ . This setup is used to give an alternative approach to the study of the module theory for canonical algebras, by first investigating the hereditary category $`\text{coh}𝒳`$. The rings described by (1.1) are $`H`$-graded factorial and it is shown in that this property characterizes those rings amongst the two-dimensional rings.
There are two main known sources of connected hereditary categories $`𝒞`$ with tilting object; the module categories of finite dimensional hereditary $`k`$-algebras and the categories $`\text{coh}𝒳`$. In addition there are the hereditary abelian categories derived equivalent to them. It is conjectured that there are no more, and in fact this is proved in for noetherian hereditary categories and more generally in under the assumption that $`𝒞`$ has at least one nonzero object of finite length, or at least one directing object, that is, an object which does not lie on a cycle of nonzero nonisomorphisms.
### 1.3. Noetherian hereditary categories with Serre functor
We recall some essential features of the classification of noetherian hereditary abelian $`\mathrm{Ext}`$-finite $`k`$-categories with Serre functor. For further details we refer to .
Assume that $`𝒞`$ is $`\mathrm{Ext}`$-finite. A Serre functor for $`𝒞`$ is an auto-equivalence $`F:D^b(𝒞)D^b(𝒞)`$ where $`D^b(𝒞)`$ denotes the bounded derived category, such that there are isomorphisms $`\mathrm{Hom}(A,B)\stackrel{}{}\mathrm{Hom}(B,FA)^{}`$ natural in $`A`$ and $`B`$ ($`()^{}`$ is the $`k`$-dual). This clearly implies that $`𝒞`$ has finite homological dimension.
For a hereditary abelian $`\mathrm{Ext}`$-finite $`k`$-category $`𝒞`$ the existence of a Serre functor implies the existence of almost split sequences, and the converse holds if $`𝒞`$ has no non-zero projective or injective objects .
Let be $`𝒞`$ a connected category. It is proved in that if $`𝒞`$ is a connected noetherian hereditary $`\mathrm{Ext}`$-finite $`k`$-category with Serre functor then $`𝒞`$ has one of the following forms.
* A category $`\text{coh}𝒪`$, where $`𝒪`$ is a sheaf of hereditary orders over a smooth projective curve.
* $`𝒞`$ is $`\mathrm{mod}\mathrm{\Lambda }`$ for a finite dimensional hereditary $`k`$-algebra $`\mathrm{\Lambda }`$.
* $`𝒞`$ is the category of finite dimensional representations of the quiver $`\stackrel{~}{A}_n`$ with cyclic orientation with $`n<\mathrm{}`$.
* $`𝒞`$ is derived equivalent to a hereditary category where all objects have finite length and having an infinite number of nonisomorphic simple objects.
The actual result in also contains a precise classification of the categories in (iv). We also recall from that for the categories in (i) there is an alternative description as follows.
* Categories of the form $`\mathrm{qgr}S`$ = $`\mathrm{gr}S`$/finite length, where $`\mathrm{gr}S`$ denotes the category of finitely generated graded modules over a commutative noetherian $``$-graded domain $`S=k+S_1+S_2+\mathrm{}+S_i+\mathrm{}`$ of Krull dimension two which is finite over its center, where the $`S_i`$ are finite dimensional over $`k`$, and $`S`$ is an isolated singularity.
The categories $`\text{coh}𝒳`$ have also as we have seen a similar description, as quotient categories starting with an $`H`$-graded ring, where $`H`$ is not necessarily $``$. But it follows from that one can assume that the rings are $``$-graded, so that the class $`\text{coh}𝒳`$ is a subclass of the $`\mathrm{qgr}S`$.
### 1.4. Classical hereditary orders.
In this section we collect some well-known properties of hereditary orders. We will loosely refer to a classical hereditary order as an order $`\mathrm{\Lambda }`$ in a central simple algebra $`A`$ over a field $`K`$ which is hereditary. Let $`R`$ be the center of $`\mathrm{\Lambda }`$. According to $`R`$ is a Dedekind ring. (This fact does not seem to be contained in exactly this form in ).
Assume that $`R`$ is a discrete valuation ring with maximal ideal $`m`$. Then according to the radical $`I`$ of $`\mathrm{\Lambda }`$ is invertible. Furthermore by there is an integer $`e`$ such that $`I^e=m\mathrm{\Lambda }`$, called the ramification index of $`\mathrm{\Lambda }/R`$. It follows from the structure theory of hereditary orders in that if $`e=1`$ then $`\mathrm{\Lambda }`$ is maximal. If the converse is true then we say that $`A/K`$ is unramified. This happens for example if $`A=M_n(K)`$.
If $`R`$ is not a discrete valuation ring then by localizing one defines ramification indices $`e_P`$ for the non-zero primes in $`R`$. By analyzing $`𝒟=\mathrm{Hom}(\mathrm{\Lambda },R)`$ it follows easily that $`\mathrm{\Lambda }/R`$ ramifies in only a finite number of primes.
## 2. Finitely generated Grothendieck groups
Let $`𝒞`$ be a noetherian $`\mathrm{Ext}`$-finite hereditary abelian $`k`$-category with Serre functor, where $`k`$ is an arbitrary field. In this section we describe which $`𝒞`$ have finitely generated Grothendieck group. For this we use the classification theorem from in the form recalled in §1.3. The main problem we need to deal with is when the category $`\text{coh}𝒪`$ of coherent modules over a sheaf $`𝒪`$ of hereditary orders over a smooth projective curve $`X`$ has finitely generated Grothendieck group.
Let $`X`$ be a regular connected curve over a field $`k`$. Let $`K`$ be the function field of $`X`$ and let $`A`$ be a central simple algebra over $`K`$. Let $`𝒪`$ be a sheaf of hereditary orders in $`A`$ over $`𝒪_X`$. Thus locally $`𝒪`$ is a hereditary order over a Dedekind ring (in the sense of ).
If $`𝒜`$ is a sheaf of rings on a topological space $`Z`$ then we use the notation $`\mathrm{coh}(𝒜)`$ for the category of coherent $`𝒜`$-modules and we write $`K_0(𝒜)`$ for $`K_0(\mathrm{coh}(𝒜))`$.
Our first aim is to give some results on $`K_0(𝒪)`$.
###### Proposition 2.1.
Let $`𝒪`$ be as above. Let $`x_1,\mathrm{},x_t`$ be the points in which $`𝒪`$ ramifies, and let $`e_1,\mathrm{},e_t`$ be the corresponding ramification indices (see §1.4).
Put
(2.1)
$$r=\underset{i}{}(e_i1).$$
Let $`\overline{𝒪}`$ be a maximal order lying over $`𝒪`$. Then
$$K_0(𝒪)K_0(\overline{𝒪})^r.$$
If $`k`$ is algebraically closed then
$$K_0(𝒪)K_0(𝒪_X)^r.$$
###### Proof.
The hereditary orders in $`A`$ containing $`𝒪`$ form a partially ordered set which we will denote by $`(𝒪)`$.
If $`𝒪^{}(𝒪)`$ lies minimally over $`𝒪`$ then one proves exactly as in \[30, Thm 1.14\] that $`K_0(𝒪)=K_0(𝒪^{})`$.
Let $`\overline{𝒪}`$ be a maximal order lying over $`𝒪`$. We deduce that $`K_0(𝒪)^rK_0(\overline{𝒪})`$, where $`r`$ is the length of a maximal chain in $`(𝒪)`$, starting in $`𝒪`$ and ending in $`\overline{𝒪}`$.
A local computation shows that $`r`$ is given by the formula (2.1), which finishes the proof of the computation of $`K_0(𝒪)`$.
If $`k`$ is algebraically closed then by Tsen’s theorem \[8, p. 374\] one has that $`\overline{𝒪}\mathrm{End}_{𝒪_X}()`$ where $``$ is a vector bundle of rank $`n`$ on $`X`$. Hence by Morita theory $`K_0(\overline{𝒪})K_0(𝒪_X)`$. ∎
###### Corollary 2.2.
Assume $`k`$ is a algebraically closed. Then $`K_0(𝒪)`$ is finitely generated if and only if $`X`$ is an open subset of $`^1`$.
###### Proof.
By the previous proposition it suffices to prove this for $`𝒪=𝒪_X`$. Let $`\overline{X}`$ be the regular projective curve associated to the function field of $`X`$ . Then $`\overline{X}`$ is a regular compactification of $`X`$. In particular $`\overline{X}X`$ is a finite number of points, whence by the localization sequence $`K_0(𝒪_X)`$ is finitely generated if and only if $`K_0(𝒪_{\overline{X}})`$ is finitely generated. Hence we may assume that $`X`$ is projective. By \[20, Ex. II.6.12, Rem. IV.4.10.4\] one has $`K_0(𝒪_X)^2J(k)`$ where $`J(k)`$ denotes the $`k`$-points of the Jacobian of $`X`$. It is well-known that $`J(k)`$ is not finitely generated if $`X^1`$ (for example because in that case $`J(k)`$ is non-trivial and divisible ). ∎
Combining with §1.3 we now get the following main result of this section.
###### Theorem 2.3.
Let $`𝒞`$ be a connected noetherian $`\mathrm{Ext}`$-finite hereditary abelian $`k`$-category with Serre functor where $`k`$ is an algebraically closed field. Then $`K_o(𝒞)`$ is finitely generated (free abelian) if and only if $`𝒞`$ has one of the following forms.
1. $`\mathrm{mod}\mathrm{\Lambda }`$ where $`\mathrm{\Lambda }`$ is an indecomposable finite dimensional hereditary $`k`$-algebra.
2. Finite dimensional representations over $`\stackrel{~}{A}_n`$ with $`n`$ finite and cyclic orientation.
3. $`\text{coh}𝒪`$ where $`𝒪`$ is a sheaf of hereditary $`𝒪_X`$-orders with $`X=^1`$.
###### Proof.
It is clear that the categories in 1. and 2. have finitely generated (free abelian) Grothendieck groups, and that $`K_o(𝒞)`$ is not finitely generated if $`𝒞`$ is derived equivalent to a hereditary category $`𝒞^{}`$ with all objects of finite length and an infinite number of nonisomorphic simple objects. In view of §1.3 the proof is completed by using Corollary 2.2. ∎
## 3. A criterion for deciding if an object is a tilting object.
The aim of this section is to give a criterion for an object of projective dimension at most one in an abelian category to be a tilting object. For this we need to recall some results on semiorthogonal pairs in triangulated categories from .
Let $`𝒜`$ be a triangulated category and let $``$, $`𝒞`$ be two strict ($`=`$ closed under isomorphisms) full triangulated subcategories of $`𝒜`$. $`(,𝒞)`$ is said to be a *semi-orthogonal pair* if $`\mathrm{Hom}_𝒜(B,C)=0`$ for $`B`$ and $`C𝒞`$. Define
$$^{}=\{A𝒜B:\mathrm{Hom}_𝒜(B,A)=0\}$$
$`{}_{}{}^{}𝒞`$ is defined similarly.
If $`𝒮`$ is a class of objects in $`𝒜`$ then the (triangulated) category generated by $`𝒮`$ is the smallest subcategory of $`𝒜`$ which is closed under shifts, cones, and isomorphisms.
The following result is a slight variation of the statement of \[4, Lemma 3.3.1\] (see also \[5, §1\]).
###### Lemma 3.1.
The following conditions are equivalent for a semi-orthogonal pair $`(,𝒞)`$.
1. $``$ and $`𝒞`$ generate $`𝒜`$.
2. For every $`A𝒜`$ there exists a distinguished triangle $`BAC`$ with $`B`$ and $`C𝒞`$.
3. $`𝒞=^{}`$ and the inclusion functor $`i_{}:𝒜`$ has a right adjoint $`i^!:𝒜`$.
4. $`={}_{}{}^{}𝒞`$ and the inclusion functor $`j_{}:𝒞𝒜`$ has a left adjoint $`j^{}:𝒜𝒞`$.
If one of these conditions holds then the functors $`i^!,j^{}`$ are exact and the triangles in 2. are (for a fixed $`A`$) unique up to unique isomorphism. They are necessarily of the form
(3.1)
$$i_{}i^!AAj_{}j^{}A$$
where the maps are obtained by adjointness from the identity maps $`i^!Ai^!A`$ and $`j^{}Aj^{}A`$. In particular triangles as in 2. are functorial.
If any of the conditions of the previous lemma holds then we say that $`(,𝒞)`$ is a semi-orthogonal decomposition of $`𝒜`$.
Now for the rest of this section let $`k`$ be a field. All categories (abelian or triangulated) will be $`k`$-linear and have finite dimensional $`\mathrm{Hom}`$’s and $`\mathrm{Ext}`$’s. We assume furthermore that for any pair $`A,B`$, there are only a finite number of non-zero $`\mathrm{Ext}^i(A,B)`$.
Let $`𝒜`$ be a triangulated or abelian category. For an object $`T`$ in $`𝒜`$ we denote by $`T^{}`$ the full subcategory of $`𝒜`$ whose objects are the $`C`$ in $`𝒜`$ with $`\mathrm{Ext}^i(T,C)=0`$ for all $`i`$. We say that an object $`T𝒜`$ is exceptional if $`\mathrm{Ext}_𝒜^i(T,T)=0`$ for $`i>0`$ and $`\mathrm{End}_𝒜(T)`$ is a (finite dimensional) division algebra. A sequence of exceptional objects $`T_1,\mathrm{},T_n`$ is an exceptional collection if $`\mathrm{Ext}_𝒜^{}(T_i,T_j)=0`$ for $`j>i`$. An exceptional collection is strongly exceptional if $`\mathrm{Ext}_𝒜^t(T_i,T_j)=0`$ for $`t>0`$ and all $`i,j`$.
The following is proved in .
###### Lemma 3.2.
Assume that $`T_1,\mathrm{},T_n`$ is an exceptional collection in a triangulated category $`𝒜`$. Let $``$ be the triangulated subcategory of $`𝒜`$ generated by $`T_1,\mathrm{},T_n`$ and put $`T=_iT_i`$. Then $`𝒜`$ has a semi-orthogonal decomposition given by $`(,^{})=(,T^{})`$.
###### Proof.
For the convenience of the reader we repeat the proof. We have to show that $`𝒜`$ is generated by $``$ and $`T^{}`$. Let $`_1`$ be the full subcategory of $`𝒜`$ consisting of objects isomorphic to finite direct sums of the form $`_jT_1[j]^{a_j}`$. Then $`_1`$ is a strict triangulated subcategory of $`𝒜`$. This can be deduced from the fact that the formation of triangles in $`𝒜`$ is compatible with direct sums \[37, Cor. II.1.2.5\]. Sending $`A𝒜`$ to $`_i\mathrm{Ext}^i(T_1,A)_DT_1[i]`$ defines a right adjoint to the inclusion $`_1𝒜`$. This yields a semi-orthogonal decomposition of $`𝒜`$ given by $`(_1,_1^{})=(_1,T_1^{})`$. In particular $`𝒜`$ is generated by $`_1`$ and $`T_1^{}`$. Now we repeat this construction with $`T_2T_1^{}`$. So if $`_2`$ is the full subcategory of $`𝒜`$ consisting of objects isomorphic to finite direct sums of the form $`_jT_2[j]^{a_j}`$ then we have that $`T_1^{}`$ is generated by $`_2`$ and $`(T_1T_2)^{}`$. Continuing this procedure we find that $`𝒜`$ is generated by $`_1,\mathrm{}_n`$ and $`T^{}`$. This finishes the proof. ∎
We point out the following consequence of Lemma 3.2.
###### Corollary 3.3.
Let $`T_1,\mathrm{},T_n`$ be a strongly exceptional collection in a triangulated category $`𝒜`$ satisfying the above assumptions. Then $`T`$ is a tilting object if and only if $`T^{}=0`$.
We shall need that semi-orthogonal decompositions behave nicely with respect to Grothendieck groups.
###### Lemma 3.4.
Assume that $`(,𝒞)`$ is a semi-orthogonal decomposition for a triangulated category $`𝒜`$. Then $`K_0(𝒜)K_0()K_0(𝒞)`$.
###### Proof.
The inclusions $`,𝒞𝒜`$ define a map $`K_0()K_0(𝒞)K_0(𝒜)`$. An inverse to this map is given by sending $`[A]`$ to $`[i^!A][j^{}A]`$ (see Lemma 3.1 for notations). ∎
###### Lemma 3.5.
Assume that $`𝒜`$ is a triangulated category, and let $`T_1,\mathrm{},T_n𝒜`$ be an exceptional collection. Put $`T=_iT_i`$. Then $`K_0(𝒜)^nK_0(T^{})`$.
###### Proof.
Using the same method as in Lemma 3.2 we find inductively using Lemma 3.4 that $`K_0(𝒜)=_iK_0(_i)K_0(T^{})`$. Now it is easy to see that sending $`_jT_i[j]^{a_j}`$ to $`_j(1)^ja_j`$ defines an isomorphism $`K_0(_i)`$. This proves what we want. ∎
If $``$ is an abelian category and $`B`$ then we will say that $`B`$ has projective dimension $`r`$ if $`\mathrm{Ext}_{}^s(B,)=0`$ for $`s>r`$.
The above results have a counterpart for abelian categories provided we work with exceptional objects of projective dimension $`1`$. This follows from the following lemma.
###### Lemma 3.6.
Assume that $``$ is an abelian category and let $`T`$ be an object in $``$ of projective dimension $`1`$. If an object in $`D^b()`$ is (right) perpendicular to $`T`$ then so is its homology. In particular $`D_T_{}^{}^b()=T_{D^b()}^{}`$.
###### Proof.
Let $`BT_{D^b()}^{}`$ and assume that $`n`$ is maximal such that $`H^n(B)0`$. Then there is a triangle
(3.2)
$$\tau _{<n}BBH^n(B)[n]$$
Applying $`\mathrm{Hom}(T,)`$ yields injections $`\mathrm{Ext}^i(T,H^n(B))\mathrm{Hom}(T,(\tau _{<n}B)[n+1+i])`$. Since for $`i0`$ the non-trivial homology of $`(\tau _{<n}B)[n+1+i]`$ occurs in degrees $`2`$ and since the projective dimension of $`T`$ is less than or equal to 1 it follows that $`\mathrm{Hom}(T,(\tau _{<n}B)[n+1+i])=0`$ for $`i0`$. In particular $`\mathrm{Ext}^i(T,H^n(B))=0`$ for $`i0`$. Since trivially $`\mathrm{Ext}^i(T,H^n(B))=0`$ for $`i<0`$ it follows that $`H^n(B)T^{}`$. But then it follows from (3.2) that $`\tau _{<n}BT^{}`$. Repeating this procedure with $`B`$ replaced by $`\tau _{<n}B`$ eventually yields that the homology of $`B`$ is in $`T^{}`$. ∎
As a corollary one obtains a proof of the following standard result.
###### Corollary 3.7.
Let $`T`$ and $``$ be as in the previous lemma. Then $`T_{}^{}`$ is an abelian category.
###### Proof.
If $`f:AB`$ is a map in $`T_{}^{}`$ then one has to show that $`\mathrm{ker}f`$, $`\mathrm{coker}fT_{}^{}`$. Since the complex represented by $`f`$ clearly lies in $`T_{D^b()}^{}`$, this follows from the previous lemma. ∎
As a consequence of Lemma 3.6 we obtain the following result on Grothendieck groups.
###### Corollary 3.8.
Assume that $``$ is an abelian category. Assume that $`T_1,\mathrm{},T_n`$ is an exceptional collection consisting of objects of projective dimension $`1`$. Let $`T=_iT_i`$. Then $`K_0()^nK_0(T^{})`$.
###### Proof.
By Lemmas 3.5 and 3.6 one has
$`K_0()`$ $`K_0(D^b())`$
$`^nK_0(T_{D^b()}^{})`$
$`^nK_0(D_T_{}^{}^b())`$
$`^nK_0(T_{}^{})`$
We now get the main result of this section.
###### Corollary 3.9.
Assume that $``$ is an abelian category of finite Krull dimension. Let $`T=_{i=1}^nT_i`$ be as in Corollary 3.8, but assume is addition that $`\mathrm{Ext}^1(T_i,T_j)=0`$ for all $`i,j`$. If $`n=\mathrm{rk}K_0()`$ then $`T`$ is a tilting object in $``$.
###### Proof.
By Lemma 3.2 we have to show $`T^{}=0`$. By Lemma 3.6 it follows that $`T^{}`$ is equal to $`D_T_{}^{}^b()`$. So it is sufficient to show that $`\stackrel{\text{def}}{=}T_{}^{}=0`$.
By Corollary 3.8 one has $`\mathrm{rk}K_0()=0`$. Since $``$ is an abelian subcategory of $``$, it also has finite Krull dimension. In particular if $`0`$ there is a quotient category $`𝒞`$ of $``$ which has finite length. Selecting a simple object in $`𝒞`$ yields a rank function on $``$ which is non-trivial. Hence $`\mathrm{rk}K_0()>0`$. This yields a contradiction. ∎
It would be interesting to know if for a nonzero hereditary abelian $`k`$-category $``$ with finite dimensional homomorphism and extension spaces we must have $`K_0()0`$.
## 4. Strongly exceptional collections for hereditary orders over discrete valuation rings
In order to construct a tilting object in the category $`\text{coh}𝒪`$ of coherent modules for a hereditary order $`𝒪`$ over $`^1`$ we need to produce some exceptional collections of modules for hereditary orders over discrete valuation rings. We start by recalling some properties for such orders. For simplicity we restrict ourselves to hereditary orders contained in a matrix ring since that is the only case we will need.
Let $`R`$ be a discrete valuation ring and let $`m`$ be its maximal ideal. Furthermore let $`K`$ be the quotient field of $`R`$ and put $`A=M_n(K)`$ . Let $`\mathrm{\Delta }`$ be a hereditary order in $`A`$ in the sense of .
Thus there exist strictly positive integers $`n_1,\mathrm{},n_t`$ such that $`n=n_1+\mathrm{}+n_t`$ and such that $`\mathrm{\Delta }`$ is isomorphic to
(4.1)
$$\left(\begin{array}{cccc}R_{n_1\times n_1}& m_{n_1\times n_2}& \mathrm{}& \\ R_{n_2\times n_1}& R_{n_2\times n_2}& \mathrm{}& \\ \mathrm{}& \mathrm{}& \mathrm{}& \\ & & & R_{n_t\times n_t}\end{array}\right)$$
Here $`()_{a\times b}`$ is a shorthand for $`M_{a\times b}()`$. Strictly speaking this is proved in only in the case that $`R`$ is complete, but as is remarked in \[29, bottom of p. 364\] the result remains valid in the case we consider.
For $`i=1,\mathrm{},t`$ put $`p_i=_{ji}n_j`$, $`q_i=_{j>i}n_j`$. Also let $`P_i`$ be the $`i`$’th indecomposable projective for $`\mathrm{\Delta }`$.
Thus by definition
$$P_i=\left(\begin{array}{c}m_{p_{i1}}\\ R_{q_{i1}}\end{array}\right)$$
where $`()_a`$ now stands for $`M_{a\times 1}()`$.
Clearly $`P_1P_2\mathrm{}P_t`$. Set $`S_i=P_1/P_i`$ for $`i=2,\mathrm{},t`$ and for the same indexes put define
$$\mathrm{\Delta }_i=\left(\begin{array}{cc}R_{p_{i1}\times p_{i1}}& m_{p_{i1}\times q_{i1}}\\ R_{q_{i1}\times p_{i1}}& R_{q_{i1}\times q_{i1}}\end{array}\right)$$
The $`\mathrm{\Delta }_i`$ are submaximal orders containing $`\mathrm{\Delta }`$ and $`S_i`$ is a simple $`\mathrm{\Delta }_i`$-module.
Using this observation one proves:
###### Lemma 4.1.
One has
$$\mathrm{\Delta }_i_\mathrm{\Delta }S_j=\{\begin{array}{cc}0\hfill & \text{if }i>j\hfill \\ S_i\hfill & \text{if }ij\hfill \end{array}$$
###### Proof.
Since one has $`\mathrm{\Delta }_i_\mathrm{\Delta }P_1=\mathrm{\Delta }_i_\mathrm{\Delta }\mathrm{\Delta }_i_{\mathrm{\Delta }_i}P_1=\mathrm{\Delta }_i_{\mathrm{\Delta }_i}P_1=P_1`$ it follows that
$$\mathrm{\Delta }_i_\mathrm{\Delta }S_j=P_1/\mathrm{\Delta }_iP_j$$
It is now easy to see that
$$\mathrm{\Delta }_iP_j=\{\begin{array}{cc}P_1\hfill & \text{if }i>j\hfill \\ P_i\hfill & \text{if }ij\hfill \end{array}$$
which yields the result. ∎
###### Corollary 4.2.
$`\mathrm{Ext}_\mathrm{\Delta }^1(S_j,S_i)=0`$, and
$$\mathrm{Hom}_\mathrm{\Delta }(S_j,S_i)=\{\begin{array}{cc}0\hfill & \text{if }i>j\hfill \\ R/m\hfill & \text{if }ij\hfill \end{array}$$
###### Proof.
One has
$$\mathrm{Ext}_\mathrm{\Delta }^p(S_j,S_i)=\mathrm{Ext}_{\mathrm{\Delta }_i}^p(\mathrm{\Delta }_i_\mathrm{\Delta }S_j,S_i)=\{\begin{array}{cc}0\hfill & \text{if }i>j\hfill \\ \mathrm{Ext}_{\mathrm{\Delta }_i}^p(S_i,S_i)\hfill & \text{if }ij\hfill \end{array}$$
It is easy to see that $`\mathrm{Hom}(S_i,S_i)=R/m`$ and one verifies directly that $`\mathrm{Ext}_{\mathrm{\Delta }_i}^1(S_i,S_i)=0`$. ∎
###### Corollary 4.3.
$`(S_i)_{i=2,\mathrm{},t}`$ is a strongly exceptional collection for $`\mathrm{\Delta }`$.
One also has
###### Lemma 4.4.
$`\mathrm{Ext}_\mathrm{\Delta }^p(S_j,P_1)=0`$ for all $`p`$.
###### Proof.
Put $`\mathrm{\Gamma }=M_n(R)`$. Then one has $`\mathrm{Ext}_\mathrm{\Delta }^p(S_j,P_1)=\mathrm{Ext}_\mathrm{\Gamma }^p(\mathrm{\Gamma }_\mathrm{\Delta }S_j,P_1)=0`$
## 5. Existence of tilting objects for hereditary orders on $`^1`$
In this section $`X`$ will be $`^1`$ for an algebraically field $`k`$. Let $`K`$ be the function field of $`X`$ and let $`A=M_n(K)`$. Let $`𝒪`$ be a sheaf of hereditary orders in $`A`$. Thus locally $`𝒪`$ is a hereditary order over a Dedekind ring (in the sense of ).
To compute global $`\mathrm{Ext}`$’s in $`\text{coh}(𝒪)`$ below we will use the fact that \[14, Prop. II.5.3\]
(5.1)
$$\mathrm{RHom}_𝒪(A,B)=R\mathrm{\Gamma }(X,\mathrm{R}\mathrm{𝑜𝑚}_𝒪(A,B))$$
together with the fact that $`\mathrm{R}\mathrm{𝑜𝑚}_𝒪(A,B)`$ can be computed locally. That is, if $`xX`$ then $`\mathrm{R}\mathrm{𝑜𝑚}_𝒪(A,B)_x=\mathrm{RHom}(A_x,B_x)`$.
Let $`\overline{𝒪}`$ be a maximal order in $`A`$ lying over $`𝒪`$. By Tsen’s theorem there exists a vector bundle $``$ of rank $`n`$ on $`X`$ such that $`\overline{𝒪}=\mathrm{𝑛𝑑}_{𝒪_X}()`$.
Fix $`i`$ and put $`R_i=𝒪_{X,x_i}`$, $`\mathrm{\Delta }_i=𝒪_{x_i}`$. Then $`\mathrm{\Delta }_i`$ is isomorphic to an order of the form (4.1) with $`t=e_i`$. We choose this isomorphism in such a way that it extends to an isomorphism between $`\overline{𝒪}_{x_i}`$ and $`M_n(R_i)`$. Let us write $`S_{ij}`$ ($`i=1,\mathrm{}t,j=2,\mathrm{},e_i`$) for the finite length $`\mathrm{\Delta }_i`$-modules which were denoted by $`S_j`$ in Section 4. We consider the $`S_{ij}`$ as $`𝒪`$-modules.
Since $`\text{coh}𝒪_X`$ has a tilting object, for example given by $`𝒪_X𝒪_X(1)`$, the same is true for $`\overline{𝒪}`$ by Morita theory. We will take $`\overline{T}=(1)`$ as a tilting object in $`\text{coh}(\overline{𝒪})`$. By the choice of the local isomorphisms, $`_{x_i}`$ will correspond to the projective $`\mathrm{\Delta }`$-module denoted by $`P_1`$ in the previous section.
###### Proposition 5.1.
$`T=_{ij}S_{ij}(1)`$ is a tilting object in $`\text{coh}(𝒪)`$.
###### Proof.
Since $`K_0(𝒪_^1)^2`$ it follows from Proposition 2.1 that the number of summands of $`T`$ is equal to the rank of $`K_0(𝒪)`$.
We want to show that the summands of $`T`$ are a strongly exceptional collection. To do this we have to compute the $`\mathrm{Ext}^{}(,)`$ between the summands of $`T`$. We first compute the $`\mathrm{R}\mathrm{𝑜𝑚}`$’s using Corollary 4.2 and lemma 4.4. The result is as follows.
(5.2)
$$\begin{array}{cccc}& & & \\ & S_{kl}& & (1)\\ & & & \\ S_{ij}& & 0& 0\\ & & & \\ & 𝒪_{x_k}& 𝒪_X& 𝒪_X(1)\\ & & & \\ (1)& 𝒪_{x_k}& 𝒪_X(1)& 𝒪_X\end{array}$$
For the square marked ‘$``$’ we have (using Corollary 4.2)
$$\mathrm{R}\mathrm{𝑜𝑚}(S_{ij},S_{kl})=\{\begin{array}{cc}0\hfill & \text{if }ik\hfill \\ 0\hfill & \text{if }i=k\text{ and }l>j\hfill \\ 𝒪_{x_i}\hfill & \text{otherwise}\hfill \end{array}$$
It now follows immediately from (5.1) that $`\mathrm{Ext}^{}`$ is zero between the summands of $`T`$. For the $`\mathrm{Hom}`$’s we find:
(5.3)
$$\begin{array}{cccc}& & & \\ & S_{kl}& & (1)\\ & & & \\ S_{ij}& & 0& 0\\ & & & \\ & k& k& 0\\ & & & \\ (1)& k& k^2& k\end{array}$$
with the ‘$``$’ entry given by
$$\mathrm{Hom}(S_{ij},S_{kl})=\{\begin{array}{cc}0\hfill & \text{if }ik\hfill \\ 0\hfill & \text{if }i=k\text{ and }l>j\hfill \\ k\hfill & \text{otherwise}\hfill \end{array}$$
So it follows in particular that $`T`$ is defined by a strongly exceptional collection. We are now done by Corolllary 3.9. ∎
## 6. Finitely generated Grothendieck groups and existence of tilting objects.
In this section we combine our previous results to get our desired connection between existence of tilting objects and the Grothendieck group being finitely generated.
The main result of this paper is the following.
###### Theorem 6.1.
Let $`𝒞`$ be a connected noetherian hereditary abelian $`\mathrm{Ext}`$-finite $`k`$-category with Serre functor, where $`k`$ is an algebraically closed field. Then the following are equivalent.
* $`K_o(𝒞)`$ is finitely generated.
* + $`𝒞`$ has a tilting object or
+ $`𝒞`$ is the category of finite dimensional representations of the quiver $`\stackrel{~}{A}_n`$ with cyclic orientation for some $`n<\mathrm{}`$.
###### Proof.
$`(b)(a)`$. We have already pointed out that $`(b)(i)`$ implies $`(a)`$ \[19, I.4.6\], and $`(b)(ii)`$ implies $`(a)`$ is obvious.
$`(a)(b)`$. Assume that $`K_o(𝒞)`$ is finitely generated. If $`𝒞=\mathrm{mod}\mathrm{\Lambda }`$ for a finite dimensional hereditary $`k`$-algebra $`\mathrm{\Lambda }`$, then $`𝒞`$ has a tilting object. If $`𝒞=\text{coh}𝒪`$ where $`𝒪`$ is a sheaf of hereditary orders over $`^1`$, it follows from Proposition 5.1 that $`𝒞`$ has a tilting object. Hence we are done using Theorem 2.3. ∎
Actually, the following related result is also of interest.
###### Theorem 6.2.
Let $`𝒞`$ be a connected noetherian $`\mathrm{Ext}`$-finite hereditary category which has no projectives or injectives and which has an object which is not of finite length. Then the following are equivalent.
1. $`𝒞`$ has a tilting object.
2. $`𝒞`$ is derived equivalent to a finite dimensional algebra.
3. $`𝒞`$ has almost split sequences and $`K_0(𝒞)`$ is finitely generated.
4. $`𝒞`$ is of the form $`\text{coh}(𝒪)`$ where $`𝒪`$ is a sheaf of hereditary $`𝒪_^1`$-orders.
5. $`𝒞`$ is of the form $`\text{coh}𝒳`$ for a weighted projective line $`𝒳`$.
###### Proof.
$`1.2`$. When $`𝒞`$ has a tilting object $`T`$, it follows from \[19, I, Th. 4.6\] that $`𝒞`$ is derived equivalent to the finite dimensional algebra End$`(T)^{\mathrm{opp}}`$.
$`2.3`$. Since the hereditary category $`𝒞`$ is derived equivalent to a finite dimensional algebra $`\mathrm{\Lambda }`$, it follows that $`\mathrm{\Lambda }`$ must have finite global dimension. Hence $`\mathrm{mod}\mathrm{\Lambda }`$, and consequently $`𝒞`$, has a Serre functor . Then it follows that $`𝒞`$ has almost split sequences .
Since $`K_0(\mathrm{mod}\mathrm{\Lambda })`$ is finitely generated, it follows that $`K_o(𝒞)`$ is finitely generated because this property is an invariant of derived equivalence.
$`3.4`$. Since $`𝒞`$ has almost split sequences and no nonzero projectives or injectives, it follows that $`𝒞`$ has a Serre functor . Since $`𝒞`$ has some object of infinite length, it follows from Theorem 2.3 that $`𝒞`$ is of the form $`\text{coh}(𝒪)`$ where $`𝒪`$ is a sheaf of hereditary $`𝒪_^1`$-orders.
$`4.1`$. This follows from Proposition 5.1.
$`5.1`$. That $`\text{coh}𝒳`$ for a weighted projective line $`𝒳`$ has a tilting object follows from , and $`1.5`$. follows from . ∎
In an appendix we give for completeness an independent proof of $`4.5`$, hence providing a proof of Theorem 6.2 without using .
## 7. Examples and comments
In this section we give some examples and comments without proofs, related to the material in this paper.
We start by pointing out how to obtain some concrete examples of categories $`\mathrm{qgr}S`$. Translation quivers $`\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is an extended Dynkin diagram occur as AR-quivers for the graded reflexive modules over invariant rings $`S=k[X,Y]^G`$ where $`G`$ is a finite group and $`k`$ is an algebraically closed field of characteristic zero (see for the definition of AR-quiver). The corresponding mesh category for $`\mathrm{\Delta }`$ is then a full subcategory of $`\mathrm{qgr}S`$ whose objects have no nonzero summands of finite length. We obtain a (finite) basis for $`K_0(\mathrm{qgr}S)`$ by considering vertices given by a ”section”. Actually such a set of vertices corresponds to a tilting object. For two-dimensional $`Z`$-graded rings $`S^{}`$ of finite (graded) representation type we have that $`\mathrm{qgr}S^{}`$ is equivalent to some $`\mathrm{qgr}k[X,Y]^G`$. The rings $`k[X,Y]`$ are Gorenstein. It is however not true for a two-dimensional isolated singularity in general that there is a commutative Gorenstein ring $`S`$ with $`\mathrm{qgr}S^{}`$ equivalent to $`\mathrm{qgr}S`$(see ).
While there is a lot of analogy with the work in , we note that there are also some differences. In the complete case, for the rings $`\mathrm{\Lambda }`$ of finite representation type, considered as orders, the rank $`n`$ of $`K_0(q(\mathrm{\Lambda }))`$ gives information on how far $`\mathrm{\Lambda }`$ is from being a maximal order (of finite representation type). In this case there was a chain $`\mathrm{\Lambda }=\mathrm{\Lambda }_1\mathrm{}\mathrm{\Lambda }_n`$ of orders with $`\mathrm{\Lambda }_n`$ maximal and such that there is no refinement of the chain. Then $`K_0(q(\mathrm{\Lambda }))`$ has rank 1, when $`\mathrm{\Lambda }`$ is a maximal order, and all commutative $`\mathrm{\Lambda }`$ of finite type are maximal orders. In the graded case however, given $`S`$ with rank $`\mathrm{qgr}(S)=n`$, there is no corresponding chain of graded orders ending up with $`k[X,Y]`$.
We point out that if $`𝒞`$ is a hereditary abelian $`k`$-category with all objects of finite length, then $`𝒞`$ does not necessarily have almost split sequences (or Serre functor). For example, this is the case if $`𝒞`$ is the category of holonomic modules over the first Weyl algebra (see ). And it follows from that it holds for the category of finite dimensional representations over $`k`$ of a finite connected quiver having oriented cycles, but which is not equal to a single oriented cycle.
## Appendix A Hereditary orders and weighted projective lines
In this appendix we will show directly that $`\text{coh}(𝒪)`$ for $`𝒪`$ a classical hereditary order over $`^1`$ is equivalent to $`\text{coh}𝒳`$ for a weighted projective line $`𝒳`$ and furthermore we will show that every weighted projective line appears in this way. As was said before this can be deduced from Theorem 6.2 together with .
We follows the methods of , except that we consider gradings by rank one abelian groups which can have torsion.
To formalize this let $`𝒟`$ be an abelian category and let $`O𝒟`$ be an object. In addition let $`(t_h)_{hH}`$ be a family of autoequivalences of $`𝒟`$ indexed by a group $`H`$ and for any pair $`h_1,h_2H`$ assume there are given natural isomorphisms $`\eta _{h_1,h_2}:t_{h_1}t_{h_2}t_{h_1h_2}`$ satisfying the cocycle condition
(A.1)
$$\eta _{h_1h_2,h_3}(\eta _{h_1,h_2}t_{h_3})=(t_{h_1}\eta _{h_2,h_3})\eta _{h_1,h_2h_3}$$
The data $`(t_h)`$, $`(\eta _{h_1,h_2})`$ can be used to put a $`H`$-graded ring structure on
$$\mathrm{\Gamma }^{}(O)=\underset{hH}{}\mathrm{Hom}(t_h^1O,O)$$
as well as a $`H`$-graded $`\mathrm{\Gamma }^{}(O)`$-module structure on
$$\mathrm{\Gamma }^{}(M)=\underset{hH}{}\mathrm{Hom}(t_h^1O,M)$$
¿¿From now on we will assume that $`H`$ is a finitely generated abelian group of rank one. We fix an element $`z`$ in $`H`$. Associated to $`H`$ there is a surjective map $`\varphi :H`$, unique up to sign. We fix the sign by imposing $`\varphi (z)>0`$. If $`U`$ is a $`H`$-graded abelian group then we say that $`U`$ has right bounded grading if $`U_h=0`$ for $`\varphi (h)0`$. If $`R`$ is a noetherian $`H`$-graded ring then we define $`\mathrm{qgr}(R)=\mathrm{gr}(R)/\mathrm{tors}(R)`$ where $`\mathrm{tors}(R)`$ consists of the right bounded modules.
The following result is an easy extension of \[1, Thm 4.5\].
###### Proposition A.1.
Let $`𝒟`$ be a noetherian $`\mathrm{Ext}`$-finite abelian category and let $`O𝒟`$. Let $`(t_h)_{hH}`$ be a system of autoequivalences as above and assume that $`(𝒟,O,t_z)`$ is an ample triple in the sense of . Then $`R=\mathrm{\Gamma }^{}(O)`$ is noetherian and the functor $`\mathrm{\Gamma }^{}`$ defines an equivalence between $`𝒟`$ and $`\mathrm{qgr}(R)`$.
Now let $`X=^1`$, $`K=k(X)`$ and let $`𝒪`$ be a sheaf of hereditary $`𝒪_X`$-orders in $`A=M_n(K)`$.
Let $`x_1,\mathrm{},x_tX`$ be the set of ramification points of $`𝒪`$. Since the analysis of the cases $`t=0`$, $`t=1`$ and $`t2`$ is somewhat different we consider the case $`t2`$ first. Afterwards we discuss the other cases. Fix an arbitrary point $`x`$ in $`X`$ distinct from $`x_1,\mathrm{},x_t`$ and let $`(f_i)_{i=1,\mathrm{},t}`$ be rational functions with divisor $`(x_i)+(x)`$.
Let $`I_i`$ be fractional $`𝒪`$-ideals in $`A`$ defined by the condition
$$(I_i)_y=\{\begin{array}{cc}(\mathrm{rad}𝒪_{x_i})^1\hfill & \text{if }y=x_i\hfill \\ 𝒪_y\hfill & \text{otherwise}\hfill \end{array}$$
¿¿From this definition we obtain canonical isomorphisms (as fractional ideals)
(A.2)
$$I_i^{e_i}I_j^{e_j}:xxf_j/f_i$$
We let $`H`$ be the abelian group of rank one generated by the elements $`h_1,\mathrm{},h_t`$, subject to the relations $`e_ih_i=e_jh_j`$ and we put $`z=h_1+\mathrm{}+h_t`$.
Every $`hH`$ has a unique representation of the form $`a_1h_1+\mathrm{}+a_th_t`$ with $`0a_i<e_i`$ for $`i>1`$. We define $`I_h=I_1^{a_1}\mathrm{}I_t^{a_t}`$. From (A.2) we obtain canonical isomorphisms $`\zeta _{h_1,h_2}:I_{h_1}I_{h_2}I_{h_1+h_2}`$.
Associated to the fractional ideals $`I_h`$ there are autoequivalences $`t_h`$ on $`\text{coh}(𝒪)`$ given by $`I_h`$. The $`\zeta _{h_1,h_2}`$ define natural isomorphisms $`\eta _{h_1,h_2}:t_{h_1}t_{h_2}t_{h_1+h_2}`$ satisfying the cocycle condition (A.1).
Now let $`\overline{𝒪}`$ be a maximal order overlying $`𝒪`$. As usual $`\overline{𝒪}=\mathrm{𝑛𝑑}()`$ for some vector bundle $``$ on $`X`$. It follows from \[31, Ch IV\] that the triple $`(\text{coh}(𝒪),,t_z)`$ is ample. Hence if we take $`O=`$ in the above notations and we put $`R=\mathrm{\Gamma }^{}()`$ then $`R`$ is a noetherian $`H`$-graded ring and $`\mathrm{\Gamma }^{}`$ defines an equivalence between $`\text{coh}(𝒪)`$ and $`\mathrm{qgr}(R)`$.
Our next aim will be to show that $`R`$ is in fact a weighted projective line. Unfortunately the autoequivalences $`\eta _{h_1,h_2}`$ clutter up our computations rather badly. Therefore we will first give a more elegant description of $`R`$.
Let $`D`$ be the graded ring defined by
$$D=A[u_1,u_1^1,\mathrm{},u_t,u_t^1]/(f_iu_i^{e_i}=f_ju_j^{e_j})]$$
$`D`$ is clearly $`H`$-graded by putting $`\mathrm{deg}u_i=h_i`$.
Let $`𝒜`$ be the graded order in $`D`$ defined by
$$𝒜=\underset{p_1,\mathrm{},p_t}{}(I_1u_1)^{p_1}\mathrm{}(I_tu_t)^{p_t}$$
Now it is not hard to see that
$$R=\mathrm{Hom}(,𝒜_𝒪)$$
We will determine the structure of $`R`$ explicitly. A local computation shows that $`R`$ is equal to
$$\underset{p_1,\mathrm{},p_t}{}\mathrm{\Gamma }(X,𝒪_X([p_1/e_1]x_1+\mathrm{}+[p_t/e_t]x_t))u_1^{p_1}\mathrm{}u_t^{p_t}$$
where $`[\alpha ]`$ denote the biggest integer not bigger than $`\alpha `$.
We first claim that $`R`$ is generated by $`u_1,\mathrm{},u_t`$. By the relations in $`D`$ it follows that
$$\begin{array}{c}𝒪_X\left(\left[p_1/e_1\right]x_1+\mathrm{}+\left[\left(p_i+e_i\right)/e_i\right]x_i+\mathrm{}+\left[p_j/e_j\right]x_j+\mathrm{}+\left[p_t/e_t\right]x_t\right)u_1^{p_1}\mathrm{}u_i^{p_i+e_i}\mathrm{}u_j^{p_j}\mathrm{}u_t^{p_t}\hfill \\ \hfill =𝒪_X\left(\left[p_1/e_1\right]x_1+\mathrm{}+\left[p_i/e_i\right]x_i+\mathrm{}+\left[\left(p_j+e_j\right)/e_j\right]x_j+\mathrm{}+\left[p_t/e_t\right]x_t\right)u_1^{p_1}\mathrm{}u_i^{p_i}\mathrm{}u_j^{p_j+e_j}\mathrm{}u_t^{p_t}\end{array}$$
as subsheaves of $`D`$. Hence to show that every section of $`𝒪_X([p_1/e_1]x_1+\mathrm{}+[p_t/e_t]x_t])u_1^{p_1}\mathrm{}u_t^{p_t}`$ is a linear combination of products of the $`u_i`$’s, it suffices to do so in the case that $`p_i<e_i`$ for $`i2`$. So below we make this assumption.
Write $`p_1=q_1+ae_1`$ where $`q_1<e_1`$. We then have $`𝒪_X([p_1/e_1]x_1+\mathrm{}+[p_t/e_t]x_t])=𝒪_X(ax_1)`$. Let $`b_1,\mathrm{},b_t`$ be such that $`b_1+\mathrm{}+b_t=a`$. Using the relations in $`D`$ we find that
$$u_1^{q_1+b_1e_1}u_2^{p_2+b_2e_2}\mathrm{}u_t^{p_t+b_te_t}=f_1^{b_2+\mathrm{}+b_t}/(f_2^{b_2}\mathrm{}f_t^{b_t})u_1^{p_1}u_2^{p_2}\mathrm{}u_t^{p_t}$$
The divisor of $`f_1^{b_2+\mathrm{}+b_t}/(f_2^{b_2}\mathrm{}f_t^{b_t})`$ is equal to $`a(x_1)+_{i=1}^tb_i(x_i)`$. Hence these rational functions clearly generate the global sections of $`𝒪_X(ax_1)`$, which is what we had to show.
We now claim that up to changing $`f_1,(f_i)_{i3}`$ by a scalar we have the following relations in $`R`$ :
(A.3)
$$u_i^{e_i}u_2^{e_2}+\lambda _iu_1^{e_1}=0(i3)$$
where the $`(\lambda _i)_{i3}`$ are suitable scalars with $`\lambda _3=1`$.
Rewriting $`u_2^{e_2}`$ and $`u_i^{e_i}`$ in terms of $`u_1`$ it follows that the relation (A.3) is equivalent to the existence of a linear dependence
(A.4)
$$f_1/f_if_1/f_2+\lambda _i=0$$
Now the divisors of $`f_1/f_i`$, $`f_1/f_2`$ and $`1`$ are respectively given by $`(x_1)+(x_i)`$, $`(x_1)+(x_2)`$ and $`0`$. In particular these three rational functions are all sections of $`𝒪_X(x_1)`$. Since $`𝒪_X(x_1)`$ has degree one, it follows that there has to be at least a linear dependence
(A.5)
$$\alpha f_1/f_i+\beta f_1/f_2+\gamma =0$$
Furthermore, inspecting divisors, it is easily seen that $`\alpha ,\beta ,\gamma `$ must all be non-zero. Dividing (A.5) by $`\beta `$ and changing $`f_i`$ by a scalar yields (A.4). To make $`\lambda _3`$ equal to $`1`$ we finish by changing $`f_1`$ by a suitable scalar.
At this point we know that $`R`$ is a quotient of the “weighted projective line”
(A.6)
$$k[u_1,\mathrm{},u_t]/(u_i^{e_i}u_2^{e_2}+\lambda _iu_1^{e_1})$$
However a straightforward computation reveals that $`R`$ and the ring defined by (A.6) have the same Hilbert series. Hence they are isomorphic. This concludes our analysis of the case $`t2`$.
We will now discuss the other cases. First let $`t=1`$. We define $`R`$ as above. Now $`dimR_i=1`$ for $`i<e_1`$ and $`dimR_e=2`$. Let $`vR_2ku_1^e`$. We leave it as an exercise to the reader to check that $`Rk[u_1,v]`$. Hence $`\text{coh}(𝒪)`$ is again described by a weighted projective line.
The case $`t=0`$ is even more trivial. In that case $`𝒪`$ is Morita equivalent to $`𝒪_^1`$. So $`\text{coh}(𝒪)`$ is in fact described by the ordinary projective line!
To finish we show that one can get all weighted projective lines from hereditary orders. It suffices to do this in the case $`t>2`$. It is convenient to choose an affine coordinate system on $`^1`$ in such a way that $`x_1=\mathrm{}`$, $`x_2=0`$, $`x_3=1`$. Then up to a scalar we have
$$f_1(z)=zx,f_i(z)=\frac{zx}{zx_i}\text{ for }i>1$$
Computing the $`\lambda _i`$ explicitly with the above procedure we find $`\lambda _i=x_i`$. This shows what we want.
## Appendix B Examples of hereditary abelian categories
In this appendix we give some sources of examples of hereditary abelian categories, which are usually not $`\mathrm{Ext}`$-finite. These are inspired by .
Let $`R`$ be a noetherian ring of Krull dimension $`n0`$, finitely generated as a module over a central subring $`C`$. Denote by $`\mathrm{Mod}R`$ the category of $`R`$-modules and as before by $`\mathrm{mod}R`$ the subcategory of finitely generated $`R`$-modules. For $`i1`$, let $`𝒞_i`$ be the subcategory of $`\mathrm{mod}R`$ whose objects have Krull dimension at most $`i`$, and let $`\stackrel{~}{𝒞}_i`$ be the subcategory of $`\mathrm{Mod}R`$ whose objects are direct limits of objects in $`𝒞_i`$. (We define $`𝒞_1=\stackrel{~}{𝒞}_1=(0)`$.) Let $`\mathrm{qmod}_i(R)=\mathrm{mod}R/𝒞_i`$ and $`\mathrm{QMod}_i(R)=\mathrm{Mod}R/\stackrel{~}{𝒞}_i`$ be the corresponding quotient categories in the sense of . These are abelian categories. Similarly we consider the case when $`S`$ is a $``$-graded noetherian ring finitely generated over $`k`$ and finitely generated as a module over a central subring $`C`$, such that $`S_i=0`$ for $`i`$ small enough. We make the similar definitions starting with the category $`\mathrm{Gr}S`$ of graded $`S`$-modules with degree zero homomorphisms, and the subcategory $`\mathrm{gr}S`$ of finitely generated modules. The corresponding quotient categories will be denoted by $`\mathrm{QGr}_i(S)`$ and $`\mathrm{qgr}_i(S)`$. *Below when we work in the graded case all objects will be implicitly considered to be graded, unless otherwise specified.*
We can now prove the following, which gives some classes of hereditary categories.
###### Proposition B.1.
Let $`R`$ be a noetherian ring of Krull dimension $`n0`$ with the above assumptions and notation. Assume that $`C`$ satisfies $`\mathrm{Kdim}C/P+\mathrm{ht}P=n`$ for every prime ideal $`P`$ in $`C`$. Then the following conditions are equivalent.
* $`\mathrm{QMod}_i(R)`$ is nonzero hereditary.
* $`\mathrm{qmod}_i(R)`$ is nonzero hereditary.
* Either $`i=n2`$ and $`\mathrm{gl}.\mathrm{dim}R_P1`$ for any prime ideal $`P`$ in $`C`$ of height at most 1 or $`i=n1`$ and $`\mathrm{gl}.\mathrm{dim}R_P1`$ for any prime ideal $`P`$ in $`C`$ of height 0.
###### Proof.
That(a) and (b) are equivalent follows from \[31, Proposition A3\]. To prove the other equivalences we first review some generalities. First of all if $`M`$ is an $`R`$-module, then by \[11, p430, Cor. 2\] the Krull dimension of $`M`$ as $`R`$-module is equal to the Krull dimension of $`M`$ as $`C`$-module.
Furthermore we claim that $`M\stackrel{~}{𝒞}_i`$ if and only if $`M_P=0`$ for all $`P\mathrm{Spec}C`$ such that $`\mathrm{Kdim}C/P<i`$ (or equivalently, if and only if $`M_P=0`$ for all $`P`$ such that $`\mathrm{ht}P>ni`$). To see this we may assume that $`M`$ is finitely generated. Then it follows from the theory of associated primes that $`M`$ has a finite filtration (as $`C`$-module) with subquotients of the form $`C/Q`$ with $`Q\mathrm{Spec}C`$. The claim is now an immediate verification.
The subcategory $`\stackrel{~}{𝒞}_i`$ of $`\mathrm{Mod}R`$ is a localizing subcategory. Since $`R`$ is a noetherian ring which is finitely generated as a module over its center, $`\stackrel{~}{𝒞}_i`$ is closed under injective envelopes \[11, p. 431\]. Denoting by $`T:\mathrm{Mod}R\mathrm{Mod}R/\stackrel{~}{C}_i`$ the associated quotient functor we have that $`T`$ preserves injective objects and injective envelopes by \[31, Proposition A4\]. So if $`0MI_0I_1\mathrm{}`$ is a minimal injective resolution in $`\mathrm{Mod}R`$, then $`0T(M)T(I_0)T(I_1)\mathrm{}`$ is a minimal injective resolution in $`\stackrel{~}{_i}`$. A similar reasoning shows that for each prime ideal $`P`$ in $`C`$, we have that $`0M_P(I_0)_P(I_1)_P\mathrm{}`$ is a minimal injective resolution in $`\mathrm{Mod}R_P`$.
(c) $``$ (a). Assume that (c) holds, and consider for $`M`$ in $`\mathrm{Mod}R`$ a minimal injective resolution $`0MI_0I_1\mathrm{}`$ in $`\mathrm{Mod}R`$, and the induced minimal injective resolution $`0M_P(I_0)_P(I_1)_P\mathrm{}`$ for a prime ideal $`P`$ in $`C`$.
If $`i=n2`$ and $`\mathrm{gl}.\mathrm{dim}R_P1`$ for $`\mathrm{ht}P1`$, we get $`(I_j)_P=0`$ for $`j2`$ and $`\mathrm{ht}P1`$, and hence $`I_j`$ is in $`\stackrel{~}{𝒞}_{n2}`$ for $`j2`$, so that $`T(I_j)=0`$. Then we have an injective resolution $`0T(M)T(I_0)T(I_1)0`$ in $`\mathrm{QMod}_{n2}(R)`$, which shows $`\text{id}_{\mathrm{QMod}_{n2}(R)}T(M)1`$. Since $`T`$ is essentially surjective, it follows that $`\mathrm{gl}.\mathrm{dim}\mathrm{QMod}_{n2}(R)1`$.
If $`i=n1`$ and $`\mathrm{gl}.\mathrm{dim}R_P1`$ for $`\mathrm{ht}P=0`$, we get that $`I_j`$ is in $`\stackrel{~}{𝒞}_{n1}`$ for $`j2`$, so that $`0T(M)T(I_0)T(I_1)0`$ is an injective resolution of $`T(M)`$ in $`\mathrm{QMod}_{n1}(R)`$. Hence we have $`\mathrm{gl}.\mathrm{dim}\mathrm{QMod}_{n1}(R)1`$.
(a) $``$ (c). Assume now that $`\mathrm{QMod}_i(R)`$ is hereditary, and let $`0MI_0I_1\mathrm{}`$ be a minimal injective resolution of some $`M`$ in $`\mathrm{Mod}R`$. Since then $`0T(M)T(I_0)T(I_1)\mathrm{}`$ is a minimal injective resolution in $`\mathrm{QMod}_i(R)`$, it follows that $`T(I_j)=0`$ for $`j>1`$. Then $`I_j`$ is in $`\stackrel{~}{𝒞}_i`$ for $`j>1`$, so that we have $`(I_j)_P=0`$ when $`\mathrm{ht}P<ni`$. Thus we we have the exact sequence $`0M_P(I_0)_P(I_1)_P0`$ in this case, and hence $`\mathrm{gl}.\mathrm{dim}R_P1`$. Since for a noetherian ring $`R`$ which is finitely generated as a module over its center we have $`\mathrm{Kdim}R_P\mathrm{gl}.\mathrm{dim}R_P`$ , and furthermore trivially $`\mathrm{Kdim}C_P\mathrm{Kdim}R_P`$, it follows that $`ni11`$, so that $`ni2`$. Hence when $`\mathrm{QMod}_i(R)0`$, we must have $`i=n2`$ or $`i=n1`$. In the first case we have $`\mathrm{gl}.\mathrm{dim}R_P1`$ for $`\mathrm{ht}P1`$, and in the second case $`\mathrm{gl}.\mathrm{dim}R_P1`$ for $`\mathrm{ht}P=0`$. ∎
In the case of graded rings $`S`$ with the assumptions listed in the above, we consider graded prime ideals $`P`$ in the center $`C`$ and graded localizations $`S_P`$ and graded global dimension. Using that also in this case $`\stackrel{~}{𝒞}_i`$ is closed under injective envelopes , the proof of Proposition B.1 is easily adapted to give the following.
###### Proposition B.2.
Let $`S`$ be a $`Z`$-graded ring of Krull dimension $`n0`$ satisfying the standard assumptions, and with the previous notation. Suppose that $`C`$ satisfies $`\mathrm{Kdim}C/P+\mathrm{ht}P=n`$ for every graded prime ideal $`P`$. Then the following are equivalent.
* $`\mathrm{QGr}_i(S)`$ is nonzero hereditary.
* $`\mathrm{qgr}_i(S)`$ is nonzero hereditary.
* Either $`i=n2`$ and $`\mathrm{gl}.\mathrm{dim}S_P1`$ for any graded prime ideal $`P`$ in $`C`$ of height at most 1 or $`i=n1`$ and $`\mathrm{gl}.\mathrm{dim}S_P1`$ for any graded prime ideal $`P`$ in $`C`$ of height 0.
We also state the following special case.
###### Corollary B.3.
Let $`S=C`$ be a $`Z`$-graded commutative domain of Krull dimension 2 satisfying the standard assumptions. Then $`\mathrm{qgr}(S)`$ is hereditary if and only if $`S`$ is an isolated singularity.
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# On the Additivity of the Entanglement of Formation
## Abstract
We study whether the entanglement of formation is additive over tensor products and derive a necessary and sufficient condition for optimality of vector states that enables us to show additivity in two special cases.
PACS: 03.67.-a
Entanglement plays a crucial role in teleportation and quantum cryptography and is currently the focus of investigations in the developing field of quantum information . Quantifying entanglement in a satisfactory way is a major issue in quantum information; there is a list of minimal desiderata and the available proposals have been proved to comply with all of them but for additivity for which only numerical support exists .
In this letter we are concerned with the entanglement of formation, defined by
$$E_f(\rho ;M_1)=\mathrm{min}\left\{\underset{iI}{}p_iS\left(\sigma _i|\text{`}_{M_1}\right)\right\},$$
(1)
where $`\rho `$ is a generic mixed state on the tensor product of two full matrix algebras $`M_1M_2`$, the minimum is computed over all decompositions $`\rho =_{iI}p_i\sigma _i`$ of $`\rho `$ into pure states $`\sigma _i`$ on $`M_1M_2`$ and $`S\left(\sigma _i|\text{`}_{M_1}\right)`$ is the von Neumann entropy of the restriction of $`\sigma _i`$ onto $`M_1`$.
The question we want to address is whether
$$E_f(\rho \rho ;M_1M_1)=2E_f(\rho ;M_1),$$
(2)
where $`M_1M_1`$ is short for $`M_11_2M_11_2`$.
Additivity or its failure will have a quantum information theoretic counterpart; there is indeed a connection between (1) and the maximal accessible information $`I(\rho )`$ of a quantum source described by a mixed state $`\rho `$ on a matrix algebra $`M`$ . If $`\rho =_\mathrm{}Lq_{\mathrm{}}\rho _{\mathrm{}}`$, then $`I(\rho ):=sup_BI_B(\rho )`$, where
$`I_B(\rho )=`$ $``$ $`{\displaystyle \underset{iI}{}}(\mathrm{Tr}(\rho b_i))\mathrm{log}(\mathrm{Tr}(\rho b_i)`$ (3)
$`+`$ $`{\displaystyle \underset{\mathrm{}L}{}}q_{\mathrm{}}{\displaystyle \underset{iI}{}}(\mathrm{Tr}(\rho _{\mathrm{}}b_i))\mathrm{log}(\mathrm{Tr}(\rho _{\mathrm{}}b_i),`$ (4)
the maximum being computed over all choices $`B=\{b_i\}`$ of positive $`b_iM`$ such that $`_{iI}b_i=1_M`$. Optimal choices correspond to optimal detection of the classical information carried by the quantum states $`\rho _{\mathrm{}}`$.
Since $`q_{\mathrm{}}\rho _{\mathrm{}}\rho `$ there exists a unique choice of operators $`0<a_{\mathrm{}}M`$, $`\mathrm{}L`$, with $`_\mathrm{}La_{\mathrm{}}=1`$, such that
$$q_{\mathrm{}}\rho _{\mathrm{}}=\sqrt{\rho }a_{\mathrm{}}\sqrt{\rho },q_{\mathrm{}}=\mathrm{Tr}(\rho a_{\mathrm{}}).$$
(5)
Let $`𝒜`$ be a commutative $`L`$ dimensional algebra with identity $`1_𝒜`$ and orthogonal projectors $`A_{\mathrm{}}`$ with $`_{\mathrm{}}A_{\mathrm{}}=1_𝒜`$. The map $`\gamma _𝒜:𝒜M`$ obtained by linear extension of $`A_{\mathrm{}}\gamma _𝒜(A_{\mathrm{}})=a_{\mathrm{}}`$ is positive and $`\gamma _𝒜(1_𝒜)=1_M`$. Therefore, given any state $`\sigma `$ on $`M`$, the linear functional $`\sigma \gamma _𝒜:𝒜𝐂`$, $`\sigma \gamma _𝒜(A_{\mathrm{}})=\mathrm{Tr}(\sigma a_{\mathrm{}})`$ defines a state on $`𝒜`$. Using (5) and the cyclicity of the trace,
$$\mathrm{Tr}(\rho _{\mathrm{}}b_i)=\frac{\mathrm{Tr}(\rho b_i)}{\mathrm{Tr}(\rho a_{\mathrm{}})}\mathrm{Tr}(\sigma _i^Ba_{\mathrm{}}),\sigma _i^B:=\frac{\sqrt{\rho }b_i\sqrt{\rho }}{\mathrm{Tr}(\rho b_i)}.$$
(6)
Setting $`p_i^B=\mathrm{Tr}(\rho b_i)`$, (4) becomes
$$I_B(\rho )=S\left(\rho \gamma _𝒜\right)\underset{iI}{}p_i^BS\left(\sigma _i^B\gamma _𝒜\right).$$
(7)
Therefore, $`I(\rho )`$ is the maximum of (7) over all possible decompositions of $`\rho `$ into pure states:
$$I(\rho )=S\left(\rho \gamma _𝒜\right)\underset{\rho =_ip_i\sigma _i}{\mathrm{min}}\underset{i}{}p_iS\left(\sigma _i\gamma _𝒜\right).$$
(8)
If $`N`$ is a subalgebra of $`M`$, substituting the restrictions $`\rho |\text{`}_N`$, $`\sigma _i|\text{`}_N`$ for $`\rho \gamma _𝒜`$, respectively $`\sigma _i\gamma _𝒜`$, we obtain the so-called entropy of a subalgebra
$$H_\rho (N):=S(\rho |\text{`}_N)\underset{\rho =_ip_i\sigma _i}{\mathrm{min}}\underset{i}{}p_iS\left(\sigma _i|\text{`}_N\right).$$
(9)
The latter quantity is the building block of an extension of the Kolmogorov-Sinai dynamical entropy (or entropy per unit time) to the quantum realm. According to the above $`E_f(\rho ;M_1)=S\left(\rho |\text{`}_{M_1}\right)H_\rho (M_1)`$.
As the von Neumann entropy is additive over tensor products, if additivity fails for the entanglement of formation, it also fails for the entropy of a subalgebra. Then, from an information-theoretic point of view, we would deem possible to extract more information about the tensor product of two states over two independent subalgebras than that obtainable from the two of them independently .
In the following we try to use some of the properties of $`H_\rho (N)`$ to investigate the general question whether
$$E_f(\rho \sigma ;M_1M_3)=E_f(\rho ;M_1)+E_f(\sigma ;M_3),$$
(10)
where $`\rho `$ and $`\sigma `$ are states on the (finite dimensional) algebras $`M_1M_2`$, respectively $`M_3M_4`$.
If $`E_f(\rho ;M_1)`$ and $`E_f(\sigma ;M_3)`$ are achieved at optimal decompositions $`\rho =_{\mathrm{}}p_{\mathrm{}}\rho _{\mathrm{}}`$ and $`\sigma =_jq_j\sigma _j`$, the factorized decomposition $`\rho \sigma =_{j,\mathrm{}}q_jp_{\mathrm{}}\rho _{\mathrm{}}\sigma _j`$ contribute to $`E_f(\rho \sigma ;M_1M_3)`$ with $`E_f(\rho _1;M_1)+E_f(\sigma ;M_3)`$. However, the latter need not be optimal and the strict inequality $`E_f(\rho \sigma ;M_1M_3)<E_f(\rho _1;M_1)+E_f(\sigma ;M_3)`$ is not excluded. In fact, a decomposition
$$\rho \sigma =\underset{i}{}\alpha _i|\psi _i\psi _i|,\alpha _i>0,\underset{i}{}\alpha _i=1,$$
(11)
might be optimal with the $`\psi _i`$ entangled states over $`M_1M_3`$. Let us consider the Schmidt decomposition
$$|\psi _i=\underset{j}{}\beta _{ij}|\varphi _{ij}^{12}|\varphi _{ij}^{34},\psi _i=1,\beta _{ij}>0,$$
(12)
where, for fixed $`i`$ the $`|\varphi _{ij}^{12}`$’s and $`|\varphi _{ij}^{34}`$’s form orthonormal bases over $`M_1M_2`$, respectively $`M_3M_4`$. If it held that
$`S\left(|\psi _i\psi _i||\text{`}_{M_1M_3}\right)`$ $``$ $`{\displaystyle \underset{j}{}}\beta _{ij}^2(S(|\varphi _{ij}^{12}\varphi _{ij}^{12}||\text{`}_{M_1})`$ (13)
$`+`$ $`S(|\varphi _{ij}^{34}\varphi _{ij}^{34}||\text{`}_{M_3}))`$ (14)
additivity would follow because tensor-product states would then never be worse than correlated ones.
Proving the sufficient condition (14) has so far escaped us; there are however particular cases where one can show additivity by using two results obtained for the entropy of a subalgebra (9). Both results concern general properties of optimal decompositions for $`H_\rho (N)`$ that we adapt to the entanglement of formation .
Proposition 1. If $`\rho `$ is a state on $`M_1M_2`$, $`E_f(\rho ;M_1)`$ is achieved at $`\rho =_{\mathrm{}}p_{\mathrm{}}\rho _{\mathrm{}}`$ and $`U`$ is a unitary operator on $`M_1M_2`$, then $`E_f(U^{}\rho U;M_1)`$ is achieved at the optimal decomposition $`U^{}\rho U=_{\mathrm{}}p_{\mathrm{}}U^{}\rho _{\mathrm{}}U`$.
Proposition 2. Let $`\rho `$ be a state on $`M_1M_2`$ and $`E_f(\rho ;M_1)`$ be achieved at $`\rho =_{\mathrm{}}p_{\mathrm{}}\rho _{\mathrm{}}`$, That is, $`E_f(\rho ;M_1)=_{\mathrm{}}p_{\mathrm{}}S\left(\rho _{\mathrm{}}|\text{`}_{M_1}\right)`$, then
$$E_f(\sigma ;M_1)=\underset{j}{}q_jS\left(\rho _j|\text{`}_{M_1}\right),$$
(15)
where $`\sigma =_jq_j\rho _j`$ is any linear convex combination of optimal states of $`\rho `$.
To the above, we add a new property. With some abuse of notation, we denote by $`E_f(\rho ;N)`$ the minimum in (9), even if there is no tensor product structure in $`N`$.
Proposition 3. Let $`|\psi _i\psi _i|`$, $`i=1,2`$, contribute to $`E_f(\rho ;N)`$ and denote
$`\sigma _i:=|\psi _i\psi _i||\text{`}_N,i=1,2;\sigma _{12}:=|\psi _1\psi _2||\text{`}_N`$
$`\sigma _{ov}(\gamma ):=\gamma \sigma _{21}+\gamma ^{}\sigma _{12}`$
$`\sigma (\gamma ):=|\gamma |^2\sigma _1+\sigma _2\sigma _{ov}(\gamma ),\widehat{\sigma }(\gamma ):={\displaystyle \frac{\sigma (\gamma )}{\mathrm{Tr}(\sigma (\gamma ))}}.`$
Then, for all complex $`\gamma `$,
$$\frac{|\gamma |^2S(\sigma _1)+S(\sigma _2)+\mathrm{Tr}\left(\sigma _{ov}(\gamma )\mathrm{log}\sigma _1\right)}{\mathrm{Tr}(\sigma (\gamma ))}S(\widehat{\sigma }(\gamma )).$$
(16)
Vice versa, if inequality (16) holds for all complex $`\gamma `$, then for all $`\rho _\lambda =\lambda |\psi _1\psi _1|+(1\lambda )|\psi _2\psi _2|`$, $`1\lambda 0`$, one gets $`E_f(\rho _\lambda ;N)=\lambda S(\sigma _1)+(1\lambda )S(\sigma _2)`$.
Proof of Necessity: Let $`\epsilon >0`$ and set
$$\rho _{\epsilon ,\gamma }=(1+\epsilon |\gamma |^2)|\psi _1\psi _1|+\epsilon (1+\epsilon |\gamma |^2)|\psi _2\psi _2|$$
be a not normalized state on $`M`$. As $`\psi _i`$, $`i=1,2`$ are optimal, Proposition 2 yields
$$E_f(\rho _{\epsilon ,\gamma };N)=(1+\epsilon |\gamma |^2)S(\sigma _1)+\epsilon (1+\epsilon |\gamma |^2)S(\sigma _2).$$
Indeed, in taking the minimum in (9) normalization is not necessary. With $`|\varphi _1:=|\psi _1+\epsilon \gamma |\psi _2`$ and $`|\varphi _2:=|\psi _1(\gamma )^1|\psi _2`$, we construct a new decomposition $`\rho _{\epsilon ,\gamma }=|\varphi _1\varphi _1|+\epsilon |\gamma |^2|\varphi _2\varphi _2|`$.
The latter cannot contribute more than $`E_f(\rho _{\epsilon ,\gamma };N)`$; therefore, $`E_f(\rho _{\epsilon ,\gamma };N)f(\epsilon )`$, where
$$f(\epsilon ):=\varphi _1^2S\left(|\varphi _1\varphi _1||\text{`}_N\right)+\epsilon |\gamma |^2\varphi _2^2S\left(|\varphi _2\varphi _2||\text{`}_N\right).$$
Inequality (16) must then hold at first order in $`\epsilon `$.
Proof of Sufficiency: By assumption, inequality (16) holds for all $`\gamma `$’s. Thus, choosing $`\alpha _i0`$ and $`\gamma _i`$ such that $`_i\alpha _i|\gamma _i|^2=\lambda `$, $`_i\alpha _i=1\lambda `$ and $`_i\alpha _i\gamma _i=0`$, we get
$$\lambda S(\sigma _1)+(1\lambda )S(\sigma _2)\underset{i}{}\alpha _i\left(\mathrm{Tr}\sigma (\gamma _i)\right)S(\widehat{\sigma }(\gamma _i)).$$
In the above, the left hand side is the contribution to $`E_f(\rho _\lambda ;N)`$ of $`\rho _\lambda =\lambda |\psi _1\psi _1|+(1\lambda )|\psi _2\psi _2|,`$ whereas the right hand side is the contribution of
$$\rho _\lambda =\underset{i}{}\alpha _i|\psi _1+\gamma _i\psi _2\psi _1+\gamma _i\psi _2|.$$
(17)
The latter are the most general decompositions of $`\rho _\lambda `$; in fact, $`\rho _\lambda `$ as an operator acts on the two-dimensional subspace spanned by the linearly independent vectors $`\psi _i`$, $`i=1,2`$. Hence, the result follows.
With the help of Propositions 1, 2 and 3 we can now prove additivity in some special cases.
Case 1: In (12) the state $`\sigma `$ factorizes over $`M_3M_4`$: $`\sigma =\sigma _3\sigma _4`$.
Let $`E_f(\rho \sigma ;M_1M_3)`$ be achieved at an optimal decomposition made of states $`|\psi _i`$ entangled over $`M_1M_3`$. Let us consider the Schmidt decompositions $`|\psi _i=_jc_{ij}|\varphi _{ij}^{13}|\varphi _{ij}^{24}`$ with $`c_{ij}0`$ and $`|\varphi _{ij}^{13}`$ and $`|\varphi _{ij}^{24}`$ forming, for each fixed $`i`$, orthonormal bases in the first and third factor, respectively. Thus,
$`\rho \sigma _3\sigma _4`$ $`=`$ $`{\displaystyle \underset{i}{}}\alpha _i|\psi _i\psi _i|`$ (18)
$`=`$ $`{\displaystyle \underset{ijk}{}}\alpha _ic_{ij}c_{ik}|\varphi _{ij}^{13}\varphi _{ik}^{13}||\varphi _{ij}^{24}\varphi _{ik}^{24}|.`$ (19)
Further, let $`|\chi _{\mathrm{}}^3`$ be eigenvectors of $`\sigma _3`$ and consider the unitary operator ($`M_3`$)
$$\widehat{U}_{\mathrm{}}=1(1i)P_{\mathrm{}},P_{\mathrm{}}:=|\chi _{\mathrm{}}^3\chi _{\mathrm{}}^3|.$$
From Proposition 1 it follows that the vectors $`U_{\mathrm{}}|\psi _i`$, $`U_{\mathrm{}}=1_11_2\widehat{U}_{\mathrm{}}1_4`$, also give $`E_f(\rho \sigma ;M_1M_3)`$. Let us concentrate on $`|\psi _1`$; together with $`U_{\mathrm{}}|\psi _1`$, they have to satisfy (16) for all $`\gamma `$. Then, according to the notation of Proposition 3,
$`\sigma _1`$ $`=`$ $`{\displaystyle \underset{j}{}}c_{1j}^2|\varphi _{1j}^{13}\varphi _{1j}^{13}||\text{`}_{M_1M_3}`$ (20)
$`\sigma _2`$ $`=`$ $`U_{\mathrm{}}\sigma _1U_{\mathrm{}}^{},\sigma _{ov}(\gamma )=\gamma U_{\mathrm{}}\sigma _1+\gamma ^{}\sigma _1U_{\mathrm{}}^{}.`$ (21)
Taking $`\gamma =1`$, it follows that $`\widehat{\sigma }(1)={\displaystyle \frac{P_{\mathrm{}}\sigma _1P_{\mathrm{}}}{\mathrm{Tr}(P_{\mathrm{}}\sigma _1)}}`$ and $`\sigma _{ov}(1)=2\sigma _1(1i)P_{\mathrm{}}\sigma _1(1+i)\sigma _1P_{\mathrm{}}`$. Inequality (16) thus becomes
$$2\mathrm{T}\mathrm{r}\left(P_{\mathrm{}}\sigma _1\mathrm{log}\sigma _1\right)\mathrm{Tr}(P_{\mathrm{}}\sigma _1\left)S\right(\widehat{\sigma }(1)).$$
(22)
We develop $`|\varphi _{1j}^{13}=_p\beta _p\mathrm{}^j|\chi _p^1|\chi _{\mathrm{}}^3`$, along an orthonormal basis for the factor $`M_1`$, then, by means of the spectral decomposition (20), setting $`\mathrm{\Delta }_j\mathrm{}:=\varphi _j^{13}|P_{\mathrm{}}|\varphi _j^{13}=_p|\beta _p\mathrm{}^j|^2`$, we get $`P_{\mathrm{}}\sigma _1P_{\mathrm{}}=\left({\displaystyle \underset{j}{}}c_{1j}^2\mathrm{\Delta }_j\mathrm{}Q_j\mathrm{}\right)P_{\mathrm{}}`$, where $`Q_j\mathrm{}:=|\widehat{\chi }_j\mathrm{}^1\widehat{\chi }_j\mathrm{}^1|`$ and $`|\widehat{\chi }_j\mathrm{}^1={\displaystyle \underset{p}{}}{\displaystyle \frac{\beta _p\mathrm{}^j}{\mathrm{\Delta }_j\mathrm{}}}|\chi _p^1`$.
Insertion in (22) leads to
$`0`$ $``$ $`{\displaystyle \underset{j}{}}c_{1j}^2\mathrm{\Delta }_j\mathrm{}\mathrm{log}{\displaystyle \frac{\mathrm{\Delta }_j\mathrm{}}{c_{1j}^2\mathrm{Tr}(P_{\mathrm{}}\sigma )}}`$
$``$ $`{\displaystyle \underset{j}{}}c_{1j}^2\left(\mathrm{\Delta }_j\mathrm{}c_{1j}^2\mathrm{Tr}(P_{\mathrm{}}\sigma )\right),`$
the latter inequality coming from $`x\mathrm{log}x/yxy`$ and holding for all orthogonal projectors $`P_{\mathrm{}}`$. Since $`_jc_{1j}^2=1`$ and $`_{\mathrm{}}\mathrm{\Delta }_j\mathrm{}=1`$, summing over $`\mathrm{}`$ we get that $`c_{1j}=1`$ for one $`j`$ and $`c_{1k}=0`$ if $`kj`$. Thus, the supposed optimal vectors $`|\psi _i`$ must be of the form $`|\psi _i=|\varphi _i^{13}|\varphi _i^{24}`$ and the supposed optimal decomposition (19) must reduce to
$$\rho \sigma _3\sigma _4=\underset{i}{}\alpha _i|\varphi _i^{13}\varphi _i^{13}||\varphi _i^{24}\varphi _i^{24}|.$$
(23)
Tracing over $`M_2M_4`$ with respect to the Schmidt decompositions $`|\varphi _i^{13}=_j\delta _{ij}^{13}|\varphi _{ij}^1|\varphi _{ij}^3`$ and $`|\varphi _i^{24}=_j\delta _i\mathrm{}^{24}|\varphi _i\mathrm{}^2|\varphi _i\mathrm{}^4`$, orthogonality yields
$$\rho =\underset{i}{}\alpha _i\underset{j\mathrm{}}{}(\delta _{ij}^{13})^2(\delta _i\mathrm{}^{24})^2|\varphi _{ij}^1\varphi _{ij}^1||\varphi _i\mathrm{}^2\varphi _i\mathrm{}^2|.$$
(24)
We thus conclude that a decomposition of $`\rho \sigma _3\sigma _4`$ as in (19) can be optimal with respect to $`M_1M_3`$ only if $`\rho `$ is not entangled over $`M_1M_2`$, in which case $`E_f(\rho \sigma _3\sigma _4;M_1M_3)=0`$ is obviously additive. If $`\rho `$ is entangled over $`M_1M_2`$, the contradiction is avoided only if the optimal decompositions have the form
$$\rho \sigma _3\sigma _4=\underset{i}{}\alpha _i|\varphi _i^{12}\varphi _i^{12}||\varphi _i^{34}\varphi _i^{34}|.$$
(25)
Thus, the optimal states cannot carry any entanglement over $`M_1M_3`$ and additivity follows.
The second case we want to discuss is somewhat the opposite of the previous one where we proved that optimal projections for the tensor products are products of optimal projectors for the factors. In the second case, we want to show that putting together couples of optimal projectors for the factors we get optimal decompositions.
Case 2: we consider the state
$$\rho _\lambda =\lambda \rho +(1\lambda )\widehat{\rho }$$
(26)
on $`M_1M_2M_3M_4`$, where $`\rho :=|\varphi ^{12}\varphi ^{12}||\varphi ^{34}\varphi ^{34}|`$ and $`\widehat{\rho }:=|\widehat{\varphi }^{12}\widehat{\varphi }^{12}||\widehat{\varphi }^{34}\widehat{\varphi }^{34}|`$.
Let $`|\varphi ^{12}`$ and $`|\widehat{\varphi }^{12}`$ be optimal vectors for some state $`\rho `$ on $`M_1M_2`$ relative to $`M_1`$ and $`|\varphi ^{34}=|\varphi ^3|\varphi ^4`$, $`|\widehat{\varphi }^{34}=|\widehat{\varphi }^3|\widehat{\varphi }^4`$ on $`M_3M_4`$ so that $`E(\rho _\lambda ;M_3)=0`$. The contribution to $`E(\rho _\lambda ;M_1M_3)`$ of the decomposition (26) is thus
$$E_\lambda :=\lambda S\left(|\varphi ^{12}\varphi ^{12}||\text{`}_{M_1}\right)+(1\lambda )S\left(|\widehat{\varphi }^{12}\widehat{\varphi }^{12}||\text{`}_{M_1}\right)$$
(27)
and we want to prove that this is the best we can have.
We proceed as follows: as for (17), a general decomposition of $`\rho _\lambda `$ is of the form $`\rho _\lambda =_i\alpha _i|\psi _i\psi _i|`$ where $`|\psi _i=|\varphi ^{12}\varphi ^{34}+\gamma _i|\widehat{\varphi }^{12}\widehat{\varphi }^{34}`$, with $`\alpha _i>0`$ and
$$\underset{i}{}\alpha _i=\lambda ,\underset{i}{}\alpha _i|\gamma _i|^2=1\lambda ,\underset{i}{}\alpha _i\gamma _i=0.$$
(28)
We now set $`b:=\varphi ^4|\widehat{\varphi }^4`$, $`a:=\sqrt{1|b|^2}`$ and construct the normalized vector state $`|\psi ^4:={\displaystyle \frac{|\widehat{\varphi }^4b|\varphi ^4}{a}}`$ such that $`\psi ^4|\varphi ^4=0`$. We can thus rewrite
$$|\psi _i=a_i|\varphi _i^{123}\varphi ^4+a\gamma _i|\widehat{\varphi }^{12}\widehat{\varphi }^3\psi ^4,$$
(29)
where
$`|\varphi _i^{123}:={\displaystyle \frac{|\varphi ^{12}\varphi ^3+b\gamma _i|\widehat{\varphi }^{12}\widehat{\varphi }^3\varphi ^4}{a_i}}`$
$`a_i^2:=1+|b|^2|\gamma _i|^2+2e\left(b\gamma _i\widehat{\varphi }^{12}|\varphi ^{12}\widehat{\varphi }^3|\varphi ^3\right).`$
With $`|\widehat{\psi }_i:={\displaystyle \frac{|\psi _i}{\sqrt{\delta _i}}}`$, $`\delta _i:=a_i^2+a^2|\gamma _i|^2`$, the decomposition (26) reads $`\rho _\lambda =_i\alpha _i\delta _i|\widehat{\psi }_i\widehat{\psi }_i|`$.
The contribution of the latter to the entanglement of formation $`E_f(\rho _\lambda ;M_1M_3)`$ is
$$E:=\underset{i}{}\alpha _i\delta _iS\left(|\widehat{\psi }_i\widehat{\psi }_i||\text{`}_{M_1M_3}\right).$$
(30)
From the orthogonality of $`\psi ^4`$ and $`\varphi ^4`$ it follows that
$$|\widehat{\psi }_i\widehat{\psi }_i||\text{`}_{M_1M_3}=\frac{a_i^2}{\delta _i}\sigma _i^{123}+\frac{a^2|\gamma _i|^2}{\delta _i}\sigma ^{123},$$
where
$`\sigma _i^{123}`$ $`:=`$ $`|\varphi _i^{123}\varphi _i^{123}||\text{`}_{M_1M_3}`$ (31)
$`\sigma ^{123}`$ $`:=`$ $`|\widehat{\varphi }^{12}\widehat{\varphi }^{12}||\text{`}_{M_1}|\widehat{\varphi }^3\widehat{\varphi }^3||\text{`}_{M_3}.`$ (32)
Concavity of the von Neumann entropy yields
$`E{\displaystyle \underset{i}{}}\alpha _i\{a_i^2S\left(\sigma _i^{123}|\text{`}_{M_1M_3}\right)`$ (33)
$`a^2|\gamma _i|^2S(|\widehat{\varphi }^{12}\widehat{\varphi }^{12}||\text{`}_{M_1})\}.`$ (34)
As done before, we construct the normalized vector $`|\psi ^3:={\displaystyle \frac{|\widehat{\varphi }^3d|\varphi ^3}{c}}`$, such that $`\psi ^3|\varphi ^3=0`$ where $`d:=\varphi ^3|\widehat{\varphi }^3`$, $`c:=\sqrt{1|d|^2}`$, and
$`|\varphi _i^{123}`$ $`:=`$ $`{\displaystyle \frac{b_i|\psi _i^{12}\varphi ^3+bc\gamma _i|\widehat{\varphi }^{12}\psi ^3}{a_i}}`$ (35)
$`|\psi _i^{12}`$ $`:=`$ $`{\displaystyle \frac{|\varphi ^{12}+bd\gamma _i|\widehat{\varphi }^{12}}{b_i}}`$ (36)
$`b_i^2`$ $`:=`$ $`1+|b|^2|d|^2|\gamma _i|^2+2e\left(bd\gamma _i\widehat{\varphi }^{12}|\varphi ^{12}\right).`$ (37)
Introducing the Schmidt decompositions over $`M_1M_2`$: $`|\psi _i^{12}=_jc_{ij}|\varphi _{ij}^1|\varphi _{ij}^2`$, $`|\widehat{\varphi }^{12}=_{\mathrm{}}d_{\mathrm{}}|\widehat{\varphi }_{\mathrm{}}^1|\widehat{\varphi }_{\mathrm{}}^2`$, and setting $`\rho _1:=|\psi _i^{12}\psi _i^{12}||\text{`}_{M_1}`$, $`\rho _2:=|\widehat{\varphi }^{12}\widehat{\varphi }^{12}||\text{`}_{M_1}`$, because of the orthogonality of $`\varphi ^3`$ and $`\psi ^3`$, the state $`\sigma _i^{123}`$ in (31) restricted to $`M_1M_3`$ can be represented as
$`\sigma _i^{123}`$ $`=`$ $`{\displaystyle \frac{1}{a_i^2}}\left(\begin{array}{cc}b_i^2\rho _1& cb_ib^{}\gamma _i^{}\sqrt{\rho _1}V\sqrt{\rho _2}\\ cb_ib\gamma _i\sqrt{\rho _2}V^{}\sqrt{\rho _1}& |b|^2c^2|\gamma _i|^2\rho _2\end{array}\right)`$
$`=`$ $`{\displaystyle \frac{1}{a_i^2}}\left(\begin{array}{cc}b_i\sqrt{\rho _1}V& 0\\ cb\gamma _i\sqrt{\rho _2}& 0\end{array}\right)\left(\begin{array}{cc}b_iV^{}\sqrt{\rho _1}& cb^{}\gamma _i^{}\sqrt{\rho _2}\\ 0& 0\end{array}\right),`$
where $`V:=_{j,\mathrm{}}\widehat{\varphi }_{\mathrm{}}^2|\varphi _{ij}^{12}|\varphi _{ij}^1\widehat{\varphi }_{\mathrm{}}^1|`$ is a unitary operator and $`\sqrt{\rho _1}V\sqrt{\rho _2}=|\psi _i^{12}\widehat{\varphi }^{12}||\text{`}_{M_1}`$.
Since $`\sigma _i^{123}=A^{}A`$ has the same entropy as $`AA^{}={\displaystyle \frac{1}{a_i^2}}\left(\begin{array}{cc}b_i^2V^{}\rho _1V+|b|^2c^2|\gamma _i|^2\rho _2& 0\\ & \\ 0& 0\end{array}\right)`$, concavity and invariance under unitary transformations of the von Neumann entropy yield
$$S\left(\sigma _i^{123}\right)=S\left(AA^{}\right)\frac{b_i^2}{a_i^2}S(\rho _1)+\frac{c^2|b|^2|\gamma _i|^2}{a_i^2}S\left(\rho _2\right),$$
whence (34) becomes
$$E\underset{i}{}\alpha _i\{b_i^2S(\rho _1)+|\gamma _i|^2(a^2+c^2|b|^2)S(\rho _2)\}.$$
(38)
Since we assumed the states $`|\varphi ^{12}`$ and $`|\widehat{\varphi }^{12}`$ in (36) to be optimal for some state on $`M_1M_2`$ when restricted to $`M_1`$, we can use the necessary condition (16). According to the notation of Proposition 3, we have $`\sigma _1=\rho _2`$, $`\sigma _2=|\varphi ^{12}\varphi ^{12}||\text{`}_{M_1}`$, $`\gamma =\gamma _i^{}b^{}d^{}`$, $`\widehat{\sigma }(\gamma )=\rho _1`$ and
$$\sigma _{ov}(\gamma )=b^{}d^{}\gamma _i^{}\sqrt{\rho _2}V^{}\sqrt{\rho _1}bd\gamma _i\sqrt{\rho _1}V\sqrt{\rho _2}.$$
From (16) and the conditions (28) it follows that
$`E`$ $``$ $`{\displaystyle \underset{i}{}}\alpha _i\{S(|\varphi ^{12}\varphi ^{12}||\text{`}_{M_1})+|\gamma _i|^2S(|\widehat{\varphi }^{12}\widehat{\varphi }^{12}||\text{`}_{M_1})`$ (40)
$`\mathrm{Tr}\left(\sigma _{ov}(b^{}d^{}\gamma _i^{})\mathrm{log}\rho _2\right)\}=E_\lambda ,`$
where $`E_\lambda `$ is the contribution (27) to the entanglement of formation $`E(\rho _\lambda ;M_1M_3)`$ of the decomposition (26), which turns out then to be already optimal.
In this letter we have derived a necessary and sufficient condition for the optimality of two vector states and showed its usefulness by proving additivity in two cases. While in case 1. additivity was rather expected because of the tensor-product state $`\rho \sigma `$, it was less so in case 2. We requested, however, additional properties on the state structure over $`M_3M_4`$: in case 1. factorization of $`\sigma =\sigma _3\sigma _4`$ and in case 2. factorization into pure states of the optimal decomposers. In both cases the state was thus separable with respect to $`M_3M_4`$.
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# B-ball Dark Matter and Baryogenesis
## Abstract
It has been recently suggested that stable, supersymmetric B-balls formed in the early universe could not only be the dark matter at the present epoch, but also be responsible for baryogenesis by their partial evaporation at high temperatures. We reinvestigate the efficiency of B-ball baryogenesis and find it to be limited by the diffusion of baryon number away from the B-balls. Successful baryogenesis may only occur for B-balls with charges $`Q{}_{}{}^{<}\mathrm{\hspace{0.17em}10}_{}^{20}5\times 10^{23}`$, which is close to the observational lower limits on the $`Q`$ of a significant B-ball dark matter component. We also present some cosmological constraints on the abundances of larger B-balls in the early universe.
PACS numbers: 95.35.+d, 98.80.Cq, 98.80.Ft
thanks: e-mail: banerjee@mpa-garching.mpg.de, jedamzik@mpa-garching.mpg.de
It is well known that particle physics models containing an unbroken U(1) symmetry allow for the existence of non-topological solitons , i.e. Q-balls, which carry a large number of a conserved global charge . If the effective potential $`U(\mathrm{\Phi })`$, of the scalar $`\mathrm{\Phi }`$ carrying the global charge grows slower than the second power of $`\mathrm{\Phi }`$, the mass of the solitonic object scales with the U(1) charge $`Q`$ as $`M_Q\stackrel{~}{M}Q^p`$ ($`0<p<1`$), where $`\stackrel{~}{M}`$ is some energy scale. The minimal supersymmetric standard model (MSSM) with supersymmetry breaking communicated at low energy scale contains an effective potential which is nearly flat $`M^4`$ at large $`\mathrm{\Phi }`$. In this case, $`M_QMQ^{3/4}`$, where $`M110`$ TeV is the SUSY breaking scale . Such Q-balls are absolutely stable at zero temperature if their mass $`M_Q`$ becomes smaller than the total mass of individual U(1) charged particles $`mQ^1`$. Moreover, large baryonic Q-balls may be efficiently produced in the early universe within a scenario of a collapsing unstable Affleck-Dine condensate . In the MSSM the role of the global charge is played by baryon or lepton number carried by squarks or sleptons respectively . Whereas L-balls (carrying leptonic charge) are not expected to survive until the present epoch for $`Q{}_{}{}^{<}\mathrm{\hspace{0.17em}10}_{}^{36}(M/\mathrm{TeV})^{4/5}`$ due to emission of massless neutrinos, B-balls (containing baryonic charge) with $`Q{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{12}(M/\mathrm{TeV})^4`$ are stable because of the largeness of the nucleon mass. It has thus been proposed that B-balls produced in the early universe are not only an attractive dark matter candidate but may also be responsible for baryogenesis via partial evaporation of B-balls in the early universe . (This is distinct from a scenario of evaporation of unstable B-balls which could be responsible for baryogenesis and the creation of neutralino dark matter.)
In this paper we reinvestigate the evaporation of B-balls. We find the efficiency of this process to be limited by the transport of baryon number away from the soliton, resulting in somewhat different conclusions than drawn in prior work . We also give some previously unmentioned cosmological limits on the abundances of B-balls.
B-balls may release baryonic charge via evaporation of squarks at temperatures $`T>m_\chi `$, where $`m_\chi `$ is the squark mass . Assuming that B-balls constitute dark matter at present (with fractional contribution to the critical density of $`\mathrm{\Omega }_Q`$), there number density $`n_Q`$ in the early universe at temperature $`T`$ has been
$$n_Q0.8\times 10^9g_{}\frac{m_p}{M}\frac{\mathrm{\Omega }_Q}{Q^p}T^3.$$
(1)
Here $`m_p`$ is proton mass, $`g_{}`$ is the number of relativistic degrees of freedom at the considered epoch ($`g_{}=3.909`$ today), and we have assumed a Hubble constant of $`H=70`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. Consider B-balls with the general properties for their mass
$$M_Q=\alpha MQ^p$$
(2)
and their radius
$$R_Q=\beta M^1Q^{p/3}.$$
(3)
For an essentially flat effective potential (e.g. $`UM^4`$) for $`\mathrm{\Phi }m_\chi `$ the coefficient $`p=3/4`$ and the numerical factors may be computed for large B-ball charge: $`\alpha (4\pi /3)\sqrt{2}`$, $`\beta 1/\sqrt{2}`$ . If, on the other side, the expression $`U(\mathrm{\Phi })/\mathrm{\Phi }^2`$ has a minimum at finite $`\mathrm{\Phi }`$ then the B-ball mass scales with $`p=1`$ .
To trace the baryon evaporation rate of a population of B-balls at high temperatures, we consider the thermodynamical properties of a volume element containing only one B-ball. Let $`Q`$ be the charge of the B-ball surrounded by $`N_B`$ baryon number carrying particles in the diffuse plasma in a volume $`V`$, such that the total charge in the volume element is $`N_{\text{tot}}=N_B+Q`$. Then the Helmholtz free energy of the system reads as follows
$$F=VT^4\frac{7\pi ^2g_B}{360}+\frac{3N_B^2}{g_BVT^2}+\alpha MQ^p,$$
(4)
where $`g_B`$ is the number of baryon charge carrying degrees of freedom and where we assumed a vanishing entropy, $`S_Q0`$, for the B-ball. Expression (4) further assumes that baryon number in the plasma is carried by fermionic relativistic degrees of freedom (cf. ). For fixed $`N_{\text{tot}}`$, $`V`$, and $`T`$ such a configuration has to minimize the Helmholtz free energy Eq. (4) in thermal and chemical equilibrium. This extremum may be obtained by varying the fraction of $`N_{\text{tot}}`$ residing in the B-ball. One thus finds that only for charge density in excess of
$$\frac{N_{\text{tot}}}{V}\left(\frac{g_B}{6}\alpha MT^2\right)^{1/(2p)}V^{(p1)/(2p)}$$
(5)
the existence of a B-ball is thermodynamically favorable. In this case the B-ball carries almost the entire charge of the volume element ($`QN_{\text{tot}}`$). Nevertheless, in chemical equilibrium a small fraction of baryon number also resides in the plasma, with baryonic density
$$n_B^{\text{eq}}=\frac{g_B}{6}p\alpha MT^2Q^{p1}.$$
(6)
We stress that by constraining our considerations to the thermodynamics of a finite domain with volume $`V`$ and baryon number $`N_{\text{tot}}`$ we have not obtained the absolute minimum for the free energy. Considering a larger domain, but with $`N_{\text{tot}}/V`$ kept fixed, one may always find a state of lower Helmholtz free energy. This is also exemplified by the odd dependence of Eq. (5) on volume (except for $`p=1`$). As already noted in , in the extreme limit the lowest free energy state is reached when all the baryon number of the universe is contained in one large B-ball.
Nevertheless, considerations of partial chemical equilibrium are appropriate for the initial conditions envisioned resulting from the breakup of an Affleck-Dine condensate. Here B-balls of typical initial charge $`Q_0`$ form, with negligible baryon number in the plasma . During the subsequent evolution of the universe the coalescence of B-balls is not possible, such that a state of even lower Helmholtz free energy is not attainable, and only partial chemical equilibrium between an individual B-ball and the plasma around it may be attained. Following similar arguments Laine & Shaposhnikov estimated the baryon evaporation rate of B-balls via squark emission at high $`T`$ by
$`\mathrm{\Gamma }_{\text{evap}}{\displaystyle \frac{dQ}{dt}}`$ $`=`$ $`\kappa (\mu _Q\mu _{\text{plasma}})T^2\mathrm{\hspace{0.17em}4}\pi R_Q^2`$ (7)
$``$ $`\kappa ^{}\mathrm{\hspace{0.17em}4}\pi R_Q^2n_B^{\text{eq}}\mathrm{for}\mu _{\text{plasma}}<<\mu _Q,`$
where $`\mu _Q`$ is the chemical potential of the B-ball. The second line of Eq. (7) is operative under the assumption that the evaporated particles can be quickly transported away from the B-ball surface in order to sustain a jump between chemical potentials of the B-ball and the plasma ($`\mu _Q\mu _{\text{plasma}}0`$). Note that $`\mu _Q=\mu _{\text{plasma}}=(p/6)\alpha MQ^{p1}`$ in chemical equilibrium. The constants $`\kappa `$ and $`\kappa ^{}`$ in Eq. (7) are $`{}_{}{}^{<}\mathrm{\hspace{0.17em}1}`$, where for $`\kappa ^{}1`$ the evaporation rate is at its physical upper limit, implying there is no dynamical suppression of the evaporation and accretion of squarks, i.e. every squark approaching a B-ball will be absorbed. This is in contrast to the evaporation of quarks from B-balls which is suppressed and thus of negligible importance. Nevertheless, baryon transport away from the B-ball is not as efficient as envisioned in Ref. . Baryons released from the B-ball surface, will establish chemical equilibrium with the B-ball, if they are not able to escape the surface layer and evaporation ceases.
The transport mechanism of baryon number is by diffusion of squarks and quarks in the hot plasma. Solving the spherical diffusion equation with diffusion constant $`D`$
$$\frac{n_B(r,t)}{t}=D\frac{1}{r}\left(\frac{^2}{r^2}rn_B(r,t)\right),$$
(8)
on the condition that the number density at the surface boundary does not change with time ($`n_B(R_Q,t)=n_B^{\text{eq}}`$), yields a steady-state solution for the density $`n_b`$. We confirmed this result numerically for all radii $`r{}_{}{}^{<}L_{\text{d}}^{}`$, where $`L_\text{d}\sqrt{Dt}`$ is the diffusion length at time $`t`$. Therefore the particle flux through the B-ball surface is constant and given by
$$\mathrm{\Gamma }_{\text{diff}}\frac{dQ}{dt}=4\pi kR_QDn_B^{\text{eq}}$$
(9)
where the diffusion constant of relativistic squarks and quarks in a hot plasma is $`DaT^1`$ with $`a6`$ for quarks and $`a4`$ for squarks, respectively . We have determined the numerical constant $`k`$ to be very close to unity, such that we will drop it in what follows. Apart from the numerical results the expression (9) can be motivated by assuming a constant flux $`dQ/dt=4\pi R_Q^2n/r`$ and approximating $`n/r`$ by $`n_B^{\text{eq}}/\mathrm{\Delta }r`$. The only time independent lengthscale in this system is the B-ball radius $`R_Q`$, so $`n_B^{\text{eq}}/\mathrm{\Delta }r=n_B^{\text{eq}}/R_Q`$.
By comparing the rates (7) and (9)
$$\mathrm{\Gamma }_{\text{diff}}/\mathrm{\Gamma }_{\text{evap}}=\frac{D}{\kappa ^{}R_Q}\frac{M}{T}Q^{p/3}$$
(10)
it is obvious that, for large B-balls, the diffusive transport is orders of magnitude less efficient than the evaporation of baryons from the B-ball. Since the evaporated baryons are still within the surface layer of the B-ball, the B-ball is at close to chemical equilibrium with the surrounding plasma and the charge emission rate (7) must be replaced by (9).
Evaporation of squarks from B-balls is only efficient for temperatures above the squark mass $`m_\chi 0.11\text{TeV}`$. For temperatures below this mass the evaporation rate is exponentially suppressed by the Boltzmann factor $`\mathrm{exp}(m_\chi /T)`$. By integrating Eq. (9) one may calculate the number of emitted baryons from a single B-ball until evaporation becomes inefficient at temperature $`T_{\text{fin}}m_\chi `$
$$\mathrm{\Delta }Qb\frac{M_0}{T_{\text{fin}}}Q_0^{\frac{4}{3}p1}.$$
(11)
Here $`Q_0`$ is the initial B-ball charge, and $`M_0=(90/32\pi ^3g_{})^{1/2}M_{\text{Pl}}3.7\times 10^{18}/\sqrt{g_{}(T_{\text{fin}})}\text{GeV}`$ is given by the time-temperature relation $`t=M_0T^2`$ during a radiation dominated universe with $`g_{}(T_{\text{fin}})200`$. Note that for the interesting case $`p=3/4`$ the numerical constant in Eq. (11) $`b=(4\pi /3)g_B\beta ap\alpha 4.7\times 10^3`$, assuming $`g_B72`$, and $`\mathrm{\Delta }Q`$ is independent of the initial charge of the B-ball. To ensure that an initially formed B-ball survives evaporation until present ($`\mathrm{\Delta }Q/Q_01`$), such a B-ball must have an initial charge of $`Q_0>10^{18}10^{19}`$.
Within a B-ball baryogenesis scenario the number density of baryonic matter $`n_B`$ and the number density of B-balls are related by
$$n_B=n_Q\mathrm{\Delta }Q.$$
(12)
Combining Eq. (11), (12) and Eq. (1) one may calculate the baryon-to-photon ratio at the present epoch
$$\eta =\frac{n_B}{n_\gamma }1.3\times 10^8b\frac{m_p}{M}\frac{\mathrm{\Omega }_Q}{Q_0^{1\frac{1}{3}p}}\frac{M_0}{T_{\text{fin}}}.$$
(13)
For a B-ball with $`p=3/4`$, $`M110\text{TeV}`$, and $`T_{\text{fin}}0.11`$TeV it is necessary to have $`Q_0{}_{}{}^{>}\mathrm{\hspace{0.17em}9}\times 10^{20}4\times 10^{23}`$ to obtain a baryon-to-photon ratio of $`\eta 3\times 10^{10}`$. This should be compared to the range $`Q_010^{22}10^{28}`$ quoted in Ref. . (Note that the above estimate has also very different dependence on the parameters $`M`$ and $`T_{\text{fin}}`$ than the estimate by Ref. .) It is intriguing that the range of B-ball charges which may yield successful baryogenesis is very close to the observational lower limits $`Q{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{21}3\times 10^{22}`$ on the charges of a significant B-ball galactic halo dark matter population .
A few comments concerning very large Q-Balls are of relevance. There are no detector limits on large B- and L-balls, either since their current flux is extremely small, for stable solitons, or since they did not survive to the present epoch, in the unstable case. There are, however, some constraints on the existence of large Q-Balls in the early universe. Consider first unstable $`(p=1)`$ B-balls. During their decay they produce baryon inhomogeneities which may lead to a scenario of inhomogeneous nucleosynthesis. It is well known , that baryon lumps with baryon number in excess of
$$N_b{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{35}(\eta _l/10^4)^{1/2}(\eta _l10^{10})$$
(14)
can not homogenize by neutron diffusion before the epoch of weak freeze-out. Here $`\eta _l`$ is the baryon-to-photon ratio in the baryon-rich region, and the value of $`10^4`$ is of particular relevance. It is expected that immediately after the B-ball decay $`\eta _l^i10^4`$ (where we assume that baryon number is in form of diffuse baryons). Neutrino heat conduction will subsequently expand the baryonic lump to an asymptotic $`\eta _l10^4`$, independent of $`\eta _l^i`$, unless the baryonic B-ball charge is in excess of $`N_b{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{44}\eta _l^i`$ . If a fraction $`f_b{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{2}`$ of all baryons resides in such baryon number enhanced regions, overproduction of <sup>4</sup>He during nucleosynthesis results. For $`N_b{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{44}\eta _l^i`$ similar constraints from <sup>4</sup>He overproduction apply, but here even more stringent constraints may be derived from a possible overproduction of heavy elements. It is interesting to note, that for $`f_b{}_{}{}^{<}\mathrm{\hspace{0.17em}10}_{}^{2}`$ B-balls with $`N_b{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{35}`$ may yield to the production of a primordial metallicity, without violating observational constraints on the light element abundances .
A more speculative constraint may apply for large stable B-balls. If their charge is in excess of $`Q{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{44}\mathrm{\Omega }_Q^{4/3}(M/\mathrm{TeV})^{4/3}`$ their mean separation ($`n_Q^{1/3}`$) at the QCD transition at temperature $`T100`$MeV is $`{}_{}{}^{>}\mathrm{\hspace{0.17em}1}`$m. In the case of a first-order QCD phase transition they may act as seeds for the nucleation of hadronic phase. Depending on the amount of supercooling which quark-gluon plasma may sustain before spontaneous nucleation is efficient, and thus on the (three) surface free energies between the participants, hadronic phase, quark-gluon phase, and B-balls, hadronic phase bubbles may only form around B-balls. If this is the case, the mean separation of baryon number enhancing quark-gluon plasma bubbles towards the end of the transition is such that baryon number inhomogeneities with $`N_b{}_{}{}^{>}\mathrm{\hspace{0.17em}10}_{}^{35}`$ of individual lumps (the baryon number within $`1`$m at $`\eta 3\times 10^{10}`$) is large enough to yield a scenario of inhomogeneous nucleosynthesis. Except for very narrow ranges in parameter space such a scenario is typically in conflict with observationally determined light element abundances.
In summary, we have reinvestigated a proposed scenario of baryogenesis by the partial evaporation of stable B-balls in the early universe. Under the assumption that the B-balls are the dark matter at the present epoch, we have found that a successful baryogenesis scenario by B-ball evaporation at high temperature requires B-balls with baryon number $`Q10^{20}5\times 10^{23}`$, which is close to the observational lower limit on the charges of a significant ($`\mathrm{\Omega }_Q1`$) galactic B-ball population. Thus, if stable B-balls are responsible for baryogenesis, and they constitute the dark matter, they could be detected in the immediate future. We have also given some limits on the existence of larger B-balls in the early universe.
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# Reasoning with Individuals for the Description Logic 𝒮ℋℐ𝒬 This paper will appear in the Proceedings of the 17th International Conference on Automated Deduction (CADE-17), Lecture Notes in Computer Science, Germany, 2000. Springer Verlag.
## 1 Motivation
A description logic (DL) knowledge base (KB) is made up of two parts, a terminological part (the terminology or Tbox) and an assertional part (the Abox), each part consisting of a set of axioms. The Tbox asserts facts about *concepts* (sets of objects) and *roles* (binary relations), usually in the form of inclusion axioms, while the Abox asserts facts about *individuals* (single objects), usually in the form of instantiation axioms. For example, a Tbox might contain an axiom asserting that Man is subsumed by Animal, while an Abox might contain axioms asserting that both Aristotle and Plato are instances of the concept Man and that the pair $`\text{Aristotle},\text{Plato}`$ is an instance of the role Pupil-of.
For logics that include full negation, all common DL reasoning tasks are reducible to deciding KB consistency, i.e., determining if a given KB admits a non-empty interpretation . There has been a great deal of work on the development of reasoning algorithms for expressive DLs , but in most cases these consider only Tbox reasoning (i.e., the Abox is assumed to be empty). With expressive DLs, determining consistency of a Tbox can often be reduced to determining the satisfiability of a single concept , and—as most DLs enjoy the tree model property (i.e., if a concept has a model, then it has a tree model)—this problem can be decided using a tableau-based decision procedure.
The relative lack of interest in Abox reasoning can also be explained by the fact that many applications only require Tbox reasoning, e.g., ontological engineering and schema integration . Of particular interest in this regard is the DL $`𝒮𝒬`$ , which is powerful enough to encode the logic $`𝒟`$ , and which can thus be used for reasoning about conceptual data models, e.g., Entity-Relationship (ER) schemas . Moreover, if we think of the Tbox as a *schema* and the Abox as (possibly incomplete) *data*, then it seems reasonable to assume that realistic Tboxes will be of limited size, whereas realistic Aboxes could be of almost unlimited size. Given the high complexity of reasoning in most DLs , this suggests that Abox reasoning could lead to severe tractability problems in realistic applications.<sup>1</sup><sup>1</sup>1Although suitably optimised algorithms may make reasoning practicable for quite large Aboxes .
However, $`𝒮𝒬`$ Abox reasoning is of particular interest as it allows $`𝒟`$ schema reasoning to be extended to reasoning about conjunctive query containment w.r.t. a schema . This is achieved by using Abox individuals to represent variables and constants in the queries, and to enforce co-references . In this context, the size of the Abox would be quite small (it is bounded by the number of variables occurring in the queries), and should not lead to severe tractability problems.
Moreover, an alternative view of the Abox is that it provides a restricted form of reasoning with *nominals*, i.e., allowing individual names to appear in concepts . Unrestricted nominals are very powerful, allowing arbitrary co-references to be enforced and thus leading to the loss of the tree model property. This makes it much harder to prove decidability and to devise decision procedures (the decidability of $`𝒮𝒬`$ with unrestricted nominals is still an open problem). An Abox, on the other hand, can be modelled by a *forest*, a set of trees whose root nodes form an arbitrarily connected graph, where number of trees is limited by the number of individual names occurring in the Abox. Even the restricted form of co-referencing provided by an Abox is quite powerful, and can extend the range of applications for the DLs reasoning services.
In this paper we present a tableaux based algorithm for deciding the satisfiability of unrestricted $`𝒮𝒬`$ KBs (i.e., ones where the Abox may be non-empty) that extends the existing consistency algorithm for Tboxes by making use of the forest model property. This should make the realisation of an efficient implementation relatively straightforward as it can be based on an existing highly optimised implementation of the Tbox algorithm (e.g., in the FaCT system ). A notable feature of the algorithm is that, instead of making a unique name assumption w.r.t. all individuals (an assumption commonly made in DLs ), increased flexibility is provided by allowing the Abox to contain axioms explicitly asserting inequalities between pairs of individual names (adding such an axiom for every pair of individual names is obviously equivalent to making a unique name assumption).
## 2 Preliminaries
In this section, we introduce the DL $`𝒮𝒬`$. This includes the definition of syntax, semantics, inference problems (concept subsumption and satisfiability, Abox consistency, and all of these problems with respect to terminologies<sup>2</sup><sup>2</sup>2We use *terminologies* instead of Tboxes to underline the fact that we allow for general concept inclusions axioms and do not disallow cycles.), and their relationships.
$`𝒮𝒬`$ is based on an extension of the well known DL $`𝒜𝒞`$ to include transitively closed primitive roles ; we call this logic $`𝒮`$ due to its relationship with the proposition (multi) modal logic $`\mathrm{𝐒𝟒}_{(𝐦)}`$ .<sup>3</sup><sup>3</sup>3The logic $`𝒮`$ has previously been called $`𝒜𝒞_{R^+}`$, but this becomes too cumbersome when adding letters to represent additional features. This basic DL is then extended with inverse roles ($``$), role hierarchies ($``$), and qualifying number restrictions ($`𝒬`$).
###### Definition 2.1
Let $`𝐂`$ be a set of *concept names* and $`𝐑`$ a set of *role names* with a subset $`𝐑_+𝐑`$ of *transitive role names*. The set of roles is $`𝐑\{R^{}R𝐑\}`$. To avoid considering roles such as $`R^{}`$, we define a function $`𝖨𝗇𝗏`$ on roles such that $`𝖨𝗇𝗏(R)=R^{}`$ if $`R`$ is a role name, and $`𝖨𝗇𝗏(R)=S`$ if $`R=S^{}`$. We also define a function $`𝖳𝗋𝖺𝗇𝗌`$ which returns $`\mathrm{true}`$ iff $`R`$ is a transitive role. More precisely, $`𝖳𝗋𝖺𝗇𝗌(R)=\mathrm{true}`$ iff $`R𝐑_+`$ or $`𝖨𝗇𝗏(R)𝐑_+`$.
A *role inclusion axiom* is an expression of the form $`RS`$, where $`R`$ and $`S`$ are roles, each of which can be inverse. A *role hierarchy* is a set of role inclusion axioms. For a role hierarchy $``$, we define the relation $`\text{*}`$ to be the transitive-reflexive closure of $``$ over $`\{𝖨𝗇𝗏(R)𝖨𝗇𝗏(S)RS\}`$. A role $`R`$ is called a *sub-role* (resp. super-role) of a role $`S`$ if $`R\text{*}S`$ (resp. $`S\text{*}R`$). A role is *simple* if it is neither transitive nor has any transitive sub-roles.
The set of $`𝒮𝒬`$-*concepts* is the smallest set such that
* every concept name is a concept, and,
* if $`C`$, $`D`$ are concepts, $`R`$ is a role, $`S`$ is a simple role, and $`n`$ is a nonnegative integer, then $`CD`$, $`CD`$, $`\neg C`$, $`R.C`$, $`R.C`$, $`nS.C`$, and $`nS.C`$ are also concepts.
A *general concept inclusion axiom* (GCI) is an expression of the form $`CD`$ for two $`𝒮𝒬`$-concepts $`C`$ and $`D`$. A *terminology* is a set of GCIs.
Let $`𝐈=\{a,b,c\mathrm{}\}`$ be a set of *individual names*. An *assertion* is of the form $`a:C`$, $`(a,b):R`$, or $`a\doteq ̸b`$ for $`a,b𝐈`$, a (possibly inverse) role $`R`$, and a $`𝒮𝒬`$-concept $`C`$. An *Abox* is a finite set of assertions.
Next, we define semantics of $`𝒮𝒬`$ and the corresponding inference problems.
###### Definition 2.2
An interpretation $`=(\mathrm{\Delta }^{},^{})`$ consists of a set $`\mathrm{\Delta }^{}`$, called the domain of $``$, and a *valuation* $`^{}`$ which maps every concept to a subset of $`\mathrm{\Delta }^{}`$ and every role to a subset of $`\mathrm{\Delta }^{}\times \mathrm{\Delta }^{}`$ such that, for all concepts $`C`$, $`D`$, roles $`R`$, $`S`$, and non-negative integers $`n`$, the following equations are satisfied, where $`\mathrm{}M`$ denotes the cardinality of a set $`M`$ and $`(R^{})^+`$ the transitive closure of $`R^{}`$:
| $`R^{}`$ | $`=`$ | $`(R^{})^+`$ | for each role $`R𝐑_+`$ |
| --- | --- | --- | --- |
| $`(R^{})^{}`$ | $`=`$ | $`\{x,yy,xR^{}\}`$ | (inverse roles) |
| $`(CD)^{}`$ | $`=`$ | $`C^{}D^{}`$ | (conjunction) |
| $`(CD)^{}`$ | $`=`$ | $`C^{}D^{}`$ | (disjunction) |
| $`(\neg C)^{}`$ | $`=`$ | $`\mathrm{\Delta }^{}C^{}`$ | (negation) |
| $`(R.C)^{}`$ | $`=`$ | $`\{xy.x,yR^{}\text{ and }yC^{}\}`$ | (exists restriction) |
| $`(R.C)^{}`$ | $`=`$ | $`\{xy.x,yR^{}\text{ implies }yC^{}\}`$ | (value restriction) |
| $`(nR.C)^{}`$ | $`=`$ | $`\{x\mathrm{}\{y.x,yR^{}\text{ and }yC^{}\}n\}`$ | ($``$-number restriction) |
| $`(nR.C)^{}`$ | $`=`$ | $`\{x\mathrm{}\{y.x,yR^{}\text{ and }yC^{}\}n\}`$ | ($``$-number restriction) |
An interpretation $``$ *satisfies* a role hierarchy $``$ iff $`R^{}S^{}`$ for each $`RS`$ in $``$. Such an interpretation is called a *model* of $``$ (written $``$).
An interpretation $``$ *satisfies* a terminology $`𝒯`$ iff $`C^{}D^{}`$ for each GCI $`CD`$ in $`𝒯`$. Such an interpretation is called a *model* of $`𝒯`$ (written $`𝒯`$).
A concept $`C`$ is called satisfiable with respect to a role hierarchy $``$ and a terminology $`𝒯`$ iff there is a model $``$ of $``$ and $`𝒯`$ with $`C^{}\mathrm{}`$. A concept $`D`$ subsumes a concept $`C`$ w.r.t. $``$ and $`𝒯`$ iff $`C^{}D^{}`$ holds for each model $``$ of $``$ and $`𝒯`$. For an interpretation $``$, an element $`x\mathrm{\Delta }^{}`$ is called an instance of a concept $`C`$ iff $`xC^{}`$.
For Aboxes, an interpretation maps, additionally, each individual $`a𝐈`$ to some element $`a^{}\mathrm{\Delta }^{}`$. An interpretation $``$ satisfies an assertion
$`\begin{array}{ccc}\hfill a:C& \text{ iff }& a^{}C^{},\hfill \\ \hfill (a,b):R& \text{ iff }& a^{},b^{}R^{},\text{ and}\hfill \\ \hfill a\doteq ̸b& \text{ iff }& a^{}b^{}\hfill \end{array}`$
An Abox $`𝒜`$ is *consistent* w.r.t. $``$ and $`𝒯`$ iff there is a model $``$ of $``$ and $`𝒯`$ that satisfies each assertion in $`𝒜`$.
For DLs that are closed under negation, subsumption and (un)satisfiability can be mutually reduced: $`CD`$ iff $`C\neg D`$ is unsatisfiable, and $`C`$ is unsatisfiable iff $`CA\neg A`$ for some concept name $`A`$. Moreover, a concept $`C`$ is satisfiable iff the Abox $`\{a:C\}`$ is consistent. It is straightforward to extend these reductions to role hierarchies, but terminologies deserve special care: In , the *internalisation* of GCIs is introduced, a technique that reduces reasoning w.r.t. a (possibly cyclic) terminology to reasoning w.r.t. the empty terminology. For $`𝒮𝒬`$, this reduction must be slightly modified. The following Lemma shows how general concept inclusion axioms can be *internalised* using a “universal” role $`U`$, that is, a transitive super-role of all roles occurring in $`𝒯`$ and their respective inverses.
###### Lemma 2.3
Let $`C,D`$ be concepts, $`𝒜`$ an Abox, $`𝒯`$ a terminology, and $``$ a role hierarchy. We define
$$C_𝒯:=\underset{C_iD_i𝒯}{\text{}}\neg C_iD_i.$$
Let $`U`$ be a transitive role that does not occur in $`𝒯`$, $`C`$, $`D`$, $`𝒜`$, or $``$. We set
$$_U:=\{RU,𝖨𝗇𝗏(R)UR\text{ occurs in }𝒯\text{}C\text{}D\text{}𝒜\text{, or }\}.$$
* $`C`$ is satisfiable w.r.t. $`𝒯`$ and $``$ iff $`CC_𝒯U.C_𝒯`$ is satisfiable w.r.t. $`_U`$.
* $`D`$ subsumes $`C`$ with respect to $`𝒯`$ and $``$ iff $`C\neg DC_𝒯U.C_𝒯`$ is unsatisfiable w.r.t. $`_U`$.
* $`𝒜`$ is consistent with respect to $``$ and $`𝒯`$ iff $`𝒜\{a:C_𝒯U.C_𝒯a\text{ occurs in }𝒜\}`$ is consistent w.r.t. $`_U`$.
The proof of Lemma 2.3 is similar to the ones that can be found in . Most importantly, it must be shown that, (a) if a $`𝒮𝒬`$-concept $`C`$ is satisfiable with respect to a terminology $`𝒯`$ and a role hierarchy $``$, then $`C,𝒯`$ have a *connected* model, i. e., a model where any two elements are connect by a role path over those roles occuring in $`C`$ and $`𝒯`$, and (b) if $`y`$ is reachable from $`x`$ via a role path (possibly involving inverse roles), then $`x,yU^{}`$. These are easy consequences of the semantics and the definition of $`U`$.
###### Theorem 2.4
Satisfiability and subsumption of $`𝒮𝒬`$-concepts w.r.t. terminologies and role hierarchies are polynomially reducible to (un)satisfiability of $`𝒮𝒬`$-concepts w.r.t. role hierarchies, and therefore to consistency of $`𝒮𝒬`$-Aboxes w.r.t. role hierarchies.
Consistency of $`𝒮𝒬`$-Aboxes w.r.t. terminologies and role hierarchies is polynomially reducible to consistency of $`𝒮𝒬`$-Aboxes w.r.t. role hierarchies.
## 3 A $`𝒮𝒬`$-Abox Tableau Algorithm
With Theorem 2.4, all standard inference problems for $`𝒮𝒬`$-concepts and Aboxes can be reduced to Abox-consistency w.r.t. a role hierarchy. In the following, we present a tableau-based algorithm that decides consistency of $`𝒮𝒬`$-Aboxes w.r.t. role hierarchies, and therefore all other $`𝒮𝒬`$ inference problems presented.
The algorithm tries to construct, for a $`𝒮𝒬`$-Abox $`𝒜`$, a tableau for $`𝒜`$, that is, an abstraction of a model of $`𝒜`$. Given the notion of a tableau, it is then quite straightforward to prove that the algorithm is a decision procedure for Abox consistency.
### 3.1 A Tableau for Aboxes
In the following, if not stated otherwise, $`C,D`$ denote $`𝒮𝒬`$-concepts, $``$ a role hierarchy, $`𝒜`$ an Abox, $`𝐑_𝒜`$ the set of roles occurring in $`𝒜`$ and $``$ together with their inverses, and $`𝐈_𝒜`$ is the set of individuals occurring in $`𝒜`$.
Without loss of generality, we assume all concepts $`C`$ occurring in assertions $`a:C𝒜`$ to be in NNF, that is, negation occurs in front of concept names only. Any $`𝒮𝒬`$-concept can easily be transformed into an equivalent one in NNF by pushing negations inwards using a combination of DeMorgan’s laws and the following equivalences:
$$\begin{array}{cccccc}\hfill \neg (R.C)& & (R.\neg C)\hfill & \hfill \neg (R.C)& & (R.\neg C)\hfill \\ \hfill \neg (nR.C)& & (n+1)R.C\hfill & \hfill \neg (nR.C)& & (n1)R.C\text{where}\hfill \\ & & & \hfill (1)R.C& :=& A\neg A\text{for some }A𝐂\hfill \end{array}$$
For a concept $`C`$ we will denote the NNF of $`\neg C`$ by $`\mathrm{}C`$. Next, for a concept $`C`$, $`\mathrm{𝖼𝗅𝗈𝗌}(C)`$ is the smallest set that contains $`C`$ and is closed under sub-concepts and $`\mathrm{}`$. We use $`\mathrm{𝖼𝗅𝗈𝗌}(𝒜):=_{a:C𝒜}\mathrm{𝖼𝗅𝗈𝗌}(C)`$ for the closure $`\mathrm{𝖼𝗅𝗈𝗌}(C)`$ of each concept $`C`$ occurring in $`𝒜`$. It is not hard to show that the size of $`\mathrm{𝖼𝗅𝗈𝗌}(𝒜)`$ is polynomial in the size of $`𝒜`$.
###### Definition 3.1
$`T=(𝐒,,,)`$ is a *tableau* for $`𝒜`$ w.r.t. $``$ iff
* $`𝐒`$ is a non-empty set,
* $`:𝐒2^{\mathrm{𝖼𝗅𝗈𝗌}(𝒜)}`$ maps each element in $`𝐒`$ to a set of concepts,
* $`:𝐑_𝒜2^{𝐒\times 𝐒}`$ maps each role to a set of pairs of elements in $`𝐒`$, and
* $`:𝐈_𝒜𝐒`$ maps individuals occurring in $`𝒜`$ to elements in $`𝐒`$.
Furthermore, for all $`s,t𝐒`$, $`C,C_1,C_2\mathrm{𝖼𝗅𝗈𝗌}(𝒜)`$, and $`R,S𝐑_𝒜`$, $`T`$ satisfies:
1. if $`C(s)`$, then $`\neg C(s)`$,
2. if $`C_1C_2(s)`$, then $`C_1(s)`$ and $`C_2(s)`$,
3. if $`C_1C_2(s)`$, then $`C_1(s)`$ or $`C_2(s)`$,
4. if $`S.C(s)`$ and $`s,t(S)`$, then $`C(t)`$,
5. if $`S.C(s)`$, then there is some $`t𝐒`$ such that $`s,t(S)`$ and $`C(t)`$,
6. if $`S.C(s)`$ and $`s,t(R)`$ for some $`R\text{*}S`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$, then $`R.C(t)`$,
7. $`x,y(R)`$ iff $`y,x(𝖨𝗇𝗏(R))`$,
8. if $`s,t(R)`$ and $`R\text{*}S`$, then $`s,t(S)`$,
9. if $`nS.C(s)`$, then $`\mathrm{}S^T(s,C)n`$,
10. if $`nS.C(s)`$, then $`\mathrm{}S^T(s,C)n`$,
11. if $`(nSC)(s)`$ and $`s,t(S)`$ then $`C(t)`$ or $`\mathrm{}C(t)`$,
12. if $`a:C𝒜`$, then $`C((a))`$,
13. if $`(a,b):R𝒜`$, then $`(a),(b)(R)`$,
14. if $`a\doteq ̸b𝒜`$, then $`(a)(b)`$,
where $``$ is a place-holder for both $``$ and $``$, and $`S^T(s,C):=\{t𝐒s,t(S)\text{and}C(t)\}`$.
###### Lemma 3.2
A $`𝒮𝒬`$-Abox $`𝒜`$ is consistent w.r.t. $``$ iff there exists a tableau for $`𝒜`$ w.r.t. $``$.
#### Proof:
For the *if* direction, if $`T=(𝐒,,,)`$ is a tableau for $`𝒜`$ w.r.t. $``$, a model $`=(\mathrm{\Delta }^{},^{})`$ of $`𝒜`$ and $``$ can be defined as follows:
$$\begin{array}{cccc}& \hfill \mathrm{\Delta }^{}& :=& 𝐒\hfill \\ \hfill \text{for concept names A in }\mathrm{𝖼𝗅𝗈𝗌}(𝒜):& \hfill A^{}& :=& \{sA(s)\}\hfill \\ \hfill \text{ for individual names }a𝐈:& \hfill a^{}& :=& (a)\hfill \\ \hfill \text{ for role names }R:& \hfill R^{}& :=& \{\begin{array}{cc}(R)^+\hfill & \text{if }𝖳𝗋𝖺𝗇𝗌(R)\hfill \\ (R)\underset{P\text{*}R,PR}{}P^{}\hfill & \text{otherwise}\hfill \end{array}\hfill \end{array}$$
where $`(R)^+`$ denotes the transitive closure of $`(R)`$. The interpretation of non-transitive roles is recursive in order to correctly interpret those non-transitive roles that have a transitive sub-role. From the definition of $`R^{}`$ and (P8), it follows that, if $`s,tS^{}`$, then either $`s,t(S)`$ or there exists a path $`s,s_1,s_1,s_2,\mathrm{},`$ $`s_n,t(R)`$ for some $`R`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$ and $`R\text{*}S`$.
Due to (P8) and by definition of $``$, we have that $``$ is a model of $``$.
To prove that $``$ is a model of $`𝒜`$, we show that $`C(s)`$ implies $`sC^{}`$ for any $`s𝐒`$. Together with (P12), (P13), and the interpretation of individuals and roles, this implies that $``$ satisfies each assertion in $`𝒜`$. This proof can be given by induction on the length $`C`$ of a concept $`C`$ in NNF, where we count neither negation nor integers in number restrictions. The only interesting case is $`C=S.E`$: let $`t𝐒`$ with $`s,tS^{}`$. There are two possibilities:
* $`s,t(S)`$. Then (P4) implies $`E(t)`$.
* $`s,t(S)`$. Then there exists a path $`s,s_1,s_1,s_2,\mathrm{},`$ $`s_n,t(R)`$ for some $`R`$ with $`𝖳𝗋𝖺𝗇𝗌(R)`$ and $`R\text{*}S`$. Then (P6) implies $`R.E(s_i)`$ for all $`1in`$, and (P4) implies $`E(t)`$.
In both cases, $`tE^{}`$ by induction and hence $`sC^{}`$.
For the converse, for $`=(\mathrm{\Delta }^{},^{})`$ a model of $`𝒜`$ w.r.t. $``$, we define a tableau $`T=(𝐒,,,)`$ for $`𝒜`$ and $``$ as follows:
$$𝐒:=\mathrm{\Delta }^{},(R):=R^{},(s):=\{C\mathrm{𝖼𝗅𝗈𝗌}(𝒜)sC^{}\},\text{ and }(a)=a^{}.$$
It is easy to demonstrate that $`T`$ is a tableau for $`D`$. ∎
### 3.2 The Tableau Algorithm
In this section, we present a completion algorithm that tries to construct, for an input Abox $`𝒜`$ and a role hierarchy $``$, a tableau for $`𝒜`$ w.r.t. $``$. We prove that this algorithm constructs a tableau for $`𝒜`$ and $``$ iff there exists a tableau for $`𝒜`$ and $``$, and thus decides consistency of $`𝒮𝒬`$ Aboxes w.r.t. role hierarchies.
Since Aboxes might involve several individuals with arbitrary role relationships between them, the completion algorithm works on a *forest* rather than on a *tree*, which is the basic data structure for those completion algorithms deciding satisfiability of a concept. Such a forest is a collection of trees whose root nodes correspond to the individuals present in the input Abox. In the presence of transitive roles, *blocking* is employed to ensure termination of the algorithm. In the additional presence of inverse roles, blocking is *dynamic*, i.e., blocked nodes (and their sub-branches) can be un-blocked and blocked again later. In the additional presence of number restrictions, *pairs* of nodes are blocked rather than single nodes.
###### Definition 3.3
A *completion forest* $``$ for a $`𝒮𝒬`$ Abox $`𝒜`$ is a collection of trees whose distinguished root nodes are possibly connected by edges in an arbitrary way. Moreover, each node $`x`$ is labelled with a set $`(x)\mathrm{𝖼𝗅𝗈𝗌}(𝒜)`$ and each edge $`x,y`$ is labelled with a set $`(x,y)_𝒜`$ of (possibly inverse) roles occurring in $`𝒜`$. Finally, completion forests come with an explicit inequality relation $`\doteq ̸`$ on nodes and an explicit equality relation $``$ which are implicitly assumed to be symmetric.
If nodes $`x`$ and $`y`$ are connected by an edge $`x,y`$ with $`R(x,y)`$ and $`R\text{*}S`$, then $`y`$ is called an $`S`$-*successor* of $`x`$ and $`x`$ is called an $`𝖨𝗇𝗏(S)`$-*predecessor* of $`y`$. If $`y`$ is an $`S`$-successor or an $`𝖨𝗇𝗏(S)`$-predecessor of $`x`$, then $`y`$ is called an $`S`$-neighbour of $`x`$. A node $`y`$ is a successor (resp. predecessor or neighbour) of $`y`$ if it is an $`S`$-successor (resp. $`S`$-predecessor or $`S`$-neighbour) of $`y`$ for some role $`S`$. Finally, *ancestor* is the transitive closure of *predecessor*.
For a role $`S`$, a concept $`C`$ and a node $`x`$ in $``$ we define $`S^{}(x,C)`$ by
$$S^{}(x,C):=\{yy\text{ is }S\text{-neighbour of }x\text{ and }C(y)\}.$$
A node is *blocked* iff it is not a root node and it is either directly or indirectly blocked. A node $`x`$ is *directly blocked* iff none of its ancestors are blocked, and it has ancestors $`x^{}`$, $`y`$ and $`y^{}`$ such that
1. $`y`$ is not a root node *and*
2. $`x`$ is a successor of $`x^{}`$ and $`y`$ is a successor of $`y^{}`$ *and*
3. $`(x)=(y)`$ and $`(x^{})=(y^{})`$ *and*
4. $`(x^{},x)=(y^{},y)`$.
In this case we will say that $`y`$ *blocks* $`x`$.
A node $`y`$ is *indirectly blocked* iff one of its ancestors is blocked, or it is a successor of a node $`x`$ and $`(x,y)=\mathrm{}`$; the latter condition avoids wasted expansions after an application of the $``$-rule.
Given a $`𝒮𝒬`$-Abox $`𝒜`$ and a role hierarchy $``$, the algorithm initialises a completion forest $`_𝒜`$ consisting only of root nodes. More precisely, $`_𝒜`$ contains a root node $`x_0^i`$ for each individual $`a_i𝐈_𝒜`$ occurring in $`𝒜`$, and an edge $`x_0^i,x_0^j`$ if $`𝒜`$ contains an assertion $`(a_i,a_j):R`$ for some $`R`$. The labels of these nodes and edges and the relations $`\doteq ̸`$ and $``$ are initialised as follows:
$$\begin{array}{ccc}\hfill (x_0^i)& :=& \{Ca_i:C𝒜\},\hfill \\ \hfill (x_0^i,x_0^j)& :=& \{R(a_i,a_j):R𝒜\},\hfill \\ \hfill x_0^i\doteq ̸x_0^j& \text{ iff }& a_i\doteq ̸a_j𝒜\text{, and}\hfill \end{array}$$
the $``$-relation is initialised to be empty. $`_𝒜`$ is then expanded by repeatedly applying the rules from Figure 1.
For a node $`x`$, $`(x)`$ is said to contain a *clash* if, for some concept name $`A𝐂`$, $`\{A,\neg A\}(x)`$, or if there is some concept $`nS.C(x)`$ and $`x`$ has $`n+1`$ $`S`$-neighbours $`y_0,\mathrm{},y_n`$ with $`C(y_i)`$ and $`y_i\doteq ̸y_j`$ for all $`0i<jn`$. A completion forest is *clash-free* if none of its nodes contains a clash, and it is *complete* if no rule from Figure 1 can be applied to it.
For a $`𝒮𝒬`$-Abox $`𝒜`$, the algorithm starts with the completion forest $`_𝒜`$. It applies the expansion rules in Figure 1, stopping when a clash occurs, and answers “$`𝒜`$ is consistent w.r.t. $``$” iff the completion rules can be applied in such a way that they yield a complete and clash-free completion forest, and “$`𝒜`$ and is inconsistent w.r.t. $``$” otherwise.
Since both the $``$-rule and the $`_r`$-rule are rather complicated, they deserve some more explanation. Both rules deal with the situation where a concept $`nR.C(x)`$ requires the identification of two $`R`$-neighbours $`y,z`$ of $`x`$ that contain $`C`$ in their labels. Of course, $`y`$ and $`z`$ may only be identified if $`y\doteq ̸z`$ is not asserted. If these conditions are met, then one of the two rules can be applied. The $``$-rule deals with the case where at least one of the nodes to be identified, namely $`y`$, is not a root node, and this can lead to one of two possible situations, both shown in Figure 2.
The upper situation occurs when both $`y`$ and $`z`$ are successors of $`x`$. In this case, we add the label of $`y`$ to that of $`z`$, and the label of the edge $`x,y`$ to the label of the edge $`x,z`$. Finally, $`z`$ inherits all inequalities from $`y`$, and $`(x,y)`$ is set to $`\mathrm{}`$, thus blocking $`y`$ and all its successors.
The second situation occurs when both $`y`$ and $`z`$ are neighbours of $`x`$, but $`z`$ is the predecessor of $`x`$. Again, $`(y)`$ is added to $`(z)`$, but in this case the inverse of $`(x,y)`$ is added to $`(z,x)`$, because the edge $`x,y`$ was pointing away from $`x`$ while $`z,x`$ points towards it. Again, $`z`$ inherits the inequalities from $`y`$ and $`(x,y)`$ is set to $`\mathrm{}`$.
The $`_r`$ rule handles the identification of two root nodes. An example of the whole procedure is given in the lower part of Figure 2. In this case, special care has to be taken to preserve the relations introduced into the completion forest due to role assertions in the Abox, and to memorise the identification of root nodes (this will be needed in order to construct a tableau from a complete and clash-free completion forest). The $`_r`$ rule includes some additional steps that deal with these issues. Firstly, as well as adding $`(y)`$ to $`(z)`$, the edges (and their respective labels) between $`y`$ and its neighbours are also added to $`z`$. Secondly, $`(y)`$ and all edges going from/to $`y`$ are removed from the forest. This will not lead to dangling trees, because all neighbours of $`y`$ became neighbours of $`z`$ in the previous step. Finally, the identification of $`y`$ and $`z`$ is recorded in the $``$ relation.
###### Lemma 3.4
Let $`𝒜`$ be a $`𝒮𝒬`$-Abox and $``$ a role hierarchy. The completion algorithm terminates when started for $`𝒜`$ and $``$.
#### Proof:
Let $`m=\mathrm{}\mathrm{𝖼𝗅𝗈𝗌}(𝒜)`$, $`n=|𝐑_𝒜|`$, and $`n_{\mathrm{max}}:=\mathrm{max}\{nnR.C\mathrm{𝖼𝗅𝗈𝗌}(𝒜)\}`$. Termination is a consequence of the following properties of the expansion rules:
1. The expansion rules never remove nodes from the forest. The only rules that remove elements from the labels of edges or nodes are the $``$\- and $`_r`$-rule, which sets them to $`\mathrm{}`$. If an edge label is set to $`\mathrm{}`$ by the $``$-rule, the node below this edge is blocked and will remain blocked forever. The $`_r`$-rule only sets the label of a root node $`x`$ to $`\mathrm{}`$, and after this, $`x`$’s label is never changed again since all edges to/from $`x`$ are removed. Since no root nodes are generated, this removal may only happen a finite number of times, and the new edges generated by the $`_r`$-rule guarantees that the resulting structure is still a completion forest.
2. Nodes are labelled with subsets of $`\mathrm{𝖼𝗅𝗈𝗌}(𝒜)`$ and edges with subsets of $`R_𝒜`$, so there are at most $`2^{2mn}`$ different possible labellings for a pair of nodes and an edge. Therefore, if a path $`p`$ is of length at least $`2^{2mn}`$, the pair-wise blocking condition implies the existence of two nodes $`x,y`$ on $`p`$ such that $`y`$ directly blocks $`y`$. Since a path on which nodes are blocked cannot become longer, paths are of length at most $`2^{2mn}`$.
3. Only the $``$\- or the $``$-rule generate new nodes, and each generation is triggered by a concept of the form $`R.C`$ or $`nR.C`$ in $`\mathrm{𝖼𝗅𝗈𝗌}(𝒜)`$. Each of these concepts triggers the generation of at most $`n_{\mathrm{max}}`$ successors $`y_i`$: note that if the $``$\- or the $`_r`$-rule subsequently causes $`(x,y_i)`$ to be changed to $`\mathrm{}`$, then $`x`$ will have some $`R`$-neighbour $`z`$ with $`(z)(y)`$. This, together with the definition of a clash, implies that the rule application which led to the generation of $`y_i`$ will not be repeated. Since $`\mathrm{𝖼𝗅𝗈𝗌}(𝒜)`$ contains a total of at most $`m`$ $`R.C`$, the out-degree of the forest is bounded by $`mn_{\mathrm{max}}n`$. ∎
###### Lemma 3.5
Let $`𝒜`$ be a $`𝒮𝒬`$-Abox and $``$ a role hierarchy. If the expansion rules can be applied to $`𝒜`$ and $``$ such that they yield a complete and clash-free completion forest, then $`𝒜`$ has a tableau w.r.t. $``$.
#### Proof:
Let $``$ be a complete and clash-free completion forest. The definition of a tableau $`T=(𝐒,,,)`$ from $``$ works as follows. Intuitively, an individual in $`𝐒`$ corresponds to a path in $``$ from some root node to some node that is not blocked, and which goes only via non-root nodes.
More precisely, a *path* is a sequence of pairs of nodes of $``$ of the form $`p=[\frac{x_0}{x_0^{}},\mathrm{},\frac{x_n}{x_n^{}}]`$. For such a path we define $`𝖳𝖺𝗂𝗅(p):=x_n`$ and $`𝖳𝖺𝗂𝗅^{}(p):=x_n^{}`$. With $`[p|\frac{x_{n+1}}{x_{n+1}^{}}]`$, we denote the path $`[\frac{x_0}{x_0^{}},\mathrm{},\frac{x_n}{x_n^{}},\frac{x_{n+1}}{x_{n+1}^{}}]`$. The set $`\mathrm{𝖯𝖺𝗍𝗁𝗌}()`$ is defined inductively as follows:
* For root nodes $`x_0^i`$ of $``$, $`[\frac{x_0^i}{x_0^i}]\mathrm{𝖯𝖺𝗍𝗁𝗌}()`$, and
* For a path $`p\mathrm{𝖯𝖺𝗍𝗁𝗌}()`$ and a node $`z`$ in $``$:
+ if $`z`$ is a successor of $`𝖳𝖺𝗂𝗅(p)`$ and $`z`$ is neither blocked nor a root node, then $`[p|\frac{z}{z}]\mathrm{𝖯𝖺𝗍𝗁𝗌}()`$, or
+ if, for some node $`y`$ in $``$, $`y`$ is a successor of $`𝖳𝖺𝗂𝗅(p)`$ and $`z`$ blocks $`y`$, then $`[p|\frac{z}{y}]\mathrm{𝖯𝖺𝗍𝗁𝗌}()`$.
Please note that, since root nodes are never blocked, nor are they blocking other nodes, the only place where they occur in a path is in the first place. Moreover, if $`p\mathrm{𝖯𝖺𝗍𝗁𝗌}()`$, then $`𝖳𝖺𝗂𝗅(p)`$ is not blocked, $`𝖳𝖺𝗂𝗅(p)=𝖳𝖺𝗂𝗅^{}(p)`$ iff $`𝖳𝖺𝗂𝗅^{}(p)`$ is not blocked, and $`(𝖳𝖺𝗂𝗅(p))=(𝖳𝖺𝗂𝗅^{}(p))`$.
We define a tableau $`T=(𝐒,,,)`$ as follows:
$$\begin{array}{ccc}\hfill 𝐒& =& \mathrm{𝖯𝖺𝗍𝗁𝗌}()\hfill \\ \hfill (p)& =& (𝖳𝖺𝗂𝗅(p))\hfill \\ \hfill (R)& =& \{p,[p|\frac{x}{x^{}}]𝐒\times 𝐒x^{}\text{ is an }R\text{-successor of }𝖳𝖺𝗂𝗅(p)\}\hfill \\ & & \{[q|\frac{x}{x^{}}],q𝐒\times 𝐒x^{}\text{ is an }𝖨𝗇𝗏(R)\text{-successor of }𝖳𝖺𝗂𝗅(q)\}\hfill \\ & & \{[\frac{x}{x}],[\frac{y}{y}]𝐒\times 𝐒x,y\text{ are root nodes, and }y\text{ is an }R\text{-neighbour of }x\}\hfill \\ \hfill (a_i)& =& \{\begin{array}{cc}[\frac{x_0^i}{x_0^i}]\hfill & \text{ if }x_0^i\text{ is a root node in }\text{ with }(x_0^i)\mathrm{}\hfill \\ [\frac{x_0^j}{x_0^j}]\hfill & \text{ if }(x_0^i)=\mathrm{},x_0^j\text{ a root node in }\text{ with }(x_0^j)\mathrm{}\text{ and }x_0^ix_0^j\hfill \end{array}\hfill \end{array}$$
Please note that $`(x)=\mathrm{}`$ implies that $`x`$ is a root node and that there is another root node $`y`$ with $`(y)\mathrm{}`$ and $`xy`$. We show that $`T`$ is a tableau for $`D`$.
* $`T`$ satisfies (P1) because $``$ is clash-free.
* (P2) and (P3) are satisfied by $`T`$ because $``$ is complete.
* For (P4), let $`p,q𝐒`$ with $`R.C(p)`$, $`p,q(R)`$. If $`q=[p|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`R`$-successor of $`𝖳𝖺𝗂𝗅(p)`$ and, due to completeness of $``$, $`C(x^{})=(x)=(q)`$. If $`p=[q|\frac{x}{x^{}}]`$, then $`x^{}`$ is an $`𝖨𝗇𝗏(R)`$-successor of $`𝖳𝖺𝗂𝗅(q)`$ and, due to completeness of $``$, $`C(𝖳𝖺𝗂𝗅(q))=(q)`$. If $`p=[\frac{x}{x}]`$ and $`q=[\frac{y}{y}]`$ for two root nodes $`x`$, $`x`$, then $`y`$ is an $`R`$-neighbour of $`x`$, and completeness of $``$ yields $`C(y)=(q)`$. (P6) and (P11) hold for similar reasons.
* For (P5), let $`R.C(p)`$ and $`𝖳𝖺𝗂𝗅(p)=x`$. Since $`x`$ is not blocked and $``$ complete, $`x`$ has some $`R`$-neighbour $`y`$ with $`C(y)`$.
+ If $`y`$ is a successor of $`x`$, then $`y`$ can either be a root node or not.
- If $`y`$ is not a root node: if $`y`$ is not blocked, then $`q:=[p|\frac{y}{y}]𝐒`$; if $`y`$ is blocked by some node $`z`$, then $`q:=[p|\frac{z}{y}]𝐒`$.
- If $`y`$ is a root node: since $`y`$ is a successor of $`x`$, $`x`$ is also a root node. This implies $`p=[\frac{x}{x}]`$ and $`q=[\frac{y}{y}]𝐒`$.
+ $`x`$ is an $`𝖨𝗇𝗏(R)`$-*successor* of $`y`$, then either
- $`p=[q|\frac{x}{x^{}}]`$ with $`𝖳𝖺𝗂𝗅(q)=y`$.
- $`p=[q|\frac{x}{x^{}}]`$ with $`𝖳𝖺𝗂𝗅(q)=uy`$. Since $`x`$ only has one predecessor, $`u`$ is not the predecessor of $`x`$. This implies $`xx^{}`$, $`x`$ blocks $`x^{}`$, and $`u`$ is the predecessor of $`x^{}`$ due to the construction of $`\mathrm{𝖯𝖺𝗍𝗁𝗌}`$. Together with the definition of the blocking condition, this implies $`(u,x^{})=(y,x)`$ as well as $`(u)=(y)`$ due to the blocking condition.
- $`p=[\frac{x}{x}]`$ with $`x`$ being a root node. Hence $`y`$ is also a root node and $`q=[\frac{y}{y}]`$.
In any of these cases, $`p,q(R)`$ and $`C(q)`$.
* (P7) holds because of the symmetric definition of the mapping $``$.
* (P8) is due to the definition of $`R`$-neighbours and $`R`$-successor.
* Suppose (P9) were not satisfied. Hence there is some $`p𝐒`$ with $`(nS.C)(p)`$ and $`\mathrm{}S^T(p,C)>n`$. We will show that this implies $`\mathrm{}S^{}(𝖳𝖺𝗂𝗅(p),C)>n`$, contradicting either clash-freeness or completeness of $``$. Let $`x:=𝖳𝖺𝗂𝗅(p)`$ and $`P:=S^T(p,C)`$. We distinguish two cases:
+ $`P`$ contains only paths of the form $`[p|\frac{y}{y^{}}]`$ and $`[\frac{x_0^i_{\mathrm{}}}{x_0^i_{\mathrm{}}}]`$. Then $`\mathrm{}P>n`$ is impossible since the function $`𝖳𝖺𝗂𝗅^{}`$ is injective on $`P`$: if we assume that there are two distinct paths $`q_1,q_2P`$ and $`𝖳𝖺𝗂𝗅^{}(q_1)=𝖳𝖺𝗂𝗅^{}(q_2)=y^{}`$, then this implies that each $`q_i`$ is of the form $`q_i=[p|\frac{y_i}{y^{}}]`$ or $`q_i=[\frac{y^{}}{y^{}}]`$. From $`q_1q_2`$, we have that $`q_i=[p|\frac{y_i}{y^{}}]`$ holds for some $`i\{1,2\}`$. Since root nodes occur only in the beginning of paths and $`q_1q_2`$, we have $`q_1=[p|(y_1,y^{})]`$ and $`q_2=[p|(y_2,y^{})]`$. If $`y^{}`$ is not blocked, then $`y_1=y^{}=y_2`$, contradicting $`q_1q_2`$. If $`y^{}`$ is blocked in $``$, then both $`y_1`$ and $`y_2`$ block $`y^{}`$, which implies $`y_1=y_2`$, again a contradiction. Hence $`𝖳𝖺𝗂𝗅^{}`$ is injective on $`P`$ and thus $`\mathrm{}P=\mathrm{}𝖳𝖺𝗂𝗅^{}(P)`$. Moreover, for each $`y^{}𝖳𝖺𝗂𝗅^{}(P)`$, $`y^{}`$ is an $`S`$-successor of $`x`$ and $`C(y^{})`$. This implies $`\mathrm{}S^{}(x,C)>n`$.
+ $`P`$ contains a path $`q`$ where $`p=[q|\frac{x}{x^{}}]`$. Obviously, $`P`$ may only contain one such path. As in the previous case, $`𝖳𝖺𝗂𝗅^{}`$ is an injective function on the set $`P^{}:=P\{q\}`$, each $`y^{}𝖳𝖺𝗂𝗅^{}(P^{})`$ is an $`S`$-successor of $`x`$, and $`C(y^{})`$ for each $`y^{}𝖳𝖺𝗂𝗅^{}(P^{})`$. Let $`z:=𝖳𝖺𝗂𝗅(q)`$. We distinguish two cases:
- $`x=x^{}`$. Hence $`x`$ is not blocked, and thus $`x`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`z`$. Since $`𝖳𝖺𝗂𝗅^{}(P^{})`$ contains only successors of $`x`$ we have that $`z𝖳𝖺𝗂𝗅^{}(P^{})`$ and, by construction, $`z`$ is an $`S`$-neighbour of $`x`$ with $`C(z)`$.
- $`xx^{}`$. This implies that $`x^{}`$ is blocked by $`x`$ and that $`x^{}`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`z`$. Due to the definition of pairwise-blocking this implies that $`x`$ is an $`𝖨𝗇𝗏(S)`$-successor of some node $`u`$ with $`(u)=(z)`$. Again, $`u𝖳𝖺𝗂𝗅^{}(P^{})`$ and, by construction, $`u`$ is an $`S`$-neighbour of $`x`$ and $`C(u)`$.
* For (P10), let $`(nS.C)(p)`$. Hence there are $`n`$ $`S`$-neighbours $`y_1,\mathrm{},y_n`$ of $`x=𝖳𝖺𝗂𝗅(p)`$ in $``$ with $`C(y_i)`$. For each $`y_i`$ there are three possibilities:
+ $`y_i`$ is an $`S`$-successor of $`x`$ and $`y_i`$ is not blocked in $``$. Then $`q_i:=[p|\frac{y_i}{y_i}]`$ or $`y_i`$ is a root node and $`q_i:=[\frac{y_i}{y_i}]`$ is in $`𝐒`$.
+ $`y_i`$ is an $`S`$-successor of $`x`$ and $`y_i`$ is blocked in $``$ by some node $`z`$. Then $`q_i=[p|\frac{z}{y_i}]`$ is in $`𝐒`$. Since the same $`z`$ may block several of the $`y_j`$s, it is indeed necessary to include $`y_i`$ explicitly into the path to make them distinct.
+ $`x`$ is an $`𝖨𝗇𝗏(S)`$-successor of $`y_i`$. There may be at most one such $`y_i`$ if $`x`$ is not a root node. Hence either $`p=[q_i|\frac{x}{x^{}}]`$ with $`𝖳𝖺𝗂𝗅(q_i)=y_i`$, or $`p=[\frac{x}{x}]`$ and $`q_i=[\frac{y_i}{y_i}]`$.
Hence for each $`y_i`$ there is a different path $`q_i`$ in $`𝐒`$ with $`S(p,q_i)`$ and $`C(q_i)`$, and thus $`\mathrm{}S^T(p,C)n`$.
* (P12) is due to the fact that, when the completion algorithm is started for an Abox $`𝒜`$, the initial completion forest $`_𝒜`$ contains, for each individual name $`a_i`$ occurring in $`𝒜`$, a root node $`x_0^i`$ with $`(x_0^i)=\{C\mathrm{𝖼𝗅𝗈𝗌}(𝒜)a_i:C𝒜\}.`$ The algorithm never blocks root individuals, and, for each root node $`x_0^i`$ whose label and edges are removed by the $`_r`$-rule, there is another root node $`x_0^j`$ with $`x_0^ix_0^j`$ and $`\{C\mathrm{𝖼𝗅𝗈𝗌}(𝒜)a_i:C𝒜\}(x_0^j)`$. Together with the definition of $``$, this yields (P12). (P13) is satisfied for similar reasons.
* (P14) is satisfied because the $`_r`$-rule does not identify two root nodes $`x_0^i,y_0^i`$ when $`x_0^i\doteq ̸y_0^i`$ holds. ∎
###### Lemma 3.6
Let $`𝒜`$ be a $`𝒮𝒬`$-Abox and $``$ a role hierarchy. If $`𝒜`$ has a tableau w.r.t. $``$, then the expansion rules can be applied to $`𝒜`$ and $``$ such that they yield a complete and clash-free completion forest.
#### Proof:
Let $`T=(𝐒,,,)`$ be a tableau for $`𝒜`$ and $``$. We use $`T`$ to trigger the application of the expansion rules such that they yield a completion forest $``$ that is both complete and clash-free. To this purpose, a function $`\pi `$ is used which maps the nodes of $``$ to elements of $`𝐒`$. The mapping $`\pi `$ is defined as follows:
* For individuals $`a_i`$ in $`𝒜`$, we define $`\pi (x_0^i):=(a_i)`$.
* If $`\pi (x)=s`$ is already defined, and a successor $`y`$ of $`x`$ was generated for $`R.C(x)`$, then $`\pi (y)=t`$ for some $`t𝐒`$ with $`C(t)`$ and $`s,t(R)`$.
* If $`\pi (x)=s`$ is already defined, and successors $`y_i`$ of $`x`$ were generated for $`nR.C(x)`$, then $`\pi (y_i)=t_i`$ for $`n`$ distinct $`t_i𝐒`$ with $`C(t_i)`$ and $`s,t_i(R)`$.
Obviously, the mapping for the initial completion forest for $`𝒜`$ and $``$ satisfies the following conditions:
$$\begin{array}{c}(x)(\pi (x)),\hfill \\ \text{if }y\text{ is an }S\text{-neighbour of }x\text{, then }\pi (x),\pi (y)(S)\text{, and }\hfill \\ x\doteq ̸y\text{ implies }\pi (x)\pi (y)\text{.}\hfill \end{array}\}$$
($``$)
It can be shown that the following claim holds:
Claim: Let $``$ be generated by the completion algorithm for $`𝒜`$ and $``$ and let $`\pi `$ satisfy $`()`$. If an expansion rule is applicable to $``$, then this rule can be applied such that it yields a completion forest $`^{}`$ and a (possibly extended) $`\pi `$ that satisfy $`()`$.
As a consequence of this claim, (P1), and (P9), if $`𝒜`$ and $``$ have a tableau, then the expansion rules can be applied to $`𝒜`$ and $``$ such that they yield a complete and clash-free completion forest. ∎
From Theorem 2.4, Lemma 3.2, 3.4 3.5, and 3.6, we thus have the following theorem:
###### Theorem 3.7
The completion algorithm is a decision procedure for the consistency of $`𝒮𝒬`$-Aboxes and the satisfiability and subumption of concepts with respect to role hierarchies and terminologies.
## 4 Conclusion
We have presented an algorithm for deciding the satisfiability of $`𝒮𝒬`$ KBs where the Abox may be non-empty and where the uniqueness of individual names is not assumed but can be asserted in the Abox. This algorithm is of particular interest as it can be used to decide the problem of conjunctive query containment w.r.t. a schema .
An implementation of the $`𝒮𝒬`$ Tbox satisfiability algorithm is already available in the FaCT system , and is able to reason efficiently with Tboxes derived from realistic ER schemas. This suggests that the algorithm presented here could form the basis of a practical decision procedure for the query containment problem. Work is already underway to test this conjecture by extending the FaCT system with an implementation of the new algorithm.
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# Estimating Hawking radiation for exotic black holes
## ACKOWLEDGEMENTS
Special Thanks to J. Koga, T. Torii and T. Tachizawa for useful discussions. T. T is thankful for financial support from the JSPS (No. 106613).
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# Geometry of River Networks II: Distributions of Component Size and Number
## I Introduction
Branching networks are an important category of all networks with river networks being a paradigmatic example. Probably as much as any other natural phenomena, river networks are a rich source of scaling laws rodriguez-iturbe97 ; rinaldo98 ; dodds2000pa . Central quantities such as drainage basin area and stream lengths are reported to closely obey power-law statistics rodriguez-iturbe97 ; rinaldo98 ; dodds2000pa ; langbein47 ; hack57 ; abrahams84 ; maritan96a ; dodds99pa . The origin of this scaling has been attributed to a variety of mechanisms including, among others: principles of optimality rodriguez-iturbe97 ; sun94 , self-organized criticality rinaldo93 , invasion percolation stark91 , and random fluctuations dodds2000pa ; leopold62 ; scheidegger67 ; manna92 . One of the difficulties in establishing any theory is that the reported values of scaling exponents show some variation abrahams84 ; maritan96a ; maritan96b .
With this variation in mind, we have in dodds2000ua extensively examined Hack’s law, the scaling relationship between basin shape and stream length. Such scaling laws are inherently broad-brushed in their descriptive content. In an effort to further improve comparisons between theory and data and, more importantly, between networks themselves, we consider here a generalization of Horton’s laws horton45 ; schumm56a . Defined fully in the following section, Horton’s laws describe how average values of network parameters change with a certain discrete renormalization of the network. The introduction of these laws by Horton may be seen as one of many examples that presaged the theory of fractal geometry mandelbrot83 . In essence, they express the relative frequency and size of network components such as stream segments and drainage basins.
Here, we extend Horton’s laws to functional relationships between probability distributions rather than simply average values. The recent work of Peckham and Gupta was the first to address this natural generalization of Horton’s laws peckham99 . Our work agrees with their findings but goes further to characterize the distributions and develop theoretical links between the distributions of several different parameters. We also present empirical studies that reveal underlying scaling functions with a focus on fluctuations and further consider deviations due to finite-size effects.
We examine continent-scale networks: the Mississippi, Amazon, Congo, Nile and Kansas river basins. As in dodds2000ua , we also examine Scheidegger’s model of directed, random networks scheidegger67 . Both real and model networks provide important tests and motivations for our generalizations of Horton’s laws.
We begin with definitions of stream ordering and Horton’s laws. Thereafter, the paper is divided into two main sections. In Section III, we first sketch the theoretical generalization of Horton’s laws. Estimates of the Horton ratios are carried out in Section IV and these provide basic parameters of the generalized laws. Empirical evidence from real continent-scale networks is then provided along with data from Scheidegger’s random network model in Section V. In Section VI we derive the higher order moments for stream length distributions and in Section VII, we consider deviations from Horton’s laws for large basins. In the Appendix A, we expand on some of the connections outlined in Section V, presenting a number of mathematical considerations on these generalized Horton distributions.
This paper is the second in a series of three on the geometry of river networks. In the first dodds2000ua we address issues of scaling and universality and provide further motivation for our general investigation. In the third article of the series dodds2000uc we extend the work of the present paper by examining how the detailed architecture of river networks, i.e., how network components fit together.
## II Stream ordering and Horton’s laws
Stream ordering was first introduced by Horton in an effort to quantify the features of river networks horton45 . The method was later improved by Strahler to give the present technique of Horton-Strahler stream ordering strahler57 . Stream ordering is a method applicable to any field where branching, hierarchical networks are important. Indeed, much use of stream ordering has been made outside of the context of river networks, a good example being the study of venous and arterial blood networks in biology zamir83 ; fung90 ; kassab93a ; kassab94a ; kassab94b ; turcotte98 ; aharinejad98 ; zamir99 . We describe two conceptions of the method and then discuss empirical laws defined with in the context of stream ordering.
A network’s constituent stream segments are ordered by an iterative pruning. An example of stream ordering for the Mississippi basin is shown in Figure 1. A source stream is defined as a section of stream that runs from a channel head to a junction with another stream (for an arboreal analogy, think of the leaves of a tree). These source streams are classified as the first order stream segments of the network. Next, remove these source streams and identify the new source streams of the remaining network. These are the second order stream segments. The process is repeated until one stream segment is left of order $`\mathrm{\Omega }`$. The order of the network is then defined to be $`\mathrm{\Omega }`$.
Once stream ordering on a network has been done, a number of natural quantities arise. These include $`n_\omega `$, the number of basins (or equivalently stream segments) for a given order $`\omega `$; $`l_\omega `$, the average main stream length; $`l_\omega ^{\text{ (s)}}`$, the average stream segment length; $`a_\omega `$, the average basin area; and the variation in these numbers from order to order. Horton horton45 and later Schumm schumm56a observed that the following ratios are generally independent of order $`\omega `$:
$$\frac{n_\omega }{n_{\omega +1}}=R_n,\frac{l_{\omega +1}}{l_\omega }=R_l,\text{and}\frac{a_{\omega +1}}{a_\omega }=R_a.$$
(1)
Since the main stream length averages $`\overline{l}_\omega `$ are combinations of stream segment lengths $`\overline{l}_\omega =_{\nu =1}^\omega \overline{l}_\omega ^{\text{ (s)}}`$ we have that the Horton ratio for stream segment lengths $`R_{l^{\text{(s)}}}`$ is equivalent to $`R_l`$. Because our theory will start with the distributions of $`l^{\text{(s)}}`$, we will generally use the ratio $`R_{l^{\text{(s)}}}`$ in place of $`R_l`$.
Horton’s laws have remained something of a mystery in geomorphology—the study of earth surface processes and form—due to their apparent robustness and hence perceived lack of physical (or geological) content. However, statements that Horton’s laws are “statistically inevitable” kirchner93 , while possibly true, have not yet been based on reasonable assumptions dodds2000pa . Furthermore, many other scaling laws can be shown to follow in part from Horton’s laws dodds99pa . Thus, Horton’s laws being without content would imply the same is true for those scaling laws that follow from them. Other sufficient assumptions include uniform drainage density (i.e., networks are space-filling) and self-affinity of single channels. The latter can be expressed as the relation maritan96a ; tarboton88 ; labarbera89 ; tarboton90
$$lL_{}^d,$$
(2)
where $`L_{}`$ is the longitudinal diameter of a basin. Scaling relations may be derived and the set of relevant scaling exponents can be reduced to just two: $`d`$ as given above and the ratio $`\mathrm{ln}R_{l^{\text{(s)}}}/\mathrm{ln}R_n`$ dodds99pa . Note that one obtains $`R_aR_n`$ so that only the two Horton ratios $`R_n`$ and $`R_{l^{\text{(s)}}}`$ are independent. Horton ratios are thus of central importance in the full theory of scaling for river networks.
## III Postulated form of Horton distributions
Horton’s laws relate quantities which are indexed by a discrete set of numbers, namely the stream orders. They also algebraically relate mean quantities such as $`\overline{a}_\omega `$. Hence we may consider a generalization to functional relationships between probability distributions. In other words, for stream lengths and drainage areas we can explore the relationships between probability distributions defined for each order.
Furthermore, as we have noted, Horton’s laws can be used to derive power laws of continuous variables such as the probability distributions of drainage area $`a`$ and main stream length $`l`$ maritan96a ; dodds99pa ; devries94 :
$$P(a)a^\tau \text{and}P(l)l^\gamma .$$
(3)
These derivations necessarily only give discrete points of power laws. In other words, the derivations give points as functions of the discrete stream order $`\omega `$ and are uniformly spaced logarithmically and we interpolate the power law from there. The distributions for stream lengths and areas must therefore have structures that when combined across orders produce smooth power laws.
For the example of the stream segment length $`l_\omega ^{\text{ (s)}}`$, Horton’s laws state that the mean $`\overline{l}_\omega ^{\text{ (s)}}`$ grows by a factor of $`R_{l^{\text{(s)}}}`$ with each integer step in order $`\omega `$. In considering $`P(l_\omega ^{\text{ (s)}},\omega )`$, the underlying probability distribution function for $`l_\omega ^{\text{ (s)}}`$, we postulate that Horton’s laws apply for every moment of the distribution and not just the mean. This generalization of Horton’s laws may be encapsulated in a statement about the distribution $`P(l_\omega ^{\text{ (s)}},\omega )`$ as
$$P(l_\omega ^{\text{ (s)}},\omega )=c_{l^{\text{(s)}}}(R_nR_{l^{\text{(s)}}})^\omega F_{l^{\text{(s)}}}(l_\omega ^{\text{ (s)}}R_{l^{\text{(s)}}}^\omega ).$$
(4)
The factor of $`(R_n)^\omega `$ indicates the that $`_{l^{\text{(s)}}=0}^{\mathrm{}}\text{d}l^{\text{(s)}}P(l_\omega ^{\text{ (s)}},\omega )(R_n)^\omega `$, i.e., the frequency of stream segments of order $`\omega `$ decays according to Horton’s law of stream number given in equation (1). Similarly, for $`l_\omega `$, $`a_\omega `$ and $`n_{\mathrm{\Omega },\omega }`$, we write
$$P(l_\omega ,\omega )=c_l(R_nR_{l^{\text{(s)}}})^\omega F_l(l_\omega R_{l^{\text{(s)}}}^\omega ),$$
(5)
$$P(a_\omega ,\omega )=c_a(R_n^2)^\omega F_a(a_\omega R_n^\omega ),$$
(6)
and
$$P(n_{\mathrm{\Omega },\omega })=c_n(R_n)^{\mathrm{\Omega }\omega }F_n(n_{\mathrm{\Omega },\omega }R_n^\omega ),$$
(7)
where constants $`c_{l^{\text{(s)}}}`$, $`c_l`$, $`c_a`$ and $`c_n`$ are appropriate normalizations. We have used the subscripted versions of the lengths and areas, $`l_\omega ^{\text{ (s)}}`$, $`l_\omega `$, and $`a_\omega `$, to reinforce that these parameters are for points at the outlets of order $`\omega `$ basins only. The quantity $`n_{\mathrm{\Omega },\omega }`$ is the number of streams of order $`\omega `$ within a basin of order $`\mathrm{\Omega }`$. This will help with some notational issues later on. The form of the distribution functions $`F_{l^{\text{(s)}}}`$, $`F_l`$, $`F_a`$ and $`F_n`$ and their interrelationships become the focus of our investigations. Since scaling is inherent in each of these postulated generalizations of Horton’s laws, we will often refer to these distribution functions as scaling functions.
We further postulate that distributions of stream segment lengths are best approximated by exponential distributions. Empirical evidence for this will be provided later on in Section V. The normalized scaling function $`F_{l^{\text{(s)}}}(u)`$ of equation (4) then has the form
$$F_{l^{\text{(s)}}}(u)=\frac{1}{\xi }e^{u/\xi }=F_{l^{\text{(s)}}}(u;\xi ),$$
(8)
where we have introduced a new length scale $`\xi `$ and stated its appearance with the notation $`F_{l^{\text{(s)}}}(u;\xi )`$. The value of $`\xi `$ is potentially network dependent. As we will show, distributions of main stream lengths, areas and stream number are all dependent on $`\xi `$ and this is the only additional parameter necessary for their description. Note that $`\xi `$ is both the mean and standard deviation of $`F_{l^{\text{(s)}}}(u;\xi )`$, i.e., for exponential distributions, fluctuations of a variable are on the order of its mean value. We may therefore think of $`\xi `$ as a fluctuation length scale. Note that the presence of exponential distributions indicates a randomness in the physical distribution of streams themselves and this is largely the topic of our third paper dodds2000uc .
Since main stream lengths are combinations of stream segment lengths, i.e. $`l_\omega =_{i=1}^\omega l_\omega ^{\text{ (s)}}`$, we have that the distributions of main stream lengths of order $`\omega `$ basins are approximated by convolutions of the stream segment length distributions. For this step, it is more appropriate to use conditional probabilities such as $`P(l_\omega ^{\text{ (s)}}|\omega )`$ where the basin order $`\omega `$ is taken to be fixed. We thus write
$$P(l_\omega |\omega )=P(l_{}^{\text{(s)}}{}_{1}{}^{}|1)P(l_{}^{\text{(s)}}{}_{2}{}^{}|2)\mathrm{}P(l_\omega ^{\text{ (s)}}|\omega ).$$
(9)
where $``$ denotes convolution. Details of the form obtained are given in Appendix A.1.
The next step takes us to the power law distribution for main stream lengths. Summing over all stream orders and integrating over $`u=l_\omega `$ we have
$$P(l)\underset{\omega =1}{\overset{\mathrm{}}{}}_{u=l}^{\mathrm{}}\text{d}uP(u,\omega ),$$
(10)
where we have returned to the joint probability for this calculation. The integral over $`u`$ is replaced by a sum when networks are considered on discrete lattices. Note that the probability of finding a main stream of length $`l`$ is independent of any sort of stream ordering since it is defined on an unordered network. The details of this calculation may be found in Appendix A.2 where it is shown that a power law $`P(l)l^\gamma `$ follows from the deduced form of the $`P(l_\omega ,\omega )`$ with $`\gamma =\mathrm{ln}R_n/\mathrm{ln}R_{l^{\text{(s)}}}`$.
## IV Estimation of Horton ratios
We now examine the usual Horton’s laws in order to estimate the Horton ratios. These ratios are seen as intrinsic parameters in the probability distribution functions given above in equations (4), (5), (6) and (7).
Figure 2(a) shows the stream order averages of $`l^{\text{(s)}}`$, $`l`$, $`a`$ and $`n`$ for the Mississippi basin. Deviations from exponential trends of Horton’s laws are evident and indicated by deviations from straight lines on the semi-logarithmic axis. Such deviations are to be expected for the smallest and largest orders within a basin dodds99pa ; dodds2000uc . For the smallest orders, the scale of the grid used becomes an issue but even with infinite resolution, the scaling of lengths, areas and number for low orders cannot all hold at the same time dodds99pa . For large orders, the decrease in sample space contributes to these fluctuations since the number of samples of order $`\omega `$ streams decays exponentially with order as $`(R_n)^{\mathrm{\Omega }\omega }`$. Furthermore, correlations with overall basin shape provide another source of deviations dodds2000uc . Nevertheless, in our theoretical investigations below we will presume exact scaling. Note also that the equivalence of $`R_n`$ and $`R_a`$ is supported by Figure 2(b) where the stream numbers $`n_w`$ have been inverted for comparison. Similar agreement is found for the Amazon and Nile as shown in Tables 1, 2, and 3 which we now discuss.
Table 1 shows the results of regression on the Mississippi data for various ranges of stream orders for stream number, area and lengths. Tables 2 and 3 show the same results carried out for the Amazon and Nile. Each table presents estimates of the four ratios $`R_n`$, $`R_a`$, $`R_l`$ and $`R_{l^{\text{(s)}}}`$. Also included are the comparisons $`R_a/R_n`$ and $`R_l/R_{l^{\text{(s)}}}`$, both of which we expect to be close to unity. For each quantity, we calculate the mean $`\mu `$, standard deviation $`\sigma `$ and normalized deviation $`\sigma /\mu `$.
Note the variation of exponents with choice of order range. This is the largest source of error in the calculation of the Horton ratios. Therefore, rather than taking a single range of stream orders for the regression, we examine a collection of ranges. Also, the deviations for high and low orders observed in Figures 2(a) and 2(b) do of course affect measurements of the Horton ratios. In all cases, we have avoided using data for the smallest and largest orders.
For the three example networks given here, the statements $`R_aR_n`$ and $`R_lR_{l^{\text{(s)}}}`$ are well supported. The majority of ranges give $`R_n/R_a`$ and $`R_l/R_{l^{\text{(s)}}}`$ very close to unity. The averages are also close to one and are different from unity mostly by within 1.0 and uniformly by within 1.5 standard deviations.
The normalized deviations, ie., $`\sigma /\mu `$, for the four ratios are all below $`0.05`$. No systematic ordering of the $`\sigma /\mu `$ is observed. Of all the data, the values for $`R_l`$ in the case of the Mississippi are the most notably uniform having $`\sigma /\mu =0.015`$. Throughout there is a slight trend for regression on lower orders to overestimate and on higher orders to underestimate the average ratios, while reasonable consistency is found at intermediate orders.
Thus, overall the ranges chosen in the tables give a reasonably even set of estimates of the Horton ratios and we will use these averages as our estimates of the ratios.
## V Empirical evidence for Horton distributions
### V.1 Stream segment length distributions
We now present Horton distributions for the Mississippi, Amazon, and Nile river basins as well as the Scheidegger model. Scheidegger networks may be thought of as collections of random-walker streams and are fully defined in dodds2000ua and extensively studied in dodds2000uc . The forms of all distributions are observed to be the same in the real data and in the model.
The first distribution is shown in Figure 3(a). This is the probability density function of $`l_{}^{\text{(s)}}{}_{4}{}^{}`$, fourth order stream segment lengths, for the Mississippi River. Distributions for different orders can be rescaled to show satisfactory agreement. This is done using the postulated Horton distribution of stream segment lengths given in equation (4). The rescaling is shown in Figure 3(b) and is for orders $`\omega =3,\mathrm{},6`$. Note the effect of the exponential decrease in number of samples with order is evident for $`\omega =6`$ since $`P(l_{}^{\text{(s)}}{}_{6}{}^{})`$ is considerably scattered. Nevertheless, the figure shows the form of these distributions to be most closely approximated by exponentials. We observe similar exponential distributions for the Amazon, the Nile and the Scheidegger model. The fluctuation length scale $`\xi `$ is found to be approximately $`800`$ meters for the Mississippi, $`1600`$ meters for the Amazon and $`1200`$ meters for the Nile.
Since $`\xi `$ is based on the definition of stream ordering, comparisons of $`\xi `$ are only sensible for networks that are measured on topographies with the same resolution. The above values of $`\xi `$ are approximate and our confidence in them would be improved with higher resolution data. Nevertheless they do suggest that fluctuations in network structure increase as we move from the Mississippi through to the Nile and then the Amazon.
### V.2 Main stream segment length distributions
The distributions of $`\omega =4`$ main stream lengths for the Amazon River is shown in Figure 4(a). Since main stream lengths are sums of stream segment lengths, their distribution has a single peak away from the origin. However, these distributions will not tend towards a Gaussian because the individual stream length distributions do not satisfy the requirements of the central limit theorem feller68I . This is because the moments of the stream segment length distributions grow exponentially with stream order. As the semi-logarithmic axes indicate, the tail may be reasonably well (but not exactly) modeled by exponentials. There is some variation in the distribution tails from region to region. For example, corresponding distributions for the Mississippi data do exhibit tails that are closer to exponentials. However, for the present work where we are attempting to characterize the basic forms of the Horton distributions, we consider these deviations to be of a higher order nature and belonging to the realm of further research.
In accord with equation (5), Figure 4(b) shows the rescaling of the main stream length distributions for $`\omega =3,\mathrm{},6`$. The ratios used, $`R_n=4.49`$ and $`R_l=2.19(R_{l^{\text{(s)}}}=2.17)`$ are taken from Table 2. Given the scatter of the distributions, it is unreasonable to perform minimization techniques on the rescaled data itself in order to estimate $`R_n`$ and $`R_l`$. This is best done by examining means, as we have done, and higher order moments which we discuss below. Furthermore, varying $`R_n`$ and $`R_l`$ from the above values by, say, $`\pm 0.05`$ does not greatly distort the visual quality of the “data collapse.”
Similar results for the Scheidegger model are shown in Figure 5. The Scheidegger model may be thought of as a network defined on a triangular lattice where at each lattice site one of two directions is chosen as the stream path dodds2000ua ; dodds2000uc . Figure 5(a) gives a single example distribution for main stream lengths of order $`\omega =6`$ basins. The tail is exponential as per the real world data. Figure 5(b) shows a collapse of main stream length distributions for orders $`\omega =4`$ through $`7`$. In contrast to the real data where an overall basin order is fixed ($`\mathrm{\Omega }`$), there is no maximum basin order here. The distributions in Figure 5(b) have an arbitrary normalization meaning the absolute values of the ordinate are also arbitrary. Otherwise, this is the same collapse as given in equation (5). For the Scheidegger model, our simulations yield $`R_n5.20`$ and $`R_{l^{\text{(s)}}}3.00`$ dodds99pa . For all distributions, we observe similar functional forms for real networks and the Scheidegger model, the only difference lying in parameters such as the Horton ratios.
### V.3 Drainage area distributions
Figure 6 shows more Horton distributions, this time for drainage area as calculated for the Nile river basin. In Figure 6, an example distribution for $`\omega =4`$ sub-basins is presented. The distribution is similar in form to those of main stream lengths of Figure 4, again showing a reasonably clear exponential tail. Rescaled drainage area distributions for $`\omega =3,\mathrm{},6`$ are presented in Figure 6(b). The rescaling now follows equation (6). Note that if $`R_n`$ and $`R_a`$ were not equivalent, the rescaling would be of the form
$$P(a_\omega ,\omega )=c_a(R_nR_a)^\omega F_a(a_\omega R_a^\omega ).$$
(11)
Since we have asserted that $`R_nR_a`$, equation (11) reduces to equation (6). The Horton ratio used here is $`R_n=4.42`$ which is in good agreement with $`R_a=4.53`$, the respective standard deviations being $`0.17`$ and $`0.10`$. Both figures are taken from the data of Table 3.
### V.4 Summing distributions to form power laws
As stated in Section III, the Horton distributions of $`a_\omega `$ and $`l_\omega `$ must combine to form power law distributions for $`a`$ and $`l`$ (see equations 3 and 10). Figure 7 provides empirical support for this observation for the example main stream lengths of the Mississippi network. The distributions for $`\omega =3`$, 4 and 5 main stream lengths are individually shown. Their combination together with the distribution of $`l_6`$ gives the reasonable approximation of a power law as shown. The area distributions combine in the same way. Note that the distributions do not greatly overlap. Each point of the power law is therefore the addition of significant contributions from only two or three of the separate distributions. The challenge here then is to understand how rescaled versions of $`F_l`$, being the basic form of the $`P(l_\omega ,\omega )`$, fit together in such a clean fashion. The details of this connection are established in Appendix A.2.
### V.5 Connecting distributions of number and area
In considering the generalized Horton distributions for number and area, we observe two main points: a calculation in the vein of what we are able to do for main stream lengths is difficult; and, the Horton distributions for area and number are equivalent.
In principle, Horton area distributions may be derived from stream segment length distributions. This follows from an assumption of statistically uniform drainage density which means that the typical drainage area drained per unit length of any stream is invariant in space. Apart from the possibility of changing with space which we will preclude by assumption, drainage density does naturally fluctuate as well dodds2000ua . Thus, we can write $`a\rho _\omega l_\omega ^{\text{ (s)}}`$ where the sum is over all orders and all stream segments and $`\rho `$ is the average drainage density.
However, we need to know for an example basin, how many instances of each stream segment occur as a function of order. For example, the number of first order streams in an order $`\mathrm{\Omega }`$ basins is $`n_{\mathrm{\Omega },1}`$. Given the distribution of this number, we can then calculate the distribution of the total contribution of drainage area due to first order streams. But the distributions of $`n_{\mathrm{\Omega },\omega }`$ are not independent so we cannot proceed in this direction.
We could potentially use the typical number of order $`\omega `$ streams, $`(R_n)^{\mathrm{\Omega }\omega }`$. Then the distribution of total area drained due to order $`\omega `$ streams would approach Gaussian because the individual distribution are identical and the central limit theorem would apply. However, because the fluctuations in total number of stream segments are so great, we lose too much information with this approach. Indeed, the distribution of area drained by order $`\omega `$ stream segments in a basin reflects variations in their number rather than length. Again, we meet up with the problem of the numbers of distinct orders of stream segment lengths being dependent.
One final way would be to use Tokunaga’s law dodds99pa ; tokunaga66 ; tokunaga78 ; tokunaga84 ; peckham95 . Tokunaga’s law states that the number of order $`\nu `$ side branches along an (absorbing) stream segment of order $`\mu `$ is given by
$$T_k=T_1(R_{l^{\text{(s)}}})^{k+1}.$$
(12)
where $`k=\mu \nu `$. The parameter $`T_1`$ is the average number of side streams having order $`\nu =\mu 1`$ for every order $`\mu `$ absorbing stream. This gives a picture of how a network fits together and may be seen to be equivalent to Horton’s laws dodds99pa . Now, even though we also understand the distributions underlying Tokunaga’s law dodds2000ua , similar technical problems arise. On descending into a network, we find the number of stream segments at each level to be dependent on all of the above.
Nevertheless, we can understand the relationship between the distributions for area and number. What follows is a generalization of the finding that $`R_nR_a`$. The postulated forms for these distributions were given in equations (6) and (7). Consider $`n_{\mathrm{\Omega },1}`$, the number of first order streams in an order $`\mathrm{\Omega }`$ basin. Assuming that, on average, first order streams are distributed evenly throughout a network, then this number is simply proportional to $`a_\mathrm{\Omega }`$. As an example, Figure 8 shows data obtained for the Scheidegger model. For the Scheidegger model, first order streams are initiated with a $`1/4`$ probability when the flow at the two upstream sites is randomly directed away, each with probability $`1/2`$. Thus, for an area $`a_\mathrm{\Omega }`$, we expect and find $`n_{\mathrm{\Omega },\omega }=a_\mathrm{\Omega }/4`$.
For higher internal orders, we can apply a simple renormalization. Assuming a system with exact scaling, the number of streams $`n_{\mathrm{\Omega },\omega }`$ is statistically equivalent to $`n_{\mathrm{\Omega }\omega +1,1}`$. Since the latter is proportional to $`a_{\mathrm{\Omega }\omega +1}`$ we have that
$$n_{\mathrm{\Omega },\omega }\rho _\omega a_{\mathrm{\Omega }\omega +1}$$
(13)
where the constant of proportionality is the density of order $`\omega `$ streams, Clearly, this equivalence improves as number increases, i.e., the difference $`\mathrm{\Omega }\omega `$ increases.
While we do not have exact forms for the area or number distributions, we note that they are similar to the main stream length distributions. Since source streams are linear basins with the width of a grid cell, the distribution of $`a_1`$ is the same as the distribution of $`l_1`$ and $`l_{}^{\text{(s)}}{}_{1}{}^{}`$, a pure exponential. Hence, $`n_{\mathrm{\Omega },\mathrm{\Omega }1}`$ is also an exponential. For increasing $`\omega `$, the distribution of $`a_\omega `$ becomes single peaked with an exponential tail, qualitatively the same as the main stream length distributions.
## VI Higher order moments
Finally, we discuss the higher order moments for the generalized Horton distributions. Figure 9 presents moments for distributions of main stream lengths for the case of the Mississippi. These moments are calculated directly from the main stream length distributions. A regular logarithmic spacing is apparent in moments for orders ranging from $`3`$ to $`7`$.
To see whether or not this is expected, we detail a few small calculations concerning moments starting from the exponential form of stream segment lengths given in equation (8). As noted previously, for an exponential distribution, $`F_{l^{\text{(s)}}}(u)=\xi ^1e^{u/\xi }`$, the mean is simply $`u=\xi `$. In general, the $`q`$th moment of an exponential distribution is
$`u^q`$ $`=`$ $`{\displaystyle _{u=0}^{\mathrm{}}}{\displaystyle \frac{u^q}{\xi }}e^{u/\xi }\text{d}u`$ (14)
$`=`$ $`\xi ^q{\displaystyle _{x=0}^{\mathrm{}}}x^qe^x\text{d}x=q!\xi ^q.`$
Assuming scaling holds exactly for across all orders, the above is precisely $`(l_{}^{\text{(s)}}{}_{1}{}^{})^q`$. Note that $`(l_{}^{\text{(s)}}{}_{1}{}^{})^q=q!l_{}^{\text{(s)}}{}_{1}{}^{}^q`$. Since the characteristic length of order $`\omega `$ streams is $`(R_{l^{\text{(s)}}})^{\omega 1}`$, we therefore have
$$(l_\omega ^{\text{ (s)}})^q=q!\xi ^q(R_{l^{\text{(s)}}})^{(\omega 1)q}=q!l_\omega ^{\text{ (s)}}^q.$$
(15)
Since main stream lengths are sums of stream segment lengths, so are their respective moments. Hence,
$`(l_\omega )^q`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{\omega }{}}}(l_{}^{\text{(s)}}{}_{k}{}^{})^q,`$ (16)
$`=`$ $`{\displaystyle \underset{k=1}{\overset{\omega }{}}}q!\xi ^q(R_{l^{\text{(s)}}})^{(k1)q},`$
$`=`$ $`q!\xi ^q{\displaystyle \underset{k=1}{\overset{\omega }{}}}(R_{l^{\text{(s)}}})^{(k1)q},`$
$`=`$ $`q!\xi ^q{\displaystyle \frac{(R_{l^{\text{(s)}}})^{q\omega }1}{R_{l^{\text{(s)}}}1}}.`$
We can now determine the log-space separation of moments of main stream length. Using Stirling’s approximation gradshteyn65 that $`\mathrm{ln}n!(n+1/2)\mathrm{ln}nn`$ we have
$$\mathrm{ln}(l_\omega )^qq\left[\xi +(R_{l^{\text{(s)}}})^\omega +\mathrm{ln}q\right]+C,$$
(17)
where $`C`$ is a constant. The $`\mathrm{ln}q`$ term inside the square brackets in equation 17 creates small deviations from linearity for $`1\omega 15`$. Thus, in agreement with Figure 9, we expect approximately linear growth of moments in log-space.
## VII Limitations on the predictive power of Horton’s laws
In this last section, we briefly examine deviations from scaling within this generalized picture of Horton’s laws. The basic question is given an approximate scaling for quantities measured at intermediate stream orders, what can we say about the features of the overall basin?
As noted in the previous section, all moments of the generalized Horton distributions grow exponentially with order. Coupling this with the fact that $`n_\omega R_n^\omega `$, i.e., the number of samples of order $`\omega `$ basins decreases exponentially with $`\omega `$, we observe that a basin’s $`a`$ and $`l`$ will potentially differ greatly from values predicted by Horton’s laws.
To illustrate this, Figure 10 specifically shows the distributions $`P(l_3)`$ and $`P(l_4)`$ scaled up to give $`P(l_{11})`$ for the Congo river. The actual Congo’s length measured at this 1000 meter resolution is represented by the solid line and is around 57% of the distribution’s mean as indicated by the dashed line. Nevertheless, we see that the measured length is within a standard deviation of the predicted value.
In Table 4, we provide a comparison of predicted versus measured main stream lengths and areas for the basins studied here. The mean for the scaled up distributions overestimates the actual values in all cases except for the Nile. Also, apart from the Nile, all values are within a standard deviation of the predicted mean. The coefficients of variation, $`\sigma _a/\overline{a}_\mathrm{\Omega }`$ and $`\sigma _l/\overline{l}_\mathrm{\Omega }`$, all indicate that fluctuations are on the order of the expected values of stream lengths and areas.
Thus, we see that by using a probabilistic point of view, this generalized notion of Horton’s laws provides a way of discerning the strength of deviations about the expected mean. In general, stronger deviations would imply that geologic conditions play a more significant role in setting the structure of the network.
## VIII Conclusion
The objective of this work has been to explore the underlying distributions of river network quantities defined with stream ordering. We have shown that functional relationships generalize all cases of Horton’s laws. We have identified the basic forms of the distributions for stream segment lengths (exponential) and main stream lengths (convolutions of exponentials) and shown a link between number and area distributions. Data from the continent-scale networks of the Mississippi, Amazon, and Nile river basins as well as from Scheidegger’s model of directed random networks provide both agreement with and inspiration for the generalizations of Horton’s laws. Finally, we have identified a fluctuation length scale $`\xi `$ which is a reinterpretation of what was previously identified as only a mean value. We see the study of the generalized Horton distributions as integral to increasing our understanding of river network structure. We also suggest that practical network analysis be extended to measurements of distributions and the length scale $`\xi `$ with the aim of refining our ability to distinguish and compare network structure.
By taking account of fluctuations inherent in network scaling laws, we are able to see how measuring Horton’s laws on low-order networks is unavoidably problematic. Moreover, as we have observed, the measurement of the Horton ratios is in general a delicate operation suggesting that many previous measurements are not without error.
The theoretical understanding of the growth and evolution of river networks requires a more thorough approach to measurement and a concurrent improvement in the statistical description of river network geometry. The present consideration of a generalization of Horton’s laws is a necessary step in this process giving rise to stronger tests of both real and synthetic data. In the following paper dodds2000uc , we round out this expanded picture of network structure by consdering the spatial distribution of network components.
## Acknowledgements
The work was supported in part by NSF grant EAR-9706220 and the Department of Energy grant DE FG02-99ER 15004. The authors would like to express their gratitude to H. Cheng for enlightening and enabling discussions.
## Appendix A Analytic connections between stream length distributions
In this appendix we consider a series of analytic calculations. These concern the connections between the distributions of stream segment lengths $`l_\omega ^{\text{ (s)}}`$, ordered basin main stream lengths $`l_\omega `$ and main stream lengths $`l`$. We will idealize the problem in places, assuming perfect scaling and infinite networks while making an occasional salubrious approximation. Also, we will treat the problem of lengths fully noting that derivations of distributions for areas follow similar but more complicated lines.
We begin by rescaling the form of stream segment length distributions
$$P(l_\omega ^{\text{ (s)}},\omega )=(R_n1)(R_nR_{l^{\text{(s)}}})^\omega F_{l^{\text{(s)}}}(lR_{l^{\text{(s)}}}^\omega ).$$
(18)
The normalization $`c_{l^{\text{(s)}}}=R_n1`$ stems from the requirement that
$$_{u=0}^{\mathrm{}}F_{l^{\text{(s)}}}(u)=1,$$
(19)
which is made purely for aesthetic purposes. As we have suggested in equation (8) and demonstrated empirical support for, $`F_{l^{\text{(s)}}}(u)`$ is well approximated by the exponential distribution $`\xi ^1e^{u/\xi }`$. For low $`u`$ and also we have noted that deviations do of course occur but they are sufficiently insubstantial as to be negligible for a first order treatment of the problem.
### A.1 Distributions of main stream lengths as a function of stream order
We now derive a form for the distribution of main stream lengths $`P(l_\omega |\omega )`$. As we have discussed, since $`l_\omega =_{i=1}^\omega l_\omega ^{\text{ (s)}}`$, we have the convolution (9). The right-hand side of equation (9) consists of exponentials as per equation (8) so we now consider the function $`K_\omega (u;\stackrel{}{a})`$ given by
$$K_\omega (u;\stackrel{}{a})=a_1e^{a_1u}a_2e^{a_2u}\mathrm{}a_\omega e^{a_\omega u},$$
(20)
where $`\stackrel{}{a}=(a_1,a_2,\mathrm{},a_\omega )`$. We are specifically interested in the case when no two of the $`a_i`$ are equal, i.e., $`a_ia_j`$ for all $`ij`$. To compute this $`\omega `$-fold convolution, we simply examine the $`K_\omega (u;\stackrel{}{a})`$ for $`\omega =2`$ and $`\omega =3`$ and identify the emerging pattern. For $`\stackrel{}{a}=(a_1,a_2)`$ we have, omitting the prefactors for the time being,
$`e^{a_1u}e^{a_2u}`$ (21)
$`=`$ $`{\displaystyle \frac{e^{a_1u}e^{a_2u}}{a_1a_2}}={\displaystyle \frac{e^{a_1u}}{a_1a_2}}+{\displaystyle \frac{e^{a_2u}}{a_2a_1}}`$
providing $`a_1a_2`$. Convolving this with $`e^{a_3u}`$ we obtain
$`e^{a_1u}e^{a_2u}e^{a_3u}=\left({\displaystyle \frac{e^{a_1u}e^{a_2u}}{a_1a_2}}\right)e^{a_3u}`$ , (22)
$`=`$ $`{\displaystyle \frac{e^{a_1u}e^{a_3u}}{(a_1a_2)(a_1a_3)}}{\displaystyle \frac{e^{a_2u}e^{a_3u}}{(a_1a_2)(a_2a_3)}},`$
$`=`$ $`{\displaystyle \frac{e^{a_1u}}{(a_1a_2)(a_1a_3)}}+`$
$`{\displaystyle \frac{e^{a_2u}}{(a_2a_1)(a_2a_3)}}+{\displaystyle \frac{e^{a_3u}}{(a_3a_1)(a_3a_2)}}.`$
Generalizing from this point, we obtain
$$K_\omega (u;\stackrel{}{a})=\left(\underset{i=1}{\overset{\omega }{}}a_i\right)\underset{i=1}{\overset{\omega }{}}\frac{e^{a_iu}}{_{j=1,ji}^\omega (a_ia_j)}.$$
(23)
Now, setting $`a_i=1/(\xi (R_{l^{\text{(s)}}})^{i1})`$ and carrying out some manipulations we obtain the following expression for $`P(l_\omega ,\omega )`$:
$`P(l_\omega ,\omega )={\displaystyle \frac{1}{(R_n)^\omega }}{\displaystyle \frac{1}{_{j=1}^\omega \xi (R_{l^{\text{(s)}}})^{i1}}}{\displaystyle \underset{i=1}{\overset{\omega }{}}}{\displaystyle \frac{e^{l_\omega /\xi (R_{l^{\text{(s)}}})^{i1}}}{_{j=1,ji}^\omega (1/\xi (R_{l^{\text{(s)}}})^{i1}1/\xi (R_{l^{\text{(s)}}})^{j1})}}`$ , (24)
$`=`$ $`{\displaystyle \frac{1}{(R_n)^\omega }}{\displaystyle \frac{1}{\xi ^\omega _{j=1}^\omega (R_{l^{\text{(s)}}})^{j1}}}{\displaystyle \underset{i=1}{\overset{\omega }{}}}e^{l_\omega /\xi (R_{l^{\text{(s)}}})^{i1}}\xi ^{\omega 1}{\displaystyle \frac{_{j=1,ji}^\omega (R_{l^{\text{(s)}}})^{i1}_{j=1,ji}^\omega (R_{l^{\text{(s)}}})^{j1}}{_{j=1,ji}^\omega (R_{l^{\text{(s)}}})^{j1}(R_{l^{\text{(s)}}})^{i1}}},`$
$`=`$ $`{\displaystyle \frac{1}{(R_n)^\omega }}{\displaystyle \frac{\xi ^{\omega 1}}{\xi ^\omega }}{\displaystyle \underset{i=1}{\overset{\omega }{}}}e^{l_\omega /\xi (R_{l^{\text{(s)}}})^{i1}}{\displaystyle \frac{(R_{l^{\text{(s)}}})^{2(i1)}_{j=1}^\omega (R_{l^{\text{(s)}}})^{i1}_{j=1}^\omega (R_{l^{\text{(s)}}})^{j1}_{k=1}^\omega (R_{l^{\text{(s)}}})^{(j1)}}{(R_{l^{\text{(s)}}})^{(\omega 1)}_{j=1,ji}^\omega (R_{l^{\text{(s)}}})^j(R_{l^{\text{(s)}}})^i}},`$
$`=`$ $`{\displaystyle \frac{1}{(R_n)^\omega }}{\displaystyle \frac{1}{\xi }}{\displaystyle \underset{i=1}{\overset{\omega }{}}}e^{l_\omega /\xi (R_{l^{\text{(s)}}})^{i1}}{\displaystyle \frac{(R_{l^{\text{(s)}}})^{(i1)(\omega 2)}(R_{l^{\text{(s)}}})^{\omega 2}/R_{l^{\text{(s)}}}}{_{j=1,ji}^\omega (R_{l^{\text{(s)}}})^j(R_{l^{\text{(s)}}})^i}},`$
$`=`$ $`{\displaystyle \frac{1}{(R_n)^\omega }}{\displaystyle \frac{1}{\xi R_{l^{\text{(s)}}}}}{\displaystyle \underset{i=1}{\overset{\omega }{}}}e^{l_\omega /\xi (R_{l^{\text{(s)}}})^{i1}}{\displaystyle \frac{(R_{l^{\text{(s)}}})^{i(\omega 2)}}{_{j=1,ji}^\omega (R_{l^{\text{(s)}}})^j(R_{l^{\text{(s)}}})^i}}`$
Note that we have added in a factor of $`1/(R_n)^\omega `$ for the appropriate normalization. In addition, one observes that $`P(0,\omega )=0`$ for all $`\omega >1`$ since all convolutions of pairs of exponentials vanish at the origin. Furthermore, the tail of the distribution is dominated by the exponential corresponding to the largest stream segment.
The next step is to connect to the power law distribution of main stream lengths, $`P(l)`$ (see Figure 7 and the accompanying discussion). On considering equation (10) we see that the problem can possibly be addressed with some form of asymptotic analysis.
Before attacking this calculation however, we will simplify the notation keeping only the important details of the $`P(l_\omega ,\omega )`$. Our main interest is to see how equation (10) gives rise to a power law. We transform the outcome of equation (24) by using $`n=\omega `$, $`u=l_\omega /\xi `$, $`r=R_{l^{\text{(s)}}}`$, and $`s=R_n`$, neglecting multiplicative constants and then summing over stream orders to obtain
$$G(u)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{s^n}\underset{i=1}{\overset{n}{}}\frac{r^{(n2)i}e^{u/r^{i1}}}{_{j=1,ji}^n(r^jr^i)}.$$
(25)
The integration over $`l_\omega `$ has been omitted meaning that the result will be a power law with one power lower than expected.
### A.2 Power law distributions of main stream lengths
We now show that this sum of exponentials $`G(u)`$ in equation (25) does in fact asymptotically tend to a power law. We first interchange the order of summation replacing $`_{n=1}^{\mathrm{}}_{i=1}^n`$ with $`_{i=1}^{\mathrm{}}_{n=i}^n`$ to give
$`G(u)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}e^{u/r^{i1}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{r^{(n2)i}}{s^n_{j=1,ji}^n(r^jr^i)}},`$ (26)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}C_ie^{u/r^{i1}}.`$
We thus simply have a sum of exponentials to contend with. The coefficients $`C_i`$ appear unwieldy at first but do yield a simple expression after some algebra which we now perform:
$`C_i={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{r^{(n2)i}}{s^n_{j=1,ji}^n(r^jr^i)}},`$ (27)
$`=`$ $`{\displaystyle \frac{1}{_{j=1}^{i1}(r^jr^i)}}{\displaystyle \underset{n=i}{\overset{\mathrm{}}{}}}{\displaystyle \frac{r^{(n2)i}}{s^n_{j=i+1}^n(r^jr^i)}},`$
$`=`$ $`{\displaystyle \frac{r^{(i2)i}}{_{j=1}^{i1}(r^jr^i)}}{\displaystyle \frac{1}{s^i}}{\displaystyle \underset{n=i}{\overset{\mathrm{}}{}}}{\displaystyle \frac{s^ir^{(n2)i}r^{(i2)i}}{s^n_{j=i+1}^n(r^jr^i)}},`$
$`=`$ $`{\displaystyle \frac{1}{_{j=1}^{i1}r^i(r^jr^i)}}{\displaystyle \frac{r^i}{s^i}}{\displaystyle \underset{n=i}{\overset{\mathrm{}}{}}}{\displaystyle \frac{r^{(ni)i}}{s^{ni}_{j=i+1}^n(r^jr^i)}},`$
$`=`$ $`{\displaystyle \frac{1}{_{j=1}^{i1}(r^{ji}1)}}{\displaystyle \frac{1}{r^is^i}}{\displaystyle \underset{n=i}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{_{j=i+1}^nsr^i(r^jr^i)}},`$
$`=`$ $`{\displaystyle \frac{1}{r^is^i}}{\displaystyle \frac{1}{_{j=1}^{i1}(r^{ji}1)}}{\displaystyle \underset{n=i}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=i+1}{\overset{n}{}}}{\displaystyle \frac{1}{s(r^{ji}1)}},`$
$`=`$ $`{\displaystyle \frac{1}{r^is^i}}{\displaystyle \frac{1}{_{j=1}^{i1}(r^{ji}1)}}{\displaystyle \underset{n=i}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=i+1}{\overset{n}{}}}{\displaystyle \frac{1}{s(r^{ji}1)}},`$
$`=`$ $`{\displaystyle \frac{1}{r^is^i}}\left({\displaystyle \frac{1}{_{k=1}^{i1}(1r^k)}}\right)\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \frac{1}{s(r^k1)}}\right).`$
In reaching the last line we have shifted the indices in several places. In the last bracketed term we have set $`k=ji`$ and then $`m=ni`$ while in the first bracketed term, we have used $`k=ji`$. Immediately of note is that the last term is independent of $`i`$ and may thus be ignored.
The first bracketed term does depend on $`i`$ but converges rapidly. Writing $`D_i=_{k=1}^{i1}(1r^k)`$ we have that $`D_i=D_m_{k=m}^{i1}(1r^k)`$. Taking $`m`$ to be fixed and large enough such that $`1r^k`$ is approximated well by $`\mathrm{exp}\{r^k\}`$ for $`km`$, we then have
$`D_i`$ $`=`$ $`D_m\mathrm{exp}\left\{{\displaystyle \underset{k=m}{\overset{i1}{}}}r^k\right\},`$ (28)
$`=`$ $`D_m\mathrm{exp}\left\{{\displaystyle \frac{r^{1m}}{(r1)}}(11/r^{im1})\right\}.`$
As $`i\mathrm{}`$, $`D_i`$ clearly approaches a product of $`D_m`$ and a constant. Therefore, the first bracketed term in equation (27) may also be neglected in an asymptotic analysis.
Hence, as $`i\mathrm{}`$, the coefficients $`C_i`$ are simply given by
$$C_i\frac{1}{s^ir_i}.$$
(29)
and we can approximate $`G(u)`$ as, boldly using the equality sign,
$$G(u)=AS(u)=A\underset{i=0}{\overset{\mathrm{}}{}}\frac{e^{u/r^i}}{r^{i(1+\gamma )}},$$
(30)
where $`A`$ comprises the constant part of the $`C_i`$ and factors picked up by shifting the lower limit of the index $`i`$ from $`1`$ to $`0`$. We have also used here the identification
$$s=r^\gamma .$$
(31)
We turn now to the asymptotic behavior of $`S(u)`$, this being the final stretch of our analysis
There are several directions one may take at this point. We will proceed by employing a transformation of $`S(u)`$ that is sometimes referred to as the Sommerfeld-Watson transformation and also as Watson’s lemma (carrier66, , p. 239). Given a sum over any set of integers $`I`$, say $`S=_{nI}f(n)`$, it can be written as the following integral
$$S=\frac{1}{2\pi i}_C\frac{\pi \mathrm{cos}\pi z}{\mathrm{sin}\pi z}f(z)𝑑z.$$
(32)
where $`C`$ is a contour that contains the points on the real axis $`n+i0`$ where $`nI`$ and none of the points of the same form with $`n/I`$. Calculation of the residues of the simple poles of the integrand return us to the original sum.
Applying the transformation to $`S(u)`$ we obtain
$$S(u)=\frac{1}{2\pi i}_C\frac{\pi \mathrm{cos}\pi z}{\mathrm{sin}\pi z}e^{ur^z}r^{z(1+\gamma )}\text{d}z.$$
(33)
The contour $`C`$ is represented in Figure 11.
We first make a change of variables, $`r^z=\rho `$. Substituting this and $`\text{d}z=\text{d}\rho /\rho \mathrm{ln}r`$ into equation (33) we have
$`S(u)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _C^{}}{\displaystyle \frac{\pi \mathrm{cos}\pi \mathrm{ln}\rho /\mathrm{ln}r}{\mathrm{sin}\pi \mathrm{ln}\rho /\mathrm{ln}r}}e^{u\rho }\rho ^{(1+\gamma )}(\text{d}\rho /\rho \mathrm{ln}r)`$ (34)
$`=`$ $`{\displaystyle \frac{1}{2i\mathrm{ln}r}}{\displaystyle _C^{}}{\displaystyle \frac{\pi \mathrm{cos}\pi \mathrm{ln}\rho /\mathrm{ln}r}{\mathrm{sin}\pi \mathrm{ln}\rho /\mathrm{ln}r}}e^{u\rho }\rho ^\gamma \text{d}\rho .`$
The transformed contour $`C^{}`$ is depicted in Figure 12.
As $`u\mathrm{}`$, the contribution to integral from the neighborhood of $`\rho =0`$ dominates. The introduction of the sin and cos terms has created an interesting oscillation that has to be handled with with some care. We now deform the integration contour $`C^{}`$ into the contour $`C^{\prime \prime }`$ of Figure 13 focusing on the interval along the imaginary axis $`[i,i]`$. Choosing this path will simplify the cos and sin expressions which at present have logs in their arguments.
The integral $`S(u)`$ is now given by $`S(u)I(u)+\text{c.c.}`$ where
$$I(u)=\frac{1}{2i\mathrm{ln}r}_0^i\frac{\pi \mathrm{cos}\pi \mathrm{ln}\rho /\mathrm{ln}r}{\mathrm{sin}\pi \mathrm{ln}\rho /\mathrm{ln}r}e^{u\rho }\rho ^\gamma \text{d}\rho .$$
(35)
Writing $`\rho =\sigma +i\tau `$ with $`\sigma =0`$, we have $`\text{d}\rho =i\text{d}\tau `$ and the following for the cos and sin terms:
$`\mathrm{cos}\pi \mathrm{ln}\rho /\mathrm{ln}r={\displaystyle \frac{\rho ^{i\pi /\mathrm{ln}r}+\rho ^{i\pi /\mathrm{ln}r}}{2}},`$ (36)
$`=`$ $`{\displaystyle \frac{\tau ^{i\pi /\mathrm{ln}r}e^{\pi ^2/2\mathrm{ln}r}+\tau ^{i\pi /\mathrm{ln}r}e^{\pi ^2/2\mathrm{ln}r}}{2}},`$
and
$`\mathrm{sin}\pi \mathrm{ln}\rho /\mathrm{ln}r={\displaystyle \frac{\rho ^{i\pi /\mathrm{ln}r}\rho ^{i\pi /\mathrm{ln}r}}{2i}},`$ (37)
$`=`$ $`{\displaystyle \frac{\tau ^{i\pi /\mathrm{ln}r}e^{\pi ^2/2\mathrm{ln}r}\tau ^{i\pi /\mathrm{ln}r}e^{\pi ^2/2\mathrm{ln}r}}{2i}}.`$
The $`\mathrm{cot}`$ term in the integrand becomes
$`{\displaystyle \frac{\mathrm{cos}\pi \mathrm{ln}\rho /\mathrm{ln}r}{\mathrm{sin}\pi \mathrm{ln}\rho /\mathrm{ln}r}}`$ $`=`$ $`i{\displaystyle \frac{1+\tau ^{2i\pi /\mathrm{ln}r}e^{\pi ^2/\mathrm{ln}r}}{1\tau ^{2i\pi /\mathrm{ln}r}e^{\pi ^2/\mathrm{ln}r}}}`$ (38)
$`=`$ $`i{\displaystyle \frac{1+\delta (\tau )}{1\delta (\tau )}},`$
where $`\delta (\tau )=\tau ^{2i\pi /\mathrm{ln}r}e^{\pi ^2/\mathrm{ln}r}`$. The integral $`I(u)`$ now becomes
$`I(u)`$ $`=`$ $`{\displaystyle \frac{i}{2\mathrm{ln}r}}{\displaystyle _0^1}{\displaystyle \frac{1+\delta (\tau )}{1\delta (\tau )}}e^{iu\tau }\tau ^\gamma e^{i\pi \gamma /2}\text{d}\tau `$ (39)
$`=`$ $`{\displaystyle \frac{e^{i\pi (1+\gamma )/2}}{2\mathrm{ln}r}}{\displaystyle _0^1}e^{iu\tau }\tau ^\gamma {\displaystyle \frac{1+\delta (\tau )}{1\delta (\tau )}}\text{d}\tau .`$
Now, since $`|\delta (\tau )|=e^{\pi ^2/\mathrm{ln}r}10^4`$ (taking $`r=R_{l^{\text{(s)}}}2.5`$), we can expand the expression as follows
$`{\displaystyle \frac{1+\delta }{1\delta }}`$ $`=`$ $`(1+\delta )(1+\delta +\delta ^2+\mathrm{})`$ (40)
$`=`$ $`1+2\delta +2\delta ^2+2\delta ^3+\mathrm{}`$
The integral in turn becomes
$`I(u)={\displaystyle \frac{i^{1+\gamma }}{2\mathrm{ln}r}}{\displaystyle _0^1}\text{d}\tau \tau ^\gamma e^{iu\tau }\times `$ (41)
$`(1+2\tau ^{2i\pi /\mathrm{ln}r}e^{\pi ^2/\mathrm{ln}r}+2\tau ^{4i\pi /\mathrm{ln}r}e^{2\pi ^2/\mathrm{ln}r}+\mathrm{}`$
$`+2\tau ^{2ni\pi /\mathrm{ln}r}e^{n\pi ^2/\mathrm{ln}r}+\mathrm{})`$
The basic $`n`$-th integral in this expansion is
$$I_n(u)=_0^1\tau ^{\gamma +2ni\pi /\mathrm{ln}r}e^{iu\tau }\text{d}\tau .$$
(42)
Substituting $`u\tau =w`$ and replacing the upper limit $`w=u`$ with $`w=\mathrm{}`$ we have
$`I_n(u)=u^{(1+\gamma +2ni\pi /\mathrm{ln}r)}{\displaystyle _0^{\mathrm{}}}\text{d}ww^{\gamma +2ni\pi /\mathrm{ln}r}e^{iw},`$ (43)
$`=`$ $`(iu)^{(1+\gamma +2ni\pi /\mathrm{ln}r)}{\displaystyle _0^{\mathrm{}}}i\text{d}w(iw)^{\gamma +2ni\pi /\mathrm{ln}r}e^{iw},`$
$`=`$ $`(iu)^{(1+\gamma +2ni\pi /\mathrm{ln}r)}{\displaystyle _0^{\mathrm{}}}\text{d}v(v)^{\gamma +2ni\pi /\mathrm{ln}r}e^v,`$
$`=`$ $`(iu)^{(1+\gamma +2ni\pi /\mathrm{ln}r)}\mathrm{\Gamma }(\gamma +2ni\pi /\mathrm{ln}r).`$
Here, we have rotated the contour along the imaginary $`iw`$-axis to the real $`v`$-axis and identified the integral with the gamma function $`\mathrm{\Gamma }`$ gradshteyn65 . The integral can now be expressed as
$$I(u)=\frac{1}{2\mathrm{ln}ru^{1+\gamma }}\left[1+2\underset{n=1}{\overset{\mathrm{}}{}}u^{2ni\pi /\mathrm{ln}r}\mathrm{\Gamma }(\gamma +2ni\pi /\mathrm{ln}r)\right].$$
(44)
We now need to show that the higher order terms are negligible. Note that their magnitudes do no vanish with increasing $`u`$ but instead are highly oscillatory terms. Using the asymptotic form of the Gamma function bender78
$$\mathrm{\Gamma }(z)=z^{z1/2}e^z\sqrt{2\pi }\left(1+O(1/z)\right),$$
(45)
we can estimate as follows for large $`n`$ that
$`|\mathrm{\Gamma }(1+\gamma +2ni\pi /\mathrm{ln}r)|`$ (46)
$``$ $`|(2i\pi n/\mathrm{ln}r+1+\gamma )^{2i\pi n/\mathrm{ln}r+1/2+\gamma }e^{\gamma 1}\sqrt{2\pi }|`$
$`=`$ $`|(e^{i\pi /2}2\pi n/\mathrm{ln}r)^{2i\pi n/\mathrm{ln}r+1/2+\gamma }e^{\gamma 1}\sqrt{2\pi }|`$
$`=`$ $`e^{\pi ^2n/\mathrm{ln}r}n^{\gamma +1/2}(2\pi /e)^{1+\gamma }(\mathrm{ln}r)^{1/2\gamma }.`$
Hence, $`\mathrm{\Gamma }(1+\gamma +2ni\pi /\mathrm{ln}r)`$ vanishes exponentially with $`n`$. For the first few values of $`n`$ taking $`\gamma =3/2`$ and $`r=2.5`$, we have $`\mathrm{\Gamma }(1+\gamma +2i\pi /\mathrm{ln}r)2.5\times 10^3`$ and $`\mathrm{\Gamma }(1+\gamma +4i\pi /\mathrm{ln}r)2.1\times 10^6`$ showing that these corrections are negligible.
Hence we are able estimate $`S(u)`$ to first order as
$$S(u)\frac{1}{\mathrm{ln}r}u^{1\gamma }.$$
(47)
Thus we have determined that a power law follows from the initial assumption that stream segment lengths follow exponential distributions.
This equivalence has been drawn as an asymptotic one, albeit one where convergences have been shown to be rapid. The calculation is clearly not the entire picture as the solution does contain small rapidly-oscillating corrections that do not vanish with increasing argument. A possible remaining problem and one for further investigation is to understand how the distributions for main stream lengths $`l_\omega `$ fit together over a range that is not to be considered asymptotic. Nevertheless, the preceding is one attempt at demonstrating this rather intriguing breakup of a smooth power law into a discrete family of functions built up from one fundamental scaling function.
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# SMALL x, SATURATION AND THE HIGH ENERGY LIMIT OF QCD
## 1 The Colored Glass Condensate
At very small values of the Bjorken $`x`$-variable, one expects QCD to enter a new regime which is caracterized by parton saturation and very high values of the QCD field strength $`F_{\mu \nu }^a1/g`$. Saturation, which is a limitation on the maximum phase-space parton density that can be reached in the hadron wavefunction, may have been already observed at HERA $`^\mathrm{?}`$, and should play an important role in the very early stages of relativistic heavy ion collisions at RHIC and LHC $`^\mathrm{?}`$ (and Refs. therein). In the saturation regime, the individual parton-parton interactions may be weak <sup>b</sup><sup>b</sup>bThis is the case if the saturation momentum $`Q_s`$ (cf. eq. (11) below) is large enough; e.g., at LHC, one expects $`Q_s23`$ GeV. (which we shall assume in what follows; i.e., we assume that $`g1`$), but the parton densities are so large that the system becomes strongly non-perturbative. Thus, at a theoretical level, understanding saturation is a challenging and fascinating problem where one has to deal with fully non-linear QCD. This is reminiscent of a similar problem in high temperature QCD where perturbation theory breaks down at the soft scale $`g^2T`$ because of the large thermal occupation numbers of the soft gluons $`^\mathrm{?}`$ (and Refs. therein).
The efforts toward understanding the region of high gluon density have uncovered a new form of matter which is formed from these gluons $`^\mathrm{?}`$. This matter is universal in that it should describe the high gluon density part of any hadron and nuclear wavefunction at small $`x`$. The combination of high density and the fact that gluons are massless bosons leads naturally to the expectation that this matter is a Bose condensate. Since the gluons carry color and local color is a gauge dependent quantity, any gauge invariant formulation will neccessarily involve an average over all colors to restore the invariance. This averaging procedure bears a formal ressemblance to the averaging over background fields done for spin glasses $`^\mathrm{?}`$. The matter is therefore called the Colored Glass Condensate (CGC).
It should be emphasized that this picture holds, strictly speaking, only in the infinite-momentum frame, where the hadron propagates almost at the speed of light, and thus appears as an infinitesimally thin two-dimensional sheet (by Lorentz contraction). In this frame, the parton interpretation makes sense and deep inelastic scattering (DIS) proceeds via the instantaneous absorbtion of the external probe (e.g., a virtual photon $`\gamma ^{}`$ with 4-momentum $`q^\mu `$) by some parton in the hadron. The Bjorken $`x_B`$ parameter is defined as $`x_BQ^2/2Pq`$, where $`Q^2q^\mu q_\mu `$, and $`P^\mu =\delta ^{\mu +}P^+`$, with large $`P^+`$, is the hadron 4-momentum <sup>c</sup><sup>c</sup>cWe use light-cone vector notations: for some arbitrary vector $`p^\mu `$, we write $`p^\mu =(p^+,p^{},𝐩_{})`$, with $`p^+(1/\sqrt{2})(p^0+p^3)`$, $`p^{}(1/\sqrt{2})(p^0p^3)`$, and $`𝐩_{}(p^1,p^2)`$.. By kinematics, $`x_B`$ coincides with the longitudinal momentum fraction $`xp^+/P^+`$ of the struck parton: $`x_B=x`$.
At $`x1`$, the gluon density increases faster, and is the driving force toward saturation <sup>d</sup><sup>d</sup>dTo directly measure the gluon density, it is convenient to consider a Gedanken experiment where the DIS is initiated by the “current” $`j\frac{1}{4}F_a^{\mu \nu }F_{\mu \nu }^a`$ which couples directly to gluons.. The dynamics that leads to this increase is the quantum evolution toward small-$`x`$ : in a parton cascade initiated by some fast parton (i.e., a hadron constituent with a relatively large longitudinal momentum $`k^+P^+`$, e.g., a valence quark), the number of radiated gluons increases exponentially with the rapidity gap $`\mathrm{\Delta }\tau \mathrm{ln}(k^+/p^+)\mathrm{ln}(1/x)`$ between the original parton and the soft (i.e., small $`p^+`$: $`p^+=xP^+P^+`$) final gluon which is struck by the external current (see Fig. 1). This (BFKL) picture—which assumes the radiated gluons to behave as free partons—ceases to be valid at very small $`x`$, where the gluon density is so large that the radiated gluons overlap each other in the transverse plane. This is the onset of saturation.
This is also the regime where the description in terms of a colored class condensate becomes appropriate: Corresponding to the strong ordering in longitudinal momenta along the cascade,
$`k^+k_0^+k_1^+k_2^+\mathrm{}k_N^+p^+,`$ (1)
there is a similar hierarchy among the lifetimes of the radiated gluons (since the latter are proportional to the respective $`k^+`$ momenta): Softer a gluon is, shorter is its lifetime.
This hierarchy in time stays at the basis of a quenched approximation $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ for the “fast” degrees of freedom: when “seen” by the soft modes with very short lifetimes, the modes with larger $`p^+`$ appear as frozen (no dynamics) and can be replaced by a static (i.e., independent of $`x^+`$) and random color charge configuration, with density $`\rho _a(x^{},x_{})`$. (This is random since the soft gluons can belong to different cascades.) The spatial correlations of the effective charge $`\rho _a(x^{},x_{})`$ reflect the quantum dynamics at large longitudinal momenta, and are encoded in a statistical weight function $`W[\rho ]`$. Because of the Lorentz contraction, we can write $`\rho _a(x^{},x_{})\delta (x^{})\rho _a(x_{})`$, and $`W`$ is a functional of the superficial charge density $`\rho _a(x_{})`$ alone.
Thus, in order to compute soft correlations in this approximation, one has to first study the (quantum) dynamics of the soft gluons in the presence of a given color charge $`\rho _a`$, and then perform a (classical) average over $`\rho _a`$, with weight function $`W[\rho ]`$. Clearly, the latter will depend upon what we call “soft” and “fast”, i.e., upon the separation scale $`\mathrm{\Lambda }`$ between fast ($`p^+>\mathrm{\Lambda }`$) and soft ($`p^+<\mathrm{\Lambda }`$) degrees of freedom. We thus write $`W[\rho ]W_\mathrm{\Lambda }[\rho ]`$.
For instance, the 2-point correlation function is obtained as (in the light-cone gauge $`A^+=0`$)
$$\mathrm{T}A^\mu (x)A^\nu (y)=𝒟\rho W_\mathrm{\Lambda }[\rho ]\left\{\frac{^\mathrm{\Lambda }𝒟AA^\mu (x)A^\nu (y)\mathrm{e}^{iS[A,\rho ]}}{^\mathrm{\Lambda }𝒟A\mathrm{e}^{iS[A,\rho ]}}\right\},$$
(2)
with the functional integral running only over soft gluon fields $`A_a^\mu (p)`$ with $`p^+<\mathrm{\Lambda }`$, and the action $`S[A,\rho ]`$ describing the dynamics of these fields in the presence of the classical source $`\rho _a`$, in the eikonal appoximation $`^\mathrm{?}`$: $`S=S_{YM}+S_W`$, where $`S_{YM}=d^4x(F_{\mu \nu }^2/4)`$ is the usual Yang-Mills action, and (with $`\stackrel{~}{x}(x^{},x_{})`$, and T denoting time-ordering of color matrices):
$$S_W\frac{i}{N_c}d^3\stackrel{~}{x}\mathrm{Tr}\left\{\rho (\stackrel{~}{x})\mathrm{T}\mathrm{exp}\left[ig𝑑x^+A^{}(x^+,\stackrel{~}{x})\right]\right\}.$$
(3)
The double averaging in eq. (2) is very similar to the one performed for spin systems in a random external magnetic field $`^\mathrm{?}`$ (which plays the role of $`\rho `$ in the above equation), and supports the physical picture of the saturation regime as a coloured glass condensate.
Of course, the final results for soft correlators must be independent of the arbitrary separation scale $`\mathrm{\Lambda }`$. The $`\mathrm{\Lambda }`$-dependence of the weight function $`W_\mathrm{\Lambda }[\rho ]`$ must cancel against the cutoff dependence of the quantum theory for the soft modes. This constraint can be formulated as a renormalization group equation for $`W_\mathrm{\Lambda }[\rho ]`$ $`^\mathrm{?}`$, to be presented in Sect. 3 below.
## 2 Saturation in the classical approximation
Consider first the simple approximation where the quantum path integral in eq. (2) is evaluated in the saddle-point (or classical) approximation $`\delta S/\delta A^\mu =0`$, and the weight function is taken simply as a Gaussian (this is the McLerran-Venugopalan model $`^\mathrm{?}`$, originally formulated for a large nucleus for which the Gaussian approximation is expected to work better) :
$$W_\mathrm{\Lambda }[\rho ]\mathrm{exp}\left\{\frac{1}{2\mu _\mathrm{\Lambda }^2}d^2x_{}\rho _a^2(x_{})\right\}.$$
(4)
Here, $`\mu _\mathrm{\Lambda }^2`$ is to the total color charge squared (per unit area) of the partons with $`p^+>\mathrm{\Lambda }`$.
The classical approximation of eq. (2) reads (with $`\stackrel{~}{x}(x^{},x_{})`$) :
$$A^\mu (x^+,\stackrel{~}{x})A^\nu (x^+,\stackrel{~}{y})_{cl}=𝒟\rho W_\mathrm{\Lambda }[\rho ]𝒜^\mu (\stackrel{~}{x})𝒜^\nu (\stackrel{~}{y}),$$
(5)
where $`𝒜^\mu (\stackrel{~}{x})`$ is the solution to the classical Yang-Mills equations with source $`\rho _a(\stackrel{~}{x})`$,
$$[D_\nu ,F^{\nu \mu }]=g\delta ^{\mu +}\delta (x^{})\rho _a(x_{}),$$
(6)
and is time-independent, like $`\rho _a`$ itself. The LC gauge solution can be written as $`^\mathrm{?}`$:
$`𝒜^+=𝒜^{}=0,𝒜^i(\stackrel{~}{x})=\theta (x^{})𝒜_+^i(x_{})+\theta (x^{})𝒜_{}^i(x_{}),`$ (7)
with $`𝒜_+^i(x_{})`$ and $`𝒜_{}^i(x_{})`$ related to $`\rho (x_{})`$ via a non-linear equation. Note that the vector potential $`𝒜^i(x^{},x_{})`$ is discontinuous at $`x^{}=0`$, so the associated electric field $`^{+i}^+𝒜^i`$ is localized at the light-cone (i.e., within the support of the source):
$$^{+i}(\stackrel{~}{x})=\delta (x^{})\left(𝒜_+^i(x_{})𝒜_{}^i(x_{})\right).$$
(8)
By using this approximation, let us compute the gluon distribution function, that is, the number of gluons per unit of $`x`$ in the hadron wavefunction having transverse momentum less than $`Q`$. This is defined as (with $`\stackrel{~}{k}^\mu (k^+,𝐤_{})`$ and $`k^+=xP^+`$) :
$$xG(x,Q^2)=\frac{1}{\pi }\frac{d^2k_{}}{(2\pi )^2}\mathrm{\Theta }(Q^2k_{}^2)F_a^{+i}(x^+,\stackrel{~}{k})F_a^{+i}(x^+,\stackrel{~}{k}).$$
(9)
In the classical approximation (8), $`F^{i+}(x^+,\stackrel{~}{k})^{i+}(k_{})`$ is independent of both $`x^+`$ and $`k^+`$, and (with $`\mathrm{\Delta }𝒜^i𝒜_+^i𝒜_{}^i`$, and the classical average defined as in eq. (5)) :
$`xG_{cl}(x,Q^2)`$ $`=`$ $`{\displaystyle \frac{A}{\pi }}{\displaystyle ^{Q^2}}{\displaystyle \frac{d^2k_{}}{(2\pi )^2}}{\displaystyle d^2x_{}\mathrm{e}^{ik_{}x_{}}\mathrm{\Delta }𝒜_a^i(0_{})\mathrm{\Delta }𝒜_a^i(x_{})_{cl}}.`$ (10)
Here, $`A`$ is the hadron transverse aria, and we have assumed transverse homogeneity. Thus, $`xG_{cl}(x,Q^2)`$ is independent of $`x`$, which reflects the absence of quantum evolution in the present, classical approximation. With the Gaussian weight function (4), and the non-linear classical solution (7), the classical gluon distribution (10) can be computed exactly $`^{\mathrm{?},\mathrm{?}}`$. One obtains:
$$\mathrm{\Delta }𝒜_a^i(0_{})\mathrm{\Delta }𝒜_a^i(x_{})_{cl}=\frac{N_c^21}{\pi \alpha _sN_c}\frac{1\mathrm{e}^{x_{}^2\mathrm{ln}(x_{}^2\mathrm{\Lambda }_{QCD}^2)Q_s^2/4}}{x_{}^2},$$
(11)
where $`Q_s\alpha _s\mu _\mathrm{\Lambda }`$ is the saturation momentum, and is a priori a function of $`\mathrm{\Lambda }`$. Remarkably, this equation displays saturation via non-linear effects in the classical solution. This interpretation can be made sharper by going to momentum space. If $`N(k_{})`$ is the Fourier transform of (11),
$$N(k_{})\alpha _s(Q_s^2/k_{}^2)\mathrm{for}k_{}^2Q_s^2,$$
(12)
which is the normal perturbative behaviour, but
$$N(k_{})\frac{1}{\alpha _s}\mathrm{ln}\frac{k_{}^2}{Q_s^2}\mathrm{for}k_{}^2Q_s^2,$$
(13)
which shows a much slower increase, i.e., saturation, at low momenta (with $`k_{}\mathrm{\Lambda }_{QCD}`$ though). According to eq. (11), saturation is also a statement about the maximum field intensity that can be reached in the system: the classical field never becomes larger than $`A^i1/g`$. This is the maximal occupation number for a classical field, since larger occupation numbers are blocked by repulsive interactions of the gluon field.
## 3 The non-linear evolution equation
In writing down the effective theory for soft gluons in eq. (2), we have assumed that the influence of the fast gluons can be reproduced by a classical color source $`\rho `$ with a peculiar structure (time-independent, and localized at $`x^{}=0`$), and some (still unspecified) weight function $`W_\mathrm{\Lambda }[\rho ]`$. In this section, we show how to construct this effective theory, step by step, by integrating quantum fluctuations in successive layers of $`p^+`$.
To this aim, it is convenient to consider a sequence of two effective theories (“Theory I” and “Theory II”) defined as in eq. (2), but with different separation scales: $`\mathrm{\Lambda }`$ in the case of Theory I, and $`b\mathrm{\Lambda }`$ for Theory II, with $`b1`$. That is, Theory II differs from Theory I in that the “semi-fast” fields with longitudinal momenta between $`b\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }`$ have been integrated out, and the associated correlations have been incorporated at tree-level, within the new weight function $`W_{b\mathrm{\Lambda }}`$. The difference $`\mathrm{\Delta }WW_{b\mathrm{\Lambda }}W_\mathrm{\Lambda }`$ can be obtained $`^\mathrm{?}`$ by matching calculations of soft ($`k^+b\mathrm{\Lambda }`$) gluon correlations in both theories. The result can be expressed as an evolution equation for $`W_\tau [\rho ]`$ (with $`\tau \mathrm{ln}(P^+/\mathrm{\Lambda })`$) with respect to variations in $`\tau `$.
The quantum corrections due to the semi-fast fields have to be computed to leading logarithmic accuracy (LLA), that is, to leading order in $`\alpha _s\mathrm{ln}(1/b)`$—indeed, it is only to this accuracy that the hierarchy of scales in eq. (1) is satisfied, and the matching is possible—, but to all orders in the classical fields and sources (since, in the saturation regime of interest, the non-linear effects are so strong that cannot be expanded in perturbation theory). This specifies the accuracy of the evolution equations to be obtained.
To this accuracy, there are only two contributions to $`\mathrm{\Delta }W`$ : the one- and two-point correlators of the fluctuating colour charge $`\delta \rho _a(x)`$ of the semi-fast gluons. Specifically, one obtains $`^{\mathrm{?},\mathrm{?}}`$
$`\delta \rho _a(x)_\rho `$ $``$ $`\alpha _s\mathrm{log}(1/b)\delta (x^{})\sigma _a(x_{}),`$
$`\delta \rho _a(x)\delta \rho _b(y)_\rho `$ $``$ $`\alpha _s\mathrm{log}(1/b)\delta (x^{})\chi _{ab}(x_{},y_{})\delta (y^{}),`$ (14)
while all the $`n`$-point correlators with $`n3`$ are of higher order in $`\alpha _s`$. In these equations, $`\mathrm{}_\rho `$ denotes a quantum expectation value over the semi-fast fields at fixed $`\rho `$. Furthermore, $`\sigma _a(x_{})`$ and $`\chi _{ab}(x_{},y_{})`$ are generally non-linear functionals of $`\rho (x_{})`$ given by one-loop diagrams within Theory I (with loop momenta restricted to the strip: $`b\mathrm{\Lambda }<|p^+|<\mathrm{\Lambda }`$).
In terms of these functions, the evolution equation for $`W_\tau [\rho ]`$ reads <sup>e</sup><sup>e</sup>eIn condensed notations where, e.g., $`\rho _x`$ stands for $`\rho _a(x_{})`$, and repeated indices are understood to be summed (integrated) over. $`^\mathrm{?}`$
$$\frac{W_\tau [\rho ]}{\tau }=\alpha _s\left\{\frac{1}{2}\frac{\delta ^2}{\delta \rho _x\delta \rho _y}[W_\tau \chi _{xy}]\frac{\delta }{\delta \rho _x}[W_\tau \sigma _x]\right\}.$$
(15)
This functional equation is equivalent to an infinite hierarchy of ordinary equations for the correlators of the charge density. For instance, by multiplying eq. (15) with $`\rho _x\rho _y`$ and functionally integrating over $`\rho `$, one obtains an evolution equation for the two-point function:
$$\frac{d}{d\tau }\rho _x\rho _y_\tau =\alpha _s\chi _{xy}+\rho _x\sigma _y+\sigma _x\rho _y_\tau ,$$
(16)
which in general, however, involves also the higher $`n`$-point functions, via $`\sigma `$ and $`\chi `$. But in the weak source approximation, i.e., with $`\sigma `$ and $`\chi `$ computed to lowest order in $`\rho `$, this becomes a closed equation for $`\rho \rho `$ which has been shown $`^\mathrm{?}`$ to be equivalent to the BFKL equation, as necessary on physical grounds. This is a highly non-trivial check of the effective theory in eq. (2).
In order to study saturation, however, one needs eq. (15) in the regime of strong background fields and sources ($`𝒜^i1/g`$ and $`\rho 1/g^2`$; cf. eqs. (11) and (6)), which requires for an exact calculation of the coefficients $`\sigma `$ and $`\chi `$. This has been done recently $`^\mathrm{?}`$, via a lenthy calculation which had to cope with difficulties related to gauge-fixing, the axial poles in the gluon propagator, and the proper definition of the singular limit $`\rho (x)\delta (x^{})\rho (x_{})`$. This makes possible to look for solutions to eq. (15). Note that, formally, this is like a functional Schrödinger equation in imaginary “time” $`\tau `$. An interesting possibility is that the “Hamiltonian” in its r.h.s. has an eigenfunction $`𝒲[\rho ]`$ of maximum eigenvalue $`\lambda `$. This would lead to an universal behavior of $`W_\tau [\rho ]`$ in the small-$`x`$, or high-energy, limit: $`W_\tau [\rho ]_\tau \mathrm{}\mathrm{e}^{\tau \lambda }𝒲[\rho ]`$, with $`\tau `$-independent $`𝒲[\rho ]`$.
## References
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# Comment on “On the Origin of the OZI Rule in QCD”, by N. Isgur and H. B. Thacker
## I Introduction
In their recent paper , Isgur and Thacker discuss an issue of paramount importance for mesonic spectroscopy, the nature of the large OZI-violating amplitude observed in the pseudoscalar nonet. The $`\pi \eta ^{}`$ splitting is the largest mass splitting among light mesons, and understanding its physical origin is clearly a key issue.
In order to clarify the physical origin of the large $`\eta ^{}`$ mass Isgur and Thacker studied its analogue in the scalar $`0^{++}`$, vector $`1^{}`$, and axialvector $`1^{++}`$ channels. They observe that the OZI-violating amplitude $`A^{0^{++}}`$ in the scalar channel is as large as the one in the pseudoscalar channel. In the first version of their paper, Isgur and Thacker claimed that the signs of $`A^{0^{++}}`$ and $`A^{0^+}`$ are the same. After being alerted to an error in their work by us and other, the second version concludes that the two amplitudes are opposite in sign. As we discuss below, this sign is crucial for phenomenology, and it is indeed what instantons predict it to be. Nevertheless, Isgur and Thacker still conclude that “… our result favors the large $`N_c`$ and not the instanton interpretation of the solution of the $`\eta ^{}`$ mass problem”. To us, this appears to be a significant misunderstanding, and we would like to clarify the issue in this comment.
We emphasize that the issue of the sign of the $`A^{0^{++}}`$ amplitude is also connected with an old misunderstanding concerning the scalar isoscalar (sigma) meson channel. Because the sigma resonance around 600 MeV is so broad and has not always been included in the Particle Data Table, it is sometimes assumed that the lightest state in this channel is located around 1.5 GeV or higher, and that the interaction must therefore be strongly repulsive. But we know that, fundamentally, the interaction in the $`0^{++}`$ channel must be very attractive, how else could chiral symmetry breaking take place? Using models of QCD it is hard to see how, after the vacuum is rearranged and chiral symmetry breaking has taken place, the $`\sigma `$ state could be pushed up to a mass much beyond 600 MeV. One might argue that there is a an attractive interaction in the $`0^{++}`$ channel, but that it respects the OZI rule. But in this case one is faced with the problem that no low mass strength is seen in the isovector channel, so the $`I=1`$ $`0^{++}`$ channel is indeed repulsive. This is clearly seen in lattice calculations of the scalar isovector correlation function.
We would also like to discuss a somewhat secondary point. Isgur and Thacker emphasize that their results are in agreement with $`1/N_c`$ arguments, even though the $`1/N_c`$ expansion makes no prediction concerning the sign of the amplitude. They also emphasize that $`1/N_c`$ arguments are generally incompatible with instantons. This is an issue on which there is a lot of confusion in the literature, and we shall comment on it below.
A general issue worth discussing in this note is the question of “conspiracies” among different hadronic amplitudes, which Isgur and Thacker mention in connection with the OZI rule in the vector channel. On this point we completely agree with Isgur and Thacker. We discuss other examples of similar conspiracies that have appeared in the study of hadronic correlation functions. Despite our criticism, we consider the paper by Isgur and Thacker to be a positive step. The problems related to OZI-violating amplitudes, chiral symmetry breaking, the $`U(1)_A`$ anomaly, and their relation to the foundations of the constituent quark model are rarely discussed in the literature on hadronic spectroscopy.
## II The sign of the OZI-violating amplitude: Scalars vs pseudoscalars
Let us first recall the setting. Isgur and Thacker consider QCD with two flavors and introduce $`4\times 4`$ mass matrices in the basis of $`u\overline{d},d\overline{u},u\overline{u},d\overline{d}`$ states of the form
$$𝐓=\left[\begin{array}{cccc}𝐒& \mathrm{𝟎}& \mathrm{𝟎}& \mathrm{𝟎}\\ \mathrm{𝟎}& 𝐒& \mathrm{𝟎}& \mathrm{𝟎}\\ \mathrm{𝟎}& \mathrm{𝟎}& 𝐒+𝐀& 𝐀\\ \mathrm{𝟎}& \mathrm{𝟎}& 𝐀& 𝐒+𝐀\end{array}\right].$$
(1)
Here, we ignore the effects of quark masses, and the eigenstates have to fall into isospin multiplets. It is then sufficient to consider the lower $`2\times 2`$ block of the mass matrix
$$𝐓\left[\begin{array}{cc}𝐃& 𝐀\\ 𝐀& 𝐃\end{array}\right].$$
(2)
where $`D=S+A`$ in the notation of Isgur and Thacker. So there are two amplitudes, $`D,A`$, and two different eigenstates in every $`J^{PC}`$ channel. In the pseudoscalar channel the physical states are the $`\eta ^{}(\overline{u}i\gamma _5u+\overline{d}i\gamma _5d)`$ and $`\pi ^0(\overline{u}i\gamma _5u\overline{d}i\gamma _5d)`$, with masses (or mass squared) $`DA`$ and $`D+A`$. We emphasize that both $`S`$ and $`A`$ depend on the quantum numbers of the current. In the case $`\mathrm{\Gamma }=i\gamma _5`$ there is no question about the sign of $`A`$: we want $`m_\eta ^{}>m_\pi `$ and thus $`A>0`$.
Similarly, there are two independent scalar channels, traditionally called $`\sigma (\overline{u}u+\overline{d}d)`$ and $`\delta ^0(\overline{u}u\overline{d}d)`$ ($`f_0`$ and $`a_0`$ in modern notation). The issue at hand is the sign of the corresponding amplitude $`A^{0++}`$.
Before we go into phenomenology, let us explain the instanton and perturbative QCD predictions. As discovered in the classical paper by t’Hooft the effect of instantons on fermionic correlation functions can be summarized in terms of an effective interaction
$$L=G\left[(\overline{\psi }\tau _a\psi )^2(\overline{\psi }\psi )^2(\overline{\psi }i\gamma _5\tau _a)^2+(\overline{\psi }i\gamma _5)^2\right],$$
(3)
where $`G`$ is an effective coupling (that depends on the instanton amplitude) and $`\tau `$ is an isospin matrix. We can directly read off the interaction in the channel characterized by the current $`\overline{\psi }\mathrm{\Gamma }\psi `$. The interaction is attractive in the pion $`\overline{\psi }i\gamma _5\stackrel{}{\tau }\psi `$ and sigma $`\overline{\psi }\psi `$ channels, and repulsive in the eta prime $`\overline{\psi }i\gamma _5\psi `$ and delta $`\overline{\psi }\stackrel{}{\tau }\psi `$ channels. So, to first order, the instanton-induced interaction corresponds to $`A^{0+}=A^{0++}`$.
It is worth repeating why this is so . The instanton interaction corresponds to the contribution of fermionic zero modes to the quark propagator. Since there is exactly one zero mode for every flavor, there are no diagonal $`(\overline{u}u)(\overline{u}u)`$ or $`(\overline{d}d)(\overline{d}d)`$ interactions. Second, since the fermion zero modes for quarks and anti-quarks have opposite chirality, the interaction is also off-diagonal in the basis spanned by right and left-handed fermions $`q_R,q_L`$
$$𝐓\left[\begin{array}{cc}\mathrm{𝟎}& 𝐀\\ 𝐀& \mathrm{𝟎}\end{array}\right].$$
(4)
This means that the instanton-induced amplitude follows very simple rules: (i) the sign flips in going from scalar to pseudoscalar, (ii) the sign flips in going from $`I=0`$ to $`I=1`$ states, (iii) to leading order there is no interaction in vector channels. This means that the OZI-violating amplitude in the vector channels receives no direct instanton contribution and is expected to be small. To summarize, we have the following prediction for the signs of the instanton contribution
$$\left[\begin{array}{cc}\eta ^{}& +\\ \pi & \\ \sigma & \\ \delta & +\end{array}\right].$$
(5)
So, to first order in the interaction, the $`\pi `$ and $`\sigma `$ form a light multiplet, and the $`\eta ^{}`$ and $`\delta `$ a heavy multiplet. The degeneracy is due to $`SU(2)`$ chiral symmetry, and the splitting between the multiplets is a manifestation of $`U(1)_A`$ violation. Of course, because the interaction in the scalar channel is so strong, the vacuum is rearranged and chiral symmetry is broken spontaneously. The pion becomes a Goldstone boson and is exactly massless in the chiral limit. The sigma corresponds to the massive excitation of the quark condensate and is pushed up to $`500600`$ MeV. But it is clear that the interaction at short distances remains attractive, so the sigma will always be lighter than the delta.
Short distance perturbative interactions cannot account for OZI-violation in the scalar channels. The reason is that chirality is strictly conserved in perturbative QCD. The only possible annihilation channels are $`\overline{q}_Lq_L`$ or $`\overline{q}_Rq_R`$, which do not contribute to the scalar or pseudoscalar channel. Perturbative effects do contribute to OZI violation in the vector channel. The OZI violating amplitudes in the vector channel (which are suppressed by more than an order of magnitude) then provide an estimate of the relative role of instanton/pQCD effects.
Phenomenologically, the situation in the scalar $`0^{++}`$ channel appears confused, because of the difficulties in classifying the observed scalar mesons. But what is important here is not whether some state is too broad to be considered a “true” resonance, or mixes strongly with $`\pi \pi `$ or $`\overline{K}K`$, etc. In fact it is pretty clear that there is strong evidence for substantial strength in the $`\sigma `$ channel at energies below 800 MeV. This strength is seen experimentally as the strong rise in the $`I=0`$ $`\pi \pi `$ phase shift and the “$`\sigma `$” meson of nuclear physics.
## III How to measure the OZI-violating amplitude
In order to measure the OZI-violating amplitudes we have to replace the mass matrices introduced above by objects more directly related to field theory. In particular, we shall consider correlation functions involving the scalar and pseudoscalar currents introduced above. The correlators are defined by
$`\mathrm{\Pi }_\pi (x,y)`$ $`=`$ $`\mathrm{Tr}\left(S(x,y)i\gamma _5S(y,x)i\gamma _5\right),`$ (6)
$`\mathrm{\Pi }_\delta (x,y)`$ $`=`$ $`\mathrm{Tr}\left(S(x,y)S(y,x)\right),`$ (7)
$`\mathrm{\Pi }_\eta ^{}(x,y)`$ $`=`$ $`\mathrm{Tr}\left(S(x,y)i\gamma _5S(y,x)i\gamma _5\right)`$ (9)
$`2\mathrm{Tr}\left(i\gamma _5S(x,x)\right)\mathrm{Tr}\left(i\gamma _5S(y,y)\right),`$
$`\mathrm{\Pi }_\sigma (x,y)`$ $`=`$ $`\mathrm{Tr}\left(S(x,y)S(y,x)\right)`$ (11)
$`2\mathrm{Tr}\left(S(x,x)\right)\mathrm{Tr}\left(S(y,y)\right),`$
where $`S(x,y)`$ is the fermion propagator and $`.`$ denotes the average over all gauge configurations. The OZI-violating difference between the $`\pi ,\eta ^{}`$ and $`\sigma ,\delta `$ channels is determined by “disconnected” (or double hairpin) contributions to the correlation functions. Note that it is important to correctly define the gamma matrices in the correlation functions in order to ensure positivity and the existence of a spectral representation. In particular, we have to use $`\mathrm{\Gamma }=i\gamma _5`$ in the pseudoscalar channel in order to guarantee $`\mathrm{\Pi }_\pi (x,y)>0`$. If this is taken care of, we can directly determine the sign of the OZI-violating amplitude from the sign of the disconnected correlation function. There is a subtlety in the scalar $`\sigma `$ channel, because we need to subtract the constant $`\overline{\psi }\psi ^2`$ term from the correlation function.
Isgur and Thacker find, in the revised version of their paper, that the disconnected correlation function is large and negative in the pseudoscalar channel, large and positive in the scalar channel, and very small in the vector channel. In Figs. 1 and 2 we compare these results to correlation functions obtained in unquenched instanton simulations . The correlation functions clearly show exactly the same pattern. There are some technical differences as compared to the work of Isgur and Thacker. We measure point-to-point rather than point-to-plane correlators, and the simulations are unquenched rather than quenched. This means that there is no need to extract amputated matrix elements, one can just measure the masses directly. We find a heavy $`\eta ^{}`$ and $`\delta `$, $`m_\delta m_\eta ^{}1`$ GeV, and a light sigma meson, $`m_\sigma 600`$ MeV.
## IV The large $`N_c`$ limit
The basis of the large $`N_c`$ approach is the assumption that $`N_c=3`$ QCD is similar to QCD in the limit $`N_c=\mathrm{}`$. In particular, it is assumed that there are no phase transitions as we go from $`N_c=3`$ to $`N_c\mathrm{}`$. Currently, the status of these assumptions is not clear, because not much is known about QCD($`N_c=\mathrm{}`$).
Let us start with what is firmly known. In the large $`N_c`$ limit gauge invariant quantities do not fluctuate. This means that all physical quantities can be obtained from a classical master field. For many years, there was no progress in obtaining the master field from QCD, except in the case of zero dimensions, or $`SU(N)`$ matrix models. Recently, Maldacena discovered a master field for $`N=4`$ supersymmetric QCD. This master field was described as a certain gravitational metric, together with a set of rules that relate gauge theory to supergravity observables. Amazingly, the same correlation functions can be obtained from an instanton master field . This configuration is a coherent superposition of many instantons with different colors orientations but the same size and position, held together by fermion exchanges. Progress was also made in understanding $`N=2`$ SUSY QCD. Seiberg and Witten determined the low energy effective action of this theory. In the semi-classical limit, their result can be expressed as a one-loop contribution plus an infinite series on $`n`$-instanton corrections, and nothing else. Witten and Seiberg’s result was generalized to arbitrary $`N_c`$ by Douglas and Shenker . They identified a special form of the large $`N_c`$ limit in which instantons (and monopoles) survive.
In practice the large $`N_c`$ expansion is used to determine the relative importance of certain classes of diagrams. In order to have a sensible large $`N_c`$ limit one requires $`g^2N_c=\mathrm{const}`$. We emphasize that the results are based on the usual perturbative expansion around the trivial vacuum. Instantons, and other non-perturbative effects, can spoil large $`N_c`$ counting rules. Applying large $`N_c`$ rules to hadronic correlation functions leads to the usual predictions
$`M_{genericmeson}`$ $``$ $`N_c^0M,`$ (12)
$`M_{genericbaryon}`$ $``$ $`N_c^1M,`$ (13)
$`M_\eta ^{}`$ $``$ $`N_c^1M,`$ (14)
where $`M`$ is some mass scale of order $`\mathrm{\Lambda }_{QCD}`$. This prediction is a little bit of an embarrassment, because in the real world the nucleon and the $`\eta ^{}`$ are very close in mass. Of course, one can argue that the numerical coefficient in the different channels could be very different. Nevertheless, the large mass scale that appears in the $`\eta ^{}`$ channel shows that large $`N_c`$ rules do not work equally well in all channels. This observation repeats itself in the scalar $`0^{++}`$ channel, where deviations from large $`N_c`$ counting are unusually charge. On the other hand, large $`N_c`$ arguments have proved to be useful in analyzing many properties of octet and decuplet baryons.
So what is the $`N_c`$ dependence of the instanton effects? In his well known paper , Witten argued that one has to choose between instantons and the large $`N_c`$ explanations of the $`U(1)_A`$ problem, because instanton effects scale as $`\mathrm{exp}(N_c)`$. However, twenty years later the dilemma does not seem that clear cut. New arguments have appeared (many by Witten himself), and instantons and the large $`N_c`$ limit may well be reconciled one day.
(i) The suppression of instantons in the large $`N_c`$ limit seems to follow from the instanton amplitude $`\mathrm{exp}(8\pi ^2/g^2)`$, together with the ’t Hooft scaling $`g^21/N_c`$. But this applies to small instantons only. For large instantons $`\rho \mathrm{\Lambda }^1`$ the action is $`O(1)`$, and there is no suppression. This is the scenario that was suggested for the $`CP^N`$ model: For small $`N`$ instantons are small and semi-classical, but for large $`N`$ instantons become strongly overlapping. The question of the instanton size distribution in QCD is a very non-trivial dynamical question. From lattice calculations we only know that the typical instanton size is about the same in $`N_c=2`$ and $`N_c=3`$ QCD, $`\rho 1/3`$ fm.
(ii) Collecting all factors of $`N_c`$ in the one-loop instanton amplitude one finds that they all exponentiate, giving $`dN/d\rho \mathrm{exp}(N_cF(\rho ))`$ . The function $`F(\rho )`$ has a non-trivial zero, so the instanton density may scale like a power of $`N_c`$, rather than an exponential, provided the instanton distribution becomes a delta function $`\delta (\rho \rho _0)`$ where $`\rho _0`$ is the zero of $`F(\rho )`$.
(iii) The second statement in Witten’s paper is that, even in the absence of instantons, the $`\eta ^{}`$ mass can be related to the topological susceptibility $`\chi _{top}`$ measured in pure gauge theory
$`{\displaystyle \frac{f_\pi ^2}{2N_f}}(m_\eta ^2+m_\eta ^{}^22m_K^2)=\chi _{top}.`$ (15)
This prediction has been checked on the lattice and it works quite well. After much effort, the lattice measurements of $`\chi _{top}`$ are stable, and it has become clear that the topological susceptibility is dominated by instantons, and not by some mysterious fluctuation.
(iv) As we emphasized above the large $`N_c`$ limit of both $`N=2`$ and $`N=4`$ SUSY QCD is consistent with the survival of instantons in this limit. The $`N=2`$ case is particularly interesting, because there are two different ways to take the large $`N_c`$ limit. In the “naive” limit, $`W`$ bosons have masses of $`O(1)`$, monopoles have masses $`O(N_c)`$, and instantons have action $`O(N_c)`$ and are not important. In this case, the one loop perturbative result is exact and nothing interesting happens. A different large $`N_c`$ limit is possible near the singular points on the space of vacua where monopoles condense. In this regime one has to tune the Higgs VEVs such that $`|v_iv_j|1/N_c`$. Then monopole masses are $`O(1)`$, and instantons are not suppressed. In this case the naive scaling relation $`g^21/N_c`$ is not satisfied.
(v) The large $`N_c`$ prediction for pure gluodynamics $`\chi _{top}\mathrm{\Lambda }_{QCD}^4N_c^2ϵ_{vac}`$ was recently verified in a string theory setting . This result is consistent with a number of scenarios for the instanton liquid. For example, the instanton ensemble might evolve from a random gas of density $`n\chi _{top}`$ at small $`N_c`$ to a strongly correlated liquid with large density $`(N/V)ϵ_{vac}`$ but very small fluctuations $`\chi _{top}=(N_IN_A)^2/VN_c^2(N/V)`$.
(vi) Instantons are not incompatible with the success of large $`N_c`$ arguments in the baryon sector. The instanton induced interaction between quarks is $`1/N_c`$ suppressed, just like one gluon exchange. Also, the topological soliton model can be derived from instantons in the large $`N_c`$ limit . This model automatically incorporates all the large $`N_c`$ predictions.
## V Smallness of OZI-violating amplitudes in channels other than $`O^\pm `$: conspiracy between hadronic amplitudes
We completely agree with Isgur and Thacker on this point. We would like to make two comments. First of all, instantons are compatible with the smallness of the OZI violating amplitude in the vector channel. The direct instanton contribution to the disconnected correlator vanishes. The next correction is suppressed also. Correlated instanton-anti-instanton pairs contribute to connected, but not to disconnected vector correlators .
Second, we would like to point out that we have compiled an extensive phenomenological analysis of hadronic correlation functions . This compilation contains a number of examples for cancellations between a large number of hadronic amplitudes. An example is the striking “superduality” phenomenon in the vector channels<sup>*</sup><sup>*</sup>*We note that this phenomenon is different from the standard parton-hadron duality. Ordinary duality holds in all channels, but usually breaks down at distances $`x>0.3`$ fm.: The full correlation function remains close to free field behavior for distances as large as $`|xy|<1.3`$ fm. In terms of the spectral representation, this behavior requires remarkable fine tuning between different hadronic amplitudes. This result can be directly verified from data, using the cross section for $`e^+e^{}`$ annihilation into hadrons.
Another important conclusions from this study is the fact that OZI violation in the scalar and pseudoscalar channel appears at very short distances. This is one more argument in favor of instantons rather than confinement as the source of OZI violation. While instantons appear at short distances, confinement is a long distance effect.
## VI Conclusions
In summary we showed that instantons explain the sign and magnitude of the OZI violating amplitude in all mesonic channels, not just for the $`\eta ^{}`$. In addition to that, instantons provide a mechanism for chiral symmetry breaking, and a successful description of light hadron spectroscopy. The role of instantons in the large $`N_c`$ limit is a complicated issue, but instantons and the large $`N_c`$ need not be incompatible. Nevertheless, naive large $`N_c`$ counting rules may not work in all channels. Instantons provide an explanation why large $`N_c`$ predictions work well in some cases, but fail badly in others.
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# 1 Introduction
## 1 Introduction
The search for Higgs bosons and supersymmetric particles is among the most important endeavors of present and future high energy physics. The novel colored particles, squarks and gluinos, and the weakly interacting gauginos can be searched for at the upgraded Tevatron, a $`p\overline{p}`$ collider with a c.m. energy of 2 TeV, and the LHC, a $`pp`$ collider with a c.m. energy of 14 TeV. Until now the search at the Tevatron has set the most stringent bounds on the colored SUSY particle masses. At the 95% CL, gluinos have to be heavier than about 180 GeV, while squarks with masses below about 180 GeV have been excluded for gluino masses below $`300`$ GeV . Stops, the scalar superpartners of the top quark, have been excluded in a significant part of the MSSM parameter space with mass less than about 80 GeV by the LEP and Tevatron experiments . Finally charginos with masses below about 90 GeV have been excluded by the LEP experiments, while the present search at the Tevatron is sensitive to chargino masses of about 60–80 GeV with a strong dependence on the specific model . Due to the negative search at LEP2 the lightest neutralino $`\stackrel{~}{\chi }_1^0`$ has to be heavier than about 30 GeV in the context of SUGRA models . In the $`R`$-parity-conserving MSSM, supersymmetric particles can only be produced in pairs. All supersymmetric particles will decay to the lightest supersymmetric particle (LSP), which is most likely to be a neutralino, stable thanks to conserved $`R`$-parity. Thus the final signatures for the production of supersymmetric particles will mainly be jets, charged leptons and missing transverse energy, which is carried away by neutrinos and the invisible neutral LSP.
In Section 2 we shall summarize the details of the calculation of the NLO QCD corrections, as described in Refs. \[??\] for the case of $`\stackrel{~}{q}\overline{\stackrel{~}{q}}`$ production. The evaluation of the full SUSY QCD corrections splits into two pieces, the virtual corrections, generated by virtual particle exchanges, and the real corrections, which originate from gluon radiation and the corresponding crossed processes with three-particle final states.
In Section 3 we shall consider the production of squarks and gluinos except stops . We assume the light-flavored squarks to be mass degenerate, which is a reasonable approximation for all squark flavors except stops, while the light quarks ($`u,d,s,c,b`$) will be treated as massless particles. The production of stop pairs requires the inclusion of mass splitting and mixing effects and will be investigated in Section 4. In Section 5 we will summarize the results for the production of charginos and neutralinos at NLO . The calculation of the LO cross sections has been performed a long time ago . Since in most of the cases the \[unphysical\] scale dependence of the LO quantities amounts up to about 50%, the determination of the NLO corrections is necessary in order to gain a reliable theoretical prediction, which can be used in present and future experimental analyses.
## 2 SUSY QCD corrections
### 2.1 Virtual corrections
The one-loop virtual corrections are built up by gluon, gluino, quark and squark exchange contributions \[see Fig. 1\]. They have to be contracted with the LO matrix elements. The calculation of the one-loop contributions has been performed in dimensional regularization, leading to the extraction of ultraviolet, infrared and collinear singularities as poles in $`ϵ=(4n)/2`$. For the chiral $`\gamma _5`$ coupling we have used the naive scheme, which is well justified in the present analysis at the one-loop levelWe have explicitly checked that the results obtained with a consistent $`\gamma _5`$ scheme are identical to the one with the naive scheme.. We have explicitly checked that after summing all virtual corrections no quadratic divergences are left over, in accordance with the general property of supersymmetric theories. The renormalization has been performed by identifying the squark and gluino masses with their pole masses, and defining the strong coupling in the $`\overline{\mathrm{MS}}`$ scheme including five light flavors in the corresponding $`\beta `$ function. The massive particles, i.e. squarks, gluinos and top quarks, have been decoupled by subtracting their contribution at vanishing momentum transfer . In dimensional regularization, there is a mismatch between the gluonic degrees of freedom \[d.o.f. = $`n2`$\] and those of the gluino \[d.o.f. = $`2`$\], so that SUSY is explicitly broken. In order to restore SUSY in the physical observables in the massless limit, an additional finite counter-term is required for the renormalization of the novel $`\stackrel{~}{q}\stackrel{~}{g}\overline{q}`$ vertex .
### 2.2 Real corrections
The real corrections originate from the radiation of a gluon in all possible ways and from the crossed processes by interchanging the gluon of the final state against a light (anti)quark in the initial state. The phase-space integration of the final-state particles has been performed in $`n=42ϵ`$ dimensions, leading to the extraction of infrared and collinear singularities as poles in $`ϵ`$. After evaluating all angular integrals and adding the virtual and real corrections, the infrared singularities cancel. The left-over collinear singularities are universal and are absorbed in the renormalization of the parton densities at NLO. We defined the parton densities in the conventional $`\overline{\mathrm{MS}}`$ scheme including five light flavors, i.e. the squark, gluino and top quark contributions are not included in the mass factorization. Finally we end up with an ultraviolet, infrared and collinear finite partonic cross section.
However, there is an additional class of physical singularities, which have to be regularized. In the second diagram of Fig. 2 an intermediate $`\stackrel{~}{q}\stackrel{~}{g}^{}`$ state is produced, before the \[off-shell\] gluino splits into a $`q\overline{\stackrel{~}{q}}`$ pair. If the gluino mass is larger than the common squark mass, and the partonic c.m. energy is larger than the sum of the squark and gluino masses, the intermediate gluino can be produced on its mass-shell. Thus the real corrections to $`\stackrel{~}{q}\overline{\stackrel{~}{q}}`$ production contain a contribution of $`\stackrel{~}{q}\stackrel{~}{g}`$ production. The residue of this part corresponds to $`\stackrel{~}{q}\stackrel{~}{g}`$ production with the subsequent gluino decay $`\stackrel{~}{g}\overline{\stackrel{~}{q}}q`$, which is already contained in the LO cross section of $`\stackrel{~}{q}\stackrel{~}{g}`$ pair production, including all final-state cascade decays. This term has to be subtracted in order to derive a well-defined production cross section. Analogous subtractions emerge in all reactions: if the gluino mass is larger than the squark mass, the contributions from $`\stackrel{~}{g}\stackrel{~}{q}\overline{q},\overline{\stackrel{~}{q}}q`$ have to be subtracted, and in the reverse case the contributions of squark decays into gluinos have to subtracted.
## 3 Production of Squarks and Gluinos
Squarks and gluinos can be produced via $`pp,p\overline{p}\stackrel{~}{q}\overline{\stackrel{~}{q}},\stackrel{~}{q}\stackrel{~}{q},\stackrel{~}{q}\stackrel{~}{g},\stackrel{~}{g}\stackrel{~}{g}`$ at hadron colliders. The hadronic squark and gluino production cross sections can be obtained from the partonic ones by convolution with the corresponding parton densities. We have performed the numerical analysis for the upgraded Tevatron and the LHC. For the natural renormalization/factorization scale choice $`Q=m`$, where $`m`$ denotes the average mass of the final-state SUSY particles, the SUSY QCD corrections are large and positive, increasing the total cross sections by 10–90% . This is shown in Fig. 3, where the K factors, defined as the ratios of the NLO and LO cross sections, are presented as a function of the corresponding SUSY particle mass for the LHC.
We have investigated the residual scale dependence in LO and NLO, which is presented in Fig. 4. The inclusion of the NLO corrections reduces the LO scale dependence by a factor 3–4 and reaches a typical level of $`15\%`$, which serves as an estimate of the remaining theoretical uncertainty. Moreover, the dependence on different sets of parton densities is rather weak and leads to an additional uncertainty of $`15\%`$. In order to quantify the effect of the NLO corrections on the search for squarks and gluinos at hadron colliders, we have extracted the SUSY particle masses corresponding to several fixed values of the production cross sections. These masses are increased by 10–30 GeV at the Tevatron and 10–50 GeV at the LHC, thus enhancing the present and future bounds on the squark and gluino masses significantly.
Finally we have evaluated the QCD-corrected transverse-momentum and rapidity distributions for all different processes. As can be inferred from Fig. 5, the modification of the normalized distributions in NLO compared to LO is less than about 15% for the transverse-momentum distributions and much less for the rapidity distributions. Thus it is a sufficient approximation to rescale the LO distributions uniformly by the K factors of the total cross sections.
## 4 Stop Pair Production
At LO only pairs of $`\stackrel{~}{t}_1`$ or pairs of $`\stackrel{~}{t}_2`$ can be produced at hadron colliders. Mixed $`\stackrel{~}{t}_1\stackrel{~}{t}_2`$ pair production is only possible at NLO and beyond. However, we have estimated that mixed stop pair production is completely suppressed by several orders of magnitude and can thus safely be neglected . The evaluation of the QCD corrections proceeds along the same lines as in the case of squarks and gluinos. The strong coupling and the parton densities have been defined in the $`\overline{\mathrm{MS}}`$ scheme with 5 light flavors contributing to their scale dependences, while the stop masses are renormalized on-shell.
The QCD corrections increase the cross sections by up to about 40% \[see Fig. 6\] and thus lead to an increase of the extracted stop masses from the measurement of the total cross section.
Moreover, as in the squark/gluino case the scale dependence is strongly reduced \[see Fig. 7\] and yields an estimate of about 15% of the remaining theoretical uncertainty at NLO. At NLO the virtual corrections depend on the stop mixing angle, the squark, gluino and second stop masses. However, it turns out that these dependences are very weak and can safely be neglected as can be inferred from Fig. 8.
## 5 Chargino and Neutralino Production
The production cross sections of charginos and neutralinos depend on several MSSM parameters, i.e. $`M_1,M_2,\mu `$ and $`\mathrm{tan}\beta `$ at LO . The cross sections are sizeable for chargino/neutralino masses below about 100 GeV at the upgraded Tevatron and less than about 200 GeV at the LHC. Due to the strong dependence on the MSSM parameters the extracted bounds on the chargino and neutralino masses depend on the specific region in the MSSM parameter space . The outline of the determination of the QCD corrections is analogous to the previous cases of squarks, gluinos and stops. The resonance contributions due to $`gq\stackrel{~}{\chi }_i\stackrel{~}{q}`$ with $`\stackrel{~}{q}q\stackrel{~}{\chi }_j`$ have to be subtracted in order to avoid double counting with the associated production of gauginos and strongly interacting squarks and gluinos. The parton densities have been defined with 5 light flavors contributing to their scale evolution in the $`\overline{\mathrm{MS}}`$ scheme, while the $`t`$-channel squark masses have been renormalized on-shell.
The QCD corrections enhance the production cross sections of charginos and neutralinos by about 10–40% \[see Fig. 9\]. The LO scale dependence is reduced to about 10% at NLO \[see Fig. 9\], which signalizes a significant stabilization of the thoretical prediction for the production cross sections . The dependence of the chargino/neutralino production cross sections on the specific set of parton densities ranges at about 15%.
## 6 Conclusions
In this work we have reviewed the status of SUSY particle production at hadron colliders at NLO. Most QCD corrections to the production processes are known, thus yielding a nearly complete theoretical status. There are especially large QCD corrections to the production of gluinos, which significantly increase the extracted bounds on the gluino mass from the negative search for these particles at the Tevatron. In all processes, which are known at NLO, the theoretical uncertainties are reduced to about 15%, which should be sufficient for the upgraded Tevatron and the LHCThe computer code PROSPINO for the production of squarks, gluinos and stops at hadron colliders is available at http://www.desy.de/$``$spira. The NLO production of gauginos and sleptons will be included soon..
Acknowledgements
I would like to thank W. Beenakker, R. Höpker, M. Krämer, M. Klasen, T. Plehn and P.M. Zerwas for their collaboration and the organizers of the Les Houches workshop for the invitation, the pleasant atmosphere and financial support.
Comparison of exact matrix element calculations with ISAJET and PYTHIA in case of degenerate spectrum in the MSSM
S. ABDULLIN, V. ILYIN and T. KON
One of the main purposes of the LHC collider is to search for physics beyond the Standard Model (SM), in particular to look for superpartners of ordinary particles expected in SUperSYmmetric extensions of the SM (SUSY). SUSY, if it exists, is expected to reveal itself at the LHC firstly via an excess of $`(multilepton+)`$ $`multijet`$ \+ $`E_T^{miss}`$ final states compared to Standard Model (SM) expectations . The bulk of these studies was carried out in the framework of the Minimal SuperSymmetric Model (MSSM), or the Supergravity-inspired subset of MSSM, minimal SUGRA (mSUGRA) with universal gaugino mass and corresponding mass hierarchy.
Unfortunately even in the framework of the MSSM one can expect scenarios with non-universal gaugino mass leading to some “pathological” SUSY cases which are difficult to observe due to a “degenerate” mass spectrum . For instance, it is quite obvious that there are no visible objects in the SUSY event if $`m_{\stackrel{~}{\chi }_1^0}m_{\stackrel{~}{g}}m_{\stackrel{~}{q}}=M_{SUSY}`$. In this case one can rely mainly on observation of jet activity which is additional to sparticle pair production. The idea is that additional high $`p_T`$ recoil partons can provide a clear $`E_T^{miss}+jet(s)`$ signature to observe the excess of SUSY events over SM expectations.
Our preliminary estimates of the MSSM signal generated with ISAJET show that the limit of visibility (at the 5 $`\sigma `$ level) in the MSSM is about 1.2 – 1.3 TeV with integrated luminosity of 100 fb<sup>-1</sup> for $`m_{\stackrel{~}{\chi }_1^0}m_{\stackrel{~}{g}}m_{\stackrel{~}{q}}`$, thus much smaller than the squark-gluino mass reach ($``$ 2.5 TeV) in mSUGRA with the same integrated luminosity .
It is well known that data bases in ISAJET and PYTHIA contain matrix elements only for $`22`$ subprocesses. So, additional partons can be produced only via initial and final state radiation (ISR/FSR parton showering). If one discusses production of hard partons, this mechanism is an approximation with a priori unknown precision. From a general point of view this approximation should underestimate the result obtained within the perturbative QCD theory through evaluation of the complete set of relevant Feynman diagrams. Thus, one might assume that there will be a substantial number of events with high $`p_T`$ of the heavy SUSY system recoiling against observable jet(s), larger than predicted by the parton showering mechanism in ISAJET and PYTHIA. We use the packages CompHEP and GRACE to evaluate $`23`$ matrix elements in leading order and calculate the distributions. Finally we compare these results with ones obtained with ISAJET and PYTHIA.
## Calculational details
To make a coherent calculation we choose a set of parameters as follows:
* gluino mass $`m_{\stackrel{~}{g}}`$ = 1.2 TeV;
* common squark mass $`m_{\stackrel{~}{q}}`$ = 100 TeV (10 TeV for PYTHIA) to neglect squark contributions;
* CTEQ4M set of structure functions;
* $`Q^2`$ = $`4m_{\stackrel{~}{g}}^2`$ ( in PYTHIA $`Q^2=(m_1^2+m_2^2)/2`$);
* NLO running strong coupling with normalization $`\alpha _S(M_Z)=0.118`$, so $`\alpha _S(2400\mathrm{GeV})=0.0817`$.
We use ISAJET 7.44 and PYTHIA 6.125 with multiple interactions switched off. We treat the jet in ISAJET/PYTHIA not as the single hardest final state parton, but use the final state parton list (excluding the two gluinos) as input to an iterative cone jet finder algorithm ($`R_{cone}`$ = 0.5) to find a group of partons forming parton-level jets.
In CompHEP and GRACE we calculate exact tree level matrix elements for the subprocesses $`:`$
$$gg\stackrel{~}{g}\stackrel{~}{g}+g;gq(Q)\stackrel{~}{g}\stackrel{~}{g}+q(Q);qQ\stackrel{~}{g}\stackrel{~}{g}+g;$$
which were convoluted with the CTEQ4m parton distributions to get LHC cross sections and distributions. The $`q`$ ($`Q`$) terms stand for quarks (antiquarks) of the first five flavours: $`u`$, $`d`$, $`c`$, $`s`$ and $`b`$.
We ask the final (additional) parton, gluon or quark, to radiate within the typical ATLAS and CMS rapidity, $`|\eta _j|<4.5`$, and to be hard enough, $`E_T^j>10`$ GeV. No cuts are applied on the gluino pair.
## Results and Conclusions
The cross section of the process $`pp\stackrel{~}{g}\stackrel{~}{g}+jet`$ within cuts on the jet turns out to be rather high, about 120 fb, with the following contributions from different channels :
$$\sigma _{gg}=78.9\text{fb};\sigma _{gq(Q)}=21.9\text{fb};\sigma _{qQ}=20.0\text{fb}.$$
We found that the shapes of the $`p_T`$ distributions are quite different in these three cases, see Fig.1. Indeed, one can expect that the $`qQ`$ channel produces less events than the $`gg`$ contribution, due to the dominance of the gluon component in the proton for small $`x`$. However, one can see that the $`gq(Q)`$ curve is higher than the gluon-gluon one for $`p_T>400`$ GeV, being much smaller even than $`qQ`$ at the soft end. Note that the set of topologies of Feynman diagrams is the same in all three channels. So, the observed difference connects surely with different interplay of contributions of vector and spinor virtual particles.
In Fig.2 the $`p_T`$ distribution for $`pp\stackrel{~}{g}\stackrel{~}{g}+jet`$ is given as a sum of the three channels discussed above. The CompHEP/GRACE result is represented by solid histograms, ISAJET by dashed histograms and PYTHIA by dot-dashed ones. One can see that the PYTHIA distribution underestimates the exact ME result smoothly. Of course this (almost constant) difference can be explained (corrected) by a different setting of some QCD switches. At the same time one can see a more pronounced deviation from the exact ME in the shape of the ISAJET curve. In particular at moderate $`p_T`$ ($`200900`$ GeV) ISAJET overestimates the exact ME result. One can also see that for very high transverse momenta ($``$ 1200 GeV for ISAJET and 1400 GeV for PYTHIA) the parton showering approximation is too crude.
Coming back to the initial idea to look at additional jet activity with high $`p_T`$ recoiling against the heavy SUSY system one can conclude that parton shower approximation works well enough up to very high values of transverse momenta. Therefore, unfortunately one can not find room for sizable improvements of parton shower predictions by using exact matrix elements in the process discussed.
Some interesting discrepancy is observed in the region of relatively small transverse momenta, $`p_T<150`$ GeV, where ISAJET and PYTHIA produce much weaker (by a factor of 2-3) jet activity than CompHEP/GRACE. However, this region is sensitive to higher order corrections (in particular due to resummation of large logarithms of $`p_T/m_{\stackrel{~}{g}}`$), thus a more careful analysis is necessary to get reliable conclusions.
The finite width effect on neutralino production
T. KON, Y. KURIHARA, M. KURODA and K. ODAGIRI
It has been proposed in that neutralino mass differences can be determined fairly accurately through the $`\mu `$-pair invariant mass distribution $`d\sigma /dM_{\mu \mu }`$ of the process
$$pp\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_2^0X\stackrel{~}{\chi }_1^+(\mu \stackrel{~}{\mu })X\stackrel{~}{\chi }_1^+\mu (\stackrel{~}{\chi }_1^0\mu )X,$$
(1)
as the distribution exhibits a sharp decrease at its kinematic endpoint. When, however, the width of the neutralino is taken into account, the smearing of the cross section might make this observation obsolete and less appealing for the determination of the neutralino mass. Therefore, it is important to know to what extent the width of neutralino makes the mass determination less accurate.
We have investigated this problem by comparing the muon-pair invariant mass distribution in the following three cases.
(case 1) Zero width approximation for $`\stackrel{~}{\chi }_2^0`$. Only those diagrams which proceed via the chain decay (1) are taken into account.
(case 2) Same as (1) but finite width of $`\stackrel{~}{\chi }_2^0`$ is used. Only those diagrams that proceed via resonances as shown in (1) are considered.
(case 3) Finite width of $`\stackrel{~}{\chi }_2^0`$ as well as the entire set of diagrams (165 diagrams in unitary gauge) that creates the final state $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^0\mu \mu `$ are taken into account. The diagrams which are considered in case 3 but not in case 2 constitute the background to this process.
In the numerical evaluation of the muon pair distribution, we have used the following set of SUSY parameters,
$$\mathrm{tan}\beta =12,M_2=200GeV,\mu =500GeV,$$
(2)
which results in the following masses for charginos and neutralinos,
$`m_{\stackrel{~}{\chi }_i^0}`$ $`=`$ $`(97,197,507,511),`$ (3)
$`m_{\stackrel{~}{\chi }_i^+}`$ $`=`$ $`(197,514).`$ (4)
In particular, in the present set of parameters we are considering, $`\stackrel{~}{\chi }_1^0`$ is almost purely bino (the SUSY partner of the $`U(1)`$ gauge boson, $`B_\mu `$), while $`\stackrel{~}{\chi }_2^0`$ is almost purely neutral wino (the SUSY partner of the neutral $`SU(2)_L`$ gauge-boson, $`W_3^\mu `$). As for the slepton masses, we have used the common mass $`m_\stackrel{~}{\mathrm{}}`$ = 150 GeV.
Assuming that other SUSY particles are heavy so that $`\stackrel{~}{\chi }_2^0`$ decays dominantly into $`\mathrm{}\stackrel{~}{\mathrm{}}`$, we have evaluated the width of $`\stackrel{~}{\chi }_2^0`$ as
$$\mathrm{\Gamma }_{\stackrel{~}{\chi }_2^0}=1.25GeV.$$
(5)
The process is described by the subprocess
$$q\overline{q}^{}\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^+(\mu \stackrel{~}{\mu })\stackrel{~}{\chi }_1^+\mu (\stackrel{~}{\chi }_1^0\mu ).$$
(6)
For the quark distribution functions, we have used QTEC-3D .
The numerical evaluation was performed by HERWIG in case 1 and case 2, while the case 3 was calculated by GRACE .
Results are shown in Fig.1(a) and (b), where the distributions of the invariant mass $`M_{\mu \mu }`$ in the second neutralino decay are shown. In Fig.1(a), the dashed line represents the zero width approximation (case 1), the dotted line represents the distribution for the finite $`\stackrel{~}{\chi }_2^0`$ width but with only the resonance diagrams (case 2), while the solid line represents the distribution for the finite $`\stackrel{~}{\chi }_2^0`$ width and full set of diagrams which give the final states $`\stackrel{~}{\chi }_1^+\stackrel{~}{\chi }_1^0\mu \mu `$ (case 3). We show a detailed comparison between case 2 and case 3 around the end point region of $`M_{\mu \mu }`$ in Fig.1(b). As one sees from the figures, the edge of the distribution at the end point of the phase space remains, although it is broadened by the finite width effect. As a result, the exact position of the end point acquires an ambiguity of about 5 GeV, although the precision of mass determination, which may additionally make use of the shape of the distribution, can presumably be made much less than this. A full analysis will need to take into account the experimental resolutions and signal selection efficiency, and this is beyond the scope of our study. The total cross section of the signal process (1) is 0.015 pb.
Acknowledgements One of the authors, K.O., would like to thank Drs. J. Kanzaki and P. Richardson for their computing help.
Width effects in slepton production $`e^+e^{}\stackrel{~}{\mu }_R^+\stackrel{~}{\mu }_R^{}`$
H.-U. MARTYN
## 1 Introduction
If supersymmetry will be discovered in nature a precise measurement of the particle spectrum will be very important in order to determine the underlying theory. The potential of the proposed Tesla Linear Collider with its high luminosity and polarization of both $`e^\pm `$ beams will allow to obtain particle masses with an accuracy of $`10^3`$ or better . At such a precision width effects of primary and secondary particles may become non-negligible.
The present case study is based on a particular $`R`$-parity conserving mSUGRA scenario, also investigated in the Ecfa/Desy Study , with parameters $`m_0=100\text{GeV},m_{1/2}=200\text{GeV},A_0=0\text{GeV},\mathrm{tan}\beta =3`$ and $`\mathrm{sgn}(\mu )>0`$. The particle spectrum is shown in fig. 1. Typical decay widths of the scalar leptons are expected to be $`\mathrm{\Gamma }0.30.5\text{GeV}`$, while the widths of the light gauginos, decaying into 3-body final states, are (experimentally) negligible.
This note presents, as an example, a simulation of right scalar muon production
$`e_R^{}e_L^+`$ $``$ $`\stackrel{~}{\mu }_R^{}\stackrel{~}{\mu }_R^+,`$
$``$ $`\mu ^{}\chi _1^0\mu ^+\chi _1^0.`$
The analysis is based on the methods and techniques described in a comprehensive study of the same Susy spectrum . The detector concept, acceptances and resolutions are taken from the Tesla Conceptual Design Report . Events are generated with the Monte Carlo program Pythia 6.115 , which includes the width of supersymmetric particles as well as QED radiation and beamstrahlung . It is assumed that both beams are polarized, right-handed electrons to a degree of $`𝒫_{e_R^{}}=0.80`$ and left-handed positrons by $`𝒫_{e_L^+}=0.60`$. A proper choice of polarizations increases the cross section by a factor of $`3`$ and reduces the background substantially, e.g. by more than an order of magnitude for Standard Model processes.
## 2 Mass determinations
Scalar muons $`\stackrel{~}{\mu }_R`$ are produced in pairs via $`s`$ channel $`\gamma `$ and $`Z`$ exchange and decay into an ordinary muon and a stable neutralino $`\chi _1^0`$ (LSP), which escapes detection. The experimental signatures are two acoplanar muons in the final state with large missing energy and nothing else in the detector. Simple selection criteria (essentially cuts on acollinearity angle and missing energy) suppress background from $`W^+W^{}`$ pairs and cascade decays of higher mass Susy particles and result in detection efficiencies around $`70\%`$.
Two methods to determine the mass of $`\stackrel{~}{\mu }_R`$ will be discussed: (i) a threshold scan of the pair production cross section, and (ii) a measurement of the energy spectrum of the decay muons, which simultaneously constrains the mass of the primary smuon and the secondary neutralino. The particle mass parameters given by the chosen Susy model are $`m_{\stackrel{~}{\mu }_R}=132.0\text{GeV}`$, $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}=0.310\text{GeV}`$ and $`m_{\chi _1^0}=71.9\text{GeV}`$.
### 2.1 Threshold scan
Cross section measurements close to production threshold are relatively simple. One essentially counts additional events with a specific signature, here two oppositely charged, almost monoenergetic muons, over a smooth background. The cross section for slepton pair production rises as $`\sigma \beta ^3`$, where $`\beta =\sqrt{14m_{\stackrel{~}{\mu }_R}^2/s}`$ is the velocity related to the $`\stackrel{~}{\mu }_R`$ mass. The excitation curve as a function of the cms energy, including effects due to QED initial state radiation and beamstrahlung, is shown in figure 2. The sensitivity to the width $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}`$ is most pronounced close to the kinematic production limit and diminishes with increasing energy. A larger width ‘softens’ the rise of the cross section with energy. Fits to various mass and/or width hypotheses are performed by simulating measurements with a total integrated luminosity of $`100\mathrm{fb}^1`$ distributed over 10 equidistant points around $`\sqrt{s}=264274\text{GeV}`$. The data may be collected within a few months of Tesla operation.
Taking the width from the model prediction $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}=310\mathrm{MeV}`$, a fit to the threshold curve gives a statistical accuracy for the smuon mass of $`\delta m_{\stackrel{~}{\mu }_R}=90\mathrm{MeV}`$. This error is considerably smaller than the expected width. A two-parameter fit yields $`m_{\stackrel{~}{\mu }_R}=132.002{}_{0.130}{}^{+0.170}\text{GeV}`$ and $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}=311{}_{225}{}^{+560}\mathrm{MeV}`$. However, both parameters are highly correlated with a correlation coefficient of $`0.95`$. Finally, if one may fix the $`\stackrel{~}{\mu }_R`$ mass from another measurement, the width can be determined to $`\delta \mathrm{\Gamma }_{\stackrel{~}{\mu }_R}=\pm 190\mathrm{MeV}`$. It should be noted that the scan procedure and choice of energy points is by no means optimized. Possibilities to reduce the correlations should be studied.
### 2.2 Energy spectrum of $`\mu ^\pm `$
For energies far above threshold, the kinematics of the decay chain of reaction (1) allows to identify and to reconstruct the masses of the primary and secondary sparticles. The isotropic decays of the scalar muons lead to a flat energy spectrum of the observed final $`\mu ^\pm `$ in the laboratory frame. The endpoints of the energy distribution are related to the masses of the $`\stackrel{~}{\mu }_R`$ and $`\chi _1^0`$ via
$`{\displaystyle \frac{m_{\stackrel{~}{\mu }}^2m_{\chi ^0}^2}{2(E_{\stackrel{~}{\mu }}+p_{\stackrel{~}{\mu }})}}`$ $`E_\mu `$ $`{\displaystyle \frac{m_{\stackrel{~}{\mu }}^2m_{\chi ^0}^2}{2(E_{\stackrel{~}{\mu }}p_{\stackrel{~}{\mu }})}}.`$ (2)
In practice the sharp edges of the energy spectrum will be smeared by effects due to detector resolution, selection criteria and in particular initial state radiation and beamstrahlung. The results of a simulation at $`\sqrt{s}=320\text{GeV}`$ assuming an integrated luminosity of $`160\mathrm{fb}^1`$ are shown in figure 3. One observes a clear signal from $`\stackrel{~}{\mu }_R`$ pair production above a small background of cascade decays $`\chi _2^0\mu ^+\mu ^{}\chi _1^0`$ from the reaction $`e_R^{}e_L^+\chi _2^0\chi _1^0`$. Contamination from chargino or $`W`$ pair production is completely negligible.
A two-parameter fit to the $`\mu `$ energy spectrum yields masses of $`m_{\stackrel{~}{\mu }_R}=132.0\pm 0.3\text{GeV}`$ and $`m_{\chi _1^0}=71.9\pm 0.2\text{GeV}`$. The statistical accuracy is of the same size as the expected width of the scalar muon. Choosing a different width $`\mathrm{\Gamma }_{\stackrel{~}{\mu }_R}`$ in the simulation modifies essentially the $`\mu `$ energy spectrum at the low endpoint and has little impact at higher energies. This is illustrated in figure 3, right part, which compares the lower part of the spectrum with the predictions of width zero and twice the expected value. The sharp rise is getting smeared out with increasing width. With the anticipated luminosity of $`160\mathrm{fb}^1`$ it may be feasible to distinguish these cases.
### 2.3 Production of other sparticles
It should be noted that estimates on the sensitivity of width effects in other slepton production channels can be obtained from the above results by scaling the cross section and taking the branching ratios into final states into account. Thus one expects e.g. a gain by a factor of $``$2 for selectron $`\stackrel{~}{e}_R`$ and sneutrino $`\stackrel{~}{\nu }_e`$ pair production. For the higher mass chargino $`\chi _2^\pm `$ and neutralinos $`\chi _{3,\mathrm{\hspace{0.17em}4}}^0`$ mass resolutions of $`0.250.50\text{GeV}`$ may be obtained from threshold scans , where the cross sections rise as $`\sigma \beta `$. The corresponding widths are expected to be $`25\text{GeV}`$ (two-body decays in a gauge boson and gaugino) and have certainly to be considered.
## 3 Conclusions
The high luminosity of Tesla allows to study the production and decays of the accessible Susy particle spectrum. Polarization of both $`e^{}`$ and $`e^+`$ beams is very important to optimize the signal and suppress backgrounds. A simulation of slepton production $`e^+e^{}\stackrel{~}{\mu }_R^+\stackrel{~}{\mu }_R^{}`$ shows that for precision mass measurements with an accuracy of $`𝒪(100\mathrm{MeV})`$ the widths of the primary particles have to be taken into account. Finally, it is worth noting that the anticipated mass resolutions from threshold scans or lepton energy spectra can only be obtained if beamstrahlung effects are well under control.
Radiative Effects on Squark Pair Production at $`e^+e^{}`$ Colliders
M. DREES, O.J.P. ÉBOLI, R.M. GODBOLE and S. KRAML
Among the advantages of $`e^+e^{}`$ colliders over hadron colliders are the fairly well defined center–of–mass energy in hard (annihilation) events, and the comparative cleanliness of the environment. Taken together, these properties make the kinematic reconstruction of $`e^+e^{}`$ events much easier than that of comparable events at hadron colliders. In particular, experiments at $`e^+e^{}`$ colliders should be able to measure the masses of new particles (with unsuppressed electroweak couplings) with an error of 1% or less , either through threshold scans, or through fitting kinematic distributions of events well above threshold. The latter method is more versatile, since the same data set can well contain several kinds of “new physics” events, which can be separated from each other (and from backgrounds) using kinematical cuts. Moreover, kinematical reconstruction allows to determine not only the masses of the new particles produced in the primary interaction, but also those of their decay products.
However, in order to correctly interpret the information contained in various distributions one needs an accurate model of the final state. In particular, if one wants to achieve errors of 1% or less, various radiative effects have to be taken into account. These are of special importance if the new particles and/or their decay products have strong interactions, since a significant fraction of all signal events will then contain hard gluons. Here we study the impact of these radiative effects on the measurement of squark masses at $`e^+e^{}`$ colliders. Squarks are the currently most plausible new particles that have both strong and electroweak interactions.
In the last few years considerable progress has been made in the accurate calculation of total cross sections for squark pair production and of squark branching ratios. In particular, one–loop corrections to these quantities from both ordinary QCD and SUSY QCD, as well as from Yukawa interactions, are now known. These have been used to estimate the error with which the squark mass can be extracted from a measurement of the total squark pair production cross section times branching ratio into a given final state. However, in such a “dynamical” determination of the squark mass one has to assume values for all the other input parameters that affect the cross section times branching ratio. These include the gluino mass and, for third generation squarks, also the masses and mixing angles of squarks, charginos, neutralinos, and Higgs bosons .
In contrast, relatively little attention has been paid to the kinematical determination of squark masses at $`e^+e^{}`$ colliders. The pioneering work by Feng and Finnell investigates the usefulness of various kinematical distributions, and concludes that experiments at a 500 GeV collider should be able to determine the mass of 200 GeV squarks with an error of $`0.5\%`$ using just 20 fb<sup>-1</sup> of data, if all squarks decay directly into a massless quark and an invisible (stable or long–lived) neutralino $`\stackrel{~}{\chi }_1^0`$. However, their study did not include any radiative effects. Here we update their analysis by including initial state radiation of photons, as well as the emission of hard gluons during the production ($`e^+e^{}\stackrel{~}{q}\stackrel{~}{q}^{}g`$) and/or decay ($`\stackrel{~}{q}q\stackrel{~}{\chi }_1^0g`$) of the squarks. We also take larger squark masses $`(m_{\stackrel{~}{q}}300`$ GeV) to account for recent experimental bounds from the Tevatron , and a correspondingly higher center–of–mass energy of 800 GeV.
For reasons of space we only briefly summarize the ingredients of our analysis; details will be given elsewhere. We treat the emission of photons off the initial state in the structure function formalism . The differential cross section is then given by
$$d\sigma =𝑑x_1𝑑x_2f_{e|e}(x_1,\sqrt{s})f_{e|e}(x_2,\sqrt{s})𝑑\widehat{\sigma }(\widehat{s}=x_1x_2s),$$
(1)
where $`\widehat{\sigma }`$ is the cross section in the absence of ISR, and
$$f_{e|e}(x,\sqrt{s})=\beta \left[(1x)^{\beta 1}\left(1+\frac{3}{4}\beta \right)\frac{\beta }{2}(1+x)\right],$$
(2)
with $`\beta =\frac{\alpha _{\mathrm{em}}}{\pi }\left(\mathrm{log}\frac{s}{m_e^2}1\right)`$, is the leading–log resummed effective $`e^\pm `$ distribution function. Note that we do not include beamstrahlung (which is expected to further smear out the peak in the $`e^+e^{}`$ luminosity at $`\widehat{s}=s`$), since it depends on details of accelerator design. Moreover, we conservatively assume that all ISR photons escape detection, even though eq.(2) is strictly valid only if there are no experimental constraints on the phase space of the emitted photons. However, we will see that the main effect of ISR is an overall reduction of the cross section by $`15\%`$; kinematical distributions are little affected even if all ISR photons are invisible.
We treat the emission of gluons during $`\stackrel{~}{q}\stackrel{~}{q}^{}`$ production as described in . In particular, we introduce a minimal gluon energy $`E_{g,\mathrm{min}}`$ to regularize IR divergences. The final results will not depend on the value of this regulator after contribution from $`\stackrel{~}{q}\stackrel{~}{q}^{}`$ and $`\stackrel{~}{q}\stackrel{~}{q}^{}g`$ events have been added, if $`E_{g,\mathrm{min}}`$ is sufficiently small.<sup>*</sup><sup>*</sup>*Of course, virtual QCD corrections to $`\stackrel{~}{q}\stackrel{~}{q}^{}`$ production have to be added for this cancellation to work. In the numerical results presented below we take $`E_{g,\mathrm{min}}=1`$ GeV. For $`\sqrt{s}=800`$ GeV and $`m_{\stackrel{~}{q}}=300`$ GeV this implies that only 18% of all squark pairs are produced together with a “hard” gluon.
The squared matrix element for $`\stackrel{~}{q}q\stackrel{~}{\chi }_1^0g`$ can be found in . We again have to include virtual QCD corrections, in this case to $`\stackrel{~}{q}q\stackrel{~}{\chi }_1^0`$ decays, to cancel IR singularities. In this case we regularize these singularities by introducing a finite gluon mass $`m_g`$, which we also set to 1 GeV in our numerical examples. For our choice $`m_{\stackrel{~}{q}}=300`$ GeV, $`m_{\stackrel{~}{\chi }}=50`$ GeV this means that nearly 90% of all squark decays produce a “hard” gluon. Altogether more than 95% of all $`e^+e^{}\stackrel{~}{q}\stackrel{~}{q}^{}q\overline{q}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0X`$ events therefore contain at least one “hard” gluon, as defined through our two IR regulators. Note that, unlike for $`e^+e^{}\stackrel{~}{q}\stackrel{~}{q}^{}`$, virtual QCD corrections to $`\stackrel{~}{q}q\stackrel{~}{\chi }_1^0`$ decays introduce UV divergences. These cancel only after including full SUSY–QCD corrections , which depend on the mass of the gluino. We take $`m_{\stackrel{~}{g}}=450`$ GeV; the shapes of the kinematical distributions we study here are almost independent of this choice.
Due to the emission of hard gluons we have up to 5 visible partons in the final stateWe allow gluon emission during the production and both decays simultaneously, i.e. we include events with 2 or 3 gluons. These contributions are formally of NNLO. However, since production and decays are independent processes, up to terms $`𝒪(\mathrm{\Gamma }_{\stackrel{~}{q}}/\sqrt{s})`$, other (as yet unknown) NNLO contributions cannot cancel these known contributions.. In the numerical results presented below we group these into exactly 2 jets (except for events with $`<2`$ partons in the acceptance region defined by $`|\mathrm{cos}(\theta )|0.9`$, which we discard), using the $`k_T`$ clustering (“Durham”) algorithm . We have checked that one obtains very similar results if one merges partons with a fixed $`y_{\mathrm{cut}}`$ parameter, rather than a fixed number of final state jets, and only uses the two hardest jets for the kinematical analysis. In order to simulate experimental resolutions, we smear the energies (but not directions) of all partons before jet merging, with a Gaussian error given by $`\delta (E)=0.3/\sqrt{E}0.01`$ ($`E`$ in GeV). Finally, we apply a set of cuts in order to suppress backgrounds. We require that the energy of each jet exceeds 15 GeV, that the missing $`p_T`$ exceeds 56 GeV (taken from after scaling up from $`\sqrt{s}=500`$ GeV to $`800`$ GeV), and that the acoplanarity angle between the two jets exceed $`30^{}`$.
In Figs. 1a,b we show the resulting distribution in the jet energy (1a) and in the variable $`m_{\stackrel{~}{q},\mathrm{min}}`$ (1b) defined below. In the absence of cuts, radiative effects and energy smearing, the jet energy distribution should be constant (flat) between two kinematical endpoints, and zero elsewhere. The cuts distort this simple shape, producing a small peak near the lower edge of the spectrum; this is shown by the solid curve in Fig. 1a. The histograms show “experimental” distributions, based on 3,500 events before cuts; for a single $`q=2/3`$ $`SU(2)`$singlet ($`\stackrel{~}{u}_R`$) squark, this would require an integrated luminosity of about 205 fb<sup>-1</sup> (including ISR and QCD corrections, and assuming all squarks decay directly into $`q\stackrel{~}{\chi }_1^0`$); the total cross section is 17.0 (11.7) fb before (after) cuts.The total cross section for $`\stackrel{~}{u}_L`$ pair production for the same mass is 26.8 fb. However, in most models one expects $`SU(2)`$ doublet squarks to predominantly decay into charginos and heavier neutralinos, rather than directly into the LSP. The dotted histogram still does not include any radiative effects, but includes a finite energy resolution as described above; clearly this has little effect on the shape of the spectrum. As mentioned earlier, including ISR (dashed histogram) reduces the cross section, but again does not distort the shape of the spectrum very much; in particular, there is only little “leakage” beyond the nominal endpoints. In contrast, including gluon emission (solid histogram) does change the shape of the distribution. In particular, there are now quite a few entries below the lower nominal endpoint. In most of these events one of the quarks falls out of the acceptance region, so that one of the jets is entirely made up of gluons. QCD corrections also increase the total cross section before (after) cuts by 24% (42%)<sup>§</sup><sup>§</sup>§Note that QCD corrections markedly increase the acceptance of the cuts. It is therefore also important to include their effect on the kinematics when trying to extract $`m_{\stackrel{~}{q}}`$ from measurements of the $`\stackrel{~}{q}\stackrel{~}{q}^{}`$ production cross section.
The variable $`m_{\stackrel{~}{q},\mathrm{min}}`$ plotted in Fig. 1b has been introduced in . It is the minimal possible squark mass, if one assumes a fixed value of $`\sqrt{s}`$ (no ISR), and if the final state consists of exactly two massless quarks and two neutralinos with equal (and known) mass.The neutralino mass is expected to be known from analyses of data at lower energy, e.g. from chargino pair events. None of these assumptions hold in our case. In fact, in some cases the reconstruction described in is impossible, i.e. some trigonometrical function acquires a value exceeding unity; we discard these events. We nevertheless find that radiative effects introduce only a modest amount of “leakage” beyond the nominal endpoint of the distribution, which is at $`m_{\stackrel{~}{q}}`$. Note, however, that these effects significantly broaden the distribution, i.e. many events have migrated to lower values of $`m_{\stackrel{~}{q},\mathrm{min}}`$ (compare the solid and dotted histograms). Since events in the peak contribute most to the determination of $`m_{\stackrel{~}{q}}`$ , this broadening is expected to increase the error on $`m_{\stackrel{~}{q}}`$.
Fig. 2 shows that a fit of the $`m_{\stackrel{~}{q},\mathrm{min}}`$ distribution nevertheless yields a smaller statistical error for $`m_{\stackrel{~}{q}}`$ than a fit of the $`E_{\mathrm{jet}}`$ distribution does. In this figure we compare a mock data set for $`m_{\stackrel{~}{q}}=300`$ GeV, based on an integrated luminosity of just 50 fb<sup>-1</sup> (852 events before cuts), with “template” distributions, which have been computed for 13 different values of $`m_{\stackrel{~}{q}}`$, leaving all other input parameters unchanged, and applying the cuts described above. This gives 13 different values of $`\chi ^2`$; the open (filled) squares have been computed from the $`E_{\mathrm{jet}}`$ ($`m_{\stackrel{~}{q},\mathrm{min}}`$) distributions. A parabola is then fitted to $`\chi ^2(m_{\stackrel{~}{q}})`$. The minimum of this parabola gives the “measured” value $`m_{\stackrel{~}{q},0}`$ of $`m_{\stackrel{~}{q}}`$, while the ($`1\sigma `$) error is computed from $`\chi ^2(m_{\stackrel{~}{q},0}\pm \delta m_{\stackrel{~}{q}})=\chi _{\mathrm{min}}^2+1`$. Note that the $`\chi ^2`$ values are only based on the shapes of the distributions, i.e. the “data” have been normalized to the template before computing $`\chi ^2`$ for a given assumed value of $`m_{\stackrel{~}{q}}`$. Fig. 2 therefore shows the results of purely kinematical determinations of $`m_{\stackrel{~}{q}}`$.
We see that even after radiative effects have been included, a fit of the $`m_{\stackrel{~}{q},\mathrm{min}}`$ distribution determines the squark mass with a statistical error of well under 1% already with 50 fb<sup>-1</sup> of data, which corresponds to only about 1 month of running time for the planned TESLA collider . However, radiative effects also introduce possible new systematic errors, e.g. due to the choice of scale in $`\alpha _s`$. (In our calculation we used $`\sqrt{\widehat{s}}`$ for the corrections to $`\stackrel{~}{q}\stackrel{~}{q}^{}`$ production, and $`m_{\stackrel{~}{q}}`$ in $`\stackrel{~}{q}`$ decays.) Moreover, we haven’t included any hadronization effects yet. Note that the lighter scalar top eigenstate $`\stackrel{~}{t}_1`$, which is likely to be the lightest, and hence most easily accessible squark, might well hadronise before it decays . In addition, the massless quarks and gluons in the final state will hadronise into jets with finite masses. Finally, the error on the assumed LSP mass also propagates into the error on $`m_{\stackrel{~}{q}}`$ . We plan to investigate these issues in a future publication.
### Acknowledgements
This work was supported in part by Conselho Nacional de Desenvolvimento Científico e Tecnológico (CNPq), by Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP), by Programa de Apoio a Núcleos de Excelência (PRONEX), by “Fonds zur Förderung der wissenschaftlichen Forschung” of Austria, project no. P13139–PHY, by the Department of Science and Technology (India), and by the US National Science Foundation under NSF grant INT-9602567.
Spin correlations and phases for SUSY particle searches at $`e^+e^{}`$ colliders
N. GHODBANE
## 1 Introduction
Since one expects high luminosities for the next generation of linear colliders (e.g. $`500fb^1`$ for the TESLA project), one can use beam polarization to reduce the standard model backgrounds and use the polarization dependence of the cross sections to study specific SUSY parameters. Moreover, as it has been stressed by several authors , spin correlations play a major role in the kinematic distributions of final particles. For this reason, we upgraded the SUSYGEN Monte Carlo generator<sup>*</sup><sup>*</sup>*A description of the Monte Carlo generator can be found in the Higgs working group report. .
Here, we firstly study the spin correlations effects in the gaugino and the stau searches. Then, in a second part, we show how the newly reconsidered CP violating phases arising from MSSM can affect the chargino searches.
## 2 Beam polarization and spin correlations
As it has been stressed by several papers the study of the angular distribution of $`e^\pm `$ produced in $`e^+e^{}\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_1^0e^+e^{}`$ will give valuable information concerning the neutralino nature and then enable MSSM parameter extraction. Here we consider the two scenarios (A: $`M_2=78`$ GeV$`/c^2`$, $`\mu =250`$ GeV$`/c^2`$, $`\mathrm{tan}\beta =2`$) and (B: $`M_2=210`$ GeV$`/c^2`$, $`\mu =60`$ GeV$`/c^2`$, $`\mathrm{tan}\beta =2`$). In model A, $`\stackrel{~}{\chi }_1^0`$ is Bino like ($`m_{\stackrel{~}{\chi }_1^0}=42.5\text{ GeV}/c^2,\stackrel{~}{\chi }_1^0=+0.98\stackrel{~}{B}+0.17\stackrel{~}{W}^30.09\stackrel{~}{H}_1^0+0.04\stackrel{~}{H}_2^0`$) whereas $`\stackrel{~}{\chi }_2^0`$ is Wino like ($`m_{\stackrel{~}{\chi }_2^0}=91.9\text{ GeV }/c^2,\stackrel{~}{\chi }_2^0=+0.14\stackrel{~}{B}0.95\stackrel{~}{W}^30.28\stackrel{~}{H}_1^0+0.05\stackrel{~}{H}_2^0`$). In model B, the two lightest neutralinos are Higgsino like $`(m_{\stackrel{~}{\chi }_1^0}=55.1\text{ GeV}/c^2,\stackrel{~}{\chi }_1^0=+0.16\stackrel{~}{B}0.09\stackrel{~}{W}^30.78\stackrel{~}{H}_1^00.60\stackrel{~}{H}_2^0)`$ and $`(m_{\stackrel{~}{\chi }_2^0}=88.9\text{ GeV }/c^2,\stackrel{~}{\chi }_2^0=+0.21\stackrel{~}{B}0.24\stackrel{~}{W}^3+0.62\stackrel{~}{H}_1^00.71\stackrel{~}{H}_2^0)`$. For each one of these two models, we considered the common sfermions mass at the GUT scale, $`m_0`$, being equal to 80 GeV$`/c^2`$ and 200 GeV/$`c^2`$.
The right side of figure 1 illustrates the effect of spin-correlations in angular distributions of $`e^{}`$, decay product of $`\stackrel{~}{\chi }_2^0`$. One can notice that the angular distribution of the final leptons depends strongly on the sfermion mass parameter $`m_0`$. For small $`m_0`$ of inclusion of spin correlations gives an effect $`20\%`$ in the angular distributions. For higher selectron masses, the total cross section is smaller, but the spin correlation effects appear to be more important $`30\%`$. If the two neutralinos are Higgsino like, the effects are negligible ($`0.7\%`$) (see figure 2). Moreover these angular distributions depend strongly on the center of mass energy.
Recent studies have shown that $`\tau `$ polarization effects yield valuable information for the MSSM parameters e.g for $`\mathrm{tan}\beta `$, the nature of the LSP and the mixing angle $`\theta _{\stackrel{~}{\tau }}`$.
Figure 3 shows the momentum distribution of pions produced by $`e^+e^{}\stackrel{~}{\tau }_1^+\stackrel{~}{\tau }_1^{}\tau ^+\stackrel{~}{\chi }_1^0\tau ^{}\stackrel{~}{\chi }_1^0\pi ^+\nu _\tau \stackrel{~}{\chi }_1^0\pi ^{}\overline{\nu }_\tau \stackrel{~}{\chi }_1^0`$. The distributions have been plotted assuming two scenarios for the stau mixing angle $`\mathrm{\Theta }_{\stackrel{~}{\tau }}=0`$ and $`\pi `$, and $`\stackrel{~}{\chi }_1^0`$ nature (Bino, Higgsino). One can see that the final particle distributions will give access to the tau polarization $`𝒫_\tau `$, to $`\mathrm{tan}\beta `$, and to $`\stackrel{~}{\chi }_1^0`$ nature and through them, to the MSSM parameters.
## 3 Phases in supersymmetry searches
In the MSSM, there are new potential sources of CP non conservation. Complex CP violating phases can arise from several parameters present in the MSSM Lagrangian: the Higgs mixing mass parameter $`\mu `$, the gauginos masses $`M_i`$, the trilinear couplings $`A_i`$. Experimental constraints on these CP violating phases come from the electric dipole moment of the electron and the neutron .
Figure 4 (left side) shows the chargino pair production cross section variation in terms of $`\varphi _\mu `$, the phase associated to the $`\mu `$ parameter, for several values of the sneutrino mass $`m_{\stackrel{~}{\nu }_e}`$ and for a value of $`\mathrm{tan}\beta =1.5`$. One sees that there is a local minimum between the two extreme values ($`\mathrm{\Phi }_\mu =0,\pi `$), which are tested at LEP searches ($`\mu >0,\mu <0`$), but also that it is not so deep as to raise doubts on the exhaustiveness of “phaseless” searches. Further the electric dipole moment experimental upper limit ($`E_e^{exp}<4.310^{27}ecm`$) will constrain these phases (figure on the right side) and rule out many scenarios for which the smallest cross section for chargino pair production is obtained for a $`\varphi _\mu `$ parameter different from $`0`$ and $`\pi `$.
At the linear collider, one cannot neglect this strong dependence of the cross section on phases.
The three-leptons signature from resonant sneutrino production at the LHC
G. MOREAU, E. PEREZ and G. POLESELLO
## 1 Introduction
In extensions of the Minimal Supersymmetric Standard Model (MSSM) where the so-called R-parity symmetry is violated, the superpotential contains some additional trilinear couplings which offer the opportunity to singly produce supersymmetric (SUSY) particles as resonances. The analysis of resonant SUSY particle production allows an easier determination of the these R-parity violating ($`\overline{)}\mathrm{R}_\mathrm{p}`$) couplings than the displaced vertex analysis for the Lightest Supersymmetric Particle (LSP) decay, which is difficult experimentally especially at hadronic colliders.
In this paper, we study the sensitivity provided by the ATLAS detector at the LHC on singly produced charginos via the $`\lambda _{211}^{}`$ coupling, the main contribution coming from the resonant process $`pp\stackrel{~}{\nu }_\mu \stackrel{~}{\chi }_1^\pm \mu ^{}`$. At hadron colliders, due to the continuous energy distribution of the colliding partons, the resonance can be probed over a wide mass range. We have chosen to concentrate on $`\lambda _{ijk}^{}L_iQ_jD_k^c`$ interactions since $`\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c`$ couplings lead to multijet final states with large QCD background. Besides, we focus on $`\lambda _{211}^{}`$ since it corresponds to first generation quarks for the colliding partons and it is not severely constrained by low energy experiments: $`\lambda _{211}^{}<0.09`$ (for $`\stackrel{~}{m}=100`$ GeV) . We consider the cascade decay leading to the three-leptons signature, namely $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0l_p^\pm \nu _p`$ (with $`l_p=e,\mu `$), $`\stackrel{~}{\chi }_1^0\mu u\overline{d},\overline{\mu }\overline{u}d`$. The main motivation lies in the low Standard Model background for this three-leptons final state. The considered branching ratios are typically of order $`B(\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0l_p^\pm \nu _p)22\%`$ (for $`m_{\stackrel{~}{l}},m_{\stackrel{~}{q}},m_{\stackrel{~}{\chi }_2^0}>m_{\stackrel{~}{\chi }_1^\pm }`$) and $`B(\stackrel{~}{\chi }_1^0\mu ud)40\%70\%`$.
## 2 Mass reconstruction
The clean final state, with only two hadronic jets, three leptons and a neutrino allows the reconstruction of the $`\stackrel{~}{\nu }`$ decay chain and the measurement of the $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\nu }_\mu `$ masses. We perform full analysis for the following point of the MSSM: $`M_1=75`$ GeV, $`M_2=150`$ GeV, $`\mu =200`$ GeV, $`\mathrm{tan}\beta =1.5`$, $`A_t=A_b=A_\tau =0`$, $`m_{\stackrel{~}{f}}=300`$ GeV and for $`\lambda _{211}^{}`$=0.09. For this set of MSSM parameters, the mass spectrum is: $`m_{\stackrel{~}{\chi }_1^0}=79.9`$ GeV, $`m_{\stackrel{~}{\chi }_1^\pm }=162.3`$ GeV and the total cross-section for the three-leptons production is 3.1 pb, corresponding to $`100000`$ events for the standard integrated luminosity of 30 fb<sup>-1</sup> expected within the first three years of LHC data taking.
The single chargino production has been calculated analytically and implemented in a version of the SUSYGEN MonteCarlo modified to include the generation of $`pp`$ processes. The generated signal events were processed through the program ATLFAST , a parameterized simulation of the ATLAS detector response.
First, we impose the following loose selection cuts in order to select the considered final state and to reduce the Standard Model (SM) background (see Section 3.1): (a) Exactly three isolated leptons with $`p_T^1>20`$ GeV, $`p_T^{2,3}>10`$ GeV and $`|\eta |<2.5`$, (b) At least two of the three leptons must be muons, (c) Exactly two jets with $`p_T>15`$ GeV, (d) The invariant mass of any $`\mu ^+\mu ^{}`$ pair must lie outside $`\pm 6.5`$ GeV of the $`Z`$ mass.
The three leptons come in the following flavor-sign configurations (+ charge conjugates): (1) $`\mu ^{}e^+\mu ^+`$ (2) $`\mu ^{}e^+\mu ^{}`$ (3) $`\mu ^{}\mu ^+\mu ^+`$ (4) $`\mu ^{}\mu ^+\mu ^{}`$, where the first lepton comes from the $`\stackrel{~}{\nu }_\mu `$, the second one from the $`W`$ and the third one from the $`\stackrel{~}{\chi }_1^0`$ decay. As a starting point for the analysis, we focus on configuration (1) where the muon produced in the $`\stackrel{~}{\chi }_1^0`$ decay is unambiguously identified as the one with the same sign as the electron. The distribution of the $`\mu `$-jet-jet invariant mass exhibits a clear peak over a combinatorial background, shown on the left side of Figure 1. After combinatorial background subtraction (right of Figure 1) an approximately Gaussian peak is left, from which the $`\stackrel{~}{\chi }_1^0`$ mass can be measured with a statistical error of $`100`$ MeV. The combinatorial background is due to events where one jet from $`\stackrel{~}{\chi }_1^0`$ decay is lost and a jet from initial state radiation is used in the combination, and its importance is reduced for heavier sneutrinos or neutralinos. Once the position of the $`\stackrel{~}{\chi }_1^0`$ mass peak is known, the reconstructed $`\stackrel{~}{\chi }_1^0`$ statistics can be increased by also considering signatures (2), (3) and (4), and by choosing as the $`\stackrel{~}{\chi }_1^0`$ candidate the muon-jet-jet combination which gives invariant mass nearest to the peak measured previously using events sample (1). For further reconstruction, we define as $`\stackrel{~}{\chi }_1^0`$ candidates the $`\mu `$-jet-jet combinations with an invariant mass within 12 GeV of the measured $`\stackrel{~}{\chi }_1^0`$ peak, yielding a total statistics of 6750 events for signatures (1) to (4) for an integrated luminosity of 30 fb<sup>-1</sup> .
For $`\stackrel{~}{\chi }_1^\pm `$ reconstruction we consider only configurations (1) and (2), for which the charged lepton from $`W`$ decay is unambiguously identified as the electron. The longitudinal momentum of the neutrino from the $`W`$ decay is calculated from the missing transverse momentum of the event ($`p_T^\nu `$) and by constraining the electron-neutrino invariant mass to the $`W`$ mass. The resulting neutrino longitudinal momentum has a twofold ambiguity. We therefore build the invariant $`W\stackrel{~}{\chi }_1^0`$ mass candidate using both solutions for the $`W`$ boson momentum. The observed peak, represented on the left side of Figure 2, can be fitted with a Gaussian shape with a width of $`6`$ GeV. Only the solution yielding the $`\stackrel{~}{\chi }_1^\pm `$ mass nearer to the measured mass peak is retained, and the $`\stackrel{~}{\chi }_1^\pm `$ candidates are defined as the combinations with an invariant mass within 15 GeV of the peak, corresponding to a statistics of 2700 events.
Finally, the sneutrino mass is reconstructed by taking the invariant mass of the $`\stackrel{~}{\chi }_1^\pm `$ candidate and the leftover muon (Figure 2, right). The $`\stackrel{~}{\nu }`$ mass peak has a width of $`10`$ GeV and 2550 events are counted within 25 GeV of the measured peak.
## 3 Analysis reach
### 3.1 Standard Model background
We consider the following SM processes for the evaluation of the background to the three-leptons signature: (1) $`\overline{t}t`$ production, followed by $`tWb`$, where the two $`W`$ and one of the $`b`$ quarks decay leptonically, (2) $`WZ`$ production, where both bosons decay leptonically, (3) $`Wt`$ production, (4) $`Wbb`$ production, (5) $`Zb`$ production. These backgrounds were generated with the PYTHIA Monte Carlo , and the ONETOP parton level generator , and passed through the ATLFAST package .
We apply to the background events the loose selection cuts described in Section 2, and in addition we reject the three same-sign muons configurations which are never generated by our signal. The background to the sneutrino decay signal is calculated by considering the events with a $`\mu `$-jet-jet invariant mass in an interval of $`\pm 15`$ GeV around the $`\stackrel{~}{\chi }_1^0`$ peak measured for the signal. In order to optimize the signal to background ratio only events containing three muons (configurations (3) and (4)), which are less likely in the Standard Model, are considered. In each event two combinations, corresponding to the two same-sign muons, can be used for the $`\stackrel{~}{\chi }_1^0`$ reconstruction. Both configurations are used when counting the number of events in the peak. In most cases, however, the difference in mass between the two combinations is such that they do not appear in the same peak region.
### 3.2 Supersymmetric background
The pair production of SUSY particles through standard $`R_p`$-conserving processes represents another source of background. A study based on the HERWIG 6.0 MonteCarlo has shown that all the SUSY events surviving the cuts described in Section 3.1 are mainly from $`pp\stackrel{~}{\chi }+X`$ reactions ($`\stackrel{~}{\chi }`$ being either a chargino or a neutralino and $`X`$ any other SUSY particle), and that the SUSY background decreases as the $`\stackrel{~}{\chi }^\pm `$ and $`\stackrel{~}{\chi }^0`$ masses increase. This behavior is due to the combination of two effects: the $`\stackrel{~}{\chi }+X`$ production cross-section decreases with increasing $`\stackrel{~}{\chi }`$ mass, and the probability of losing two of the four jets from the decays of the two $`\stackrel{~}{\chi }_1^0`$ in the event becomes smaller as the $`\stackrel{~}{\chi }^\pm `$ and $`\stackrel{~}{\chi }_1^0`$ masses increase. The SUSY background is only significant for $`\stackrel{~}{\chi }_1^\pm `$ masses lower than $`200`$ GeV.
Besides, it can be assumed that the $`\stackrel{~}{\chi }_1^0`$ mass will be derived from inclusive $`\stackrel{~}{\chi }_1^0`$ reconstruction in SUSY pair production as shown in and . Hence, even in the cases where a significant $`\stackrel{~}{\chi }_1^0`$ peak can not be observed above the SUSY background, we can proceed to the further steps in the kinematic reconstruction. The strong kinematic constraint obtained by requiring both the correct $`\stackrel{~}{\chi }_1^0`$ mass and a peak structure in the $`\stackrel{~}{\chi }_1^0W`$ invariant mass will then allow to separate the single sneutrino production from other SUSY processes.
Therefore, only the Standard Model background is considered in the evaluation of the analysis reach presented below.
### 3.3 Reach in the mSUGRA parameter space
In Figure 3, we show the regions of the $`m_0m_{1/2}`$ plane where the signal significance exceeds 5 $`\sigma `$ ($`\frac{S}{\sqrt{B}}>5`$ with $`S=Signal`$ and $`B=SMBackground`$) after the set of cuts described in Section 3.1 has been applied, within the mSUGRA model. The full mass reconstruction analysis of Section 2 is possible only above the dashed line parallel to the $`m_0`$ axis. Below this line the decay $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0W^\pm `$ is kinematically closed, and the $`W`$ mass constraint can not be applied to reconstruct the neutrino longitudinal momentum.
The basic feature in Figure 3 is a decrease of the sensitivity on $`\lambda _{211}^{}`$ as $`m_0`$ increases. This is due to a decrease of the partonic luminosity as $`m_{\stackrel{~}{\nu }}`$ increases. The sensitivity on $`\lambda _{211}^{}`$ is also observed to decrease as $`m_{\stackrel{~}{\chi }_1^\pm }`$ approaches $`m_{\stackrel{~}{\nu }}`$. There are two reasons. First, in this region the phase space factor of the decay $`\stackrel{~}{\nu }\stackrel{~}{\chi }_1^\pm \mu ^{}`$ following the resonant sneutrino production is suppressed, thus reducing the branching fraction. Secondly, as the $`\stackrel{~}{\nu }_\mu `$ and the $`\stackrel{~}{\chi }_1^\pm `$ become nearly degenerate the muon from the decay becomes on average softer, and its $`p_T`$ can fall below the analysis requirements. In the region $`m_{\stackrel{~}{\chi }_1^\pm }>m_{\stackrel{~}{\nu }}`$, shown as a hatched region in the upper left of the plots, the resonant sneutrino production contribution vanishes and there is essentially no sensitivity to $`\lambda _{211}^{}`$. Finally, the the sensitivity vanishes for low values of $`m_{1/2}`$. This region, below the LEP 200 kinematic limit for $`\stackrel{~}{\chi }_1^\pm `$ detection, corresponds to low values of the $`\stackrel{~}{\chi }_1^0`$ mass. In this situation the two jets from the $`\stackrel{~}{\chi }_1^0`$ decay are soft, and one of them is often below the transverse momentum requirement, or they are reconstructed as a single jet.
For high $`\mathrm{tan}\beta `$, the three-lepton signature is still present, but it may be produced through the decay chain $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\tau }_1\nu _\tau `$, followed by $`\stackrel{~}{\tau }_1\tau \stackrel{~}{\chi }_1^0`$. The full kinematic reconstruction becomes very difficult, but the signal efficiency is essentially unaffected, as long as the mass difference between the lightest $`\stackrel{~}{\tau }`$ and the $`\stackrel{~}{\chi }_1^0`$ is larger than $`50`$ GeV. For a smaller mass difference the charged lepton coming from the $`\tau `$ decay is often rejected by the analysis cuts.
## 4 Conclusion
In conclusion we have shown that if minimal supersymmetry with R-parity violation is realized in Nature, the three-leptons signature from resonant sneutrino production will be a privileged channel for precisely measuring sparticle masses in a model-independent way as well as for testing a broad region of the mSUGRA parameter space.
This signature can lead to a high sensitivity to the $`\lambda _{211}^{}`$ coupling and should also allow to probe an unexplored range of values for many other $`\overline{)}\mathrm{R}_\mathrm{p}`$ couplings of the type $`\lambda _{1jk}^{}`$ and $`\lambda _{2jk}^{}`$.
Resonant slepton production at the LHC
H. DREINER<sup>*</sup><sup>*</sup>*E-mail address: dreiner@v2.rl.ac.uk, P. RICHARDSONE-mail address: p.richardson1@physics.ox.ac.uk, and M. H. SEYMOURE-mail address: M.Seymour@rl.ac.uk
## 1 Introduction
In R-parity violating ($`\overline{)}\mathrm{R}_\mathrm{p}`$) models the single resonant production of charged sleptons in hadron-hadron collisions is possible. The most promising channels for the discovery of these processes, at least with small $`\overline{)}\mathrm{R}_\mathrm{p}`$ couplings, involve the gauge decays of these resonant sleptons. In particular if we consider the production of a charged slepton, this can then decay to give a neutralino and a charged lepton, i.e. the process
$$\mathrm{u}+\overline{\mathrm{d}}\stackrel{~}{\mathrm{}}^+\mathrm{}^++\stackrel{~}{\chi }^0.$$
(1)
In addition to this $`s`$-channel process there are $`t`$-channel processes involving squark exchange. The neutralino decays via the crossed process to give a charged lepton, which due to the Majorana nature of the neutralino can have the same charge as the lepton from the slepton decay. We therefore have a like-sign dilepton signature which we expect to have a low Standard Model background.
## 2 Backgrounds
The dominant Standard Model backgrounds to this process come from
* Gauge boson pair production, i.e. production of ZZ or WZ followed by leptonic decays of the gauge bosons with some of the leptons not being detected.
* $`\mathrm{t}\overline{\mathrm{t}}`$ production. Either the t or $`\overline{\mathrm{t}}`$ decays semi-leptonically, giving one lepton. The second top decays hadronically. A second lepton with the same charge can be produced in a semi-leptonic decay of the bottom hadron formed in the hadronic decay of the second top, i.e.
$`\mathrm{t}`$ $``$ $`\mathrm{W}^+\mathrm{b}\mathrm{e}^+\overline{\nu _\mathrm{e}}\mathrm{b},`$
$`\overline{\mathrm{t}}`$ $``$ $`\mathrm{W}^{}\overline{\mathrm{b}}\mathrm{q}\overline{\mathrm{q}}\overline{\mathrm{b}},\overline{\mathrm{b}}\mathrm{e}^+\overline{\nu _\mathrm{e}}\overline{\mathrm{c}}.`$ (2)
* $`\mathrm{b}\overline{\mathrm{b}}`$ production. If either of these quarks hadronizes to form a $`\mathrm{B}_{\mathrm{d},\mathrm{s}}^0`$ meson this can mix to give a $`\overline{\mathrm{B}}_{\mathrm{d},\mathrm{s}}^0`$. This means that if both the bottom hadrons decay semi-leptonically the leptons will have the same charge as they are both coming from either b or $`\overline{\mathrm{b}}`$ decays.
* Single top production. A single top quark can be produced together with a $`\overline{\mathrm{b}}`$ quark by either an $`s`$\- or $`t`$-channel W exchange. This can then give one charged lepton from the top decay, and a second lepton with the same charge from the decay of the meson formed after the b quark hadronizes.
* Non-physics backgrounds. There are two major sources: (i) from misidentifying the charge of a lepton, e.g. in Drell-Yan production, and (ii) from incorrectly identifying an isolated hadron as a lepton. This means that there is a major source of background from W production with an additional jet faking a lepton.
Early studies of like-sign dileptons at the LHC only studied the backgrounds from heavy quark production. It was found that by imposing cuts on the transverse momentum and isolation of the leptons the heavy quark backgrounds could be significantly reduced. However more recent studies of the like-sign dilepton production at the LHC and the Tevatron suggest that a major source of background to like-sign dilepton production is from gauge boson pair production and from fake leptons. Here we will consider the backgrounds from gauge boson pair production as well as heavy quark production. The study of the non-physics backgrounds (e.g. fake leptons) requires a full simulation of the detector and it is therefore beyond the scope of our study. In particular the background from fake leptons cannot be reliably calculated from Monte Carlo simulations and must be extracted from data . We can use the differences between the $`\overline{)}\mathrm{R}_\mathrm{p}`$ signature we are considering and the MSSM signatures considered in to reduced the background from gauge boson pair production.
We impose the following cuts
* A cut on the transverse momentum of the like-sign leptons $`p_T>40`$ GeV.
* An isolation cut on the like-sign leptons so that the transverse energy in a cone of radius $`R=\sqrt{\mathrm{\Delta }\varphi ^2+\mathrm{\Delta }\eta ^2}=0.4`$ about the direction of each lepton is less than $`5`$ GeV.
* A cut on the transverse mass, $`M_T^2=2|p_T_{\mathrm{}}||p_{T_\nu }|(1\mathrm{cos}\mathrm{\Delta }\varphi _\mathrm{}\nu )`$, where $`p_T_{\mathrm{}}`$ is the transverse momentum of the charged lepton, $`p_{T_\nu }`$ is the transverse momentum of the neutrino, assumed to be all the missing transverse momentum in the event, and $`\mathrm{\Delta }\varphi _\mathrm{}\nu `$ is the azimuthal angle between the lepton and the neutrino, i.e. the missing momentum in the event. We cut out the region where $`60\text{GeV}<M_T<85\text{GeV}`$.
* A veto on the presence of a lepton in the event with the same flavor but opposite charge (OSSF) as either of the leptons in the like-sign pair if the lepton has $`p_T>10`$ GeV and which passes the same isolation cut as the like-sign leptons.
* A cut on the missing transverse energy, $`E_{\mathrm{𝑚𝑖𝑠𝑠}}^T<20`$ GeV .
While these cuts were chosen to reduce the background we have not attempted to optimize them. The first two cuts are designed to reduce the background from heavy quark production. As can be seen in Fig. 1, these cuts reduce this background by several orders of magnitude. The remaining cuts are designed to reduce the background from gauge boson pair, in particular WZ, production which is the major source of background after the imposition of the isolation and $`p_T`$ cuts. The transverse mass cut is designed to remove events with leptonic W decays as can be seen in Fig. 2a. The veto on the presence of OSSF leptons is designed to remove events where one lepton from the dilepton pair comes from the leptonic decay of a Z boson. The missing transverse energy cut again removes events with leptonic W decays, this is mainly to reduce the background from WZ production, as seen in Fig. 2b. The effect of these cuts on the heavy quark and gauge boson pair backgrounds are shown in Figs. 1 and 3, respectively.
The backgrounds from the various processes are summarized in Table 1. The simulations of the $`\mathrm{b}\overline{\mathrm{b}}`$, $`\mathrm{t}\overline{\mathrm{t}}`$ and single top production were performed using HERWIG6.1 . The simulations of gauge boson pair production used PYTHIA6.1 . The major contribution to the background comes from WZ production; the major contribution to the error comes from $`\mathrm{b}\overline{\mathrm{b}}`$. For the $`\mathrm{b}\overline{\mathrm{b}}`$ simulation we have required a parton-level cut of $`40`$ GeV on the transverse momentum of the bottom quarks. This should not affect the results provided we impose a cut of at least $`40`$ GeV on the $`p_T`$ of the leptons. We also forced the B meson produced to decay semi-leptonically. In events where there was one $`\mathrm{B}_{\mathrm{d},\mathrm{s}}^0`$ meson this meson was forced to mix, if there was more than one $`\mathrm{B}_{\mathrm{d},\mathrm{s}}^0`$ then one of the mesons was forced to mix and the others forced to not mix. Even with these cuts it is impossible to simulate the full luminosity with the resources available, due to the large cross section for $`b\overline{b}`$ production. This gives the large error on the estimate of this background.
## 3 Signal
We used HERWIG6.1 to simulate the signal. This version includes the resonant slepton production, including the $`t`$-channel diagrams, and the R-parity violating decay of the neutralino including a matrix element for the decay . We will only consider first generation quarks as the cross sections for processes with higher generation quarks are suppressed by the parton distributions. There are upper bounds on the $`\overline{)}\mathrm{R}_\mathrm{p}`$ couplings from low energy experiments. The bound on $`\lambda _{}^{}{}_{111}{}^{}`$ from neutrino-less double beta decay is very strict so we consider muon production via the coupling $`\lambda _{}^{}{}_{211}{}^{}`$, which has a much weaker bound,
$$\lambda _{}^{}{}_{211}{}^{}<0.059\times \left(\frac{M_{\stackrel{~}{d}_R}}{100\mathrm{G}\mathrm{e}\mathrm{V}}\right),$$
(3)
from the ratio $`R_\pi =\mathrm{\Gamma }(\pi e\nu )/\mathrm{\Gamma }(\pi \mu \nu )`$ .
We have performed a scan in $`M_0`$ using HERWIG with the following SUGRA parameters, $`M_{1/2}=300`$GeV, $`A_0=300`$GeV, $`\mathrm{tan}\beta =2`$, $`\mathrm{sgn}\mu =+`$, and with the $`\overline{)}\mathrm{R}_\mathrm{p}`$ coupling $`\lambda _{}^{}{}_{211}{}^{}=0.01`$. The number of events which pass the cuts given in Section 2 are shown in Fig. 4a, while the efficiency of the cuts, i.e. the fraction of the signal events which have a like-sign dilepton pair passing the cuts, is shown in Fig. 4b. The dip in the efficiency between $`140\text{GeV}<M_0<180\text{GeV}`$ is due to the resonant production of the second lightest neutralino becoming accessible. Just above threshold the efficiency for this channel is low due to the low $`p_T`$ of the lepton produced in the slepton decay.
If we conservatively take a background of $`7.6`$ events, i.e. 1$`\sigma `$ above the central value of our calculation, a 5$`\sigma `$ fluctuation of the background would correspond to 20 events, using Poisson statistics. This is given as a dashed line in Fig. 4a. As can be seen for a large range of values of $`M_0`$ resonant slepton production can be discovered at the LHC, for $`\lambda _{211}^{}=0.01`$. The production cross section depends quadratically on the $`\overline{)}\mathrm{R}_\mathrm{p}`$ Yukawa coupling and hence it should be possible to probe much smaller couplings for small values of $`M_0`$.
As can be seen in Fig. 5, at this SUGRA point the sdown mass varies between $`622`$GeV at $`M_0=50`$GeV and $`784`$GeV at $`M_0=500`$GeV. The corresponding limit on the coupling $`\lambda _{}^{}{}_{211}{}^{}`$ varies between 0.37 and 0.46. We can probe couplings of $`\lambda _{}^{}{}_{211}{}^{}=2\times 10^3`$ for $`M_0=50`$GeV which corresponds to a smuon mass of $`223`$GeV, and at couplings of $`\lambda _{}^{}{}_{211}{}^{}=10^2`$ we can probe values of $`M_0`$ up to $`500`$GeV, i.e. a smuon mass of $`540`$GeV. This is more than an order of magnitude smaller than the current upper bounds on the $`\overline{)}\mathrm{R}_\mathrm{p}`$ coupling given above for these values of $`M_0`$. This is a greater range of couplings and smuon masses than can be probed at the Tevatron . The backgrounds are higher at the LHC but this is compensated by the higher energy and luminosity leading to significantly more signal events.
## 4 Conclusions
We have considered the backgrounds to like-sign dilepton production at the LHC and find a background after cuts of $`5.1\pm 2.5`$ events for an integrated luminosity of $`10fb^1`$. This means, taking a conservative estimate of the background of 7.6 events, that 20 events would correspond to a $`5\sigma `$ discovery. For a full analysis however, non-physics backgrounds must also be considered.
A preliminary study of the signal suggests that an efficiency for detecting the signal in excess of 20% can be achieved over a range of points in SUGRA parameter space. At the SUGRA point studied this means we can probe $`\overline{)}\mathrm{R}_\mathrm{p}`$ couplings of $`2\times 10^3`$ for a smuon mass of $`223`$GeV and up to smuon masses of $`540`$GeV for couplings of $`10^2`$, and higher masses for larger couplings.
A more detailed scan of SUGRA parameter space for this signal remains to be performed.
$`\stackrel{~}{\chi }_1^0`$ reconstruction in mSUGRA models with $`R`$parity breaking $`LQD`$ term
L. MEGNER and G. POLESELLO
## 1 Introduction
The $`R`$parity violating ($`\overline{)}\mathrm{R}_\mathrm{p}`$) extension of the MSSM contains the following additional terms in the superpotential, which are trilinear in the quark and lepton superfields,
$`W_{Rodd}={\displaystyle \underset{i,j,k}{}}\left({\displaystyle \frac{1}{2}}\lambda _{ijk}L_iL_jE_k^c+\lambda _{ijk}^{}L_iQ_jD_k^c+{\displaystyle \frac{1}{2}}\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c\right),`$ (1)
where $`i,j,k`$ are flavor indices. The presence of R-violating terms will mainly manifest itself in two ways: by production of single sparticles, or by the fact that the LSP is not stable and decays to standard model particles. See for a recent review on the subject.
The reconstruction of the $`\stackrel{~}{\chi }_1^0`$ from its decay products in SUSY particle pair production offers the possibility of detecting $`R`$parity violation for values of the $`R`$-violating couplings too small to yield a significant cross-section for single sparticle production. The $`\stackrel{~}{\chi }_1^0`$ decay can be studied for all three coupling types, $`\lambda `$, $`\lambda ^{}`$ and $`\lambda ^{\prime \prime }`$ . The hypothesis of a dominating $`\lambda ^{}`$ is the most favorable case for the kinematic reconstruction of SUSY cascade decays, if the decay of the $`\stackrel{~}{\chi }_1^0`$ into a charged lepton has a sizeable branching fraction. In this case, the charged lepton can be used as the initiator of the mass reconstruction, yielding a reduced combinatorial background with respect to three-jet case, and all three particles from the $`\stackrel{~}{\chi }_1^0`$ decay are reconstructed in the detector, allowing to fully reconstruct the $`\stackrel{~}{\chi }_1^0`$ mass peak. In a pioneering work it was explicitly demonstrated that with appropriate selection criteria on jets and leptons a clear $`\stackrel{~}{\chi }_1^0`$ peak could be observed in the lepton-jet-jet invariant mass over an acceptable combinatorial background. That work was based on $`\stackrel{~}{\chi }_1^0`$ decaying with 100$`\%`$ branching fraction to $`eqq^{}`$, with the momenta of the decay products distributed according to the phase space. The present study is based on a detailed implementation of the $`R`$parity violating decays in the HERWIG 6.0 Monte Carlo including matrix elements for the $`\stackrel{~}{\chi }_1^0`$ decay . We aim to to verify if the $`\stackrel{~}{\chi }_1^0`$ peak can be observed for a broad range of SUSY models, and to develop techniques for evaluating and subtracting the combinatorial background, as a starting point for performing precision measurements of the parameters of the underlying model.
## 2 The SUSY model
We work in the framework of the minimal Supergravity inspired model (mSUGRA) as implemented in the ISASUSY Monte Carlo . The SUSY events were generated with version 6.0 of the HERWIG Monte Carlo , which implements all the $`R`$parity conserving two-to-two SUSY processes. This version of HERWIG also includes the simulation of all $`R`$parity violating sparticle decays . The produced events were then passed through ATLFAST, a fast simulation of the ATLAS detector .
We assumed for this study a single dominant R-violating coupling $`\lambda _{111}^{}`$=0.01. The results obtained are independent of the precise value of the $`\lambda ^{}`$ coupling, as long as it induces a prompt $`\stackrel{~}{\chi }_1^0`$ decay. The results are valid for all the couplings of the type $`\lambda _{ijk}^{}`$ with $`i=1,2`$. They should be taken with some care if $`j`$ or $`k`$ is equal to three, since in this case a $`b`$ quark is among the $`\stackrel{~}{\chi }_1^0`$ decay products.
In models with a single dominant $`\lambda ^{}`$ term $`\stackrel{~}{\chi }_1^0`$ has two decay modes: $`\stackrel{~}{\chi }_1^0l^\pm qq^{}`$ and $`\stackrel{~}{\chi }_1^0\nu q\overline{q}`$. The relative branching ratio of the two modes is a function of the parameters of the model. The mode with a charged lepton is the better mode for reconstructing the $`\stackrel{~}{\chi }_1^0`$ mass, because there is no neutrino which escapes detection. We have verified that the charged lepton mode has a significant branching fraction over all of the mSUGRA parameters space allowed by present experimental searches.
## 3 $`\stackrel{~}{\chi }_1^0`$ reconstruction in mSUGRA models
### 3.1 Introduction
The possibility to extract a $`\stackrel{~}{\chi }_1^0`$ mass peak is determined by:
* the number of produced $`\stackrel{~}{\chi }_1^0`$ for which the decay $`\stackrel{~}{\chi }_1^0l`$-jet-jet with $`l=e,\mu `$ can be completely reconstructed in the detector;
* the combinatorial background, both from SUSY events and from Standard Model processes.
These factors depend on the SUSY model considered, in particular the $`\stackrel{~}{\chi }_1^0`$ mass, the squark and gluino masses, and the jet and lepton multiplicity in SUSY events. Our aim is to determine if the reconstruction can be performed for the full parameter space. For this reason we chose to perform a detailed study for a few representative points in parameter space, trying to span a range as broad as possible of the parameters listed above. A convenient choice are the models already studied in detail for R-conserving mSUGRA , which present very varied phenomenologies, and also give the possibility to cross-check the results of our analysis with previous detailed studies. We have therefore selected points 1, 4 and 5, and we have added a low mass point, which we call 7, in a region with low branching fraction ($`27\%`$) for the decay $`\stackrel{~}{\chi }_1^0l^\pm qq^{}`$. The mSUGRA parameters of the studied points are given in Table 1.
### 3.2 Selection criteria for jets
In order to extract a significant $`\stackrel{~}{\chi }_1^0`$ peak the selection criteria need to be optimized separately for each addressed model. At this level, the important parameter is not the purity of the obtained $`\stackrel{~}{\chi }_1^0`$ peak, but its significance, and the possibility of precisely measuring the peak position.
An initial event selection was first made in order to reject the Standard model background. The events were required to have:
* at least 6 jets with p$`{}_{\mathrm{t}}{}^{}>`$ 15 GeV
* at least one lepton with p$`{}_{\mathrm{t}}{}^{}>`$ 20 GeV in point 4 and point 7, and at least two leptons with p$`{}_{\mathrm{t}}{}^{}>`$ 10 GeV in point 1 and point 5.
The reason for requiring two leptons in point 1 and 5 is that a stronger $`t\overline{t}`$ background rejection is needed, since the SUSY cross-section for these points is lower. These are preliminary cuts and will be further optimized in the final analysis . For the $`\stackrel{~}{\chi }_1^0`$ reconstruction the jets with the highest momentum were not considered, since there are likely to come from the start of the decay chain. The 3rd to the 8th jet in decreasing momentum scale were considered in point 1 and 5 and the 5th to the 8th jets in point 4 and 7. The $`\stackrel{~}{\chi }_1^0`$ was reconstructed starting from an identified electron, and the invariant mass was calculated for the electron-jet1-jet2 combinations such that $`\mathrm{\Delta }R`$(electron-j1)$`<2`$, $`\mathrm{\Delta }R`$(electron-j2)$`<2`$, $`\mathrm{\Delta }R`$(j1-j2)$`<2`$ with $`\mathrm{\Delta }R=\sqrt{\mathrm{\Delta }\eta ^2+\mathrm{\Delta }\varphi ^2}`$, where $`\eta `$ is the pseudorapidity, and $`\varphi `$ the angle in the plane perpendicular to the beam axis. In figure 1 the invariant mass distribution obtained in this way is shown in solid lines for the four points. A statistically significant peak, corresponding to the $`\stackrel{~}{\chi }_1^0`$ mass is seen in all four plots. The hatched histograms superimposed on the plots are the $`t\overline{t}`$ background. This important background can easily be evaluated and subtracted by performing the $`\stackrel{~}{\chi }_1^0`$ reconstruction procedure using the wrong flavor of the lepton as described in the ‘wrong flavor’ technique below. In the next section we concentrate on the more difficult combinatorics from SUSY events.
### 3.3 Subtraction of combinatorial background
In order to estimate the shape of the combinatorial background two different techniques were combined.
In the ‘wrong flavor’ technique, the combinatorial background is evaluated using all of the lepton-jet-jet combinations satisfying the cuts for the $`\stackrel{~}{\chi }_1^0`$ candidate, but using the leptons with the wrong flavor, i.e. electrons if the $`\stackrel{~}{\chi }_1^0`$ decays into muons and vice versa.
In the ‘mixed events’ technique, the idea is to construct lepton-jet-jet combinations and make sure that not all objects in such a combination come from the same $`\stackrel{~}{\chi }_1^0`$, i.e. the combination does not contribute to the signal. In order to make sure that this is the case we use two different events (respectively 1 and 2). The directions ($`\eta `$,$`\varphi `$) of the three objects (lepton, jet, jet) and the momenta of two of the objects (lepton-jet or jet-jet) are taken from one of the candidate $`\stackrel{~}{\chi }_1^0`$ combinations in event 1. The momentum of the third object is taken from a combination in event 2.
These two approaches have complementary strengths and weaknesses. The ‘wrong flavor’ approach correctly subtracts all of the combinatorial background not involving one of the actual leptons from $`\stackrel{~}{\chi }_1^0`$ decay, whereas the ‘mixed events’ background approximately accounts for events where only two of the three particles come from a $`\stackrel{~}{\chi }_1^0`$ decay.
In order to combine the strengths of the two methods the background is estimated through a linear combination of the backgrounds calculated with the two previous methods. The relative weight of the two components is a function both of the characteristics of the SUSY model and of the applied cuts. In each case it is therefore necessary to optimize the parameters of the linear combination. A prescription which was found to give reasonable results is:
* Normalize the backgrounds estimated with the wrong lepton technique and with mixing events technique to the measured mass spectrum separately. This normalization is done in a region outside of the $`\stackrel{~}{\chi }_1^0`$ peak.
* Evaluate the relative weight of the two different background components by performing a least-square fit to the observed $`m_{ljj}`$ mass distribution outside of the peak.
The relative weight factors obtained exhibit some sensitivity to the choice of the mass windows over which the least square fit is performed. The ‘real’ background shape is best reproduced when the lower limit of the mass window is as near as possible to the peak position. Thus the estimate of the peak purity is affected by a systematic uncertainty, whereas the peak position has been checked to be independent of the choice of the subtraction technique. The sum of the background calculated with this combined method and the $`t\overline{t}`$ background is shown for all the four points as dashed lines in figure 1. In all four cases the distribution for the combinatorial background is nicely reproduced.
We show in Figure 2 the results for Point 5 in order to illustrate the power of this method. In the upper plot the invariant mass distribution for all the accepted electron-jet-jet combinations in SUSY events (solid line) is shown. The combinatorial background evaluated with the ‘combined’ method is superimposed as a hatched histogram. The lower left plot shows the signal peak after background subtraction (full line); the actual signal from $`\stackrel{~}{\chi }_1^0`$ decay is superimposed (dashed line). Likewise, the full line in the lower right plot is the background distribution estimate obtained with the ‘combined method’, and the dashed line is the mass distribution for actual background combinatorics. Distributions showing the same excellent agreement are obtained for all four considered models.
## 4 Conclusions
We have performed a reconstruction of $`\stackrel{~}{\chi }_1^0`$ for the decay $`\stackrel{~}{\chi }_1^0l^\pm qq^{}`$ for different scenarios within the mSUGRA model with an $`R`$parity violating $`\lambda ^{}`$ term. We have shown that it is possible to define simple cuts on the jet and lepton multiplicity and topology which allow to observe a statistically significant peak in the lepton-jet-jet invariant mass distribution, at the position of the $`\stackrel{~}{\chi }_1^0`$ mass. We have discussed a method to subtract the combinatorial background. Using this subtraction the $`\stackrel{~}{\chi }_1^0`$ mass and the number of reconstructed $`\stackrel{~}{\chi }_1^0`$s can be precisely measured. This technique works effectively for all four points considered in this analysis. The mass of the reconstructed $`\stackrel{~}{\chi }_1^0`$ can be used as a starting point for the reconstruction of particles further up in the decay chain.
Supersymmetry with $`R`$ parity violation at the LHC: discovery potential from single top production
P. CHIAPPETTA, A. DEANDREA, E. NAGY, S. NEGRONI, G. POLESELLO and J.M. VIREY
The feasibility of single top quark production via squark and slepton exchanges to probe several combinations of $`R`$ parity violating couplings at hadron colliders has been studied . According to those studies, the LHC is better at probing the $`B`$ violating couplings $`\lambda ^{\prime \prime }`$ whereas the Tevatron and the LHC have a similar sensitivity to $`\lambda ^{}`$ couplings. We perform a complete and detailed study including all signal channels using a Monte Carlo generator based on Pythia 6.1 , taking into account all the backgrounds and including the ATLAS detector response using ATLFAST 2.0 .
The $`R`$-parity violating parts of the Lagrangian that contribute to single top production are:
$$L_\overline{)}R=\lambda _{ijk}^{}\stackrel{~}{e}_L^i\overline{d}_R^ku_L^j\lambda _{ijk}^{^{\prime \prime }}(\stackrel{~}{d}_R^k\overline{u}_L^id_L^j+\stackrel{~}{d}_R^j(\overline{d}_L^k)^cu_L^i)+h.c.$$
(1)
The superscript $`c`$ corresponds to charge conjugation. There are altogether 27 and 9 $`\lambda _{ijk}^{}`$ and $`\lambda _{ijk}^{\prime \prime }`$ Yukawa couplings, respectively. The most suppressed couplings are $`\lambda _{111}^{}`$, $`\lambda _{133}^{}`$, $`\lambda _{112}^{^{\prime \prime }}`$, $`\lambda _{113}^{^{\prime \prime }}`$ (see for detailed up to date reviews of the existing bounds). In order to fix the kinematical variables, the reaction we consider is
$$u_i(p_1)+d_j(p_2)t(p_3)+b(p_4),$$
(2)
the $`p_k`$ being the 4-momenta of the particles and the indices $`i`$ and $`j`$ refer to the generations of the $`u`$ and $`d`$-type quarks.
For valence-valence (VV) or sea-sea (SS) subprocesses, the scalar slepton exchange in the $`\widehat{u}`$-channel is taken into account but appears to be suppressed within our assumptions about the $`\lambda ^{}`$ couplings and sfermion masses. The down type squark exchange in the $`\widehat{s}`$-channel squared amplitude is dominant.
Let us now consider the subprocesses involving valence-sea (VS) quarks. Concerning $`R`$ parity violating terms, slepton exchange in the $`\widehat{s}`$-channel and down type squark exchange in the $`\widehat{u}`$-channel contribute. The dominant terms are the squared amplitude due to $`\stackrel{~}{e}`$ exchange, and for initial quarks of the same generation ($`i=j`$), the interference between $`W`$ and $`\stackrel{~}{d}`$. The result is sensitive to the interference term only if the product of $`\lambda ^{\prime \prime }`$ couplings is large (around $`10^1`$). For subprocesses involving quarks of different generations in the initial state the situation is more complex and all amplitudes have to be taken into account.
We have carried out a feasibility study to detect single top production through $`R`$-parity violation at the LHC by measuring the $`l\nu bb`$ final state using the following procedure.
First, we have implemented the partonic $`22`$ cross sections in the PYTHIA event generator. Providing PYTHIA with the flavors and momenta of the initial partons using a given parton distribution function (p.d.f.)<sup>*</sup><sup>*</sup>* We have used the CTEQ3L p.d.f., complete final states including initial and final state radiations and hadronization are generated.
The events were processed through ATLFAST to simulate the response of the ATLAS detector. In particular the energy of electrons, photons and muons was smeared according to the resolution of the relevant detector element in the pseudorapidity range $`|\eta |<2.5`$. Finally, a simple fixed cone algorithm (of radius $`R=0.4`$) was used to reconstruct the parton jets. The minimum transverse energy of a jet was set at 15 GeV. According to the expected b-tagging performance of the ATLAS detector for low luminosity at the LHC we have assumed a $`60\%`$ b-tag efficiency for a factor 100 of rejection against light jets.
The same procedure was applied to the SM background with the exception that we used besides PYTHIA also the ONETOP event generator.
The integrated luminosity for one year at low luminosity at the LHC is taken to be 10 fb<sup>-1</sup>. In Table 1 we display the total cross section values for different initial parton flavors in the case of exchanged squarks of mass of 600 GeV and of $`R`$-parity conserving width $`\mathrm{\Gamma }_R`$ = 0.5 GeV. We took for all $`\lambda ^{\prime \prime }=10^1`$, which yields a natural width of the squark which is smaller than the experimental resolution. Table 4 contains the same information for slepton exchange ( $`\lambda ^{}=10^1`$, for a slepton of mass of 250 GeV and a width of $`\mathrm{\Gamma }_R`$ = 0.5 GeV). Other processes are not quoted because the small value of the limits of their couplings prevents their detection.
In order to study the dependence of the signal on the mass and the width of the exchanged particle we have fixed the couplings to $`10^1`$ and have chosen three different masses for the exchanged squarks: 300, 600 and 900 GeV, respectively. For each mass value we have chosen two different $`\mathrm{\Gamma }_R`$: 0.5 and 20 GeV, respectively. For the first case $`\mathrm{\Gamma }_{\overline{)}\mathrm{R}_\mathrm{p}}`$ dominates, whereas in the last one, when $`\mathrm{\Gamma }_{tot}\mathrm{\Gamma }_R`$, the single top-production cross section decreases by a factor $`10`$. We have considered here the $`ub`$ parton initial state, since this has the highest cross section value.
In order to study the dependence on the parton initial state we have fixed the mass of the exchanged squark to 600 GeV and its width with $`\mathrm{\Gamma }_R`$ = 0.5 GeV and varied the initial state according to the first line of Table 1.
Finally, for the exchanged sleptons we have studied only one case, namely the $`u\overline{d}`$ initial state with a mass and width of the exchanged slepton of 250 GeV and 0.5 GeV, respectively. In each case we have generated about $`10^5`$ signal events.
The irreducible backgrounds are single top production through a virtual $`W`$ (noted $`W^{}`$), or through $`W`$-gluon fusion. $`W`$-gluon fusion is the dominant process (for a detailed study see ). A $`Wbb`$ final state can be obtained either in direct production or through $`Wt`$ or $`t\overline{t}`$ production. Finally, the reducible background consists of $`W+`$n$`j`$ events where two of the jets are misidentified as $`b`$-jets.
For the $`tb`$ final state first we reconstruct the top quark. The top quark can be reconstructed from the $`W`$ and from one of the $`b`$-quarks in the final state, requiring that their invariant mass satisfy $`150M_{Wb}200`$ GeV. The $`W`$ can be in turn reconstructed from either of the two decay channels: $`Wu\overline{d}`$, $`Wl\nu `$. Here we have considered only the latter case which gives a better signature due to the presence of a high $`p_t`$ lepton and missing energy. The former case suffers from multi-jets event backgrounds. As we have only one neutrino, its longitudinal momentum can be reconstructed by using the $`W`$ and top mass constraints. The procedure used is the following :
\- we select events with two b-jets of $`p_t40`$ GeV, with one lepton of $`p_t25`$ GeV, with $`E_t^{miss}35`$ GeV and with a jet multiplicity $`3`$,
\- we reconstruct the longitudinal component ($`p_z`$) of the neutrino by requiring $`M_{l\nu }`$ = $`M_W`$. This leads to an equation with twofold ambiguity on $`p_z`$.
\- More than 80$`\%`$ of the events have at least one solution for $`p_z`$. In case of two solutions, we calculate $`M_{l\nu b}`$ for each of the two b-jets and we keep the $`p_z`$ that minimizes $`|M_{top}M_{l\nu b}|`$.
\- we keep only events where $`150M_{l\nu b}200`$ GeV.
Next, the reconstructed top quark is combined with the $`b`$ quark not taking part in the top reconstruction. In order to reduce the $`t\overline{t}`$ background to a manageable level, we need to apply a strong jet veto on the third jet by requiring that its $`p_t`$ should be $`20`$ GeV.
The invariant mass distribution of the $`tb`$ final state after the cuts described above is shown in Fig. 1. Once an indication for a signal is found, we count the number of signal ($`N_s`$) and background ($`N_b`$) events in an interval corresponding to 2 standard deviations around the signal peak for an integrated luminosity of 30 fb<sup>-1</sup>. Then we rescale the signal peak by a factor $`\alpha `$ such that $`N_s/\sqrt{N_b}=5`$. By definition the scale-factor $`\alpha `$ determines the limit of sensitivity for the lowest value of the $`\lambda ^{\prime \prime }`$ ($`\lambda ^{}`$) coupling we can test with the LHC: $`\lambda _{ijk}^{\prime \prime }\lambda _{lmn}^{\prime \prime }0.01\sqrt{\alpha }`$.
In Table 2 we show the limits obtained for the combinations of $`\lambda _{132}^{\prime \prime }\lambda _{332}^{\prime \prime }`$ for different masses and widths of the exchanged $`\stackrel{~}{s}`$-quark. Also shown are the current limits assuming a mass for $`\stackrel{~}{m}_f`$ = 100 GeV, the number of signal and background events, as well as the experimentally observable widths of the peak ($`\mathrm{\Gamma }_{exp}`$).
In Table 3 we compile the sensitivity limit of the bilinear combination of the different Yukawa couplings one can obtain after 3 years of LHC run with low luminosity, if the exchanged squark has a mass of 600 GeV.
For the exchanged sleptons (cf Table 4) we have calculated the sensitivity limit of the bilinear combination of the different Yukawa couplings only for the most favorable case, i.e. for the $`u\overline{d}`$ partonic initial state. We obtain 4.63$`\times 10^3`$ for the limits on $`\lambda _{11k}^{}\lambda _{k33}^{}`$ (in comparison with the limit of 2.8$`\times 10^3`$ obtained by Oakes et al.).
For those cases where the exchanged squark (slepton) can be discovered at the LHC we have made an estimate on the precision with which one can determine its mass. For this purpose, we have subtracted the background under the mass peak and fitted a Gaussian curve on the remaining signal. For the assumed value of the coupling constant, the error on the mass determination is dominated by the $`1\%`$ systematic uncertainty on the jet energy scale in ATLAS.
Probing $`\overline{)}\mathrm{R}_\mathrm{p}`$ couplings through indirect effects on Drell-Yan production at the LHC
D. CHOUDHURY and R.M. GODBOLE
Within the Standard Model, electroweak gauge invariance ensures both lepton number and baryon number conservation (at least in the perturbative context). However, this is not so within the Minimal Supersymmetric Standard Model (MSSM). In most standard treatments, such interactions are avoided by the imposition of a discrete symmetry called $`R`$-parity, which implies a conserved multiplicative quantum number, $`R(1)^{3B+L+S}`$, where $`B`$ is baryon number, $`L`$ is lepton number, and $`S`$ is spin . All ordinary particles are $`R`$-parity even, while all superpartners are $`R`$-parity odd. If $`R`$-parity is conserved, superpartners must be produced in pairs and the lightest superpartner, or LSP, is absolutely stable. Unfortunately, $`R`$-parity invariance of the SUSY Lagrangian is an ad hoc assumption and not derived from any known fundamental principle. Hence, it is of interest to consider $`R`$-parity–violating extensions of the MSSM, especially since the experimental signatures of low energy supersymmetry would then be radically different. Curiously, $`\overline{)}\mathrm{R}_\mathrm{p}`$ interactions can improve the agreement between theory and precision electroweak measurements, and also offer explanations for experimental anomalies such as the high-$`Q^2`$ excess reported at HERA .
The most general $`\overline{)}\mathrm{R}_\mathrm{p}`$ terms in the superpotential consistent with Lorentz invariance, gauge symmetry, and supersymmetry are<sup>*</sup><sup>*</sup>*We do not consider here the bilinear terms that mix lepton and Higgs superfields .
$$W_{\mathit{}}=\lambda _{ijk}L_iL_jE_k^c+\lambda _{ijk}^{}L_iQ_jD_k^c+\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c,$$
(1)
where $`L_i`$ and $`Q_i`$ are the $`SU(2)`$ doublet lepton and quark fields and $`E_i^c,U_i^c,D_i^c`$ are the singlet superfields. The $`UDD`$ couplings violate baryon number while each of the other two sets violate lepton number. As our aim is to explore dilepton final states in $`pp`$ collisions, we shall explicitly forbid the $`UDD`$ interactions as an economical way to avoid unacceptably rapid proton decay.
The remarkable agreement between low-energy experimental data and SM expectations imply quite severe bounds on the strength of many $`\overline{)}\mathrm{R}_\mathrm{p}`$ operators . Thus, within the context of colliders it is natural to consider the production of supersymmetric particles to be governed mainly by gauge interactions with $`\overline{)}\mathrm{R}_\mathrm{p}`$ being important only in the subsequent decay. However, such analyses are neither very sensitive to the actual size of the $`\overline{)}\mathrm{R}_\mathrm{p}`$ coupling nor are they particularly useful for the case of heavy superpartners (when the pair-production cross sections are smaller). A complementary approach is afforded by processes like Drell-Yan production where the exchange of a heavy squark is governed by the relevant $`\overline{)}\mathrm{R}_\mathrm{p}`$ coupling. In this note, we refine the analysis of Ref. in the context of the LHC, thereby extending its reach.
Expanding the superfield components in (1), we obtain the interaction Lagrangian that connects quarks to leptons:
$$_{LQD}=\lambda _{ijk}^{}\left\{\stackrel{~}{\nu }_{i\mathrm{L}}\overline{d}_{k\mathrm{R}}d_{j\mathrm{L}}\stackrel{~}{e}_{i\mathrm{L}}\overline{d}_{k\mathrm{R}}u_{j\mathrm{L}}+\stackrel{~}{d}_{j\mathrm{L}}\overline{d}_{k\mathrm{R}}\nu _{i\mathrm{L}}\stackrel{~}{u}_{j\mathrm{L}}\overline{d}_{k\mathrm{R}}e_{i\mathrm{L}}+\stackrel{~}{d}_{k\mathrm{R}}^c\nu _{i\mathrm{L}}d_{j\mathrm{L}}\stackrel{~}{d}_{k\mathrm{R}}^ce_{i\mathrm{L}}u_{j\mathrm{L}}\right\}+\mathrm{h}.\mathrm{c}.$$
(2)
With only one $`\lambda _{ijk}^{}`$ being nonzero, the simplest processes (with an observable final state) to which $`_{LQD}`$ would contribute are
* $`u_j\overline{u}_je_i^{}e_i^+`$ ($`\stackrel{~}{d}_{k\mathrm{R}}`$)
* $`d_k\overline{d}_ke_i^{}e_i^+`$ ($`\stackrel{~}{u}_{j\mathrm{L}}`$)
* $`u_j\overline{d}_je_i^+\nu _i`$ ($`\stackrel{~}{d}_{k\mathrm{R}}`$)
* $`q_j\overline{q}_j^{}q_j\overline{q}_j^{}`$ ($`\stackrel{~}{\nu }_{i\mathrm{L}},e_{i\mathrm{L}}`$)
with the particle in parentheses being exchanged in the $`t`$-channel. If we allow more than one $`\overline{)}\mathrm{R}_\mathrm{p}`$ coupling to be non-zero, more exotic final states would be possible. However, as simultaneous existence of more than one such coupling is disfavored from the data on flavor-changing neutral current processes, we shall not consider this possibility. Furthermore, the limits on $`LQD`$ couplings of muons with the first-generation quarks found in nucleon targets are weaker than those for electrons. Hence, in this note, we shall restrict ourselves to perhaps the most optimistic case, namely the study of a $`\mu ^+\mu ^{}`$ final state in the presence of a non-zero $`\lambda _{211}^{}`$. The corresponding details for the other cases will be presented elsewhere.
The signature we focus on is a $`\mu ^+\mu ^{}`$ without any missing transverse momentum. To ensure detectability, we demand that
$$p_T(\mu ^\pm )50\text{GeV}\mathrm{and}\left|\eta (\mu ^\pm )\right|<3.$$
(3)
The lowest order SM process leading to such a final state is the Drell–Yan mechanism i.e. $`q\overline{q}(\gamma ^{},Z^{})\mu ^+\mu ^{}`$. A non-zero $`\lambda _{211}^{}`$ induces an additional $`t`$–channel diagram and the expressions for the total cross sections can be found in Ref.. QCD corrections to the Drell–Yan process have been calculated to the next-to-leading order and are a function of the c.m. energy ($`\sqrt{s}`$) of the collider, the structure functions used and the subprocess scale $`M`$. The dependence on $`M`$ is marginal though and one may approximate it by a scale–independent $`K`$-factor $`1.12`$ (for the CTEQ4M structure functions that we use). The same factor approximates the K-factor for the dimuon mass distributions also for the range of dimuon masses and rapidities that we use. In the absence of a calculation of the higher-order corrections to the $`\overline{)}\mathrm{R}_\mathrm{p}`$ contribution, we assume that $`K`$-factor for the full theory is the same as that within the SM.
To maximize the sensitivity, we need to consider the differential distributions especially since the event topology of the $`\overline{)}\mathrm{R}_\mathrm{p}`$ “signal” is quite different from that of the SM “background”. A convenient set of independent kinematical variables is given by the dimuon invariant mass $`M`$, the rapidity of the $`\mu ^+\mu ^{}`$ pair $`\eta _{\mathrm{pair}}`$ and the difference of the individual rapidities $`\mathrm{\Delta }\eta `$.
It is obvious that the contribution due to a heavy squark exchange would be non-negligible only for M
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[-0.07cm] $\sim$}}~{}m_{\rm sq}. Thus the relative deviations are pronounced only for large values of $`M`$ (see Fig.1). It has been demonstrated that the parton luminosities can be measured with high accuracy at the LHC. Hence a change in the absolute value of the cross-section can indeed be used to look at the sensitivity of the signal to $`\overline{)}\mathrm{R}_\mathrm{p}`$ couplings. As a matter of fact such an analysis of the differential distribution in the dimuon mass of the dimuon pair cross-section has been used by the D0 collaboration to put limits on the Quark-Electron Compositeness scale, using the Tevatron data. Recall that the analysis here does include the effect of the higher order corrections on the SM contribution to the $`\mu ^+\mu ^{}`$ pairs.
Note that $`\mathrm{\Delta }\eta (=\eta _+\eta _{}`$), the difference of the lepton rapidities is directly related to the scattering angle in the subprocess center-of-mass frame. Hence one expects to see the difference between an $`s`$–channel and a $`t`$–channel process in this distribution. This is borne out strikingly in the second of Fig. 1. The symmetry about $`\mathrm{\Delta }\eta =0`$ is dictated by the fact that the initial state is symmetric. And finally, one considers $`\eta _{\mathrm{pair}}`$ which is a measure of the Lorentz boost of the center-of-mass frame, or in other words, the mismatch of the quark and antiquark energies.
To quantify our comparison of the differential distributions in the two cases (SM vs. $`\overline{)}\mathrm{R}_\mathrm{p}`$), we devise a $`\chi ^2`$ test. We divide the $`M`$$`\eta _{\mathrm{pair}}`$$`\mathrm{\Delta }\eta `$ hypervolume into equal sized bins and compare the number of signal ($`N_n^{\mathrm{SM}+\mathit{}}`$) and background ($`N_n^{\mathrm{SM}}`$) events in each bin for a given integrated luminosity. We then define a $`\chi ^2`$ test of discrimination
$$\chi ^2=\underset{n}{}\frac{(N_n^{\mathrm{SM}+\mathit{}}N_n^{\mathrm{SM}})^2}{N_n^{\mathrm{SM}}+(ϵN_n^{\mathrm{SM}})^2}$$
(4)
where $`ϵ`$ is a measure of the systematic errors, arising mainly from the uncertainties in luminosity measurement and quark densities. To be specific, we use a uniform grid of ($`20\text{GeV},0.15,0.15`$) and choose two representative values for $`ϵ`$, namely $`5\%`$ and $`15\%`$.
In Fig. 2, we present the exclusion plots in the ($`m_{\mathrm{sq}},\lambda _{211}^{}`$) plane that can be achieved at the LHC for a given luminosity. The interpretation is simple. The area above the respective curves can be ruled out at 95% C.L. As expected, for $`ϵ=0.15`$ we do slightly worse than for the case of $`ϵ=0.05`$. And similarly, we do better with larger luminosity. What is most interesting, however, is that we do significantly better than the bounds obtained from low-energy experiments, viz.
$$\lambda _{211}^{}<0.059(100\text{GeV}/m_{\stackrel{~}{d}_R}).$$
The LHC bound clearly does not scale as the mass of the squark. This is easily understood from a perusal of the first of Figs. 1. For smaller squark masses, the $`\overline{)}\mathrm{R}_\mathrm{p}`$ contribution is also peaked at smaller values of $`M`$. However, the SM contribution is also much bigger at such $`M`$. Consequently, the relative deviations are smaller and so are the contributions to the $`\chi ^2`$. For larger squark masses, the $`\overline{)}\mathrm{R}_\mathrm{p}`$ contributions concentrate at larger $`M`$ and hence the relative deviations are larger too. This is also the reason why our analysis shows much more sensitivity to the $`\overline{)}\mathrm{R}_\mathrm{p}`$ couplings than the analysis of Ref. .
It is clear that a similar analysis can also be performed to study the sensitivity of LHC experiments to effects of Leptoquarks/compositeness etc. which would give rise to four-fermion currents with unusual chiral structure.
Neutrino masses, R-parity violating supersymmetry and collider signals
A.K. DATTA, B. MUKHOPADHYAYA, S. ROY and F. VISSANI
In a scenario with R-parity violating supersymmetry one can have masses for the neutrinos. The seed of this lies in the superpotential, where, upon admission of lepton number violation, one can write the following terms over and above those of the minimal SUSY standard model (MSSM):
$$W_{\mathit{}}=\lambda _{ijk}L_iL_jE_k^c+\lambda _{ijk}^{}L_iQ_jD_k^c+\lambda _{ijk}^{\prime \prime }U_i^cD_j^cD_k^c+ϵ_iL_iH_2$$
(1)
If one considers the framework where R-parity is violated through a term of the type $`ϵ_iL_iH_2`$, then one neutrino acquires a mass at the tree level through mixing with neutralinos. We have argued that this may give rise to the mass splitting corresponding to the $`\nu _\mu \nu _\tau `$ oscillation, which provides an explanation of the atmospheric muon neutrino deficit recorded by the superkamiokande (SK) experiment . The other two neutrinos remain massless at the tree level. However, the $`\lambda `$-and $`\lambda ^{^{}}`$-type terms can at the same time induce mass terms for all the three generations at the one-loop level, which can be responsible for a smaller mass splitting between the the first two neutrino generation, thereby explaining the solar neutrino problem.
The relevant parameters in this scenario, in addition to those contained in the MSSM, are the $`ϵ_i`$’s mentioned above and vacuum expectation values (vev) for the scalar neutrinos, which are unavoidable consequences of the former. These can be lumped into one ‘basis-independent’ parameter, namely, the vev of the sneutrino corresponding to the state which acquires the tree-level mass. The SK results restrict this vev to be of the order of 100 keV. It is also possible to constrain the parameter space of the soft bilinear SUSY breaking terms of this theory from large angle mixing and neutrino mass-squared difference as demanded by the atmospheric $`\nu _\mu `$ data. In addition, non-trivial restrictions on the hierarchy of the $`ϵ`$-parameters of different flavors or between the $`ϵ`$’s and $`\mu `$ follow if we demand that flavor-changing neutral currents be suppressed . Moreover, if one wants to solve the solar neutrino problem in this scenario, either through MSW solution or vacuum oscillation, then the $`\lambda `$ and $`\lambda ^{^{}}`$ couplings also get restricted to values $`10^5`$ for the MSW solution, and to much smaller values for the ‘just-so’ solution.
Apart from the neutrino-neutralino mixing mentioned above, such a scenario induces mixing in the chargino-charged lepton sector as well. These mixings result in some couplings like $`\stackrel{~}{\chi }_1^0\tau W`$, $`\stackrel{~}{\chi }_1^0\nu Z`$ etc., which are characteristic of bilinear R-parity violation.
If the lightest neutralino ($`\stackrel{~}{\chi }_1^0`$) is heavier than the $`Z`$, then we have the following additional decay channels of $`\stackrel{~}{\chi }_1^0`$ :
$$\stackrel{~}{\chi }_1^0\nu _lZ,l=e,\mu ,\tau $$
(2)
$$\stackrel{~}{\chi }_1^0lW,l=e,\mu ,\tau $$
(3)
This kind of decay will not be possible with only the trilinear R-violating interactions. Also, if the $`\lambda `$ and $`\lambda ^{^{}}`$ terms have to be responsible for loop-induced solutions to the solar neutrino problem, then they are restricted to such values that three-body decays of the lightest neutralino triggered by them are dominated by the two-body decays as and when they are allowed. In such a case, demanding maximal mixing between the second and third generations (as required by the SK data), one would expect to get equal numbers of muons and taus in decays in the channel $`\stackrel{~}{\chi }_1^0lW`$ . This can be a remarkable feature in collider searches of the lightest neutralino. Another important observation is that, thanks to the small R-violating couplings necessitated by the small tree-level neutrino mass, the decay length of the lightest neutralino can be as large as 1 - 10 mm or so. This implies a decay gap that can be observed in collider experiments.
Among other signals that have been investigated , the signal in the form of *equal* number of *like sign* $`\mu `$ and $`\tau `$ events accompanied by *like sign on-shell* $`W`$’s may be seen at the upgraded Tevatron, LHC or future generation $`e^+e^{}`$ colliders when a pair of lightest neutralinos ($`\stackrel{~}{\chi }_1^0`$) is produced directly or via cascades. Each neutralino may decay directly to $`\mu `$ or $`\tau `$ (with BR $`40\%`$) along with *on-shell* $`W`$’s. We emphasize that this is a consequence of neutralino-neutrino mixing triggered exclusively by the bilinear R-parity violating terms in the superpotential. A very interesting possibility, essentially due to the Majorana nature of neutralinos, is that of having like-sign dilepton events accompanied by two W’s of identical charge . Of course, for such a signal to be distinct, one requires the associated $`W`$’s to be *on-shell*. In other words, such a signal is typical of a situation where $`\chi _1^0`$ is heavier than the W. Otherwise, 3-body decays of the lightest neutralino take place, and these can have contributions from trilinear R-violating couplings as well.
The interesting feature is that there is almost no Standard Model background for such a signal. Taming the MSSM background, too, is unlikely to pose any serious difficulty, so long as the W’s can be identified . One may further exploit the large decay gap ($``$ few mm) in the lightest neutralino decay for strengthening the signal.
Introduction to GMSB phenomenology at TeV colliders
S. AMBROSANIO<sup>*</sup><sup>*</sup>*E-mail: ambros@mail.cern.ch
## 1 Introduction to Gauge-Mediated SUSY Breaking
Since no superpartners have been detected at collider experiments so far, supersymmetry (SUSY) cannot be an exact symmetry of Nature. The requirement of “soft” supersymmetry breaking alone is not sufficient to reduce the free parameters to a number suitable for predictive phenomenological studies. Hence, motivated theoretical hypotheses on the nature of SUSY breaking and the mechanism through which it is transmitted to the visible sector of the theory \[here assumed to be the one predicted by the minimal SUSY extension of the standard model (MSSM)\] are highly desirable. If SUSY is broken at energies of the order of the Planck mass and the SUSY breaking sector communicates with the MSSM sector through gravitational interactions only, one falls in the supergravity-inspired (SUGRA) scheme. The most recognized alternative to SUGRA is based instead on the hypothesis that SUSY breaking occurs at relatively low energy scales and is mediated mainly by gauge interactions (GMSB) . A good theoretical reason to consider such a possibility is that it provides a natural, automatic suppression of the SUSY contributions to flavor-changing neutral currents and CP-violating processes. A pleasant consequence is that, at least in the simplest versions of GMSB, the MSSM spectrum and other observables depend on just a handful of parameters, typically
$$M_{\mathrm{mess}},N_{\mathrm{mess}},\mathrm{\Lambda },\mathrm{tan}\beta ,\mathrm{sign}(\mu ),$$
(1)
where $`M_{\mathrm{mess}}`$ is the overall messenger scale; $`N_{\mathrm{mess}}`$ is the so-called messenger index, parameterizing the structure of the messenger sector; $`\mathrm{\Lambda }`$ is the universal soft SUSY breaking scale felt by the low-energy sector; $`\mathrm{tan}\beta `$ is the ratio of the vacuum expectation values of the two Higgs doublets; sign($`\mu `$) is the ambiguity left for the SUSY higgsino mass after conditions for correct electroweak symmetry breaking (EWSB) are imposed (see e.g. Refs. ).
The phenomenology of GMSB (and, more generally, of any theory with low-energy SUSY breaking) is characterized by the presence of a very light gravitino $`\stackrel{~}{G}`$ ,
$$m_{3/2}m_{\stackrel{~}{G}}=\frac{F}{\sqrt{3}M_P^{}}\left(\frac{\sqrt{F}}{100\mathrm{TeV}}\right)^22.37\mathrm{eV},$$
(2)
where $`\sqrt{F}`$ is the fundamental scale of SUSY breaking, 100 TeV is a typical value for it, and $`M_P^{}=2.44\times 10^{18}`$ GeV is the reduced Planck mass. Hence, the $`\stackrel{~}{G}`$ is always the lightest SUSY particle (LSP) in these theories. If $`R`$-parity is assumed to be conserved, any produced MSSM particle will finally decay into the gravitino. Depending on $`\sqrt{F}`$, the interactions of the gravitino, although much weaker than gauge and Yukawa interactions, can still be strong enough to be of relevance for collider physics. As a result, in most cases the last step of any SUSY decay chain is the decay of the next-to-lightest SUSY particle (NLSP), which can occur outside or inside a typical detector or even close to the interaction point. The pattern of the resulting spectacular signatures is determined by the identity of the NLSP and its lifetime before decaying into the $`\stackrel{~}{G}`$,
$$c\tau _{\mathrm{NLSP}}\frac{1}{100}\left(\frac{\sqrt{F}}{100\mathrm{TeV}}\right)^4\left(\frac{m_{\mathrm{NLSP}}}{100\mathrm{GeV}}\right)^5,$$
(3)
where $``$ is a number of order unity depending on the nature of the NLSP.
The identity of the NLSP \[or, to be more precise, the identity of the sparticle(s) having a large branching ratio (BR) for decaying into the gravitino and the relevant SM partner\] determines four main scenarios giving rise to qualitatively different phenomenology:
Occurs whenever $`m_{\stackrel{~}{N}_1}<(m_{\stackrel{~}{\tau }_1}m_\tau )`$. Here typically a decay of the $`\stackrel{~}{N}_1`$ to $`\stackrel{~}{G}\gamma `$ is the final step of decay chains following any SUSY production process. As a consequence, the main inclusive signature at colliders is prompt or displaced photon pairs + X + missing energy. $`\stackrel{~}{N}_1`$ decays to $`\stackrel{~}{G}Z^0`$ and other minor channels may also be relevant at TeV colliders.
Defined by $`m_{\stackrel{~}{\tau }_1}<\mathrm{Min}[m_{\stackrel{~}{N}_1},m_{\stackrel{~}{\mathrm{}}_R}]m_\tau `$, features $`\stackrel{~}{\tau }_1\stackrel{~}{G}\tau `$ decays, producing $`\tau `$ pairs or charged semi-stable $`\stackrel{~}{\tau }_1`$ tracks or decay kinks + X + missing energy. Here and in the following, $`\mathrm{}`$ stands for $`e`$ or $`\mu `$.
When $`m_{\stackrel{~}{\mathrm{}}_R}<\mathrm{Min}[m_{\stackrel{~}{N}_1},m_{\stackrel{~}{\tau }_1}+m_\tau ]`$, $`\stackrel{~}{\mathrm{}}_R\stackrel{~}{G}\mathrm{}`$ decays are also open with large BR. In addition to the signatures of the stau NLSP scenario, one also gets $`\mathrm{}^+\mathrm{}^{}`$ pairs or $`\stackrel{~}{\mathrm{}}_R`$ tracks or decay kinks.
If $`|m_{\stackrel{~}{\tau }_1}m_{\stackrel{~}{N}_1}|<m_\tau `$ and $`m_{\stackrel{~}{N}_1}<m_{\stackrel{~}{\mathrm{}}_R}`$, both signatures of the neutralino NLSP and stau NLSP scenario are present at the same time, since $`\stackrel{~}{N}_1\stackrel{~}{\tau }_1`$ 2–body decays are not allowed by phase space.
Note that in the GMSB parameter space the relation $`m_{\stackrel{~}{\mathrm{}}_R}>m_{\stackrel{~}{\tau }_1}`$ always holds. Also, one should keep in mind that the classification above is only valid as an indicative scheme in the limit $`m_e`$, $`m_\mu 0`$, neglecting also those cases where a fine-tuned choice of $`\sqrt{F}`$ and the sparticle masses may give rise to competition between phase-space suppressed decay channels from one ordinary sparticle to another and sparticle decays to the gravitino .
In this report, two important aspects of the GMSB phenomenology at TeV colliders will be treated:
The consequences of the GMSB hypothesis on the light Higgs spectrum using the most accurate tools available today for model generation and $`m_h`$ calculation;
Studies and possible measurements at the LHC with the ATLAS detector in the stau NLSP or slepton co-NLSP scenarios, with focus on determining the fundamental SUSY breaking scale $`\sqrt{F}`$.
For this purpose, we generated about 30000 GMSB models under well defined hypotheses, using the program SUSYFIRE , as described in the following section.
## 2 GMSB Models
In the GMSB framework, the pattern of the MSSM spectrum is simple, as all sparticle masses are generated in the same way and scale approximately with a single parameter $`\mathrm{\Lambda }`$, which sets the amount of soft SUSY breaking felt by the visible sector. As a consequence, scalar and gaugino masses are related to each other at a high energy scale, which is not the case in other SUSY frameworks, e.g. SUGRA. Also, it is possible to impose other conditions at a lower scale to achieve EWSB and further reduce the dimension of the parameter space.
To build our GMSB models, we adopt the usual phenomenological approach, in particular following Ref. , where problems relevant for GMSB physics at TeV colliders were also approached. We do not specify the origin of the SUSY higgsino mass $`\mu `$, nor do we assume that the analogous soft SUSY breaking parameter $`B\mu `$ vanishes at the messenger scale. Instead, we impose correct EWSB to trade $`\mu `$ and $`B\mu `$ for $`M_Z`$ and $`\mathrm{tan}\beta `$, leaving the sign of $`\mu `$ undetermined. However, we are aware that to build a satisfactory GMSB model one should also solve the latter problem in a more fundamental way, perhaps by providing a dynamical mechanism to generate $`\mu `$ and $`B\mu `$, possibly with values of the same order of magnitude. This might be accomplished radiatively through some new interactions. However, in this case the other soft terms in the Higgs potential, namely $`m_{H_{1,2}}^2`$, will be also affected and this will in turn change the values of $`|\mu |`$ and $`B\mu `$ coming from EWSB conditions . Within the study (A), we are currently considering some “non-minimal” possibilities for GMSB models that to some extent take this problem into account, and we are trying to assess the impact on the light Higgs mass. We do not treat this topic here, but refer to for further details.
To determine the MSSM spectrum and low-energy parameters, we solve the renormalization group (RG) evolution with the following boundary conditions at the $`M_{\mathrm{mess}}`$ scale,
$`M_a`$ $`=`$ $`N_{\mathrm{mess}}\mathrm{\Lambda }g\left({\displaystyle \frac{\mathrm{\Lambda }}{M_{\mathrm{mess}}}}\right){\displaystyle \frac{\alpha _a}{4\pi }},(a=1,2,3)`$
$`\stackrel{~}{m}^2`$ $`=`$ $`2N_{\mathrm{mess}}\mathrm{\Lambda }^2f\left({\displaystyle \frac{\mathrm{\Lambda }}{M_{\mathrm{mess}}}}\right){\displaystyle \underset{a}{}}\left({\displaystyle \frac{\alpha _a}{4\pi }}\right)^2C_a,`$ (4)
respectively for the gaugino and the scalar masses. In Eq. (4), $`g`$ and $`f`$ are the one-loop and two-loop functions whose exact expressions can be found e.g. in Ref. , and $`C_a`$ are the quadratic Casimir invariants for the scalar fields. As usual, the scalar trilinear couplings $`A_f`$ are assumed to vanish at the messenger scale, as suggested by the fact that they (and not their squares) are generated via gauge interactions with the messenger fields at the two loop-level only.
To single out the interesting region of the GMSB parameter space, we proceed as follows. Barring the case where a neutralino is the NLSP and decays outside the detector (large $`\sqrt{F}`$), the GMSB signatures are very spectacular and are generally free from SM backgrounds. Keeping this in mind and being interested in GMSB phenomenology at future TeV colliders, we consider only models where the NLSP mass is larger than 100 GeV, assuming that searches at LEP and the Tevatron, if unsuccessful, will in the end exclude a softer spectrum in most cases. We require that $`M_{\mathrm{mess}}>1.01\mathrm{\Lambda }`$, to prevent an excess of fine-tuning of the messenger masses, and that the mass of the lightest messenger scalar be at least 10 TeV. We also impose $`M_{\mathrm{mess}}>M_{\mathrm{GUT}}\mathrm{exp}(125/N_{\mathrm{mess}})`$, to ensure perturbativity of gauge interactions up to the GUT scale. Further, we do not consider models with $`M_{\mathrm{mess}}>\mathrm{\hspace{0.33em}10}^5\mathrm{\Lambda }`$. As a result of this and other constraints, the messenger index $`N_{\mathrm{mess}}`$, which we assume to be an integer independent of the gauge group, cannot be larger than 8. To prevent the top Yukawa coupling from blowing up below the GUT scale, we require $`\mathrm{tan}\beta >1.2`$ (and in some cases $`>1.5`$). This is also motivated by the current bounds from SUSY Higgs searches at LEP II . Models with $`\mathrm{tan}\beta >\mathrm{\hspace{0.33em}55}`$ (with a mild dependence on $`\mathrm{\Lambda }`$) are forbidden by the EWSB requirement and typically fail to give $`m_A^2>0`$.
To calculate the NLSP lifetime relevant to our study (B), one needs to specify the value of the fundamental SUSY breaking scale $`\sqrt{F}`$ on a model-by-model basis. Using perturbativity arguments, for each given set of GMSB parameters, it is possible to determine a lower bound according to Ref. ,
$$\sqrt{F}\sqrt{F_{\mathrm{mess}}}\sqrt{\mathrm{\Lambda }M_{\mathrm{mess}}}>\mathrm{\Lambda }.$$
(5)
On the contrary, no solid arguments can be used to set an upper limit on $`\sqrt{F}`$ of relevance for collider physics, although some semi-qualitative cosmological arguments are sometimes evoked.
In order to generate our model samples using SUSYFIRE, we used logarithmic steps for $`\mathrm{\Lambda }`$ (between about 45 TeV/$`N_{\mathrm{mess}}`$ and about 220 TeV/$`\sqrt{N_{\mathrm{mess}}}`$, which corresponds to excluding models with sparticle masses above $`4`$ TeV), $`M_{\mathrm{mess}}/\mathrm{\Lambda }`$ (between 1.01 and $`10^5`$) and $`\mathrm{tan}\beta `$ (between 1.2 and about 60), subject to the constraints described above. SUSYFIRE starts from the values of SM particle masses and gauge couplings at the weak scale and then evolves up to the messenger scale through RGEs. At the messenger scale, it imposes the boundary conditions (4) for the soft sparticle masses and then evolves the full set of RGEs back to the weak scale. The decoupling of each sparticle at the proper threshold is taken into account. Two-loop RGEs are used for gauge couplings, third generation Yukawa couplings and gaugino soft masses. The other RGEs are taken at the one-loop level. At the scale $`\sqrt{m_{\stackrel{~}{t}_1}m_{\stackrel{~}{t}_2}}`$, EWSB conditions are imposed by means of the one-loop effective potential approach, including corrections from stops, sbottoms and staus. The program then evolves up again to $`M_{\mathrm{mess}}`$ and so on. Three or four iterations are usually enough to get a good approximation for the MSSM spectrum.
The light Higgs boson spectrum in GMSB models
S. AMBROSANIO, S. HEINEMEYER and G. WEIGLEIN
## 1 Introduction
Within the MSSM, the masses of the $`𝒞𝒫`$-even neutral Higgs bosons are calculable in terms of the other low-energy parameters. The mass of the lightest Higgs boson, $`m_h`$, has been of particular interest, as it is bounded to be smaller than the $`Z^0`$ boson mass at the tree level. The one-loop results for $`m_h`$ have been supplemented in the last years with the leading two-loop corrections, performed in the renormalization group (RG) approach , in the effective potential approach and most recently in the Feynman-diagrammatic (FD) approach . The two-loop corrections have turned out to be sizeable. They can lower the one-loop results by up to 20%. These calculations predict an upper bound on $`m_h`$ of about $`m_h130`$ GeV for an unconstrained MSSM with $`m_t=175`$ GeV and a common SUSY mass scale $`M_{\mathrm{SUSY}}1`$ TeV.
As discussed in the Introduction , the GMSB scenario provides a relatively simple set of constraints and thus constitutes a very predictive and readily testable realization of the MSSM. The main goal of the present analysis is to study the spectrum of the lightest neutral $`𝒞𝒫`$-even Higgs boson, $`m_h`$, within the GMSB framework. Particular emphasis is given to the maximal value of $`m_h`$ achievable in GMSB after an exhaustive scanning of the parameter space. Our results are discussed in terms of the GMSB constraints on the low-energy parameters and compared to the cases of a SUGRA-inspired or an unconstrained MSSM.
## 2 Calculation of $`m_h`$
We employ the currently most accurate calculation to evaluate $`m_h`$, based on the FD approach as given in Refs. . The most important radiative corrections to $`m_h`$ arise from the top and scalar top sector of the MSSM, with the input parameters $`m_t`$, the masses of the scalar top quarks, $`m_{\stackrel{~}{t}_1}`$, $`m_{\stackrel{~}{t}_2}`$, and the $`\stackrel{~}{t}`$-mixing angle, $`\theta _{\stackrel{~}{t}}`$. Here we adopt the conventions introduced in Ref. . The complete diagrammatic one-loop result has been combined with the dominant two-loop corrections of $`𝒪(\alpha \alpha _s)`$ and with the subdominant corrections of $`𝒪(G_F^2m_t^6)`$ . GMSB models are generated with the program SUSYFIRE, according to the discussion of . For this study, we consider only models with $`\mathrm{tan}\beta >1.5`$ and $`m_A>80`$ GeV . In addition, we always use $`m_t=175`$ GeV. A change of 1 GeV in $`m_t`$ translates roughly into a shift of 1 GeV (with the same sign) in $`m_h`$ as well. Thus, changing $`m_t`$ affects our results on $`m_h`$ in an easily predictable way.
The results of the $`m_h`$ calculation have been implemented in the Fortran code FeynHiggs . This program has been combined with SUSYFIRE, which has been used to calculate the low energy parameters $`m_{\stackrel{~}{t}_1}`$, $`m_{\stackrel{~}{t}_2}`$, $`\theta _{\stackrel{~}{t}}`$, $`\mu `$, $`M_1`$, $`M_2`$, $`m_{\stackrel{~}{g}}`$, $`\mathrm{}`$ for each of the $``$30000 GMSB models generated. These have then been passed to FeynHiggs for the $`m_h`$ evaluation in a coherent way. Indeed, we transfer the $`\overline{\mathrm{MS}}`$ parameters in the SUSYFIRE output to on-shell parameters before feeding them into FeynHiggs. Compact expressions for the relevant transition formulas can be found in Refs. .
Compared to an existing analysis in the GMSB framework , we use a more complete evaluation of $`m_h`$. This leads in particular to smaller values of $`m_h`$ for a given set of input parameters in our analysis. Also, in Ref. although some GMSB scenarios with generalized messenger sectors were considered, the parameter space for the “minimal” case with a unique, integer messenger index $`N_{\mathrm{mess}}=N_1=N_2=N_3`$ was not fully explored. Indeed, $`\mathrm{\Lambda }`$ was in most cases limited to values smaller than 100 TeV and $`M_{\mathrm{mess}}`$ was fixed to $`10^5`$ TeV. Furthermore, partly as a consequence of the above assumptions, the authors did not consider models with $`N_{\mathrm{mess}}>4`$, i.e. their requirements for perturbativity of the MSSM gauge couplings up to the GUT scale were stronger than ours. We will see in the following section that maximal $`m_h`$ values in our analysis are instead obtained for larger values of the messenger scale and the messenger index.
## 3 The Light Higgs Spectrum in GMSB
In the following, we give some results in the form of scatter plots showing the pattern in GMSB for $`m_h`$, $`m_A`$ as well as other low-energy parameters of relevance for the light Higgs spectrum.
In Fig. 1(a), we show the dependence of $`m_h`$ on $`\mathrm{tan}\beta `$, where only models with $`\mathrm{tan}\beta >1.5`$, $`m_A>80`$ GeV and $`m_{\mathrm{NLSP}}>100`$ GeV are considered, while $`m_t`$ is fixed to 175 GeV. The dependence is strong for small $`\mathrm{tan}\beta <\mathrm{\hspace{0.33em}10}`$, while for larger $`\mathrm{tan}\beta `$ the increase of the lightest Higgs mass is rather mild. The maximum values for $`m_h124`$ GeV are achieved for $`\mathrm{tan}\beta >50`$. It should be noted that for very large $`\mathrm{tan}\beta >\mathrm{\hspace{0.33em}52}`$, we also find a few models with relatively small $`m_h<\mathrm{\hspace{0.33em}100}`$ GeV. This is due to the fact that in this case EWSB conditions tend to drive $`m_A`$ toward very small values . This is made visible by the scatter plot in Fig. 1(b), where the pseudoscalar Higgs mass is shown as a function of $`\mathrm{tan}\beta `$. For such small values of $`m_A`$ and for large $`\mathrm{tan}\beta `$, the relation $`m_hm_A`$ holds. Thus small $`m_h`$ values are quite natural in this region of the parameter space. On the other hand, one can see that extremely large values of $`m_A>\mathrm{\hspace{0.33em}2}`$ TeV can only be obtained for small or moderate $`\mathrm{tan}\beta <\mathrm{\hspace{0.33em}10}`$ GeV. A comparison between Fig 1(a) and (b) reveals that the largest $`m_h`$ values $`>\mathrm{\hspace{0.33em}123}`$ GeV correspond in GMSB to $`m_A`$ values in the 300–800 GeV range. Indeed, it has been checked that such large $`m_h`$ values are in general obtained in the FD calculation for $`300<m_A<\mathrm{\hspace{0.33em}1000}`$ GeV, see Ref. .
In Fig. 2, we show the dependence of the lightest Higgs boson mass on the stop mixing parameter $`x_{\mathrm{top}}`$ defined by
$$x_{\mathrm{top}}\frac{A_{\mathrm{top}}\mu /\mathrm{tan}\beta }{m_S},\mathrm{where}m_S=\sqrt{(m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{t}_2}^2)/2}.$$
(1)
For equal soft SUSY breaking parameters in the stop sector with the $`D`$-terms neglected, $`x_{\mathrm{top}}`$ corresponds to the ratio $`X_t/M_S`$ of the off-diagonal and diagonal entries in the stop mixing matrix, see e.g. Ref. .
Maximal $`m_h`$ values are obtained for $`x_{\mathrm{top}}\pm 2`$, a minimum is reached around $`x_{\mathrm{top}}0`$. Thus, for large $`m_h`$ values a large numerator in Eq. (1) is required. From Fig. 3(a), one can see that in GMSB only negative values of $`A_{\mathrm{top}}`$ are allowed at the electroweak scale, as a consequence of the fact that the trilinear couplings are negligible at the messenger scale. Due to the logarithmic dependence of $`m_h`$ on the stop masses, relatively large values of $`|A_{\mathrm{top}}|`$ are needed for large $`m_h`$. In addition, large $`\mathrm{tan}\beta `$ is also required. From Fig. 3(b) one can check that this leads to values of $`x_{\mathrm{top}}0.95`$, which can only be achieved for positive $`\mu `$. Fig. 4(a) shows the dependence of $`A_{\mathrm{top}}`$ on $`\mu `$. Large values of $`|A_{\mathrm{top}}|`$ are only reached for large $`|\mu |`$ values. Therefore maximal $`h`$ masses are obtained for relatively large and positive $`\mu `$, as can be seen in Fig. 4(b).In general, for large values of $`|\mu |`$ and $`\mathrm{tan}\beta `$ the effects of the corrections from the $`b`$$`\stackrel{~}{b}`$ sector can become important, leading to a decrease in $`m_h`$. For the GMSB models under consideration, however, this is not the case as a consequence of the relatively large $`\stackrel{~}{b}`$ masses.
All these arguments about the combination of low energy parameters needed for large $`m_h`$ in GMSB are summarized in Tab. 1. where we report the 10 models in our sample that give rise to the highest $`m_h`$ values. Together with $`m_h`$, Tab. 1 shows the corresponding input GMSB parameters (Eq. (1) in ) as well as the values of the low energy parameters mentioned above.
It is interesting to note that all the models shown in Tab. 1 feature a large messenger index and values of the messenger scale not far from the maximum we allowed while generating GMSB models. We could not construct a single model with $`m_h>\mathrm{\hspace{0.33em}122.5}`$ GeV having $`N_{\mathrm{mess}}<6`$ or $`M_{\mathrm{mess}}<10^5`$ TeV, for $`m_t=175`$ GeV. It is hence worth mentioning here that our choice of imposing $`M_{\mathrm{mess}}/\mathrm{\Lambda }<10^5M_{\mathrm{mess}}<\mathrm{\hspace{0.33em}2}\times 10^{10}`$ GeV does not correspond to any solid theoretical prejudice. On the other hand it is true that $`M_{\mathrm{mess}}>\mathrm{\hspace{0.33em}3}\times 10^8`$ GeV always corresponds to gravitino masses larger than $`1`$ keV, due to Eqs. (2) and (5) in . The latter circumstance might be disfavored by cosmological arguments . A curious consequence is that the GMSB models with the highest $`m_h`$ belong always to the stau NLSP or slepton co-NLSP scenarios.
Note also that restricting ourselves to GMSB models with $`\mathrm{\Lambda }<100`$ TeV, $`M_{\mathrm{mess}}<10^5`$ TeV and $`N_{\mathrm{mess}}4`$, we find a maximal $`m_h`$ value of 122.2 GeV, for $`m_t=175`$ GeV and $`\mathrm{tan}\beta 52`$. This is to be compared with the one-loop result of Ref. , $`m_h`$(max) = 131.7, for $`\mathrm{tan}\beta `$ around 30 (the assumed value of $`m_t`$ is not quoted).
Values for $`m_h`$ slightly larger than those we found here may also arise from non-minimal contributions to the Higgs potential, in connection with a dynamical generation of $`\mu `$ and $`B\mu `$ . A treatment of this problem can be found in Ref. .
One should also keep in mind that our analysis still suffers from uncertainties due to unknown higher order corrections both in the RGEs for GMSB model generation and in the evaluation of $`m_h`$ from low energy parameters. A rough estimate of these effects leads to shifts in $`m_h`$ not larger than 3 to 5 GeV.
## 4 Conclusions
We conclude that in the minimal GMSB framework described above, values of $`m_h>\mathrm{\hspace{0.33em}124.2}`$ GeV are not allowed for $`m_t=175`$ GeV. This is almost 6 GeV smaller than the maximum value for $`m_h`$ one can achieve in the MSSM without any constraints or assumptions about the structure of the theory at high energy scales . On the other hand, the alternative mSUGRA framework allows values of $`m_h`$ that are $`3`$ GeV larger than in GMSB . This makes the GMSB scenario slightly easier to explore via Higgs boson search. This result was expected in the light of the rather strong GMSB requirements, such as the presence of a unique soft SUSY breaking scale, the relative heaviness of the squarks and the gluino compared to non-strongly interacting sparticles, and the fact that the soft SUSY breaking trilinear couplings $`A_f`$ get nonzero values at the electroweak scale only by RGE evolution. Nevertheless, once the whole parameter space is explored, it is not true that mGMSB gives rise to $`m_h`$ values that are considerably smaller than in mSUGRA. Even smaller differences in the maximal $`m_h`$ might be present when considering non-minimal, complex messenger sectors or additional contributions to the Higgs potential . In any case, as for mSUGRA, current LEP II or Tevatron data on Higgs boson searches are far from excluding mGMSB, and the upgraded Tevatron and the LHC will certainly be needed to deeply test any realistic SUSY model.
Measuring the SUSY breaking scale at the LHC in the slepton NLSP scenario of GMSB models
S. AMBROSANIO, B. MELE, S. PETRARCA, G. POLESELLO and A. RIMOLDI
## 1 Introduction
The fundamental scale of SUSY breaking $`\sqrt{F}`$ is perhaps the most important quantity to determine from phenomenology in a SUSY theory. In the mSUGRA framework, the gravitino mass sets the scale of the soft SUSY breaking masses in the MSSM ($`0.11`$ TeV), so that $`\sqrt{F}`$ is typically large $`10^{1011}`$ GeV (Eq. (2) in ). As a consequence, the interactions of the $`\stackrel{~}{G}`$ with the other MSSM particles $`F^1`$ are too weak for the gravitino to be of relevance in collider physics and there is no direct way to access $`\sqrt{F}`$ experimentally. In GMSB theories, the situation is completely different. The soft SUSY breaking scale of the MSSM and the sparticle masses are set by gauge interactions between the messenger and low energy sectors to be $`\alpha _{\mathrm{SM}}\mathrm{\Lambda }`$ (Eq. (4) in ), so that typical $`\mathrm{\Lambda }`$ values are $`10100`$ TeV. On the other hand, $`\sqrt{F}`$ is subject to the lower bound (5) in only, which tells us that values well below $`10^{10}`$ GeV and even as low as several tens of TeV are perfectly reasonable. The $`\stackrel{~}{G}`$ is in this case the LSP and its interactions are strong enough to allow NLSP decays to the $`\stackrel{~}{G}`$ inside a typical detector size. The latter circumstance gives us a chance for extracting $`\sqrt{F}`$ experimentally through a measurement of the NLSP mass and lifetime (Eq. (3) in ).
Furthermore, the possibility of determining $`\sqrt{F}`$ with good precision opens a window on the physics of the SUSY breaking (“secluded”) sector and the way this SUSY breaking is transmitted to the messenger sector. Indeed, the characteristic scale of SUSY breaking felt by the messengers (and hence the MSSM sector) given by $`\sqrt{F_{\mathrm{mess}}}`$ in Eq. (5) can also be determined once the MSSM spectrum is known. By comparing the measured values of $`\sqrt{F}`$ and $`\sqrt{F_{\mathrm{mess}}}`$ it might well be possible to get information on the way the secluded and messenger sectors communicate with each other. For instance, if it turns out that $`\sqrt{F_{\mathrm{mess}}}\sqrt{F}`$, then it is very likely that the communication occurs radiatively and the ratio $`\sqrt{F_{\mathrm{mess}}/F}`$ is given by some loop factor. On the contrary, if the communication occurs via a direct interaction, this ratio is just given by a Yukawa-type coupling constant, with values $`<\mathrm{\hspace{0.33em}1}`$, see Refs. .
An experimental method to determine $`\sqrt{F}`$ at a TeV scale $`e^+e^{}`$ collider through the measurement of the NLSP mass and lifetime was presented in Ref. , in the neutralino NLSP scenario. Here, we are concerned with the same problem, but at a hadron collider, the LHC, and in the stau NLSP or slepton co-NLSP scenarios. These scenarios provide a great opportunity at the LHC, since the characteristic signatures with semi-stable charged tracks are muon-like, but come from massive sleptons with $`\beta `$ significantly smaller than 1. In particular, we perform our simulations in the ATLAS muon detector, whose large size and excellent time resolution allow a precision measurement of the slepton time of flight from the production vertex out to the muon chambers and hence of the slepton velocity. Moreover, in the stau NLSP or slepton co-NLSP scenarios, the knowledge of the NLSP mass and lifetime is sufficient to determine $`\sqrt{F}`$, since the factor $``$ in Eq.(3) of is exactly equal to 1. This is not the case in the neutralino NLSP scenario where $``$ depends at least on the neutralino physical composition and more information and measurements are needed for extracting a precise value of $`\sqrt{F}`$.
## 2 Choice of the Sample Models and Event Simulation
The two main parameters affecting the experimental measurement at the LHC of the slepton NLSP properties are the slepton mass and momentum distribution. Indeed, at a hadron collider most of the NLSPs come from squark and gluino production followed by cascade decays. Thus, the momentum distribution is in general a function of the whole MSSM spectrum. However, one can approximately assume that most of the information on the NLSP momentum distribution is provided by the squark mass scale $`m_{\stackrel{~}{q}}`$ only (in the stau NLSP scenario or slepton co-NLSP scenarios of GMSB, one generally finds $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{q}}`$). To perform detailed simulations, we select a representative set of GMSB models generated by SUSYFIRE. We limit ourselves to models with $`m_{\mathrm{NLSP}}>100`$ GeV, motivated by the discussion in , and $`m_{\stackrel{~}{q}}<2`$ TeV, in order to yield an adequate event statistics after a three-year low-luminosity run (corresponding to 30 fb<sup>-1</sup>) at the LHC. Within these ranges, we choose eight extreme points (four in the stau NLSP scenario and four in the slepton co-NLSP scenario) allowed by GMSB in the ($`m_{\mathrm{NLSP}}`$, $`m_{\stackrel{~}{q}}`$) plane, in order to cover the various possibilities.
In Tab. 1, we list the input GMSB parameters we used, while in Tab. 2 we report the corresponding values of the stau mass, the squark mass scale and the gluino mass. The “NLSP” column indicates whether the model belongs to the stau NLSP or slepton co-NLSP scenario. The last column gives the total cross section in pb for producing any pairs of SUSY particles at the LHC.
For each model, the events were generated with the ISAJET Monte Carlo that incorporates the calculation of the SUSY mass spectrum and branching fraction using the GMSB parameters as input. We have checked that for the eight model points considered the sparticle masses calculated with ISAJET are in good agreement with the output of SUSYFIRE.
The generated events were then passed through ATLFAST , a fast particle-level simulation of the ATLAS detector. The ATLFAST package, however, was only used to evaluate the efficiency of the calorimetric trigger that selects the GMSB events. The detailed response of the detector to the slepton NLSP has been parameterized for this work using the results of a full simulation study, as described in the next section.
## 3 Slepton detection
The experimental signatures of heavy long-lived charged particles at a hadron collider have already been studied both in the framework of GMSB and in more general scenarios . The two main observables one can use to separate these particles from muons are the high specific ionization and the time of flight in the detector.
We concentrate here on the measurement of the time of flight, made possible by the timing precision ($`<\mathrm{\hspace{0.33em}1}`$ ns) and the size of the ATLAS muon spectrometer.
It was demonstrated with a full simulation of the ATLAS muon detector that the $`\beta `$ of a particle can be measured with a resolution that can be approximately parameterized as $`\sigma (\beta )/\beta ^2=0.028`$. The resolution on the transverse momentum measurement for heavy particles is found to be comparable to the one expected for muons. We have therefore simulated the detector response to NLSP sleptons by smearing the slepton momentum and $`\beta `$ according to the parameterizations in Ref. .
An important issue is the online selection of the SUSY events. We have not made any attempt to evaluate whether the heavy sleptons can be selected using the muon trigger. For the event selection, we rely on the calorimetric $`E_T^{\mathrm{miss}}`$ trigger, consisting in the requirement of at least one hadronic jet with $`p_T>50`$ GeV, and a transverse momentum imbalance calculated only from the energy deposition in the calorimeter larger than 50 GeV. We checked that this trigger has an efficiency in excess of 80% for all the considered models.
A detailed discussion of the experimental assumptions underlying the results presented here is given in Ref. .
## 4 Event Selection and Slepton Mass Measurement
In order to select a clean sample of sleptons, we apply the following requirements:
* at least one hadronic jet with $`P_T>50`$ GeV and a calorimetric $`E_T^{\mathrm{miss}}>50`$ GeV (trigger requirement);
* at least one candidate slepton satisfying the following cuts:
+ $`|\eta |<`$2.4 to ensure that the particle is in the acceptance of the muon trigger chamber, and therefore both coordinates can be measured;
+ $`\beta _{\mathrm{meas}}<0.91`$, where $`\beta _{\mathrm{meas}}`$ is the $`\beta `$ of the particle measured with the time of flight in the precision chambers;
+ The $`P_T`$ of the slepton candidate, after the energy loss in the calorimeters has been taken into account, must be larger than 10 GeV, to ensure that the particle traverse all of the muon stations.
If we consider an integrated luminosity of 30 fb<sup>-1</sup>, a number of events ranging from a few hundred for the models with 2 TeV squark mass scale to a few hundred thousand for a 500 GeV mass scale survive these cuts and can be used for measuring the NLSP properties.
From the measurements of the slepton momentum and of particle $`\beta `$, the mass can be determined using the standard relation $`m=p\frac{\sqrt{1\beta ^2}}{\beta }`$. For each value of $`\beta `$ and momentum, the measurement error is known and it is given by the parameterizations in Ref. . Therefore, the most straightforward way of measuring the mass is just to use the weighted average of all the masses calculated with the above formula.
In order to perform this calculation, the particle momentum is needed, which implies measuring the $`\eta `$ coordinate. In fact, with the precision chambers only one can only measure the momentum components transverse to the beam axis.
The measurement of the second coordinate must be provided by the trigger chambers, for which only a limited time window around the beam crossing is read out, therefore restricting the $`\beta `$ range where this measurement is available. Hence, we have evaluated the achieved measurement precision for two different $`\beta `$ intervals: $`0.6<\beta <0.91`$ and $`0.8<\beta <0.91`$ for the eight sample points. We found a statistical error well below the 0.1% level for those model points having $`m_{\stackrel{~}{q}}<1300`$ GeV. Even for the three models (B5, B7, B8) with lower statistics ($`m_{\stackrel{~}{q}}2`$ TeV), the error stays below the 0.4% level.
Many more details, tables and figures about this part of our study can be found in Ref. .
## 5 Slepton Lifetime Measurement
The measurement of the NLSP lifetime at a high energy $`e^+e^{}`$ collider was studied in detail in Ref. for the neutralino NLSP case. Similar to that study, the measurement of the slepton NLSP lifetime we are interested in here can be performed by exploiting the fact that two NLSPs are produced in each event. One can therefore select $`N_1`$ events where a slepton is detected through the time-of-flight measurement described above, count the number of times $`N_2`$ when a second slepton is observed and use this information to measure the lifetime. Although in principle very simple, in practice this method requires an excellent control on all possible sources of inefficiency for detecting the second slepton.
We give here the basis of the method, without mentioning the experimental details. We provide an estimate of the achievable statistical error for the models considered and a parameterization of the effect on the lifetime measurement of a generic systematic uncertainty on the slepton efficiency. In case the sparticle spectrum and BRs can be measured from the SUSY events, as e.g. shown in Ref. , an accurate simulation of all the SUSY production processes can be performed, and the results from this section are representative of the measurement precision achievable in a real experiment.
Another method based on the same principles, but assuming minimal knowledge of the SUSY spectrum, is described in Ref. , where a detailed estimate of the achievable systematic precision is given.
We define $`N_1`$ starting from the event sample defined by the cuts discussed in Sec. 4, with the additional requirement that, for a given value of the slepton lifetime, at least one of the produced sleptons decays at at a distance from the interaction vertex $`>10`$ m, and is therefore reconstructed in the muon system. For the events thus selected, we define $`N_2`$ as the subsample where a second particle with a transverse momentum $`>10`$ GeV is identified in the muon system. The search for the second particle should be as inclusive as possible, in order to minimize the corrections to the ratio. In particular, the cut $`\beta _{\mathrm{meas}}<0.91`$ is not applied, but particles with a mass measured from $`\beta `$ and momentum incompatible with the measured slepton mass are rejected. This leaves a background of high momentum muons in the sample that can be statistically subtracted using the momentum distribution of electrons. The ratio
$$R=\frac{N_2}{N_1}$$
(1)
is a function of the slepton lifetime. Its dependence on the NLSP lifetime $`c\tau `$ in meters in shown in Fig. 1 for four among our eight sample models. The curves for the model points not shown are either very similar to one of the curves we show or are mostly included between the external curves corresponding to points B1 and B8, thus providing no essential additional information. Note that the curve for model 6 starts from $`c\tau =2.5`$ m and not from $`c\tau =50`$ cm, as for the other models. This is due to the large value of $`M_{\mathrm{mess}}`$ (cfr. Tab. 1), determining a minimum NLSP lifetime allowed by theory which is macroscopic in this case (Eqs. (3) and (5) in ).
The probability for a particle of mass $`m`$, momentum $`p`$ and proper lifetime $`\tau `$ to travel a distance $`L`$ before decaying is given by the expression
$$P(L)=e^{mL/pc\tau }.$$
(2)
$`N_2`$ is therefore a function of the momentum distribution of the slepton, which is determined by the details of the SUSY spectrum. One therefore needs to be able to simulate the full SUSY cascade decays in order to construct the $`c\tau `$$`R`$ relationship.
The statistical error on $`R`$ can be evaluated as
$$\sigma (R)=\sqrt{\frac{R(1R)}{N_1}}.$$
(3)
Relevant for the precision with which the SUSY breaking scale can be measured is instead the error on the measured $`c\tau `$. This can be extracted from the curves shown in Fig. 1 and can be evaluated as
$$\sigma (\mathrm{c}\tau )=\sigma (R)/\left[\frac{R(c\tau )}{c\tau }\right].$$
(4)
The measurement precision calculated according to this formula is shown in Figs. 2 and 3 for the eight sample points, for an integrated luminosity of 30 fb<sup>-1</sup>. The full line in the plots is the error on $`c\tau `$ considering the statistical error on $`R`$ only. The available statistics is a function of the strongly interacting sparticles’ mass scale. Even if a precise $`R`$$`c\tau `$ relation can be built from the knowledge of the model details, there will be a systematic uncertainty in the evaluation of the losses in $`N_2`$, because of sleptons produced outside the $`\eta `$ acceptance, or absorbed in the calorimeters, or escaping the calorimeter with a transverse momentum below the cuts. The full study of these uncertainties is in progress.
At this level, we just parameterize the systematic error as a term proportional to $`R`$, added in quadrature to the statistical error. We choose two values, $`1\%R`$ and $`5\%R`$, and propagate the error to the $`c\tau `$ measurement. The results are represented by the dashed and dotted lines in Figs. 2 and 3.
For the models with squark mass scales up to 1200 GeV, assuming a 1% systematic error on the measured ratio, a precision better than 10% on the $`c\tau `$ measurement can be obtained for lifetimes between 0.5–1 m and 50–80 m. If the systematic uncertainty grows to 5%, the 10% precision can only be achieved in the range 1–10 m. If the mass scale goes up to 2 TeV, even considering a pure statistical error only, a 10% precision is not achievable. However a 20% precision is possible over $`c\tau `$ ranges between 5 and 100 m, assuming a 1% systematic error.
Note that the curves corresponding to the model points B2, B6 and B7 do not start from $`c\tau =50`$ cm, but from the theoretical lower limit on $`c\tau `$ of 1.8, 2.5 and 6.1 meters, respectively.
## 6 Determining the SUSY Breaking Scale $`\sqrt{F}`$
Using the measured values of $`c\tau `$ and the NLSP mass, the SUSY breaking scale $`\sqrt{F}`$ can be calculated from Eq.(3) in , where $`=1`$ for the case where the NLSP is a slepton. From simple error propagation, the fractional uncertainty on the $`\sqrt{F}`$ measurement can be obtained adding in quadrature one fourth of the fractional error in $`c\tau `$ and five fourths of the fractional error on the slepton mass.
In Figs. 4 and 5, we show the fractional error on the $`\sqrt{F}`$ measurement as a function of $`\sqrt{F}`$ for our three different assumptions on the $`c\tau `$ error. The uncertainty is dominated by $`c\tau `$ for the higher part of the $`\sqrt{F}`$ range and grows quickly when approaching the lower limit on $`\sqrt{F}`$. This is because very few sleptons survive and the statistical error on both $`m_\stackrel{~}{\mathrm{}}`$ and $`c\tau `$ gets very large. If we assume a 1% systematic error on the ratio $`R`$ from which $`c\tau `$ is measured (dashed lines in Figs. 4 and 5), the error on $`\sqrt{F}`$ is better than 10% for $`1000<\sqrt{F}<\mathrm{\hspace{0.33em}4000}`$ TeV for model points B1–B4 with higher statistics. For points B5–B8, in general one can explore a range of higher $`\sqrt{F}`$ values with a small relative error, essentially due to the heaviness of the decaying NLSP in these models. Note also that the theoretical lower limit (5) in on $`\sqrt{F}`$ is equal to about 1200, 1500, 3900, 8900 TeV respectively in model points B2, B5, B6, B7, while it stays well below 1000 TeV for the other models.
## 7 Conclusions
We have discussed a simple method to measure at the LHC with the ATLAS detector the fundamental SUSY breaking scale $`\sqrt{F}`$ in the GMSB scenarios where a slepton is the NLSP and decays to the gravitino with a lifetime in the range 0.5 m $`<c\tau _{\mathrm{NLSP}}<\mathrm{\hspace{0.33em}1}`$ km. This method requires the measurement of the time of flight of long lived sleptons and is based on counting events with one or two identified NLSPs. It relies on the assumptions that a good knowledge of the MSSM sparticle spectrum and BRs can be extracted from the observation of the SUSY events and that the systematic error in evaluating the slepton losses can be kept below the few percent level. We performed detailed, particle level simulations for eight representative GMSB models, some of them being particularly hard due to low statistics. We found that a level of precision of a few 10’s % on the SUSY breaking scale measurement can be achieved in significant parts of the $`1000<\sqrt{F}<\mathrm{\hspace{0.33em}30000}`$ TeV range, for all models considered. More details as well as a full study of the systematics associated with this procedure and another less “model-dependent” method to measure $`\sqrt{F}`$ is presented in detail in Ref. .
Anomaly mediated SUSY breaking at the LHC
F.E. PAIGE and J. WELLS
## 1 Introduction
The signatures for SUSY at the LHC depend very much on the SUSY masses, which presumably result from spontaneous SUSY breaking. It is not possible to break SUSY spontaneously using just the MSSM fields; instead one must do so in a hidden sector and then communicate the breaking through some interaction. In supergravity models, the communication is through gravity. In gauge mediated models it is through gauge interactions; the gravitino is then very light and can play an important role. Simple examples of both have been discussed previously. A third possibility is that the hidden sector does not have the right structure to provide masses through either mechanism; then the leading contributions come from a combination of gravity and anomalies. This is known as Anomaly Mediated SUSY Breaking (AMSB), and it predicts a different pattern of masses and signatures.
## 2 Anomaly-Mediated Supersymmetry Breaking
In the supersymmetric standard model there exist AMSB contributions to the soft mass parameters that arise via the superconformal anomaly . The effect can be understood by recognizing several important features of supersymmetric theories. First, supersymmetry breaking can be represented by a chiral superfield $`\mathrm{\Phi }=1+m_{3/2}\theta ^2`$ which also acts as a compensator for super-Weyl transformations. Treating $`\mathrm{\Phi }`$ as a spurion, one can transform a theory into a super-conformally invariant theory. Even if a theory is superconformal at the outset (i.e., no dimensionful couplings), the spurion $`\mathrm{\Phi }`$ is employed since the quantum field theory requires a regulator that implies scale dependence (Pauli-Villars mass, renormalization scale in dimensional reduction, etc.). To preserve scale invariance the renormalization scale parameter $`\mu `$ in a quantum theory then becomes $`\mu /\sqrt{\mathrm{\Phi }^{}\mathrm{\Phi }}`$. It is the dependence of the regulator on $`\mathrm{\Phi }`$ that induces supersymmetry breaking contributions to the scalars and gauginos.
The anomaly induced masses can be derived straightforwardly for the scalar masses. The Kähler kinetic terms depend on wave function renormalization as in the following superfield operator,
$$d^2\theta d^2\overline{\theta }Z_Q\left(\frac{\mu }{\sqrt{\mathrm{\Phi }^{}\mathrm{\Phi }}}\right)Q^{}Q.$$
(1)
Taylor expanding $`Z`$ around $`\mu `$ and projecting out the $`FF^{}`$ terms yields a supersymmetry breaking mass for the scalar field $`\stackrel{~}{Q}`$:
$$m_{\stackrel{~}{Q}}^2=\frac{1}{4}\frac{d^2\mathrm{ln}Z_Q}{d(\mathrm{ln}\mu )^2}m_{3/2}^2=\frac{1}{4}\left(\frac{\gamma _Q}{g}\beta _g+\frac{\gamma _Q}{y}\beta _y\right)m_{3/2}^2.$$
(2)
Similar calculations can be done for the gauginos and the $`A`$ terms:
$`M_i`$ $`=`$ $`{\displaystyle \frac{g_i^2}{2}}{\displaystyle \frac{dg_i^2}{d\mathrm{ln}\mu }}m_{3/2}={\displaystyle \frac{\beta _{g_i}}{g_i}}m_{3/2},`$ (3)
$`A_y`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}{\displaystyle \frac{d\mathrm{ln}Z_{Q_a}}{d\mathrm{ln}\mu }}m_{3/2}={\displaystyle \frac{\beta _y}{y}}m_{3/2}`$ (4)
where the sum over $`a`$ includes all fields associated with the Yukawa coupling $`y`$ in the superpotential.
There are several important characteristics of the AMSB spectrum to note. First, the equations for the supersymmetry breaking contributions are scale invariant. That is, the value of the soft masses at any scale is obtained by simply plugging in the gauge couplings and Yukawa couplings at that scale into the above formulas. Second, the masses are related to the gravitino mass by a one loop suppression. In AMSB $`M_im_{3/2}\alpha _i/4\pi `$, whereas in SUGRA $`M_im_{3/2}`$. While the AMSB contributions are always present in a theory independent of how supersymmetry breaking is accomplished, they may be highly suppressed compared to standard hidden sector models. Therefore, for AMSB to be the primary source of scalar masses, one needs to assume or arrange that supersymmetry breaking is not directly communicated from a hidden sector. This can be accomplished, for example, by assuming supersymmetry breaking on a distant brane . Finally, the squared masses of the sleptons are negative (tachyonic) because $`\beta _g>0`$ for $`U(1)`$ and $`SU(2)`$ gauge groups. This problem rules out the simplest AMSB model based solely on eqs. 2-4.
Given the tachyonic slepton problem, it might seem most rational to view AMSB as a good idea that did not quite work out. However, there are many reasons to reflect more carefully on AMSB. As already mentioned above, AMSB contributions to scalar masses are always present if supersymmetry is broken. Soft masses in the MSSM come for free, whereas in all other successful theories of supersymmetry breaking a communication mechanism must be detailed. In particle, hidden sector models require singlets to give the gauginos an acceptable mass. In AMSB, singlets are not necessary. Also, there may be small variations on the AMSB idea that can produce a realistic spectrum and can have important phenomenological consequences. This is our motivation for writing this note.
## 3 Two realistic minimal models of AMSB: mAMSB and DAMSB
As we discussed in the introduction, the pure AMSB model gives negative squared masses for the sleptons, thus breaking electromagnetic gauge invariance, so some additional contributions must be included. The simplest assumption that solves this problem is to add at the GUT scale a single universal scalar mass $`m_0^2`$ to all the sfermions’ squared masses. We will call this model mAMSB. The description and many phenomenological implications of this model are given in Refs. . The parameters of the model after the usual radiative electroweak symmetry breaking are then
$$m_0,m_{3/2},\mathrm{tan}\beta ,sgn\mu =\pm .$$
This model has been implemented in ISAJET 7.48 ; a pre-release version of ISAJET has been used to generate the events for this analysis.
For this note the AMSB parameters were chosen to be
$$m_0=200\text{GeV},m_{3/2}=35\text{TeV},\mathrm{tan}\beta =3,sgn\mu =+$$
For this choice of parameters the slepton squared masses are positive at the weak scale, but they are still negative at the GUT scale. This means that charge and color might be broken (CCB) at high temperatures in the early universe. However, at these high energies there are also large finite temperature effects on the mass, which are positive (symmetry restoration occurs at higher $`T`$). In fact, a large class of SUSY models with CCB minima naturally fall into the correct SM minimum when you carefully follow the evolution of the theory from high T to today. If CCB minima are excluded at all scales, then the value of $`m_0`$ must be substantially larger, so the sleptons must be quite heavy.
The masses from ISAJET 7.48 for this point are listed in Table 1. The mass spectrum has some similarity to that for SUGRA Point 5 studied previously : the gluino and squark masses are similar, and the decays $`\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}\mathrm{}`$ and $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0h`$ are allowed. Thus, many of the techniques developed for Point 5 are applicable here. But there are also important differences. In particular, the $`\stackrel{~}{\chi }_1^\pm `$ is nearly degenerate with the $`\stackrel{~}{\chi }_1^0`$, not with the $`\stackrel{~}{\chi }_2^0`$. The mass splitting between the lightest chargino and the lightest neutralino must be calculated as the difference between the lightest eigenvalues of the full one-loop neutralino and chargino mass matrices. The mass splitting is always above $`m_{\pi ^\pm }`$, thereby allowing the two-body decays $`\chi _1^\pm \chi _1^0+\pi ^\pm `$ . Decay lifetimes of $`\chi ^\pm `$ are always less than 10 cm over mAMSB parameter space, and are often less than 1 cm.
Another unique feature of the spectrum is the near degeneracy of the $`\stackrel{~}{\mathrm{}}_L`$ and $`\stackrel{~}{\mathrm{}}_R`$ sleptons. The mass splitting is
$`m_{\stackrel{~}{\mathrm{}}_L}^2m_{\stackrel{~}{\mathrm{}}_R}^20.037\left(m_Z^2\mathrm{cos}2\beta +M_2^2\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{\mathrm{}}_R}}{m_Z}}\right).`$ (5)
There is no symmetry requiring this degeneracy, but rather it is an astonishing accident and prediction of the mAMSB model.
It is instructive to compare the masses from ISAJET with those calculated in Ref. to provide weak-scale input to ISAJET. These masses are listed in the right hand side of Table 1. Since the agreement is clearly adequate for the purposes of the present study, no attempt has been made to understand or resolve the differences. It is clear, however, that if SUSY is discovered at the LHC and if masses or combinations of masses are measured with the expected precision, then more work is needed to compare the LHC results with theoretical models in a sufficiently reliable way.
Another variation on AMSB is deflected AMSB (DAMSB). The idea is based on Ref. who demonstrated that realistic sparticle spectrums with non-tachyonic sleptons can be induced if a light modulus field $`X`$ (SM singlet) is coupled to heavy, non-singlet vector-like messenger fields $`\mathrm{\Psi }_i`$ and $`\overline{\mathrm{\Psi }}_i`$:
$$W_{\mathrm{mess}}=\lambda _\mathrm{\Psi }X\mathrm{\Psi }_i\overline{\mathrm{\Psi }}_i.$$
To ensure gauge coupling unification we identify $`\mathrm{\Psi }_i`$ and $`\overline{\mathrm{\Psi }}_i`$ as $`5+\overline{5}`$ representations of $`SU(5)`$. When the messengers are integrated out at some scale $`M_0`$, the beta functions do not match the AMSB masses, and the masses are deflected from the AMSB renormalization group trajectory. The subsequent evolution of the masses below $`M_0`$ induces positive mass squared for the sleptons, and a reasonable spectrum can result. Although there may be additional significant parameters associated with the generation of the $`\mu `$ and $`B_\mu `$ term in the model, we assume for this discussion that they do not affect the spectra of the MSSM fields. The values of $`\mu `$ and $`B_\mu `$ are then obtained by requiring that the conditions for EWSB work out properly.
The parameters of DAMSB are
$$m_{3/2},n,M_0,\mathrm{tan}\beta ,sgn\mu =\pm $$
where $`n`$ is the number of $`5+\overline{5}`$ messenger multiplets, and $`M_0`$ is the scale at which the messengers are integrated out. Practically, the spectrum is obtained by imposing the boundary conditions at $`M_0`$, and then using SUSY soft mass renormalization group equations to evolve these masses down to the weak scale. Expressions for the boundary conditions can be found in Refs. , and details on how to generate the low-energy spectrum are given in . The resulting spectrum of superpartners is substantially different from that of mAMSB. The most characteristic feature of the DAMSB spectrum is the near proximity of all superpartner masses. In Table 2 we show the spectrum of a model with $`n=5`$, $`M_0=10^{15}\mathrm{GeV}`$, and $`\mathrm{tan}\beta =4`$ as given in . The LSP is the lightest neutralino, which is a Higgsino. (Actually, the LSP is the fermionic component of the modulus $`X`$, but the decay of $`\chi _1^0`$ to it is much greater than collider time scales.) All the gauginos and squarks are between $`300\mathrm{GeV}`$ and $`500\mathrm{GeV}`$, while the sleptons and higgsinos are a bit lighter ($`150\mathrm{GeV}`$ to $`250\mathrm{GeV}`$) in this case.
In summary, we have outlined two interesting directions to pursue in modifying AMSB to make a realistic spectrum. The first direction we call mAMSB, and is constructed by adding a common scalar mass to the sfermions at the GUT scale to solve the negative squared slepton mass problem of pure AMSB. The other direction that we outlined is deflected anomaly mediation that is based on throwing the scalar masses off the pure AMSB renormalization group trajectory by integrating out heavy messenger states coupled to a modulus. The spectra of the two approaches are significantly different, and we should expect the LHC signatures to be different as well. In this note, we study the mAMSB carefully in a few observables to demonstrate how it is distinctive from other, standard approaches to supersymmetry breaking, such as mSUGRA and GMSB.
## 4 LHC studies of the example mAMSB model point
We now turn to a study of the example mAMSB spectra presented in Table 1. A sample of $`10^5`$ signal events was generated; since the total signal cross section is $`16\mathrm{nb}`$, this corresponds to an integrated LHC luminosity of $`6\mathrm{fb}^1`$. All distributions shown in this note are normalized to $`10\mathrm{fb}^1`$, corresponding to one year at low luminosity at the LHC. Events were selected by requiring
* At least four jets with $`p_T>100,50,50,50\text{GeV}`$;
* $`\text{ / }E_T>\mathrm{min}(100\text{GeV},0.2M_{\mathrm{eff}})`$;
* Transverse sphericity $`S_T>0.2`$;
* $`M_{\mathrm{eff}}>600\text{GeV}`$;
where the “effective mass” $`M_{\mathrm{eff}}`$ is given by the scalar sum of the missing $`E_T`$ and the $`p_T`$’s of the four hardest jets,
$$M_{\mathrm{eff}}=\text{ / }E_T+p_{T,1}+p_{T,2}+p_{T,3}+p_{T,4}.$$
Standard model backgrounds from gluon and light quark jets, $`t\overline{t}`$, $`W+\mathrm{jets}`$, $`Z+\mathrm{jets}`$, and $`WW`$ have also been generated, generally with much less equivalent luminosity. The $`M_{\mathrm{eff}}`$ distributions for the signal and the sum of all backgrounds with all except the last cut are shown in Figure 1. The ISAJET IDENT codes for the SUSY events contributing to this plot are also shown. It is clear from this plot that the Standard Model backgrounds are small with these cuts, as would be expected from previous studies .
The mass distribution for $`\mathrm{}^+\mathrm{}^{}`$ pairs with the same and opposite flavor is shown in Figure 2. The opposite-flavor distribution is small, and there is a clear endpoint in the same-flavor distribution at
$$M_{\mathrm{}\mathrm{}}^{\mathrm{max}}=\sqrt{\frac{(M_{\stackrel{~}{\chi }_2^0}^2M_\stackrel{~}{\mathrm{}}^2)(M_\stackrel{~}{\mathrm{}}^2M_{\stackrel{~}{\chi }_1^0}^2)}{M_\stackrel{~}{\mathrm{}}^2}}=213.6,215.3\text{GeV}$$
corresponding to the endpoints for the decays $`\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}_{L,R}^\pm \mathrm{}^{}\stackrel{~}{\chi }_1^0\mathrm{}^+\mathrm{}^{}`$. This is similar to what is seen in SUGRA Point 5, but in that case only one slepton contributes. It is clear from the $`e^+e^{}+\mu ^+\mu ^{}e^\pm \mu ^{}`$ dilepton distribution with finer bins shown in the same figure that the endpoints for $`\stackrel{~}{\mathrm{}}_R`$ and $`\stackrel{~}{\mathrm{}}_L`$ cannot be resolved with the expected ATLAS dilepton mass resolution. More work is needed to see if the presence of two different endpoints could be inferred from the shape of the edge of the dilepton distribution.
Since the main source for $`\stackrel{~}{\chi }_2^0`$ is $`\stackrel{~}{q}_R\stackrel{~}{\chi }_2^0q`$, information on the squark masses can be obtained by combining the leptons from $`\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}\mathrm{}`$ decays with one of the two hardest jets in the event, since the hardest jets are generally products of the squark decays. Figure 3 shows the distribution for the smaller of the two $`\mathrm{}^+\mathrm{}^{}j`$ masses formed with the two leptons and each of the two hardest jets in the event. The dashed curve in this figure shows the same distribution for $`M_{\mathrm{}\mathrm{}}>175\text{GeV}`$, for which the backgrounds are smaller. Both distributions should have endpoints at the kinematic limit for $`\stackrel{~}{q}_R\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}\mathrm{}\stackrel{~}{\chi }_1^0\mathrm{}\mathrm{}`$,
$$\left[\frac{(M_{\stackrel{~}{q}_R}^2M_{\stackrel{~}{\chi }_2^0}^2)(M_{\stackrel{~}{\chi }_2^0}^2M_{\stackrel{~}{\chi }_1^0}^2)}{M_{\stackrel{~}{\chi }_2^0}^2}\right]^{1/2}=652.9\text{GeV}.$$
Figure 3 also shows the $`\mathrm{}^\pm j`$ mass distribution formed with each of the two leptons combined with the jet that gives the smaller of the two $`\mathrm{}\mathrm{}j`$ masses. This should have a 3-body endpoint at
$$\left[\frac{(M_{\stackrel{~}{q}_R}^2M_{\stackrel{~}{\chi }_2^0}^2)(M_{\stackrel{~}{\chi }_2^0}^2M_\stackrel{~}{\mathrm{}}^2)}{M_{\stackrel{~}{\chi }_2^0}^2}\right]^{1/2}=605.4\text{GeV}.$$
The branching ratio for $`\stackrel{~}{b}_1\stackrel{~}{\chi }_2^0b`$ is very small, so the same distributions with $`b`$-tagged jets contain only a handful of events and cannot be used to determine the $`\stackrel{~}{b}_1`$ mass.
The decay chain $`\stackrel{~}{q}_R\stackrel{~}{\chi }_2^0q\stackrel{~}{\mathrm{}}_{L,R}^\pm \mathrm{}^{}q\stackrel{~}{\chi }_1^0\mathrm{}^+\mathrm{}^{}q`$ also implies a lower limit on the $`\mathrm{}\mathrm{}q`$ mass for a given limit on $`z=\mathrm{cos}\theta ^{}`$ or equivalently on the $`\mathrm{}\mathrm{}`$ mass. For $`z>0`$ (or equivalently $`M_{\mathrm{}\mathrm{}}>M_{\mathrm{}\mathrm{}}^{\mathrm{max}}/\sqrt{2}`$) this lower limit is
$$\begin{array}{ccc}(M_\mathrm{}\mathrm{}q^{\mathrm{min}})^2\hfill & =& \frac{1}{4M_2^2M_e^2}\times \hfill \\ & & [M_1^2M_2^4+3M_1^2M_2^2M_e^2M_2^4M_e^2M_2^2M_e^4M_1^2M_2^2M_q^2\hfill \\ & & M_1^2M_e^2M_q^2+3M_2^2M_e^2M_q^2M_e^4M_q^2+(M_2^2M_q^2)\times \hfill \\ & & \sqrt{(M_1^4+M_e^4)(M_2^2+M_e^2)^2+2M_1^2M_e^2(M_2^46M_2^2M_e^2+M_e^4)}]\hfill \\ M_\mathrm{}\mathrm{}q^{\mathrm{min}}\hfill & =& 376.6\text{GeV}\hfill \end{array}$$
where $`M_q`$, $`M_2`$, $`M_e`$, and $`M_1`$ are the (average) squark, $`\stackrel{~}{\chi }_2^0`$, (average) slepton, and $`\stackrel{~}{\chi }_1^0`$ masses. To determine this lower edge, the larger of the two $`\mathrm{}\mathrm{}j`$ masses formed from two opposite-sign leptons and one of the two hardest jets is plotted in Figure 4. An endpoint at about the right value can clearly be seen.
The $`\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{}^+\mathrm{}^{}q`$, $`\mathrm{}^\pm q`$, and lower $`\mathrm{}^+\mathrm{}^{}q`$ edges provide four constraints on the four masses involved. Since the cross sections are similar to those for SUGRA Point 5, we take the errors at high luminosity to be negligible on the $`\mathrm{}^+\mathrm{}^{}`$ edge, 1% on the $`\mathrm{}^+\mathrm{}^{}q`$ and $`\mathrm{}^\pm q`$ upper edges, and 2% on the $`\mathrm{}^+\mathrm{}^{}q`$ lower edge. Random masses were generated within $`\pm 50\%`$ of their nominal values, and the $`\chi ^2`$ for the four measurements with these errors were used to determine the probability for each set of masses. The resulting distribution for the $`\stackrel{~}{\chi }_1^0`$ mass, also shown in Figure 4, has a width of $`\pm 11.7\%`$, about the same as for Point 5; the errors for the other masses are also comparable. Of course, the masses being measured in this case are different: for example the squark mass is the average of the $`\stackrel{~}{q}_R`$ rather than the $`\stackrel{~}{q}_L`$ masses.
The leptons from $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0\mathrm{}^\pm \nu `$ are very soft. This implies that the rate for events with one or three leptons or for two leptons with opposite flavor are all suppressed. Figure 5 shows as a solid histogram the multiplicity of leptons with $`p_T>10\text{GeV}`$ and $`|\eta |<2.5`$ for the AMSB signal with a veto on hadronic $`\tau `$ decays. The same figure shows the distribution for a model with the same weak-scale mass parameters except that the gaugino masses $`M_1`$ and $`M_2`$ are interchanged. This model has a wino $`\stackrel{~}{\chi }_1^\pm `$ approximately degenerate with the $`\stackrel{~}{\chi }_2^0`$ rather than with the $`\stackrel{~}{\chi }_1^0`$. Clearly the AMSB model has a much smaller rate for single leptons and a somewhat smaller rate for three leptons; these rates can be used to distinguish AMSB and SUGRA-like models.
While the decay $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0h`$ is kinematically allowed, the branching ratio is only about 0.03%. Other sources of $`h`$ in SUSY events are also quite small, so in contrast to SUGRA Point 5 there is no strong $`hb\overline{b}`$ signal. However, there is a fairly large branching ratio for $`\stackrel{~}{g}\stackrel{~}{b}\overline{b},\stackrel{~}{t}\overline{t}`$ with $`\stackrel{~}{b}\stackrel{~}{\chi }_1^0b`$, $`\stackrel{~}{t}\chi _1^+b`$, giving two hard $`b`$ jets and hence structure in the $`M_{bb}`$ distribution. For this analysis $`b`$ jets were tagged by assuming that any $`B`$ hadron with $`p_{T,B}>10\text{GeV}`$ and $`|\eta _B|<2`$ is tagged with an efficiency $`ϵ_B=60\%`$; the jet with the smallest
$$R=\sqrt{(\mathrm{\Delta }\eta )^2+(\mathrm{\Delta }\varphi )^2}$$
was then taken to be $`b`$ jets. The two hardest jets generally come from the squarks. To reconstruct $`\stackrel{~}{g}\stackrel{~}{b}\overline{b}`$ one of the two hardest jets, tagged as a $`b`$, was combined with any remaining jet, also tagged as a $`b`$. In addition to the standard multijet and $`\text{ / }E_T`$ cuts, a cut $`M_{\mathrm{eff}}>1200\text{GeV}`$ was made to reduce the Standard Model background. The resulting distributions for the $`b`$ jet multiplicity and for the smallest $`bb`$ dijet mass are shown in Figure 6. The dijet mass should have an endpoint at the kinematic limit for $`\stackrel{~}{g}\stackrel{~}{b}_1\overline{b}\stackrel{~}{\chi }_1^0b\overline{b}`$,
$$M_{bb}^{\mathrm{max}}=\sqrt{\frac{(M_{\stackrel{~}{g}}^2M_{\stackrel{~}{b}}^2)(M_{\stackrel{~}{b}}^2M_{\stackrel{~}{\chi }_1^0}^2)}{M_{\stackrel{~}{b}}^2}}=418.7\text{GeV}.$$
While the figure is roughly consistent with this, the endpoint is not very sharp; more work is needed to assign an error and to understand the high mass tail. There should also be a $`b\overline{t}`$ endpoint resulting from $`\stackrel{~}{g}\stackrel{~}{t}\overline{t}`$, $`\stackrel{~}{t}\stackrel{~}{\chi }_1^+b`$, with $`M_{\stackrel{~}{\chi }_1^+}M_{\stackrel{~}{\chi }_1^0}`$ and essentially invisible. Of course $`m_t`$ has to be kept in the formula. This would be an apparent strong flavor violation in gluino decays and so quite characteristic of these models. Reconstructing the top is more complicated, so this has not yet been studied.
The splitting between the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\chi }_1^0`$ is very small in AMSB models. ISAJET gives a splitting of $`0.189\text{GeV}`$ for this point and $`c\tau =2.8\mathrm{cm}`$, with the dominant decay being the two-body mode $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0\pi `$ via a virtual $`W`$. Ref. gives a somewhat smaller value of $`\mu `$ and so a smaller splitting. The lifetime is of course quite sensitive to the exact splitting. Since the pion or electron is soft and so difficult to reconstruct, it seems better to look for the tail of long-lived winos. The signature is an isolated stiff track in a fraction of the events that ends in the tracking volume and produces no signal in the calorimeter or muon system. Figure 7 shows the radial track length $`R_T`$ distribution in units of $`c\tau `$ for winos with $`|\eta |<1`$ and the (generated) momentum distribution for those with $`R_T>10c\tau `$. Note that the ATLAS detector has three layers of pixels with very low occupancy at radii of 4, 11, and 14 cm and four double layers of silicon strips between 30 and 50 cm. It seems likely that the background for tracks that end after the pixel layers would be small.
It is instructive to compare this signature to that for GMSB models with an NLSP slepton. Both models predict long-lived charged particles with $`\beta <1`$. In the GMSB models, two NLSP sleptons occur in every SUSY event, and they decay into a hard $`e`$’s, $`\mu `$’s, or $`\tau `$’s plus nearly massless $`\stackrel{~}{G}`$’s. In the AMSB models, only a fraction of the SUSY events contain long-lived charged tracks, and these decay into a soft pion or electron plus an invisible particle. A detailed tracking simulation should be done for both cases.
Acknowledgements: This work was supported in part by the U.S. Department of Energy under Contract DE-AC02-98CH10886. We also acknowledge the support of the Les Houches Physics Center, where part of this work was done.
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# 1 Introduction
## 1 Introduction
Supersymmetric gauge theories in low dimensions have been shown to be related to non-perturbative objects in M/string theory , and are therefore of particular interest nowadays. More dramatically, there is a growing body of evidence suggesting that gauged matrix models in $`0+1`$ and $`1+1`$ dimensions may offer a non-perturbative formulation of string theory . There is also a suggestion that large $`N`$ gauge theories in various dimensions may be related to theories with gravity .
It is therefore interesting to study directly the non-perturbative properties of a class of supersymmetric matrix models at finite and large $`N_c`$, where $`N_c`$ is the number of gauge colors. Supersymmetric Discrete Light Cone Quantization (SDLCQ) is a unique non-perturbative numerical approximation that is manifestly supersymmetric at each stage of the calculation. In simplest terms SDLCQ is the approximation that arises when a theory is confined to a box of length $`L`$ in the spatial light-cone directions. This leads to a discrete Fock space basis and the supercharges are finite-dimensional matrices in this basis. This combination of the DLCQ method and supersymmetry is well defined and has proven to be a powerful tool which has allowed to solve a large class of problem that have not been solved previously . It appears that supersymmetric theories are completely well defined, when formulated on a compact space . All of the models that have been addressed so far have had sufficient supersymmetry to make them completely finite so that no renormalization was necessary. It is important to note that numerical approximations that do break supersymmetry will still be faced with a renormalization problem.
In the light-cone formulation the supercharges $`Q_\alpha ^+`$ and $`Q_\alpha ^{}`$ have several interesting and unique properties. Consider for example a pure Yang-Mills theory in $`D`$ dimensions, which has a boson multiplet and a fermion multiplet. Since the longitudinal momentum operator is a kinematic operator, the supercharge $`Q_\alpha ^+`$ must be quadratic in the fields, while the supercharge $`Q_\alpha ^{}`$, whose square is the Hamiltonian, in general has both quadratic and cubic terms in the fields. The dynamics are carried by the cubic terms in the supercharge in the sense that a theory with only quadratic terms in $`Q_\alpha ^{}`$ will be a non-interacting theory. The supercharges contain an odd number of fermion fields therefore the fermion fields must be periodic<sup>1</sup><sup>1</sup>1We will not consider the possibility of twisted boundary conditions here, but clearly this is an interesting direction to explore..
We have formulated a number of supersymmetric theories imposing periodic boundary conditions on the fields and compared our results with other numerical results. We find excellent agreement and we find the SDLCQ converges extremely fast, much faster than any version of standard DLCQ . Among the most interesting results we found were a number of exactly massless bound states in some theories. These are states that are destroyed by one supercharge, $`Q_\alpha ^{}`$, but not the other, $`Q_\alpha ^+`$, since they have finite momentum. These massless states persist at all values of the coupling and it is clear that these states are BPS states which saturate the bound $`|Z|=M`$ where the central charge $`Z`$ is zero. It would be extremely interesting to find BPS states with non-zero masses numerically.
## 2 Formulation of the Theory
In a light-cone formulation the algebra of the supercharges with a central charge extension takes the form
$`\{Q_\alpha ^{},Q_\beta ^{}\}`$ $`=`$ $`P^{}\delta _{\alpha ,\beta }`$ (1)
$`\{Q_\alpha ^+,Q_\beta ^+\}`$ $`=`$ $`P^+\delta _{\alpha ,\beta }`$
$`\{Q_\alpha ^+,Q_\beta ^{}\}`$ $`=`$ $`P_{}\gamma _{\alpha ,\beta }^{}+Z\delta _{\alpha ,\beta },`$
where $`P^+`$ is the longitudinal momentum, $`P_{}`$ is the transverse momentum and $`P^{}`$ is the Hamiltonian. $`Z`$ is the central charge extension and we have suppressed the spinor indices. In light-cone quantized field theories one has transverse boost invariance so one can always work in a frame where the total transverse momentum is zero, thus $`P_{}`$ on all physical states can be taken to be zero.
It is well known that the central charge extension $`Z`$ can be written as an integral over the boundary of the space, and it will therefore vanish if one uses periodic boundary conditions. Therefore the last anti-commutator in Eq. (1) always takes the form
$$\{Q_\alpha ^+,Q_\beta ^{}\}=0.$$
(2)
Without a central charge extension the BPS states of a theory will simply be massless states. The mass of these states is protected by the BPS symmetry and they will remain massless at all couplings. We have seen these states in a recent SDLCQ calculation in 2+1 dimensions . This theory is easily extended to an $`𝒩=2`$ supersymmetry, but without a central charge there is no hope of seeing massive BPS states. A massive BPS state would be a very striking feature in the spectrum the theory, since it would have a fixed mass as a function of the coupling.
We propose therefore to extend the definition of the supercharge in order to include anti-periodic boundary conditions. This extension breaks the supersymmetry at finite resolution, but at infinite resolution it will be restored. We define momentum-shifted analogues of the standard supercharges, $`Q_{\pm \frac{1}{2}}^{}`$, which carry momentum $`\pm \frac{\pi }{2L}`$. The Hamiltonian $`P^{}`$ is defined by
$$\{Q_{+\frac{1}{2},\alpha }^{},Q_{\frac{1}{2},\beta }^{}\}=2\sqrt{2}P^{}\delta _{\alpha ,\beta },$$
(3)
with the normalization of Ref. . To test this method in praxi, we study the eigenvalue problem
$$2P^+P^{}|\psi =M^2|\psi ,$$
(4)
defined by the Hamiltonian, Eq. (3), within the context of the supersymmetric theory of adjoint fermions in $`1+1`$ dimensions. This is one of the simplest possible supersymmetric theories and all the low energy bound states of this theory are well known . The supercharge density is $`q^{}(x)`$
$$q^{}(x)=ig2^{7/4}Tr\left[\psi (x)\psi (x)\frac{1}{_{}}\psi (x)\right].$$
(5)
where $`\psi (x)`$ is a real adjoint fermion field. We expand the field into its modes which will carry half-integer momenta, because we impose anti-periodic boundary conditions
$$\psi (x)_{ij}=\sqrt{\frac{1}{4L}}\underset{n=\frac{1}{2},\frac{3}{2},\frac{5}{2},\mathrm{}}{}\left(B_{ij}(n)e^{in\pi x^{}/L}+B_{ji}^{}(n)e^{in\pi x^{}/L}\right).$$
(6)
We define $`Q_{\pm \frac{1}{2}}^{}`$ to be
$`Q_{\pm \frac{1}{2}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2L}}{\displaystyle _L^L}𝑑x^{}e^{\pi x^{}/2L}q^{}(x)`$
$`=`$ $`{\displaystyle \frac{ig}{2^{1/4}}}\sqrt{{\displaystyle \frac{L}{\pi }}}{\displaystyle \underset{n_1,n_2=\frac{1}{2},\frac{3}{2},\frac{5}{2},\mathrm{}}{}}({\displaystyle \frac{1}{n_1}}{\displaystyle \frac{1}{n_2}}{\displaystyle \frac{1}{n_1+n_2\frac{1}{2}}})\times `$
$`\times `$ $`\left(B_{ik}^{}(n_1)B_{kj}^{}(n_2)B_{ij}(n_1+n_2{\displaystyle \frac{1}{2}})+B_{ij}^{}(n_1+n_2\pm {\displaystyle \frac{1}{2}})B_{ik}(n_1)B_{kj}(n_2)\right).`$
Since we are using a real representation for the fermions it is easy to verify that the Fock space matrices $`Q_{\pm \frac{1}{2}}^{}`$ satisfy,
$$Q_{+\frac{1}{2}}^{}=\left[Q_{\frac{1}{2}}^{}\right]^t.$$
(8)
If all the wavefunctions of the physical states vanish sufficiently fast for small momentum, then for large resolution one would expect that the half unit of momentum which the supercharges carry is negligible compared to momentum the partons carry. Thus, $`Q_{\pm \frac{1}{2}}^{}Q^{}`$ in the continuum limit $`K\mathrm{}`$. However, since we have broken the supersymmetry one might face a mass renormalization which is generally required in these two-dimensional theories. In spite of this potential difficulty we proceed to calculate the spectrum of this theory, and will find that the theory indeed need not be renormalized.
## 3 Numerical Results
In order to obtain the mass squared eigenvalues $`M^2`$ of the theory, we have to solve the eigenvalue problem, Eq. (4). We use the standard large $`N_c`$ discrete Fock basis. Its states are of the form
$$|p_1,p_2,\mathrm{},p_n=\frac{1}{N_{c}^{}{}_{}{}^{n/2}\sqrt{s}}Tr\left[B^{}(p_1)B^{}(p_2)\mathrm{}B^{}(p_n)\right]|0,$$
where the symmetry factor $`s`$ counts the number of times the set of momenta $`(p_1,\mathrm{},p_n)`$ is mapped onto itself under cyclic permutations. When using anti-periodic boundary conditions, the fermion bound states will lie in the sector of half-integer harmonic resolution $`K`$ and boson bound states will have integer $`K`$. The action of the operators, Eq. (2), can be readily evaluated
$`Q_{\pm \frac{1}{2}}^{}|p_1,p_2,\mathrm{},p_n={\displaystyle \frac{ig\sqrt{L}}{2^{7/4}}}\sqrt{{\displaystyle \frac{N_c}{\pi }}}{\displaystyle \underset{i=1}{\overset{n}{}}}()^{i+1}`$
$`\{()^{n(i+1)}{\displaystyle \underset{k=\frac{1}{2}}{\overset{p_i\pm 1}{}}}\left({\displaystyle \frac{1}{k}}+{\displaystyle \frac{1}{p_ik\pm \frac{1}{2}}}{\displaystyle \frac{1}{p_i}}\right)|k,p_ik\pm {\displaystyle \frac{1}{2}},L_n^{(i1)}(p_i)`$
$`+()^{ni}({\displaystyle \frac{1}{p_{i1}}}+{\displaystyle \frac{1}{p_i}}{\displaystyle \frac{1}{p_i+p_{i1}\pm \frac{1}{2}}})|p_i+p_{i1}\pm {\displaystyle \frac{1}{2}},L_n^{(i1)}(p_i,p_{i1})\},`$
where $`L_n^{(j)}`$ is a permutation of $`nk`$ momenta
$$L_n^{(j)}(p_{i_1},\mathrm{},p_{i_k})=\{p_{1+j},p_{2+j},\mathrm{},p_{i_11+j},p_{i_1+1+j},\mathrm{},p_{i_k1+j},p_{i_k+1+j},\mathrm{}p_{n+j}\}.$$
This compact form allows for an easy computer implementation. Since the supercharges change fermions to bosons and vice versa, they change the resolution $`K`$ by $`\pm \frac{1}{2}`$. Thus in this discrete representation the supercharge matrix,
$$K|Q_{\pm \frac{1}{2}}^{}|K\frac{1}{2},$$
will in general not be a square matrix. The Hamiltonian, however, is constructed as the anti-commutator of the of the two momentum-shifted supercharge matrices, Eq. (3). From Eq. (8) it follows then that the Hamiltonian is a hermitian matrix. In this matrix representation the Hamiltonian at resolution $`K`$ will receive contributions from intermediate states with resolution $`K\pm \frac{1}{2}`$. The supercharge matrices in a $`(fermion,boson)^t`$ basis have the structure
$`Q_{\pm \frac{1}{2}}^{}=\left(\begin{array}{cc}\text{0}& A_{\pm \frac{1}{2}}\\ B_{\pm \frac{1}{2}}& \text{0}\end{array}\right)`$ (11)
The Hamiltonian for the fermions is therefore,
$$P_{ferm}^{}=A_{+\frac{1}{2}}B_{\frac{1}{2}}+A_{\frac{1}{2}}B_{+\frac{1}{2}},$$
(12)
and the Hamiltonian for the bosons is,
$$P_{bose}^{}=B_{+\frac{1}{2}}A_{\frac{1}{2}}+B_{\frac{1}{2}}A_{+\frac{1}{2}}.$$
(13)
These two matrices are subsequently diagonalized to yield the fermionic the bosonic masses. We note that these matrices have in general different dimensions.
In a first step, we calculated the fermionic and bosonic spectrum up to harmonic resolution $`K=10`$ by solving the eigenvalue problem, Eq. (4). We then focused on the six lowest eigenvalues whose eigenfunctions for resolutions up to $`K=10`$ are solely built out of states with at most four and five particles for the bosonic and fermionic spectrum, respectively. A truncation to four particles allowed us to go to resolution $`K=25`$ in the bosonic sector with relative ease. In the fermionic sector of the spectrum, the larger dimension of the fermionic matrix, Eq. (12), prevented us from going quite as high.
The resulting bound state masses in units of $`g^2N_c/\pi `$ are shown in Fig. 1. We fitted all six curves to the function
$$M^2+a\sqrt{\frac{1}{K}}+b\frac{1}{K}.$$
We find the continuum masses to be
$$M_{aSDLCQ}^2=22.7,43.0,57.7,58.0,63.4,67.0,$$
(14)
compared to the values
$$M_{SDLCQ}^2=23.8,47.9,62.1,62.6,63.8,64.7,$$
(15)
of ordinary SDLCQ at resolution $`K=8`$. While the agreement is far from perfect, it is clear that for this model the anti-periodic SDLCQ method works, in the sense that it correctly reproduces the spectrum of the theory. It is fair to say, however, that the exact continuum values are quite sensitive to the fitting function.
In Ref. we compared the lowest eigenvalue as calculated in SDLCQ and standard DLCQ as a function of the resolution $`K`$. We found that the convergence of standard DLCQ is slow and clearly nonlinear as a function of $`1/K`$. Unfortunately, we see from Fig. 2 that the convergence for anti-periodic SDLCQ is, roughly speaking, about as bad as the standard Hamiltonian DLCQ method<sup>2</sup><sup>2</sup>2For comparison, we used the same fitting function that was used in the DLCQ calculations of Ref. for our data. .
In Fig. 2 we show the lowest fermion bound state mass and the lowest boson bound state mass as a function of $`1/K`$ for the anti-periodic SDLCQ approximation. While the fermions and bosons appear at different resolutions, they can be fitted by the same curve for large enough $`K`$. Thus the supersymmetry of the spectrum appears to be recovered by the approximation already at low harmonic resolution. Unfortunately, although the potential perils of supersymmetry breaking seem to be absent already at $`K=10`$ for the lowest eigenstates, the convergence remains unpleasantly slow. In principle, however, the method appears to be working, and we anticipate room for numerical improvements.
## 4 Conclusions
In the present note we introduced a novel version of SDLCQ that allows for the use of anti-periodic boundary conditions for all fields. Consequently, the formalism can support for the first time boundary integral contributions. We provided evidence that the approximation converges to the standard supersymmetric results. The new approach inherits the very advantageous feature of absence of mass renormalization from SDLCQ. Unfortunately however, the approximation does not enjoy the same rapid convergence as SDLCQ. The lack of rapid convergence is a severe problem, since supersymmetric theories with non-trivial central charges ($`𝒩2`$) have many species of particles and therefore a very large Fock space. Thus going to sufficiently high resolution to obtain accurate results for the spectrum will be quite difficult. Of course, the standard DLCQ approach can be used with anti-periodic boundary conditions. But, as we saw in Fig. 2, its convergence is no better than that of the anti-periodic SDLCQ approximation. Moreover, the standard DLCQ approach will certainly have renormalization problems in higher dimensions.
## Acknowledgments
This work was supported in part by the US Department of Energy. The authors would like to thank Oleg Lunin for many helpful discussions.
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# Enlarged quintessence cosmology
## I Introduction
In the wake of the recent measurements of distant exploding stars, supernovae, the existence of negative-pressure dark energy has begun to gain broad consideration. Using Type Ia supernovae as standard candles to gauge the expansion of the Universe, observers have found evidence that the Universe is accelerating . A new component with significant negative pressure, called quintessence matter (Q-matter for short) will in fact cause the cosmic expansion to speed up, so the supernovae observations provide empirical support for a new form of energy with strong negative pressure .
Different forms for the quintessence energy have been proposed. They include a cosmological constant (or more generaly a variable cosmological term), a scalar field , a frustrated network of non–Abelian cosmic strings, and a frustrated network of domain walls , . All these proposals assume the Q-matter behaves as a perfect fluid with a linear baryotropic equation of state, and so some effort has been invested in determinating its adiabatic index at the present epoch -see e.g. . This new energy is to be added to the more familiar components: i.e., normal matter (luminous and dark) plus radiation. The contribution of the radiation component is known to be negligible at the current epoch whereas the main contribution to the former comes from cold dark matter (CDM) .
However, QCDM models (including those in which the quintessence energy is just the energy of the quantum vacuum) find dificulties in explaining why the energy densities of the CDM and Q–matter should be comparable today. Since both energies redshift at different rates the conditions at the early universe have to be set very carefully for both energy densities to be of the same order today, though one may always invoke some version or other of the anthropic principle to soften a bit the problem. This is the coincidence problem . Recently, a promising solution, for spatially–flat metrics, based in the notion of “tracker fields”, fields that roll down their potential according to an attractorlike solution to the equations of motion , has been proposed. Unfortunately it is not clear to what extent these fields have the ability to solve the coincidence problem and, at the same time, drive the Universe to the current phase of accelerated expansion.
In any case, all these models overlook the fact that since the cosmic fluid consists in a mixture of different fluids a dissipative pressure may naturally arise which, for expanding universes, is bound to further decrease the total pressure . Recently it has been proposed that the CDM must self–interact to explain the structure of the halos of the galaxies -see however . This self–interaction leads naturally to a viscous pressure whose magnitud will depend on the mean free path of the CDM particles. On the other hand it has been suggested that dissipative fluids (or equivalently particle production processes) can drive a phase of accelerated expansion -see and references therein.
This paper investigates how the combined action of dissipative normal matter and a quintessence scalar field may lead to the current accelerated expansion stage, and at the same time provide a solution to the coincidence problem different from that relying on a tracker field. It is shown in Sec. II that the flat coincident solution is an attractor. Section III deepens in the consequences brought about by a dissipative pressure in the strees–energy tensor, i.e., it explores the dynamics of quintessence–dissipative dark matter (QDDM) models, while Sec. IV briefly outlines some specific models. Finally, Sec. V summarizes our findings. It should be understood that this work mainly deals with the present and late Universe since it is precisely at these stages where the dynamic effects of quintessence and dissipative pressure become important. Units in which $`c=8\pi G=k_B=1`$ are used throughout.
## II Cosmological problems
This section shows that the Friedmann-Lamaitre-Robertson-Walker (FLRW) universe filled with perfect normal matter plus quintessence fluid, corresponding to some scalar field governed by Klein–Gordon equation, cannot at the same time drive an accelerated expansion and solve the coincidence problem. To solve it, without abandoning the FLRW geometry, some additional contribution to the stress–energy tensor, such as a bulk dissipative pressure, is needed.
The overall stress–energy tensor of the cosmic fluid without the dissipative pressure reads
$$T_{ab}=\rho u_au_b+ph_{ab},(h_{ab}=g_{ab}+u_au_bu^au_a=1),$$
(1)
where $`\rho =\rho _m+\rho _\varphi `$ and $`p=p_m+p_\varphi `$. Here $`\rho _m`$ and $`p_m`$ are the energy density and pressure of the matter whose equation of state is $`p_m=(\gamma _m1)\rho _m`$ with adiabatic index in the interval $`1\gamma _m2`$. Likewise $`\rho _\varphi `$ and $`p_\varphi `$, the energy density and pressure of the minimally coupled self–interacting Q-matter field $`\varphi `$, i.e.,
$$\rho _\varphi =\frac{1}{2}\dot{\varphi }^2+V(\varphi ),p_\varphi =\frac{1}{2}\dot{\varphi }^2V(\varphi ),$$
(2)
are related by an equation of state similar to that of the matter, viz. $`p_\varphi =(\gamma _\varphi 1)\rho _\varphi `$, so that its adiabatic index is given by
$$\gamma _\varphi =\frac{\dot{\varphi }^2}{\dot{\varphi }^2/2+V(\varphi )},$$
(3)
where for non–negative potentials $`V(\varphi )`$ one has $`0\gamma _\varphi 2`$. The scalar field can be properly interpreted as Q–matter provided $`\gamma _\varphi <1`$ -see e.g. . As usual an overdot means derivative with respect to cosmic time. In general $`\gamma _\varphi `$ varies as the Universe expands, and the same is true for $`\gamma _m`$ since the massive and massless components of the matter fluid redshift at different rates.
The Friedmann equation together with the energy conservation of the normal matter fluid and quintessence (Klein-Gordon equation) are
$$H^2+\frac{k}{a^2}=\frac{1}{3}(\rho _m+\rho _\varphi )(k=1,0,1),$$
(4)
$$\dot{\rho _m}+3H\gamma _m\rho _m=0,$$
(5)
$$\ddot{\varphi }+3H\dot{\varphi }+V^{}=0,$$
(6)
where $`H\dot{a}/a`$ denotes the Hubble factor and the prime means derivative with respect to $`\varphi `$. Introducing $`\mathrm{\Omega }_m\rho _m/\rho _c`$, $`\mathrm{\Omega }_\varphi ,\rho _\varphi /\rho _c`$, with $`\rho _c3H^2`$ the critical density, and $`\mathrm{\Omega }_kk/(aH)^2`$ plus the definition $`\mathrm{\Omega }\mathrm{\Omega }_m+\mathrm{\Omega }_\varphi `$ the set of equations (46) can be recast as (cf. )
$$\mathrm{\Omega }_m+\mathrm{\Omega }_\varphi +\mathrm{\Omega }_k=1,$$
(7)
$$\dot{\mathrm{\Omega }}=\mathrm{\Omega }\left(\mathrm{\Omega }1\right)\left(3\gamma 2\right)H,$$
(8)
$$\dot{\mathrm{\Omega }}_\varphi =\left[2+\left(3\gamma 2\right)\mathrm{\Omega }3\gamma _\varphi \right]\mathrm{\Omega }_\varphi H,$$
(9)
where $`\gamma `$ is the average adiabatic index given by
$$\gamma \mathrm{\Omega }=\gamma _m\mathrm{\Omega }_m+\gamma _\varphi \mathrm{\Omega }_\varphi .$$
(10)
Next subsections investigate the flatness and coincidence problems. The former will be solved if the solution $`\mathrm{\Omega }=1`$ to equation (8) becomes an attractor at late time. In its turn, the coincidence problem will be solved if the ratio $`\mathrm{\Omega }_\varphi /\mathrm{\Omega }_m`$ becomes asymptotically a constant. We shall therefore explore the possibility of constant stable solutions in the $`(\mathrm{\Omega },\mathrm{\Omega }_m,\mathrm{\Omega }_\varphi )`$ space.
### A The flatness problem and accelerated expansion
The combined measurements of the cosmic microwave background temperature fluctuations and the distribution of galaxies on large scales seem to imply that the Universe may be flat or nearly flat . Hence the interesting solution at late times of (8) is $`\mathrm{\Omega }=1`$ (i.e., $`k=0`$), and so we discard the solution $`\mathrm{\Omega }=0`$ as incompatible with observation. The solution $`\mathrm{\Omega }=1`$ is asymptotically stable for expanding universes ($`H>0`$) provided that the condition $`\dot{\mathrm{\Omega }}/\mathrm{\Omega }<0`$ holds in a neighborhood of $`\mathrm{\Omega }=1`$ and this implies $`\gamma <2/3`$. Hence the matter stress violates the strong energy condition (SEC) $`\rho +3p0`$ and as a consequence the Universe accelerates its expansion, i.e., $`\ddot{a}/a=(\rho +3p)/6>0`$.
Let us examine more closely the implications of the current accelerated expansion for the QCDM model. Since the mixture of Q–matter and perfect dark matter fluid violates the SEC, $`\gamma _\varphi `$ must be low enough. Namely, since $`\gamma <2/3`$, $`\gamma _m1`$, and $`\gamma _\varphi <\gamma _m`$, equation (10) implies $`\gamma _\varphi <\gamma `$. Then, introducing $`\mathrm{\Omega }=1`$ in equation (9) we obtain
$$\dot{\mathrm{\Omega }}_\varphi =3(\gamma \gamma _\varphi )\mathrm{\Omega }_\varphi H,$$
(11)
and therefore $`\dot{\mathrm{\Omega }}_\varphi >0`$, i.e., $`\mathrm{\Omega }_\varphi `$ will grow until the constraint (7) is saturated, giving $`\mathrm{\Omega }_\varphi =1`$ in the asymptotic regime. Thus the matter fluid yields a vanishing contribution to the energy density of the Universe at large time. This implies that a flat FLRW universe driven by a mixture of normal perfect fluid and quintessence matter cannot both drive an accelerated expansion and solve the coincidence problem. Therefore some other contribution must enter the stress–energy tensor of the cosmic fluid, i.e., it must be “enlarged”.
### B The coincidence problem
As shown above we cannot have both the current accelerated expansion and the coincidence problem solved within a model that assumes a perfect matter fluid and a quintessential scalar field. Things fare differently when a dissipative pressure enters the play. Because of the FLRW metric, velocity gradients causing shear viscosity and temperature gradients leading to heat transport must be absent. Therefore the only admissible dissipative term corresponds to a bulk dissipative pressure $`\pi `$. This quantity is always negative for expanding fluids (i.e., $`\pi <0`$ so long as $`H>0`$) and may be understood either as a viscous pressure or as the effect of particle production. On general grounds the former possibility is usually thought to give just a small contribution to the overall pressure, however the impact of the latter is not so much limited. The expression of the bulk stress when interpreted that way is $`\pi =(\rho _m+p_m)\mathrm{\Gamma }/3H`$ where $`\mathrm{\Gamma }`$ denotes the particle production rate. This process is dissipative in the sense that the produced particles imply an augment of the phase space volume. A recent discussion about the interplay between dissipative bulk pressure and cosmological particle production can be found in .
So the total stress–energy tensor of the cosmic medium, made up of a dissipative but otherwise normal fluid plus the Q–matter fluid, reads
$$T_{ab}=(\rho _m+\rho _\varphi +p_m+p_\varphi +\pi )u_au_b+(p_m+p_\varphi +\pi )g_{ab}.$$
(12)
A parallel calculation to that of above leads to the corresponding Einstein-Klein-Gordon field equations
$$\dot{\mathrm{\Omega }}=\mathrm{\Omega }\left(\mathrm{\Omega }1\right)\left[3\left(\gamma +\frac{\pi }{\rho }\right)2\right]H,$$
(13)
and
$$\dot{\mathrm{\Omega }}_\varphi =\left\{2+\left[3\left(\gamma +\frac{\pi }{\rho }\right)2\right]\mathrm{\Omega }3\gamma _\varphi \right\}H\mathrm{\Omega }_\varphi ,$$
(14)
instead of equations (8) and (9). The energy conservation of the normal matter is
$$\dot{\rho _m}+3\left(\gamma _m+\frac{\pi }{\rho _m}\right)\rho _mH=0.$$
(15)
Owing to the presence of the dissipative bulk stress the constraint $`\gamma <2/3`$ does not longer have to be fulfilled for the solution $`\mathrm{\Omega }=1`$ of equation (13) to be stable. Likewise, inspection of (14) shows that when $`\mathrm{\Omega }=1`$ one can have $`\dot{\mathrm{\Omega }}_\varphi <0`$ just by choosing the ratio $`\pi /\rho `$ sufficiently negative . Thereby the constraint (7) allows a nonvanishing $`\mathrm{\Omega }_m`$ at large times. By contrast tracker fields based models (valid only when $`\mathrm{\Omega }_k=0`$) predict that $`\mathrm{\Omega }_m0`$ asymptotically . A fixed point solution of equation (13) is $`\mathrm{\Omega }=1`$. Note that equations (7) and (14) have fixed point solutions $`\mathrm{\Omega }_m=\mathrm{\Omega }_{m0}`$ and $`\mathrm{\Omega }_\varphi =\mathrm{\Omega }_{\varphi 0}`$, respectively, when the partial adiabatic indices and the dissipative pressure are related by
$$\gamma _m+\frac{\pi }{\rho _m}=\gamma _\varphi =\frac{2\dot{H}}{3H^2}.$$
(16)
Then smaller $`\gamma _\varphi `$, the larger the dissipative effects. Let us investigate the requirements imposed by the stability of these solutions. From (13) we see that $`\gamma +\pi /\rho <2/3`$ must be fulfilled if the solution $`\mathrm{\Omega }=1`$ is to be asymptotically stable. This condition, together with (16), leads to the additional constraint on the viscosity pressure
$$\pi <\left(\frac{2}{3}\gamma _m\right)\rho _m,$$
(17)
which is negative for ordinary matter fluids. Also by virtue of (11) and the first equality in (16) we obtain from (17) that $`\gamma _\varphi <2/3`$.
In the special case of a spatially flat universe ($`\mathrm{\Omega }=1`$), the stability of the solutions $`\mathrm{\Omega }_{m0}`$ and $`\mathrm{\Omega }_{\varphi 0}`$ may be studied directly from (14). Namely, setting $`\mathrm{\Omega }_\varphi =\mathrm{\Omega }_{\varphi 0}+\omega `$ and using (10) it follows that
$$\dot{\omega }=3\mathrm{\Omega }_m\left(\gamma _m\gamma _\varphi +\frac{\pi }{\rho _m}\right)H\left(\mathrm{\Omega }_{\varphi 0}+\omega \right).$$
(18)
Accordingly the solution $`\mathrm{\Omega }=1`$, $`\mathrm{\Omega }_\varphi =\mathrm{\Omega }_{\varphi 0}`$ is stable for the class of models that satisfies $`\psi \gamma _m\gamma _\varphi +(\pi /\rho _m)<0`$ and $`\psi 0`$ for $`t\mathrm{}`$. Note that this coincides with the attractor condition (16).
In order to study the stability of the solutions $`\mathrm{\Omega }_{m0}`$ and $`\mathrm{\Omega }_{\varphi 0}`$ when $`k0`$ it is advisable to derive a dynamic equation for the density ratio parameter
$`ϵ{\displaystyle \frac{\mathrm{\Omega }_m}{\mathrm{\Omega }_\varphi }}.`$
For this purpose we combine the logarithmic derivative of $`ϵ`$ with the definitions of $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\varphi `$ and the energy conservation equations (15) and (6) -the latter written in terms of $`\rho _\varphi `$. It yields
$$\dot{ϵ}=3\left(\gamma _\varphi \gamma _m\frac{\pi }{\rho _m}\right)Hϵ.$$
(19)
To calculate $`\gamma _\varphi `$ we use (15) together with (3), (4) and (6), obtaining
$$\gamma _\varphi =\gamma _m+\frac{\pi }{\rho _m}\frac{1}{\mathrm{\Omega }_\varphi }\left[\frac{2\dot{H}}{3H^2}+\gamma _m+\frac{\pi }{\rho _m}+\left(\frac{2}{3}\gamma _m\frac{\pi }{\rho _m}\right)\mathrm{\Omega }_k\right].$$
(20)
Introducing (20) in (19) we get
$$\dot{ϵ}=\frac{3Hϵ}{\mathrm{\Omega }_\varphi }\left[\frac{2\dot{H}}{3H^2}+\gamma _m+\frac{\pi }{\rho _m}+\left(\frac{2}{3}\gamma _m\frac{\pi }{\rho _m}\right)\mathrm{\Omega }_k\right],$$
(21)
and perturbating this expression about the solution $`ϵ_0𝒪(1)`$, (i.e., using the ansatz $`ϵ=ϵ_0+\delta `$ with $`|\delta |1`$) we obtain with the help of (16)
$$\dot{\delta }=\frac{3}{\mathrm{\Omega }_\varphi }\left(\frac{2}{3}\gamma _\varphi \right)\mathrm{\Omega }_kH\left(ϵ_0+\delta \right)$$
(22)
near the attractor. For $`\mathrm{\Omega }_k>0`$ (negatively spatially curved universes) it follows that $`\delta `$ decreases, i.e., the ratio $`(\mathrm{\Omega }_m/\mathrm{\Omega }_\varphi )_0`$ is a stable solution. For $`\mathrm{\Omega }_k<0`$ one has to go beyond the linear perturbative regime and/or restrict the class of models as in the spatially flat case to determine the stability of the solution. We defer this to a future research.
As we mentioned above, recently there have been some claims that CDM must not be a perfect fluid because it ought to self–interact (with a mean free path in the range $`1\text{kpc}l1\text{ Mpc}`$) if one wish to explain the structure of the halos of galaxies . In this light it is not unreasonable to think that this same interaction is the origin of the dissipative pressure $`\pi `$ at cosmological scales. Bearing in mind that $`l=1/n\sigma `$, with $`n`$ the number density of CDM particles and $`\sigma `$ the interaction cross section, a simple estimation reveals that at such scales $`l`$ is lower than the Hubble distance $`H^1`$ and accordingly the fluid approximation we are using is valid.
## III QDDM asymptotic era
Bulk viscosity arises typically in mixtures–either of different particles species, as in a radiative fluid, or of the same species but with different energies, as in a Maxwell–Boltzmann gas. Physically, we can think of $`\pi `$ as the internal “friction” that sets in due to the different cooling rates in the expanding mixture.
Any dissipation in exact FLRW universes have to be scalar in nature, and in principle it may be modelled as a bulk viscosity effect within a nonequilibrium thermodynamic theory such as the Israel–Stewart’s , . In that formulation, the transport equation for the bulk viscous pressure takes the form
$$\pi +\tau \dot{\pi }=3\zeta H\frac{1}{2}\pi \tau \left[3H+\frac{\dot{\tau }}{\tau }\frac{\dot{\zeta }}{\zeta }\frac{\dot{T}}{T}\right],$$
(23)
where the positive–definite quantity $`\zeta `$ stands for the phenomenological coefficient of bulk viscosity, $`T`$ the temperature of the fluid, and $`\tau `$ the relaxation time associated to the dissipative pressure -i.e., the time the system would take to reach the thermodynamic equilibrium state if the velocity divergence were suddenly turned off . Usually $`\zeta `$ is given by the kinetic theory of gases or a fluctuation-dissipation theorem or both .
Provided the factor within the square bracket in (23) is small it can be approximated by the more manageable truncated transport equation
$$\pi +\tau \dot{\pi }=3\zeta H,$$
(24)
widely used in the literature. This as well as (23) meets the requirements of causality and stability to be fulfilled by any physically acceptable transport equation .
### A The quasiperfect regime
Here we obtain an explicit expression for the leading behavior of the attractor solution at late time. We begin by writing the equation of motion for the Hubble factor that follows from combining (13) and (15) with (24)
$$\ddot{H}+3\gamma H\dot{H}+\tau ^1\left[\dot{H}+\frac{3}{2}\left(\gamma +\tau \dot{\gamma }\right)H^2\frac{3}{2}\zeta H\right]$$
$$+\frac{k}{a^2}\left[(1\frac{3}{2}\gamma )\left(2H\tau ^1\right)+\frac{3}{2}\dot{\gamma }\right]=0.$$
(25)
We next evaluate (10) on the attractor and insert it together with (16) in (25) to obtain
$$\nu ^1\left(\frac{\ddot{H}}{H}+3\gamma _m\dot{H}\right)+\dot{H}+\frac{3\gamma _m}{2}H^2\frac{3\zeta }{2\mathrm{\Omega }_{m0}}H=0.$$
(26)
Observation seems to rule out huge entropy production processes on large scales, otherwise the flux of gamma-rays we witness should be much higher . Hence we shall assume that the viscous effects are not as large as that, but however not altogether negligible. If $`\tau `$ is the relaxation time, then $`\nu =\left(\tau H\right)^1`$ is the number of relaxation times in a Hubble time -for quasistatic expansions $`\nu `$ is proportional to the number of particle interactions in a Hubble time. Perfect fluid behavior occurs in the limit $`\nu \mathrm{}`$, and a consistent hydrodynamical description of the fluids requires $`\nu >1`$. Thus we are lead to assume that $`\tau H`$ is small and we propose a “quasiperfect” expansion in powers of $`\nu ^1`$.
Let us show that the attractor solution of leading behavior $`at^\sigma `$ when $`t\mathrm{}`$, with $`\sigma `$ a positive–definite constant, is consistent in the quasiperfect regime. Indeed, by virtue of (16) it implies $`\gamma _\varphi 2/3\sigma `$ and
$$\frac{\pi }{\rho _m}\left(\gamma _m\frac{2}{3\sigma }\right).$$
(27)
For approximately constant $`\gamma _m`$ ($`\gamma _m=1`$ for CDM), we get from (27)
$$\frac{\dot{\pi }}{\pi }\frac{\dot{\rho }_m}{\rho _m}2\frac{H}{\sigma }.$$
(28)
Hence (24) becomes
$$\pi \left(1\frac{2}{\nu \sigma }\right)3\zeta H,$$
(29)
and we get to leading order in $`\nu ^1`$
$$\zeta \mathrm{\Omega }_{m0}\left(\gamma _m\frac{2}{3\sigma }\right)H.$$
(30)
Then integration of (15) yields $`\rho _ma^{2/\sigma }`$, the same scaling law as $`\rho _\varphi `$. Finally its insertion in (4) leads back to $`at^\sigma `$, showing the consistency of our assumptions. These results correspond to the lowest order in the quasiperfect expansion.
To go a step further we introduce the expansion of $`H`$ in powers of $`\nu ^1`$
$$H=H_0\left(1+h_1\nu ^1+\mathrm{}\right)$$
(31)
in (26), and assuming that $`|\dot{\tau }|\nu ^1`$, it follows the approximated solution
$$H\frac{\sigma }{t}\left[1+\left(\frac{3\gamma _m\sigma 2}{\sigma }+\frac{\theta }{t}\right)\frac{1}{\nu }\right],$$
(32)
where $`\theta `$ is an arbitrary integration constant. This expression reveals that the power law is an attractor solution and that for $`\sigma >1`$ (deceleration parameter $`q=(\sigma 1)/\sigma <0`$), CDM viscosity provides an accelerated expansion scenario that also solves the coincidence problem. It can be shown that $`\gamma _\varphi `$ does not pick any correction of order $`t^1`$ from the subdominant term in (32). Instead the first correction appears to the order $`\nu ^2t^2`$, and this fact shows the high degree of correction of the approximation that $`\gamma _\varphi `$ takes a constant value in the late time regime.
We note that this attractor solution works for any viscosity coefficient with leading behavior (30). In particular the case $`\zeta \sqrt{\rho _m}`$, investigated in , satisfies this requirement.
### B Full causal corrections
Here we gauge the changes bring about by the the full transport equation (23) on the expansion exponent obtained in the previous section. Using that equation and the viscosity coefficient found in (30) we get
$$\nu ^1\left\{\frac{\ddot{H}}{H}\frac{1+2r}{2}\frac{\dot{H}^2}{H^2}+\frac{3}{2}\left[\gamma _m\left(\frac{3}{2}r\right)+1\right]\dot{H}+\frac{9}{4}\gamma _mH^3\right\}+\dot{H}+\frac{1}{\sigma }H^2=0,$$
(33)
where, to estimate the corrections, we have assumed that in the asymptotic regime $`T\rho ^r`$, with $`r`$ a positive–definite constant, and we have used that $`\rho _mϵ_0\rho /\left(1+ϵ_0\right)`$ in this regime. This power law relationship is the simplest way to guarantee a positive heat capacity. Usually $`p`$, $`\rho `$, $`T`$ and the particle number density $`n`$ are equilibrium magnitudes related by equations of state of the form $`\rho =\rho (T,n)`$ and $`p=p(T,n)`$. Further the thermodynamic relation
$$\left(\frac{\rho }{n}\right)_T=\frac{\rho +p}{n}\frac{T}{n}\left(\frac{p}{T}\right)_n$$
(34)
holds. This directly follows from the requirement that the entropy is a state function . In the particular case of a material fluid with $`\rho =\rho (T)`$ and constant adiabatic index, this relation imposes the constraint $`r=\left(\gamma 1\right)/\gamma `$, so that $`0r1/2`$ for $`1\gamma 2`$. Inserting (31) in (33), and assuming that $`|\dot{\tau }|\nu ^1`$, we obtain the approximate solution
$$H\frac{\sigma }{t}\left\{1+\left[\frac{2}{\sigma }+\frac{3}{2}\left(\gamma _m\left(\frac{1}{2}r\frac{3}{2}\sigma \right)+1+\frac{2r+1}{3\sigma }\right)+\frac{C}{t}\right]\frac{1}{\nu }\right\}.$$
(35)
Comparison with (32) shows that except when $`\sigma 1`$, the use of the complete transport equation leads a to slighty slower rate of expansion at late time. Also, in this regime, the equilibrium temperature decreases as $`Tt^{2r}a^{2r/\sigma }`$.
### C Late Q-matter dynamics
In virtue of (10) the density parameter ratio can be written in terms of the adiabatic indices
$$ϵ=\frac{\gamma \gamma _\varphi }{\gamma _m\gamma }.$$
(36)
Since $`\gamma _m`$ is approximately constant, if $`\gamma \gamma _0`$ in the asymptotic regime when $`ϵϵ_0`$, then $`\gamma _\varphi `$ must also approach a constant value. Hence from (3) get the constraint
$$V(\varphi )/\dot{\varphi }^2C$$
(37)
with $`C>1`$ as $`\gamma _\varphi <2/3`$. Potentials that satisfy this constraint have been investigated in and . Then (6) becomes
$$\frac{\ddot{\varphi }}{\dot{\varphi }}+\frac{3H}{1+2C}0.$$
(38)
An interesting potential that meets this constraint is the exponential potential
$$V(\varphi )=V_0\mathrm{exp}(A\varphi ),$$
(39)
where $`A`$ and $`V_0`$ are constants. Now, integrating (38) and using (37) we get
$$\varphi (t)\frac{1}{A}\left[\mathrm{ln}\frac{V_0A^2\gamma _\varphi }{2\left(2\gamma _\varphi \right)}+2\mathrm{ln}t\right].$$
(40)
Hence $`\varphi `$ slowly rolls down the exponential potential as $`\dot{\varphi }1/t`$ when $`t\mathrm{}`$. Also we find that $`C(3\sigma 1)/2`$ with $`\gamma _\varphi 2/3\sigma <2/3`$ on the attractor, irrespective of $`V_0`$ and $`A`$.
Perfect fluid QCDM models based on the exponential potential are ruled out by observations . However we shall demonstrate in the next secion that in the realm of QDDM models the exponential potential yields satisfactory results without any fine tuning of the parameters.
## IV QDDM models
This section explores the dynamical evolution of a universe filled with a viscous material fluid and a quintessence scalar field by resorting to models based on simple relationships for the nonequilibrium quantities. This allows us to explore the large dissipative regime where the nonequilibrium pressure has a magnitude comparable with the energy density. Recently tracker–field models with inverse power potentials have attracted much interest . Here we will investigate some QDDM models with exponential potentials (39) such that for a wide range of initial conditions the scalar field settles into an attractor solution that depends only upon a few nonequilibrium thermodynamical parameters, addressing the coincidence problem.
### A Linear dissipative regime
The linear regime $`\zeta =\alpha H`$, with $`\alpha `$ a constant in the interval $`0<\alpha <1`$, arises for instance when the coefficient of bulk viscosity takes the form of a radiating fluid. We further assume that the number of interactions of a generic CDM particle in a Hubble time is larger than unity so that the hydrodynamic regime is respected. Here we investigate models with the limiting behavior $`\gamma 2/3`$ in the asymptotic regime. We begin by inserting the ansatz
$$\gamma =\frac{2}{3}\left(1+\chi \right)$$
(41)
in equation (25). It is immediately seen that the latter splits in two equations, namely
$$\ddot{H}+(2+\nu )H\dot{H}+\nu \left(1\frac{3}{2}\alpha \right)H^3=0,$$
(42)
and
$$\left[H^2+\frac{k}{a^2}\right]\dot{\chi }+\left[2H\left(\dot{H}\frac{k}{a^2}\right)+\tau ^1\left(H^2+\frac{k}{a^2}\right)\right]\chi =0.$$
(43)
Replacing the solution of (43) in (41) it follows,
$$\gamma =\frac{2}{3}\left(1+b\frac{a^{2\nu }}{k+\dot{a}^2}\right),$$
(44)
where $`b`$ is an arbitrary integration constant. This expression for $`\gamma `$ will be sensible only if it meets the restriction $`0\gamma 2`$. Then (42) can be transformed into a linear differential equation of second order whose general solution in parametrized form is already known , . In particular $`at^\sigma `$, where $`\sigma `$ is the largest root of $`\nu (13\alpha /2)\sigma ^2(2+\nu )\sigma +2=0`$, is an asymptotic stable solution in the limit $`t\mathrm{}`$. Then, using (10) together with the attractor conditions $`\mathrm{\Omega }=1`$, $`\gamma =2/3`$ and (16), we find
$$\sigma =\frac{2\left(1\mathrm{\Omega }_{m0}\right)}{23\mathrm{\Omega }_{m0}\gamma _m}.$$
(45)
Hence the quintessence adiabatic index $`\gamma _\varphi =2/3\sigma `$ depends solely on the dark matter parameters $`\mathrm{\Omega }_m`$ and $`\gamma _m`$ in the asymptotic regime. Moreover, a relationship between the dissipative parameters $`\alpha `$ and $`\nu `$ follows from (29) and (16), namely
$$\alpha =\mathrm{\Omega }_{m0}\left(1\frac{2}{\nu \sigma }\right)\left(\gamma _m\frac{2}{3\sigma }\right),$$
(46)
and the requirement $`\alpha >0`$ implies $`\sigma \nu >2`$ or $`\nu >\nu _{min}=3\gamma _\varphi `$. The same condition arises from the requirement that $`\gamma 2/3`$ when $`t\mathrm{}`$. Equations (45) and (46) show that $`\alpha `$ grows with $`\nu `$ attaining $`\alpha _{max}=(3\gamma _m2)\mathrm{\Omega }_m/3(1\mathrm{\Omega }_m)`$ in the limit $`\nu \mathrm{}`$.
As it follows from (45) there will be accelerated expansion (i.e., $`\sigma >1`$) if $`\mathrm{\Omega }_{m0}<2/3\gamma _m<2/3`$ and accordingly we obtain a family of exact solutions describing a QDDM scenario that solves the coincidence problem regardless of the value of the spatial curvature. Figure 1 depicts the dependence of $`\gamma _\varphi `$ on $`\mathrm{\Omega }_m`$ when $`\gamma _m=1`$. To make a rough estimate of the cosmological parameters in the late time era we assume that our Universe is currently close to the asymptotic attractor regime and use the current observational bounds. After the combination of low redshift, type Ia supernovae and COBE measurements determines (for a spatially flat universe) the range $`\mathrm{\Omega }_m0.30.4`$ and $`\gamma _\varphi <0.6`$. From figure 1 it is seen that our linear dissipative model satisfies comfortably these constraints. For $`\mathrm{\Omega }_m=0.3`$ we get from (45) $`\sigma 1.27`$. This is fully consistent with current estimations of $`Ht`$ today , as they provide a lower bound for $`\sigma `$ in a universe that started only recently a phase of accelerated expansion and approaches asymptotically the attractor regime.
### B Viscous speed regime
This scenario is somewhat more general than the previous one. It arises when the bulk viscosity coefficient is given in terms of the speed of the bulk viscous signal $`v`$ by
$$\frac{\zeta }{\tau }=v^2\gamma _m\rho _m.$$
(47)
We further assume $`\tau `$ related to $`H`$ by the same expression as before, only that to simplify the calculations we now take $`\nu `$ constant. Hence (25) becomes
$$\nu ^1\left[h^{\prime \prime }+3\gamma h^{}+3\gamma ^{}h9\frac{v^2\gamma _mϵ}{1+ϵ}h\right]+h^{}+3\gamma h=0,$$
(48)
where
$$hH^2+\frac{k}{a^2},$$
(49)
and the prime indicates derivative with respect to $`\eta \mathrm{ln}a`$. Here we have used the scale factor $`a(t)`$ as a coordinate instead of the cosmological time $`t`$, i.e., $`a`$ is assumed to be a monotonic function of $`t`$.
Equation (48) is useful to study the asymptotic stability of FLRW expansions at late time because it can be rewritten in terms of the derivative of a Lyapunov function
$$\frac{d}{d\eta }\left\{\frac{1}{2}h^2+\frac{3}{2}\left[\gamma ^{}+\gamma \nu 3\frac{v^2\gamma _mϵ}{1+ϵ}\right]h^2\right\}=\left(3\gamma +\nu \right)h^2$$
$$+\frac{3}{2}\left[\gamma ^{}+\gamma \nu 3\frac{v^2\gamma _mϵ}{1+ϵ}\right]^{}h^2.$$
(50)
If the adiabatic index does not decreases too fast in the attractor era, a sufficient condition for the Lyapunov function to have a minimum at the phase space point $`(h,h^{})=(0,0)`$ is that $`\nu >3v^2`$. Within this scenario the parameters of the matter fluid can be taken as quasistatic, yielding an asymptotically stable minimum. The leading behavior of the solutions in this quasistatic regime is given by
$$h^2=c_1a^{\lambda _1}+c_2a^{\lambda _2},$$
(51)
where
$$\lambda _{1,2}=\frac{1}{2}\left\{\left(3\gamma +\nu \right)\pm \left[\left(3\gamma \nu \right)^2+36\gamma _mv^2\mathrm{\Omega }_m\right]^{1/2}\right\}$$
(52)
and the parameters are evaluated at the asymptotic attractor era. For large scale factor, equation (51) reduces to $`h^2c_1a^{\lambda _1}`$ when $`2<\lambda _1<0`$, and we have once again a power-law accelerated cosmic expansion with $`\sigma =2/\lambda _1`$. Likewise equation (51) reduces to $`h^2c_1a^2`$ for $`\lambda _1=2`$, and $`h^20`$ for $`\lambda _1<2`$. These two latter cases correspond to linear evolutions at late time.
Combining (29) and (16) we get
$$\lambda ^2+\left(3\gamma _m+\nu \right)\lambda +3\gamma _m(\nu 3v^2)=0$$
(53)
and using (10) together with the attractor constraints $`\mathrm{\Omega }=1`$ and (16) we find that $`\lambda _1=3\gamma _\varphi `$ also satisfies (53). Hence we obtain
$$\gamma _\varphi =\frac{1}{6}\left\{3\gamma _m+\nu \left[\left(3\gamma _m\nu \right)^2+36\gamma _mv^2\right]^{1/2}\right\}.$$
(54)
In this model the quintessence adiabatic index does depend on the parameters $`\nu `$ and $`v`$ while it is independent of the density parameter $`\mathrm{\Omega }_m`$. We plot $`\gamma _\varphi `$ in figure 2 for $`\gamma _m=1`$. As it can be seen a wide range of the parameter space $`(\nu ,v)`$ is consistent with a spatially–flat accelerated universe and such that $`\gamma _\varphi <0.6`$ for any value of $`\mathrm{\Omega }_m`$. This shows another solution to the coincidence problem for any value of the spatial curvature and compatible with accelerated expansion. It is also seen that dissipative effects enlarge the parameter space where observational data has to be fitted, but global dynamic information alone cannot determine the specific values of $`\nu `$ and $`v`$. Note that the smaller $`\gamma _\varphi `$, the larger the dissipative contribution to the sound speed, and the smaller the interaction rate.
## V Discussion
We have proved that the coincidence problem and an accelerated expansion phase of FLRW cosmologies cannot be simultaneously addressed by the combined effect of a perfect fluid and Q-matter. Nonetheless, if nonbaryonic dark matter behaves as a dissipative fluid rather than a perfect one, both problems may find a simultaneous solution. This is so because an imperfect (i.e., dissipative fluid) expanding in a FLRW background possess a negative pressure $`\pi `$ that enters the conservation equations of general relativity. The models presented here are compatible with a negative deceleration parameter at present time. In consequence, the quintessence scenario becomes more robust when the dissipative effect of the nonequilibrium pressure arising in the CDM gas is allowed into the picture.
Recently attempts have been made to constraint the state equation of the cosmic fluids (Q–matter included) by considering gravitational lensing effects, the mass power spectrum and the anisotropies of the cosmic backgroud radiation ,. We have shown specific models with an ample region in the space of out–of–equilibrium thermodynamic parameters satisfying this constraint in the asymptotic attractor regime which our Universe may well be approaching. We would like to point out that the parameter space should be enlarged by adding these out–of–equilibrium parameters when fitting the observational data. Unfortunately there is some degenerancy in the determination of these parameters from constraints arising from the cosmological dynamics alone. We hope, however, that simulations of structure formation that include dissipative effects will ultimately prove instrumental in discriminating between different models.
On the attractor asymptotic regime the dissipative matter fluid and the scalar field contribute in a fixed ratio to the pressure and energy density along the QDDM era. This scenario ameliorates the self–adjusting model as it allows for $`0<\gamma <2`$ for a wide range of initial conditions. It also improves on the tracking models as it solves the coincidence problem in the late accelerated expansion phase. While keeping a finite difference between the adiabatic indices of quintessence and matter fluid, this difference arises in the viscous pressure.
As it has been noted the Q–matter proposal may entail some undesirable effects such as the variation of the constant of nature and the presence of unobserved long range forces. Efforts to solve this difficulties by coupling the quintessence field to the electromagnetic field , and to the curvature of the metric have been made. Further, a time dependent but otherwise smooth scalar field such as those studied so far, are somehow unphysical as they violate the principle of equivalence -the Q–matter must experience clustering in some degree and so it cannot be entirely smooth. Therefore to respect the equivalence principle one should assume that $`\varphi `$ varies with position as well. Accordingly one should be led to forsake the FLRW metric an take up some inhomogeneous one instead, only that in such a case the computational effort is bound to be enormous and most likely no exact solution will emerge.
Despite that the Q–matter proposal cannot be regarded at this moment with unreserved confidence, we feel this idea is still worth exploring in the hope that the aforesaid difficulties may soon find a satisfactory answer.
## Acknowledgments
This work has ben patially supported by the Spanish Ministry of Education under grant PB94-0718. LPC and ASJ thank the University of Buenos Aires for partial support under project TX-93.
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# University of Wisconsin - Madison MADPH-00-1168March 2000 OVERVIEW OF NEUTRINO OSCILLATION PHYSICS11footnote 1Talk presented at the 7th International Symposium on Particles, Strings and Cosmology (PASCOS 99), Granlibakken, California, December 1999
## 1 Neutrino Masses
Tree-level mass generation occurs through the Higgs mechanism. The Dirac mass $`m_D`$ arises in a lepton conserving ($`\mathrm{\Delta }L=0`$) interaction and requires a right-handed neutrino. A Majorana mass $`m_M`$ occurs through a $`\mathrm{\Delta }L=2`$ process with only a left-handed light neutrino field and a heavy isosinglet intermediate field $`N^c`$. Then the see-saw mechanism with $`m_D10^2`$ GeV and $`m_M>10^{12}`$ GeV generates light neutrinos
$$m_\nu =m_D^2/m_M$$
(1)
that are nearly Majorana. In the case that $`m_Dm_M\mathrm{eV}`$, as can be realized in some models, active–sterile neutrino oscillations can take place.
Neutrino mass can alternatively be generated radiatively by new interactions, such as the $`R`$-parity violating $`\nu b\stackrel{~}{b}`$ interaction. Recent reviews of theoretical models for neutrino mass generation are given in Ref. .
## 2 Three-Neutrino Oscillations in Vacuum and Matter
The relation of the three-neutrino flavor eigenstates to the mass eigenstates is
$$\left(\begin{array}{c}\nu _e\\ \nu _\mu \\ \mu _\tau \end{array}\right)=U\left(\begin{array}{c}\nu _1\\ \nu _2\\ \nu _3\end{array}\right),$$
(2)
where $`U`$ is the $`3\times 3`$ Maki-Nakagawa-Sakata (MNS) mixing matrix. It can be parametrized by
$`U=\left(\begin{array}{ccc}c_{13}c_{12}& c_{13}s_{12}& s_{13}e^{i\delta }\\ c_{23}s_{12}s_{13}s_{23}c_{12}e^{i\delta }& c_{23}c_{12}s_{13}s_{23}s_{12}e^{i\delta }& c_{13}s_{23}\\ s_{23}s_{12}s_{13}c_{23}c_{12}e^{i\delta }& s_{23}c_{12}s_{13}c_{23}s_{12}e^{i\delta }& c_{13}c_{23}\end{array}\right)\left(\begin{array}{ccc}1& 0& 0\\ 0& e^{i\varphi _2}& 0\\ 0& 0& e^{i(\varphi _3+\delta )}\end{array}\right)`$ (9)
where $`c_{ij}=\mathrm{cos}\theta _{ij}`$ and $`s_{ij}=\mathrm{sin}\theta _{ij}`$. The extra diagonal phases are present for Majorana neutrinos but do not affect oscillation phenomena.
With three neutrinos there are two independent $`\delta m^2`$ and $`\delta m_a^2\delta m_b^2`$ is indicated by the oscillation evidence. The vacuum oscillation probabilities are
$$P(\nu _\alpha \nu _\beta )=A_{\alpha \beta }\mathrm{sin}^2\mathrm{\Delta }_aB_{\alpha \beta }\mathrm{sin}^2\mathrm{\Delta }_b+ϵ_{\alpha \beta }J\mathrm{sin}2\mathrm{\Delta }_b,$$
(10)
where $`\mathrm{\Delta }_a\delta m_a^2L/4E_\nu `$. $`A_{\alpha \beta }`$ is the amplitude of the leading oscillation, $`B_{\alpha \beta }`$ the amplitude of the sub-leading oscillation and $`J`$ the CP-violating amplitude; all are determined by the $`U`$ matrix elements. The physical variable is $`L/E_\nu `$, where $`L`$ is the baseline from source to detector and $`E_\nu `$ is the neutrino energy.
In matter, $`\nu _e`$ scatter differently from $`\nu _\mu `$ and $`\nu _\tau `$, and the effective neutrino mixing amplitude in matter can be very different from the vacuum amplitude. For the leading oscillation the matter and vacuum oscillation mixings are related in the approximation of constant matter density by
$$\mathrm{sin}^22\theta _{13}^m=\frac{\mathrm{sin}^22\theta _{13}}{\left(\mathrm{cos}2\theta _{13}A/\delta m_a^2\right)^2+\mathrm{sin}^22\theta _{13}},$$
(11)
where
$$A=2\sqrt{2}G_FY_e\rho (x)E_\nu .$$
(12)
Here $`Y_e`$ is the electron fraction and $`\rho (x)`$ is the density at path-length $`x`$. The $`\nu _e\nu _\mu `$ (or $`\nu _e\nu _\tau `$) oscillation argument in matter of constant density is
$$\mathrm{\Delta }_a^m=\frac{1.27\delta m_a^2(\mathrm{eV}^2)L(\mathrm{km})}{E_\nu (\mathrm{GeV})}\sqrt{\left(\frac{A}{\delta m_a^2}\mathrm{cos}2\theta _{13}\right)^2+\mathrm{sin}^22\theta _{13}}.$$
(13)
Resonance enhancements in matter are possible for $`\delta m_a^2>0`$, while suppression occurs for $`\delta m_a^2<0`$. It is significant that the resonant energies correspond to neutrino energies relevant to the atmospheric and solar anomalies.
$`\mathrm{Earth}:`$ $`E_\nu 15\mathrm{GeV}\left({\displaystyle \frac{\delta m_a^2}{3.5\times 10^3\mathrm{eV}^2}}\right)\left({\displaystyle \frac{1.5\mathrm{gm}/\mathrm{cm}^3}{\rho Y_e}}\right)`$ (14)
$`\mathrm{Sun}:`$ $`E_\nu 10\mathrm{MeV}\left({\displaystyle \frac{\delta m_b^2}{10^5\mathrm{eV}^2}}\right)\left({\displaystyle \frac{10\mathrm{g}/\mathrm{cm}^3}{\rho Y_e}}\right).`$ (15)
### 2.1 Atmospheric Neutrino Oscillations
The Kamiokande, SuperKamiokande (SuperK), Macro and Soudan atmospheric neutrino measurements show a $`\mu /e`$ ratio that is about 0.6 of expectations. The SuperK experiment has established the dependence of the $`e`$ and $`\mu `$ event rates on zenith angle, or equivalently the baseline $`L`$. $`N(e)`$ is independent of $`L`$ and validates the $`\nu _e`$ flux calculation (within 20%). $`N(\mu )`$ depletion increases with $`L`$. The muon event distributions are will described by $`\nu _\mu \nu _\tau `$ vacuum oscillations with a $`\nu _\mu `$ survival probability
$$P(\nu _\mu \nu _\mu )=1A_{\mu \tau }\mathrm{sin}^2(1.27\delta m_a^2L/E_\nu )$$
(16)
with $`\delta m_a^2=3.5\times 10^3\mathrm{eV}^2`$ and maximal or near maximal amplitude
$$A_{\mu \tau }=1_{0.2}^{+0.0}(i.e.,|\theta _{32}45^{}|<13^{}).$$
(17)
With further data accumulation a slightly lower central value is now indicated ($`\delta m_a^2=2.8\times 10^3\mathrm{eV}^2`$). The $`\nu _\mu \nu _\tau `$ oscillations are not resolved due to smearing of $`L`$ and inferred $`E_\nu `$ values, and equally good fits to the SuperK data are found with oscillation and neutrino decay ($`\nu _2\overline{\nu }_4+J`$) models. For now we assume the simplest interpretation of the SuperK data, namely oscillations. We note that in the CHOOZ reactor experiment $`\overline{\nu }_e`$ disappearance is not observed at the $`\delta m_a^2`$ scale and the corresponding constraint on 3-neutrino mixing for $`\delta m_a^2=3.5\times 10^3\mathrm{eV}^2`$ is
$$A_{\mu e}<0.2,|U_{e3}|<0.23,\theta _{13}<13^{}.$$
(18)
### 2.2 Solar Neutrino Oscillations
The solar neutrino experiments sample different $`\nu _e`$ energy ranges and find different flux deficits compared to the Standard Solar Model (SSM) as follows:
$$\begin{array}{ccc}\nu _e{}_{}{}^{71}\mathrm{Ga}{}_{}{}^{71}\mathrm{Ge}e\hfill & \begin{array}{c}\mathrm{GALLEX}\hfill \\ \mathrm{SAGE}\hfill \end{array}\hfill & \begin{array}{c}0.60\pm 0.06\hfill \\ 0.52\pm 0.06\hfill \end{array}\hfill \\ \nu _e{}_{}{}^{37}\mathrm{Cl}{}_{}{}^{37}\mathrm{Ar}e\hfill & \mathrm{Homestake}\hfill & 0.33\pm 0.03\hfill \\ \nu e\nu e\hfill & \mathrm{SuperK}\hfill & 0.47\pm 0.02\hfill \end{array}$$
(19)
Thus the $`\nu _e`$ survival probability is inferred to be energy dependent.
Global oscillation fits have been made using floating <sup>8</sup>B and hep flux normalizations which are somewhat uncertain in the SSM. The relative normalizations from the fits range from 0.5 to 1.2 for the <sup>8</sup>B flux and 1 to 25 for the hep flux. These global fits include the data on (i) total rates, assuming all the experiments are okay — the different Cl suppression ratio plays a vital role; (ii) the night-day asymmetry, which is observed at the 2$`\sigma `$ level — the large angle matter solution gives night rates $`>`$ day rates; (iii) seasonal dependence beyond $`1/r^2`$, which can occur for vacuum solutions. The oscillation analyses generally agree on the allowed $`\delta m_{21}^2`$ and $`\mathrm{sin}^22\theta _{12}`$ regions for acceptable solutions. Typical candidate solar solutions are given in Table 1. In the case of vacuum oscillations (VO) several discrete regions of $`\delta m_{21}^2`$ are possible.
### 2.3 3-Neutrino Mixing Matrix
Once the solar oscillation solution is pinned down, and $`\theta _{12}`$ is thus determined, we will have approximate knowledge of the mixing angles of the 3-neutrino matrix, with $`\theta _{23}\pi /4`$ and $`\theta _{13}0`$ from the atmospheric and CHOOZ data. Upcoming experiments are expected to shed light on the solar solution. In the SNO experiment, which is now taking data, and the forthcoming ICARUS experiment, the high energy $`\nu _e`$ CC events may distinguish LAM, SAM, and LOW solutions with large hep flux contributions from the VO or the SAM sterile neutrino solutions. Also, the neutral-current to charged-current ratio will distinguish active from sterile oscillations. The Borexino experiment can measure the VO seasonal variation of the <sup>7</sup>Be line flux. The KamLand reactor experiment to measure the $`\overline{\nu }_e`$ survival probability will be sensitive to the LAM and LOW solar solutions.
The CP phase $`\delta `$ may be measurable at a neutrino factory if the solar solution is LAM. The CP violation comes in only at the sub-leading oscillation scale. An apparent CP-odd asymmetry is induced by matter.
### 2.4 Models
For maximal mixing in both atmospheric and solar sectors, there is an unique mixing matrix
$$U=\left(\begin{array}{ccc}1/\sqrt{2}& 1/\sqrt{2}& 0\\ 1/2& 1/2& 1/\sqrt{2}\\ 1/2& 1/2& 1/\sqrt{2}\end{array}\right).$$
(20)
In this bimaximal mixing model, there would be no CP-violating effects. However, because $`U_{e3}=0`$, long-baseline experiments would have some sensitivity to the sub-leading LAM solar scale oscillations.
Many unification models predict that the neutrino masses are Majorana and hierarchical, there is no cosmologically significant dark matter, and the SAM solar solution (small $`\theta _{12}`$ mixing) obtains.
### 2.5 Beyond 3 Neutrinos
The LSND evidence for $`\nu _\mu \nu _e`$ oscillations with $`\delta m^21\mathrm{eV}^2`$, $`\mathrm{sin}^22\theta 10^2`$ requires a $`\delta m^2`$ scale distinct from the atmospheric and solar oscillation scales, and thus a sterile neutrino state would be needed to explain all the oscillation phenomena. Then to also satisfy limits from CDHS accelerator and Bugey reactor experiments, the mass hierarchy must be two separated pairs. Such a scenario would allow even more interesting effects at a neutrino factory, such as large CP violation, since both the leading and sub-leading oscillation scales would be accessible. The MiniBooNE experiment will settle whether the LSND evidence is real. Other interest in sterile neutrinos comes from $`r`$-process nucleosynthesis if it occurs in supernovae.
## 3 Long-Baseline Experiments
Long-baseline experiments are needed to (i) confirm the atmospheric evidence for $`P(\nu _\mu \nu _\mu )`$ at accelerators; (ii) resolve the leading $`\nu _\mu \nu _\mu `$ oscillation and exclude the neutrino decay possibility; (iii) precisely measure $`|\delta m_a^2|`$; (iv) exclude $`\nu _\mu \nu _s`$ disappearance, although SuperK has now shown that oscillations to sterile neutrinos are excluded at 99% CL; (v) measure $`|U_{e3}|`$ from $`\nu _e\nu _\mu `$ appearance, which requires a muon decay source for the neutrino beam; (vi) determine the sign of $`\delta m_a^2`$ from matter effects in the Earth’s crust; and (vii) search for CP violation.
The first long-baseline experiments will measure the energy dependence of the produced muons and measure the neutral-current to charged-current ratio, to partially address the first four issues listed above. The K2K experiment from KEK to SuperK is in operation, with a baseline $`L=250`$ km and mean neutrino energy $`E_\nu =1.4`$ GeV. The MINOS experiment from Fermilab to Soudan, with $`L=732`$ km and possible energies of $`E_\nu =3,6,12`$ GeV will begin in 2002. A 10% precision on $`|\delta m_a^2|`$ may ultimately be possible at MINOS. The ICANOE and OPERA long-baseline experiments from CERN to Gran Sasso with $`L743`$ km have been approved.
Muon storage rings could provide intense neutrino beams ($`10^{19}\text{}10^{21}`$ per year) that would yield thousands of charged-current neutrino interactions in a reasonably sized detector (10–50 kt) anywhere on Earth. These neutrino factories would have pure neutrino beams ($`\nu _e,\overline{\nu }_\mu `$ from stored $`\mu ^+`$ and $`\overline{\nu }_e,\nu _\mu `$ from stored $`\mu ^{}`$) with 50% $`\nu _e`$ or $`\overline{\nu }_e`$ components. Detection of wrong-sign muons (the muons with opposite sign to the charge current from the beam muon neutrino) would signal $`\nu _e\nu _\mu `$ or $`\overline{\nu }_e\overline{\nu }_\mu `$ appearance oscillations. We now discuss the capability of a neutrino factory with $`2\times 10^{20}`$ muons a year and a 10 kt detector to resolve the issues raised in the preceding section.
With an $`E_\mu =30`$ GeV storage ring at a baseline of $`L=2800`$ km, a statistical precision of a few % on $`\mathrm{sin}^22\theta _{23}`$ is possible in $`\nu _\mu `$ survival measurements. This accuracy in measuring $`\mathrm{sin}^22\theta _{23}`$ would differentiate the bimaximal model prediction of $`\mathrm{sin}^22\theta _{23}=1`$ from the democratic model prediction of $`\mathrm{sin}^22\theta _{23}=8/9`$.
With stored muon energies $`E_\mu =10`$ to 50 GeV and baselines of $`L=732`$ to 7332 km and 1 kt detector, there would be hundreds of events per year from $`\nu _\mu \nu _\tau `$ oscillations for an intensity of $`2\times 10^{20}`$ neutrinos.
The wrong-sign muon event rates are approximately proportional to $`\mathrm{sin}^22\theta _{13}`$. For non-zero $`\mathrm{sin}^22\theta `$ the observation of $`\nu _e\nu _\mu `$ and $`\overline{\nu }_e\overline{\nu }_\mu `$ appearance oscillations at baselines long enough to have significant matter effects will allow a determination of the sign of $`\delta m_{32}^2`$, and thus determine the pattern of the masses (a $`1+2`$ mass hierarchy versus a $`2+1`$ hierarchy for three neutrinos). A proof of the principle that the sign of $`\delta m^2`$ can be so determined has been given for a baseline $`L=2800`$ km. In $`\mu ^+`$ appearance, $`\delta m_{32}^2>0`$ gives a smaller rate and harder spectrum than $`\delta m_{32}^2<0`$, while the results are opposite in $`\mu ^{}`$ appearance.
In optimizing $`E_\mu `$ and $`L`$ for long-baseline experiments to find the sign of $`\delta m_{32}^2`$, $`L=732`$ km is too short (matter effects are small) and $`L=7332`$ km is too far (event rates are low). The sensitivity to determine the sign of $`\delta m_{32}^2`$ improves linearly with $`E_\mu `$. There is a tradeoff between energy, detector size and muon beam intensity.
## 4 Absolute Neutrino Masses
Oscillation phenomena determine only mass-squared differences, leaving the absolute mass scale unknown. However, because the atmospheric and solar $`\delta m^2`$ values are $`(1\mathrm{eV})^2`$, all mass eigenvalues are approximately degenerate if at the $`1`$ eV scale. Thus all neutrino mass eigenvalues are bounded by the tritium limit from the Troitsk and Mainz experiments,
$$m_j<3\mathrm{eV}\mathrm{for}j=1,2,3.$$
(21)
Neutrinoless double-beta decay (0$`\nu \beta \beta `$) provides a probe of Majorana neutrino mass. The rate is proportional to the $`\nu _e\nu _e`$ element of the neutrino mass matrix. The present limit from the Heidelberg experiment is
$$M_{\nu _e\nu _e}<0.2\mathrm{eV}\times f,$$
(22)
where the factor $`f`$ represents uncertainty in the nuclear matrix elements, which might be as large as a factor of 3. The $`0\nu \beta \beta `$ limit translates to a bound on the summed neutrino Majorana masses of
$$m_\nu <0.75\mathrm{eV}\times f$$
(23)
in the SAM solar solution. No similar constraints apply to the LAM, LOW or VO solutions where the bound can be satisfied by having opposite CP parity of $`\nu _1`$ and $`\nu _2`$ mass eigenstates. Future experiments may probe to
$$|M_{\nu _e\nu _e}|=0.01\mathrm{eV},$$
(24)
which would provide sensitivity down to
$$m_\nu =0.08\mathrm{eV}\times f$$
(25)
in the SAM solution.
Measurements of the power spectrum by the MAP and PLANCK satellites may determine $`m_\nu `$ down to $`0.4`$ eV . The heights of the acoustic peaks can also decide how the mass is distributed among the neutrino eigenstates.
## 5 Summary
We have entered an exciting new era in the study of neutrino masses and mixing. From the SuperK evidence on atmospheric neutrino oscillations, we already have a surprising amount of information about the neutrino mixing matrix (near maximal $`\mathrm{sin}^22\theta _{23}`$ and near minimal $`\mathrm{sin}^22\theta _{13}`$). The SuperK, SNO, Borexino, KamLand, and ICARUS experiments are expected to differentiate among the candidate solar oscillation possibilities and determine $`\mathrm{sin}^22\theta _{12}`$. MiniBooNE will tell us whether a sterile neutrino is mandated. Neutrino factories will study the leading oscillations, determine the sign of $`\delta m_a^2`$, measure $`U_{e3}`$, and possibly detect CP violation. The GENIUS $`0\nu \beta \beta `$ experiment and the MAP and PLANCK satellite measurements of the power spectrum will probe the absolute scale of neutrino masses. There is a synergy of particle, physics, nuclear physics, and cosmology occurring in establishing the fundamental properties of neutrinos. A theoretical synthesis should emerge from these experimental pillars.
A more complete version of this review, including figures and more extensive references, can be found in Ref. .
## Acknowledgments
I thank my collaborators K. Whisnant, T. Weiler, S. Pakvasa, S. Geer, R. Raja, J. Learned, P. Lipari, and M. Lusignoli. This research was supported in part by the U.S. Department of Energy under Grant No. DE-FG02-95ER40896 and in part by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation.
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# Sonic black holes in dilute Bose-Einstein condensates
## I Introduction
Many investigations of dilute gas Bose-Einstein condensates are directed towards experimentally creating nontrivial configurations of the semiclassical mean field, or to predicting the properties of such configurations in the presence of quantum fluctuations. Such problems are hardly peculiar to condensates: the quantum neighborhoods of interesting classical backgrounds are important areas of research in most fields of physics. But ultracold dilute gases are so easy to manipulate and control, both experimentally and theoretically , that they may allow us to decipher less amenable systems by analogy. As an essay in such an application of condensates, in this paper we discuss the theoretical framework and propose an experiment to create the analog of a black hole in the laboratory and simulate its radiative instabilities.
It is now commonly believed that, even in the context of elementary particle physics, quantum field theory arises from a still unknown underlying structure: it is an effective dynamical theory, describing the low energy limit of collective phenomena of the underlying microscopic theory. From this viewpoint, our description of (more) fundamental phenomena, such as gravity or electromagnetism, is actually similar to the theoretical descriptions of many phenomena of condensed matter. To understand superfluidity, superconductivity or dilute Bose-Einstein condensation, we describe the dynamics of the system in terms of collective modes (quasiparticles) whose typical size is much larger than the distances between the particles that constitute the underlying medium; but even electrons and photons must be considered as the quasiparticles of the deeper theory we do not yet know. In this sense we may say that the major difference between our fundamental theories and those we use in condensed matter is that in the latter case the next microscopic level of description is actually well understood.
With this fundamental background in mind, it is not so surprising that condensed matter analogs of nontrivial configurations appearing in relativistic quantum field theories and gravitation can be constructed. For example, <sup>3</sup>He has been proposed as a laboratory counterpart of high-energy particle physics. It has been argued that, under appropriate conditions, excitations around the ground state of the system may resemble the particle-spectrum of gauge theories of high energy physics . These condensed matter systems have also been used to simulate topological defects characteristic of gauge theories and which are considered to have played a cosmological role in the early stages of the evolution of the universe such as monopoles and cosmic strings .
The past decade has witnessed an increasing interest in simulating gravitational configurations and processes in condensed matter systems in the laboratory. The key observation was originally made by Unruh and further analyzed by Visser : phononic propagation in a fluid is described by a wave equation which, under appropriate conditions, can be interpreted as propagation in an effective relativistic curved spacetime background, the spacetime metric being entirely determined by the physical properties of the fluid under study, namely, its density and flow velocity. Unruh urged a specific motivation for examining the hydrodynamic analogue of an event horizon , namely that as an experimentally and theoretically accessible phenomenon it might shed some light on the Hawking effect (thermal radiation from black holes, stationary insofar as the back reaction is negligible). In particular, one would like to gain insight into the role in the Hawking process of ultrahigh frequencies .
An event horizon for sound waves appears in principle wherever there is a surface through which a fluid flows at the speed of sound, the flow being subsonic on one side of the surface and supersonic on the other. There is a close analogy between sound propagation on a background hydrodynamic flow, and field propagation in a curved spacetime; and although hydrodynamics is only a long-wavelength effective theory for physical (super)fluids, so also field theory in curved spacetime is to be considered a long-wavelength approximation to quantum gravity . Determining whether and how sonic black holes radiate sound, in a full calculation beyond the hydrodynamic approximation or in an actual experiment, can thus offer some suggestions about black hole radiance and its sensitivity to high frequency physics (beyond the Planck scale). The possibility that such high frequencies might have consequences for observably low frequency phenomena is one of the main reasons that black holes have deserved much attention: there is reason to expect that an event horizon can act as a microscope, giving us a view into physics on scales below the Planck length. This is because modes coming from an event horizon are redshifted into the low-energy regime as they propagate out to be observed far away from the black hole. Conversely, if we imagine tracking the observed signal back towards its source, the closer we get to the horizon, the shorter the wave length of the signal must be, until at the very horizon we must either reach infinite energy scales, or encounter a breakdown in general relativity and quantum field theory in curved spacetime .
In understanding this problem, hydrodynamic and condensed matter analogs of black holes may offer some of the experimental guidance otherwise difficult to obtain in the case of gravity. Under appropriate conditions and approximations (which can be basically summarized in the requirement the wave lengths of the perturbations be sufficiently large), the propagation of collective fluctuations (phonons) admits an effective general relativistic description, in terms of a spacetime metric. This long wavelength regime would correspond analogically to quantum field theory in curved spacetime. The effective phonon metric may describe black holes, as in general relativity, and so a phonon Hawking effect may be possible; and certainly the problem of arbitrarily high frequencies at the horizon is also present. But in this case, when at short wave lengths the metric approximation is no longer valid and a more microscopic theory must be used instead, the accurate microscopic theory is actually known. And if the hydrodynamic system is a dilute Bose-Einstein condensate, the microscopic theory is actually tractable enough that we can make reliable calculations from first principles. As we will argue, trapped bosons at ultralow temperature can indeed provide an analogue to a black-hole spacetime. Similar analogues have been proposed in other contexts, such as superfluid helium , solid state physics , and optics ; but the outstanding recent experimental progress in cooling, manipulating and controlling atoms make Bose-Einstein condensates an especially powerful tool for this kind of investigation.
The basic challenge of our proposal is to keep the trapped Bose-Einstein gas sufficiently cold and well isolated to maintain a locally supersonic flow long enough to observe its intrinsic dynamics. Detecting thermal phonons radiating from the horizons would obviously be a difficult additional problem, since such radiation would be indistinguishable from many other possible heating effects. This further difficulty does not arise in our proposal, however, because the black-hole radiation we predict is not quasistationary, but grows exponentially under appropriate conditions. It should therefore be observable in the next generation of atom traps, and may also raise new issues in the theory of gravitational black holes.
In this paper, we extend and generalize the results of Ref. , including a more detailed analysis of the model considered there as well as a new qualitatively different case. The paper is organized as follows. In Sec. II, we show how sonic horizons in dilute Bose-Einstein condensates may appear in the hydrodynamic approximation, discuss the regime of validity of such approximation, as well the validity of one-dimensional models. Section III is devoted to the study of sonic horizons in condensates subject to tight ring-shaped external potentials. We present numerical results showing that both stable and unstable black holes may be created under realistically attainable conditions in current or near-future laboratories. We also study the nature of the dynamical instabilities that appear for certain configurations. In Sec. IV, we discuss a different configuration, namely that of a sink-generated black hole in an infinite one-dimensional condensate, and show that there also exist black-hole configurations, although they are not stable. We summarize and conclude in Sec. V. The Appendix is devoted to the issues of redundancy and normalization of the dynamically unstable Bogoliubov modes (associated with complex eigenfrequencies).
## II Sonic black holes in condensates
A Bose-Einstein condensate is the ground state of a second quantized many body Hamiltonian for $`N`$ interacting bosons trapped by an external potential $`V_{\mathrm{ext}}(𝐱)`$ . At zero temperature, when the number of atoms is large and the atomic interactions are sufficiently small, almost all the atoms are in the same single-particle quantum state $`\mathrm{\Psi }(𝐱,t)`$, even if the system is slightly perturbed. The evolution of $`\mathrm{\Psi }`$ is then given by the well-known Gross-Pitaevskii equation, which in appropriate units can be written as
$$i\mathrm{}_t\mathrm{\Psi }=\left(\frac{\mathrm{}^2}{2m}^2+V_{\mathrm{ext}}+\frac{4\pi a\mathrm{}^2}{m}|\mathrm{\Psi }|^2\right)\mathrm{\Psi },$$
where $`m`$ is the mass of the individual atoms and $`a`$ is the scattering length. The wave function of the condensate is normalized to the total number of atoms $`d^3𝐱|\mathrm{\Psi }(𝐱,t)|^2=N`$.
Our purposes do not require solving the Gross-Pitaevskii equation with some given external potential $`V_{\mathrm{ext}}(𝐱)`$; our concern is the propagation of small collective perturbations of the condensate, around a background stationary state
$$\mathrm{\Psi }_s(𝐱,t)=\sqrt{\rho (𝐱)}e^{i\vartheta (𝐱)}e^{i\mu t/\mathrm{}},$$
where $`\mu `$ is the chemical potential. Thus it is only necessary that it be possible, in any external potential that can be generated, to create a condensate in this state. Indeed, many realistic techniques for “quantum state engineering,” to create designer potentials and bring condensates into specific states, have been proposed, and even implemented successfully ; and our simulations indicate that currently known techniques should suffice to generate the condensate states that we propose.
Perturbations about the stationary state $`\mathrm{\Psi }_s(𝐱,t)`$ obey the Bogoliubov system of two coupled second order differential equations. Within the regime of validity of the hydrodynamic (Thomas-Fermi) approximation , these two equations for the density perturbation $`\varrho `$ and the phase perturbation $`\varphi `$ in terms of the local speed of sound
$$c(𝐱)\frac{\mathrm{}}{m}\sqrt{4\pi a\rho (𝐱)},$$
and the background stationary velocity
$$𝐯\frac{\mathrm{}}{m}\vartheta $$
read
$$\dot{\varrho }=\left(\frac{m}{4\pi a\mathrm{}}c^2\varphi +𝐯\varrho \right)\dot{\varphi }=𝐯\varphi \frac{4\pi a\mathrm{}}{m}\varrho .$$
Furthermore, low frequency perturbations are essentially just waves of (zero) sound. Indeed, the Bogoliubov equations may be reduced to a single second order equation for the condensate phase perturbation $`\varphi `$. This differential equation has the form of a relativistic wave equation $`_\mu (\sqrt{g}g^{\mu \nu }_\nu \varphi )=0`$, with $`g=detg_{\mu \nu }`$, in an effective curved spacetime with the metric $`g_{\mu \nu }`$ being entirely determined by the local speed of sound $`c`$ and the background stationary velocity $`𝐯`$. Up to a conformal factor, this effective metric has the form
$$(g_{\mu \nu })=\left(\begin{array}{cc}(c^2𝐯^2)& 𝐯^\mathrm{T}\\ 𝐯& \mathrm{𝟏}\end{array}\right).$$
This class of metrics can possess event horizons. For instance, if an effective sink for atoms is generated at the center of a spherical trap (such as by an atom laser out-coupling technique ), and if the radial potential profile is suitably arranged, we can produce densities $`\rho (r)`$ and flow velocities $`𝐯(𝐱)=v(r)𝐫/r`$ such that the quantity $`c^2𝐯^2`$ vanishes at a radius $`r=r_h`$, being negative inside and positive outside. The sphere at radius $`r_h`$ is a sonic event horizon completely analogous to those appearing in general relativistic black holes, in the sense that sonic perturbations cannot propagate through this surface in the outward direction . This can be seen explicitly by writing the equation for the radial null geodesics of the metric $`g_{\mu \nu }`$:
$$\dot{r}_\pm =v\pm c,$$
which can be obtained from setting the proper interval $`ds^2=g_{\mu \nu }dx^\mu dx^\nu `$ equal to zero and restricting the allowed motion to the radial direction, so that
$$(c^2v^2)+2v\dot{r}+\dot{r}^2=0.$$
The ingoing null geodesic $`r_{}(t)`$ is not affected by the presence of the horizon and can cross it in a finite coordinate time $`t`$. The outgoing null geodesic $`r_+(t)`$ on the other hand needs an infinite amount of time to leave the horizon since $`\dot{r}_+=0`$ at the horizon. The physical mechanism of the sonic black hole is quite simple: inside the horizon, the background flow speed $`v`$ is larger than the local speed of sound $`c`$, and so sound waves are inexorably dragged inwards.
In fact there are two conditions which must hold for this dragged sound picture to be accurate. Wavelengths larger than the black hole itself will of course not be dragged in, but merely diffracted around it. And perturbations must have wavelengths
$$\lambda \frac{\pi \mathrm{}}{mc},\frac{\pi \mathrm{}}{mc\sqrt{|1v/c|}}.$$
Otherwise they do not behave as sound waves since they lie outside the regime of validity of the hydrodynamic approximation. These short-wavelength modes must be described by the full Bogoliubov equations, which allow signals to propagate faster than the local sound speed, and thus permit escape from sonic black holes. So, to identify a condensate state $`\mathrm{\Psi }_s`$ as a sonic black hole, there must exist modes with wavelengths larger than these lower limits (which in terms of the local healing length $`\xi (𝐱)\mathrm{}/[mc(𝐱)]`$ read $`\lambda 2\pi \xi ,2\pi \xi /\sqrt{|1v/c|}`$), but also smaller than the black hole size. Even if such an intermediate range does exist, the modes outside it may still affect the stability of the black hole as discussed below.
As it stands, this description is incomplete. The condensate flows continually inwards and therefore at $`r=0`$ there must be a sink that takes atoms out of the condensate. Otherwise, the continuity equation $`(\rho 𝐯)=0`$, which must hold for stationary configurations will be violated. From the physical point of view, such a sink can be accomplished by means of an outcoupler laser beam at the origin. (Such outcouplers are the basic mechanisms for making trapped condensates into “atom lasers,” and they have already been demonstrated experimentally by several groups. A tightly focused laser pulse changes the internal state of the atoms at a particular point in the trap, and can also be made to give them a large momentum impulse. This ejects them so rapidly through the always dilute condensate cloud that they do not significantly disturb it; effectively, they simply disappear.)
We have analyzed several specific systems which may be suitable theoretical models for future experiments, and have found that the qualitative behavior is analogous in all of them. Black holes which require atom sinks are both theoretically and experimentally more involved, however; moreover, maintaining a steady transonic flow into a sink may require either a very large condensate or some means of replenishment. We will therefore first discuss an alternative configuration which may be experimentally more accessible and whose description is particularly simple: a condensate in a very thin ring that effectively behaves as a periodic one-dimensional system (Fig. 1). Under conditions that we will discuss, the supersonic region in a ring may be bounded by two horizons: a black hole horizon through which phonons cannot exit, and a ‘white hole’ horizon through which they cannot enter. Then we will analyze another simple one-dimensional model, of a long, straight condensate with an atom sink at the center (Fig. 2).
The existence of instabilities that do not show up in the one-dimensional approximation is an important question in condensate physics, which is under active theoretical and experimental investigation. The essential principles have long been clear, inasmuch as the current dilute condensates really are the weakly interacting Bose gases that have been used as toy models for superfluidity for several decades. The fact that actual critical velocities in liquid helium are generally far below the Landau critical velocity is understood to be due partly to the roton feature of the helium dispersion relation, but this is not present in the dilute condensates. Viscosity also arises due to surface effects, however, and these may indeed afflict dilute condensates as well. The point here is that in addition to the bulk phonon modes considered by Landau, and quite adequately represented in our one-dimensional analysis, there may in principle be surface modes, with a different (and generally lower) dispersion curve. If such modes exist and are unstable, it is very often the case that, as they grow beyond the perturbative regime, they turn into quantized vortices, which can cut through the supercurrent and so lower it.
Whether or not such unstable surface modes actually exist in the Bogoliubov spectrum of a dilute condensate is an issue that has recently been analyzed both numerically and analytically, and it is quite clear that such surface modes exist only if the confining potential is quite rough (which is not only easy to avoid with a magnetic or optical trapping field, but very hard to achieve) , or if the condensate dynamics in the directions perpendicular to the flow is hydrodynamic. That is, the condensate must be at least a few healing lengths thick, so that surface modes decaying on the healing length scale can satisfy all the required boundary conditions . By saying that we are considering an effectively one-dimensional condensate, we mean precisely that this is not the case. For instance, for the tight-ring model, in this regime, the radial trap scale is the shortest length scale in the problem, and the radial trap frequency is the highest frequency; this effectively means that excitations of nontrivial radial modes, including surface modes, are energetically frozen out. (In the limit of radial confinement within the scattering length, our model breaks down for other reasons — but the scattering length can easily be two orders of magnitude smaller than the healing length). The issue of the supercurrent stability in tightly confined ring shaped traps has been addressed in Ref. where the authors arrive at a positive conclusion and also clarify the role of finite temperature and possible trap anisotropy.
## III Sonic black/white holes in a ring
In a sufficiently tight ring-shaped external potential of radius $`R`$, motion in radial ($`r`$) and axial ($`z`$) cylindrical coordinates is effectively frozen. We can then write the wave function as $`\mathrm{\Psi }(z,r,\theta ,\tau )=f(z,r)\mathrm{\Phi }(\theta ,\tau )`$ and normalize $`\mathrm{\Phi }`$ to the number of atoms in the condensate $`_0^{2\pi }𝑑\theta |\mathrm{\Phi }(\theta )|^2=N`$, where with the azimuthal coordinate $`\theta `$ we have introduced the dimensionless time $`\tau =(\mathrm{}/mR^2)t`$. The Gross-Pitaevskii equation thus becomes effectively one-dimensional:
$$i_\tau \mathrm{\Phi }=\left(\frac{1}{2}_\theta ^2+𝒱_{\mathrm{ext}}+\frac{𝒰}{N}|\mathrm{\Phi }|^2\right)\mathrm{\Phi },$$
(1)
where $`𝒰4\pi aNR^2𝑑z𝑑rr|f(z,r)|^4`$ and $`𝒱_{\mathrm{ext}}(\theta )`$ is the dimensionless effective potential (in which we have already included the chemical potential) that results from the dimensional reduction. The stationary solution can then be written as $`\mathrm{\Phi }_s(\theta ,\tau )=\sqrt{\rho (\theta )}e^{i{\scriptscriptstyle 𝑑\theta v(\theta )}}`$ and the local dimensionless angular speed of sound as $`c(\theta )=\sqrt{𝒰\rho (\theta )/N}`$. Periodic boundary conditions around the ring require the “winding number” $`w(1/2\pi )_0^{2\pi }𝑑\theta v(\theta )`$ to be an integer.
The qualitative behavior of horizons in this system is well represented by the two-parameter family of condensate densities
$$\rho (\theta )=\frac{N}{2\pi }(1+b\mathrm{cos}\theta ),$$
where $`b[0,1]`$. Continuity, $`_\theta (\rho v)=0`$, then determines the dimensionless flow-velocity field
$$v(\theta )=\frac{𝒰w\sqrt{1b^2}}{2\pi c(\theta )^2},$$
which depends on $`w`$ as a third discrete independent parameter. Requiring that $`\mathrm{\Phi }_s(\theta ,\tau )`$ be a stationary solution to Gross-Pitaevskii equation then determines how the trapping potential must be modulated as a function of $`\theta `$. All the properties of the condensate, including whether and where it has sonic horizons, and whether or not they are stable, are thus functions of $`𝒰`$, $`b`$ and $`w`$. For instance, if we require that the horizons be located at $`\theta _h=\pm \pi /2`$, which imposes the relation $`𝒰=2\pi w^2(1b^2)`$, then we must have $`c^2v^2`$ positive for $`\theta (\pi /2,\pi /2)`$, zero at $`\theta _h=\pm \pi /2`$, and negative otherwise, provided that $`𝒰<2\pi w^2`$. The further requirement that perturbations on wavelengths shorter than the inner and the outer regions are indeed phononic implies $`𝒰2\pi `$, which in turn requires $`w1`$ and $`1b1/w^2`$. In fact, detailed analysis shows that $`w5`$ is sufficient.
### A Stability
The mere existence of a black hole solution does not necessarily mean that it is physically realizable: it should also be stable over sufficiently long time scales. Since stability must be checked for perturbations on all wavelengths, the full Bogoliubov spectrum must be determined. For large black holes within large, slowly varying condensates, this Bogoliubov problem may be solved using WKB methods that closely resemble those used for solving relativistic field theories in true black hole spacetimes . A detailed adaptation of these methods to the Bogoliubov problem will be presented elsewhere . The results are qualitatively similar to those we have found for black holes in finite traps with low winding number, where we have resorted to numerical methods because, in these cases, WKB techniques may fail for just those modes which threaten to be unstable.
Our numerical approach for our three-parameter family of black/white holes in the ring-shaped condensate has been to write the Bogoliubov equations in discrete Fourier space, and then truncate the resulting infinite-dimensional eigenvalue problem. Indeed, writing the wave funtion as $`\mathrm{\Phi }=\mathrm{\Phi }_s+\phi e^{i{\scriptscriptstyle 𝑑\theta v(\theta )}}`$, decomposing the perturbation $`\phi `$ in discrete modes
$$\phi (\theta ,\tau )=\underset{\omega ,n}{}e^{i\omega \tau }e^{in\theta }A_{\omega ,n}u_{\omega ,n}(\theta )+e^{i\omega ^{}\tau }e^{in\theta }A_{\omega ,n}^{}v_{\omega ,n}^{}(\theta ),$$
and substituting into the Gross-Pitaevskii equation we obtain the following equation for the modes $`u_{\omega ,n}`$ and $`v_{\omega ,n}`$:
$$\omega \left(\begin{array}{c}u_{\omega ,n}\\ v_{\omega ,n}\end{array}\right)=\underset{p}{}\left(\begin{array}{cc}h_{np}^+& f_{np}\\ f_{np}& h_{np}^{}\end{array}\right)\left(\begin{array}{c}u_{\omega ,p}\\ v_{\omega ,p}\end{array}\right).$$
In this equation,
$`f_{np}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}𝑑\theta e^{i(np)\theta }c(\theta )^2,`$ (2)
$`h_{np}^\pm `$ $`=`$ $`\pm {\displaystyle \frac{n^2}{2}}\delta _{np}+{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}𝑑\theta e^{i(np)\theta }`$ (4)
$`\times \left[pv(\theta ){\displaystyle \frac{1}{2}}v^{}(\theta )\pm \left(c(\theta )^2+{\displaystyle \frac{1}{2}}{\displaystyle \frac{c^{\prime \prime }(\theta )}{c(\theta )}}\right)\right],`$
which, after some lengthy calculations, can be written as
$`f_{np}`$ $`=`$ $`{\displaystyle \frac{𝒰}{2\pi }}\left(\delta _{n,p}+{\displaystyle \frac{b}{2}}\delta _{n,p+1}+{\displaystyle \frac{b}{2}}\delta _{n,p1}\right),`$ (5)
$`h_{np}^\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}(n+p)w\sqrt{1b^2}\alpha _{np}`$ (6)
$`\pm `$ $`\left(f_{np}+{\displaystyle \frac{4n^21}{8}}\delta _{n,p}+{\displaystyle \frac{1b^2}{8}}\beta _{np}\right),`$ (7)
where
$`\alpha _i`$ $`=`$ $`{\displaystyle \underset{j|i|,i+j\mathrm{even}}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{b}{2}}\right)^j\left(\begin{array}{c}j\\ (i+j)/2\end{array}\right),`$ (10)
$`\beta _i`$ $`=`$ $`{\displaystyle \underset{j|i|,i+j\mathrm{even}}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{b}{2}}\right)^j\left(\begin{array}{c}j\\ (i+j)/2\end{array}\right)(j+1).`$ (13)
Eliminating Fourier components above a sufficiently high cutoff $`Q`$ has negligible effect on possible instabilities, which can be shown to occur at relatively long wavelengths. We face then an eigenvalue problem for the $`2(Q+1)\times 2(Q+1)`$ matrix built out of blocks of the form
$$\left(\begin{array}{cc}h_{np}^+& f_{np}\\ f_{np}& h_{np}^{}\end{array}\right).$$
The numerical solution to this eigenvalue equation, together with the normalization condition $`𝑑\theta (u_{\omega ^{},n}^{}u_{\omega ^{},n^{}}v_{\omega ^{},n}^{}v_{\omega ^{},n^{}})=\delta _{nn^{}}\delta _{\omega \omega ^{}}`$, provides the allowed frequencies. Real negative eigenfrequencies for modes of positive norm are always present, which means that black hole configurations are energetically unstable, as expected. This feature is inherent in supersonic flow, since the speed of sound is also the Landau critical velocity. In a sufficiently cold and dilute condensate, however, the time scale for dissipation may in principle be made very long, and so these energetic instabilities need not be problematic .
More serious are dynamical instabilities, which occur for modes with complex eigenfrequencies. Since the Bogoliubov theory is based on a quantized Hamiltonian that is Hermitian, there are certainly no complex energy eigenvalues; but the natural frequencies of normal modes can indeed be complex \[in which case the usual rule, that energy eigenvalues are $`\mathrm{}(n+1/2)`$ times the mode frequencies, simply breaks down\]. A detailed discussion of the quantum mechanics of dynamical instability is presented in the Appendix; for the purposes of our main discussion it suffices to note that complex (mode) eigenfrequencies are indeed genuine physical phenomena, and by no means a numerical artifact. For sufficiently high values of the cutoff (e.g., $`Q25`$ in our calculations), the complex eigenfrequencies obtained from the truncated eigenvalue problem become independent of the cutoff within the numerical error. The existence and rapidity of dynamical instabilities depend sensitively on $`(𝒰,b,w)`$. For instance, see Fig. 3 for a contour plot of the maximum of the absolute values of the imaginary part of all eigenfrequencies for $`w=7`$, showing that the regions of instability are long, thin fingers in the $`(𝒰,b)`$ plane. Figure 4 shows the size of the largest absolute value of the instabilities for each point on the dashed curve of Fig. 3. It illustrates the important fact that the size of the imaginary parts, which gives the rate of the instabilities, increases starting from zero, quite rapidly with $`b`$, although they remain small as compared with the real parts.
### B Creation of a black/white hole
The stability diagram of Fig. 3 suggests a strategy for creating a sonic black hole from an initial stable state. Within the upper subsonic region, the vertical axis $`b=0`$
corresponds to a homogeneous persistent current in a ring, which can in principle be created using different techniques . Gradually changing $`𝒰`$ and $`b`$, it is possible to move from such an initial state to a black/white hole state, along a path lying almost entirely within the stable region, and only passing briefly through instabilities where they are sufficiently small to cause no difficulty.
Indeed, we have simulated this process of adiabatic creation of a sonic black/white hole by solving numerically (using the split operator method) the time-dependent Gross-Pitaevskii equation (1) that provides the evolution of the condensate when the parameters of the trapping potential change so as to move the condensate state along various paths in parameter space. One of these paths is shown in Fig. 3 (light-grey solid line): we start with a current at $`w=7`$, $`b=0`$, and sufficiently high $`𝒰`$ \[Fig. 5(a)\]; we then increase $`b`$ adiabatically keeping $`𝒰`$ fixed until an appropriate value is reached \[Fig. 5(b)\]; finally, keeping $`b`$ constant, we decrease $`𝒰`$ adiabatically (which can be physically implemented by decreasing the radius of the ring trap), until we meet the dashed contour for black holes of comfortable size \[Fig. 5(c)\]. Our simulations confirm that the small instabilities which briefly appear in the process of creation do not disrupt the adiabatic evolution. The final quantum state of the condensate, obtained by this procedure, indeed represents a stable black/white hole. We have further checked the stability of this final configuration by numerically solving the Gross-Pitaevskii equation (1) for very long periods of time (as compared with any characteristic time scale of the condensate) and for fixed values of the trap parameters. This evolution reflects the fact that no imaginary frequencies are present, as predicted from the mode analysis, and that the final state is indeed stationary \[Fig. 5(d)\]. Once the black/white hole has been created, one could further change the parameters $`(𝒰,b)`$ so as to move between the unstable ‘fingers’ into a stable region of higher $`b`$ (a deeper hole).
### C Quasiparticle pair creation
Instead of navigating the stable region of parameter space, one could deliberately enter an unstable region \[Fig. 5(e)–5(i)\]. In the this case, the black hole should disappear in an explosion of phonons, which may be easy to detect experimentally. Such an event might be related to the evaporation process suggested for real black holes, in the sense that pairs of quasiparticles are created near the horizon in both positive and negative energy modes. We will explain this point briefly; a more detailed exposition is included in the Appendix.
In the language of second quantization, the perturbation field operator $`\phi `$ satisfies the linear equation
$$i\dot{\phi }=\frac{1}{2}\phi ^{\prime \prime }iv\phi ^{}+\left(\frac{1}{2}\frac{c^{\prime \prime }}{c}\frac{i}{2}v^{}+c^2\right)\phi +c^2\phi ^{},$$
which, taking into account that $`[\phi (\theta ),\phi ^{}(\theta ^{})]=\delta (\theta \theta ^{})`$, can be written as
$$i\dot{\phi }=[\phi ,H],$$
where the Bogoliubov Hamiltonian is
$`H={\displaystyle 𝑑\theta }`$ $`[{\displaystyle \frac{1}{2}}\phi ^{}\phi ^{\prime \prime }iv\phi ^{}\phi ^{}+({\displaystyle \frac{1}{2}}{\displaystyle \frac{c^{\prime \prime }}{c}}{\displaystyle \frac{i}{2}}v^{}+c^2)\phi ^{}\phi `$ (15)
$`+{\displaystyle \frac{1}{2}}(\phi ^{}\phi ^{}+\phi \phi )].`$
The Hermiticity of the Bogoliubov linearized Hamiltonian implies that eigenmodes with complex frequencies appear always in dual pairs, whose frequencies are complex conjugate. In the language of second quantization, the linearized Hamiltonian for each such pair has the form
$$H=\underset{n}{}(\omega A_{\omega ^{},n}^{}A_{\omega ,n}+\omega ^{}A_{\omega ,n}^{}A_{\omega ^{},n}),$$
and the only nonvanishing commutators among these operators are $`[A_{\omega ,n},A_{\omega ^{},n^{}}^{}]=\delta _{nn^{}}`$. The asterisk on the subscript is important: the mode with frequency $`\omega ^{}`$ is a different mode from the one with frequency $`\omega `$, and $`A_{\omega ^{},n}^{}`$ is not the Hermitian conjugate of $`A_{\omega ,n}^{}`$. It is therefore clear that none of these operators is actually a harmonic oscillator creation or annihilation operator in the usual sense. However, the linear combinations
$$a_n=\frac{1}{\sqrt{2}}(A_{\omega ,n}+A_{\omega ^{},n}),b_n=\frac{i}{\sqrt{2}}(A_{\omega ,n}^{}+A_{\omega ^{},n}^{})$$
and their Hermitian conjugates are true annhilation and creation operators, with the standard commutation relations, and in terms of these the Bogoliubov Hamiltonian becomes
$$H=\underset{n}{}\left[\mathrm{Re}(\omega )(a_n^{}a_nb_n^{}b_n)\mathrm{Im}(\omega )(a_n^{}b_n^{}+a_nb_n)\right].$$
This interaction obviously leads to self-amplifying creation of positive and negative frequency pairs. Evaporation through an exponentially self-amplifying instability is not equivalent, however, to the usual kind of Hawking radiation ; this issue will be discussed in detail elsewhere.
## IV Sink-generated black holes
Condensates which develop black hole behaviors by means of flows generated by laser-driven sinks also present regions of stability and instability in parameter space and, in this sense, their behavior is analogous to that in a ring-shaped trapping potential. Here we present a simple model that exhibits the main qualitative features of more general situations and that can be studied analytically. Although in this model, we study a condensate of infinite size, in more realistic models or experiments, it will suffice to take condensates which are sufficiently large, since the stability pattern is not significantly affected by the (large but finite) size of the condensate.
### A The model
Let us consider a tight cigar-shaped condensate of infinite size such that the motion in the $`(y,z)`$ plane is effectively frozen. In appropriate dimensionless units, the effectively one-dimensional Gross-Pitaevskii equation thus becomes:
$$i_\tau \mathrm{\Phi }=\left(\frac{1}{2}_x^2+𝒱_{\mathrm{ext}}+𝒰|\mathrm{\Phi }|^2\right)\mathrm{\Phi },$$
with the normalization condition
$$\underset{D\mathrm{}}{lim}\frac{1}{2D}_D^D𝑑x|\mathrm{\Phi }(x)|^2=n.$$
In this equation, $`𝒱_{\mathrm{ext}}`$ is the dimensionless effective potential that results form the dimensional reduction, which already includes the chemical potential.
In order to obtain a black hole configuration, let us choose the potential $`𝒱_{\mathrm{ext}}`$ so that it produces a profile for the speed of sound $`c(x)=\sqrt{𝒰\rho (x)}`$ of the form
$$c(x)=\{\begin{array}{cc}c_0,\hfill & |x|<L\hfill \\ c_0[1+(\sigma 1)x/ϵ],\hfill & L<|x|<L+ϵ\hfill \\ \sigma c_0,\hfill & L+ϵ<|x|\hfill \end{array},$$
with $`\sigma >1`$, and a flow velocity in the inward direction. The continuity equation then provides the flow velocity profile
$$v(x)=\frac{v_0c_0^2}{c(x)^2}\frac{x}{|x|},$$
where $`v_0`$ is the absolute value of the flow velocity in the inner region.
As it stand this model fails to fulfill the continuity equation at $`x=0`$. In order to take this into account, we will also introduce a sink of atoms at $`x=0`$ that takes atoms out of the condensate (this can be physically implemented by means of a laser). From the mathematical point of view, it can be modeled by an additional term in the equation of the form $`iE\delta (x)`$ which indeed induces loss of atoms at $`x=0`$. Equivalently, it can be represented by boundary conditions of the form
$`\mathrm{\Phi }(0^+,\tau )\mathrm{\Phi }(0^{},\tau )`$ $`=`$ $`0,`$ (16)
$`\mathrm{\Phi }^{}(0^+,\tau )\mathrm{\Phi }^{}(0^{},\tau )`$ $`=`$ $`2iE\mathrm{\Phi }(0,\tau ),`$ (17)
which determine the flow velocity inside in terms of the characteristics of the outcoupler laser, namely, $`v_0=E`$.
Perturbations $`\varphi `$ around this stationary state $`\mathrm{\Phi }_s=\sqrt{\rho }e^{i{\scriptscriptstyle v(x)𝑑x}}`$, such that $`\mathrm{\Phi }=\mathrm{\Phi }_s+\varphi `$ (note that for convenience we have chosen a different convention as compared with the ring in which $`\mathrm{\Phi }=\mathrm{\Phi }_s+\phi e^{i{\scriptscriptstyle v}}`$) must satisfy the boundary conditions (17) and the equation
$$i\dot{\varphi }=\frac{1}{2}\varphi ^{\prime \prime }+(c^2v^2/2+c^{\prime \prime }/2c)\varphi +c^2e^{2i^xv}\varphi ^{}$$
where
$$\frac{c^{\prime \prime }}{c}=\frac{\sigma 1}{ϵ}[\delta (|x|L)\frac{1}{\sigma }\delta (|x|Lϵ)].$$
As a further simplifying assumption, we will assume that $`v_0ϵ1`$ so that
$$\left|_L^x𝑑x^{}v(x^{})\right|_L^{L+ϵ}𝑑x\frac{v_0}{[1+(\sigma 1)x/ϵ]^2}v_0ϵ1.$$
Let us now expand the perturbation $`\varphi `$ in modes
$$\varphi =\underset{\omega ,k}{}\left[A_{\omega ,k}u_{\omega ,k}(x)e^{i\omega \tau }+A_{\omega ,k}^{}v_{\omega ,k}(x)^{}e^{i\omega ^{}\tau }\right].$$
Then, the modes $`u_{\omega ,k}(x)`$ and $`v_{\omega ,k}(x)`$ satisfy, in each region, the Bogoliubov equations
$`\omega u_{\omega ,k}`$ $`=`$ $`{\displaystyle \frac{1}{2}}u_{\omega ,k}^{\prime \prime }+(c^2v^2/2)u_{\omega ,k}+c^2e^{2i^xv}v_{\omega ,k},`$ (18)
$`\omega v_{\omega ,k}`$ $`=`$ $`{\displaystyle \frac{1}{2}}v_{\omega ,k}^{\prime \prime }(c^2v^2/2)v_{\omega ,k}c^2e^{2i^xv}u_{\omega ,k}.`$ (19)
### B Matching conditions
The intermediate regions $`L<|x|<L+ϵ`$, provide the connection between the perturbation modes in the inner and outer regions. Once these connection formulas have been established, in the limit of small $`ϵ`$, we will only need to study the inside and outside modes and their relation through such formulas. The case of an abrupt horizon, in which the background condensate velocity is steeply and linearly ramped within a very short interval, is obviously quite special; and it does not particularly resemble the horizon of a large black hole in Einsteinian gravity. But the connection formula that we derive for this case will qualitatively resemble those that are obtained, with considerably more technical effort, for smoother horizons. And the results we will obtain for the global Bogoliubov spectrum of the condensate black hole will indeed be representative of more generic cases.
In the intermediate regions $`L<|x|<L+ϵ`$, the factors $`e^{\pm 2i^xv}`$ in the last terms of Eqs. (19) become $`1+𝒪(ϵ)`$. Then the solution of these equations is
$`u_{\omega ,k}`$ $`=`$ $`\alpha _{\omega ,k}+\beta _{\omega ,k}x/ϵ+𝒪(ϵ^2),`$ (20)
$`v_{\omega ,k}`$ $`=`$ $`\gamma _{\omega ,k}+\kappa _{\omega ,k}x/ϵ+𝒪(ϵ^2)`$ (21)
as can be easily seen by defining the variable $`q=x/ϵ`$ so that the equations become
$$_q^2u_{\omega ,k}=_q^2v_{\omega ,k}=𝒪(ϵ^2).$$
The singular character of $`c^{\prime \prime }/c`$ at $`|x|=L,L+ϵ`$ can be substituted by matching conditions at $`|x|=L,L+ϵ`$ which will in turn provide the connection formulas between the modes outside ($`|x|>L`$) and the modes inside ($`|x|<L`$). Furthermore, the symmetry of the problem allows us to study the region $`x>0`$.
These matching conditions are
$`\varphi (L^+)\varphi (L^{})`$ $`=`$ $`0,`$ (22)
$`\varphi ^{}(L^+)\varphi ^{}(L^{})`$ $`=`$ $`{\displaystyle \frac{\sigma 1}{ϵ}}\varphi (L),`$ (23)
$`\varphi (L+ϵ^+)\varphi (L+ϵ^{})`$ $`=`$ $`0,`$ (24)
$`\varphi ^{}(L+ϵ^+)\varphi ^{}(L+ϵ^{})`$ $`=`$ $`{\displaystyle \frac{\sigma 1}{\sigma ϵ}}\varphi (L+ϵ).`$ (25)
These equations together with the form of the modes in the region $`L<x<L+ϵ`$, provide the connection formulas between the inside and outside modes (from now on we will drop the subindex $`\omega `$):
$`u_{\mathrm{in},k}(L)`$ $`=`$ $`ϵu_{\mathrm{out},k}^{}(L)+{\displaystyle \frac{1}{\sigma }}u_{\mathrm{out},k}(L),`$ (26)
$`u_{\mathrm{in},k}^{}(L)`$ $`=`$ $`\sigma u_{\mathrm{out},k}^{}(L),`$ (27)
and likewise for the modes $`v_{\mathrm{in},\mathrm{out}}`$.
### C Dispersion relation and boundary conditions for large $`x`$
In each of the regions (inside and outside), we can write
$$u_k(x)=u_ke^{i(k|v|)(xL)},v_k(x)=v_ke^{i(k+|v|)(xL)}.$$
Upon substitution of this expansion into the Bogoliubov equations (19), we obtain, for each region, the following set of algebraic equations:
$$h_k^{}u_k+c^2v_k=0,c^2u_k+h_k^+v_k=0,$$
where $`h_k^\pm =k^2/2+c^2\pm (k|v|+\omega )`$. For these equations to have a solution, the determinant must vanish thus providing the dispersion relation
$$k^4/4+(c^2v^2)k^22\omega |v|k\omega ^2=0$$
(28)
which, for fixed $`\omega `$, is a fourth-order equation for $`k`$. For each of the four solutions $`u_k`$ and $`v_k`$ must be related by
$$v_k=h_ku_k\text{with}h_k=\frac{1}{c^2}(k^2/2+c^2k|v|\omega ).$$
The constant coefficients $`u_k`$ can be regarded as normalization constants and will be set to unity. Let us study the possible solutions to the dispersion relation depending on whether $`\omega `$ is real or complex.
#### Complex frequencies.—
In this case all four solutions are pure complex, two of them with positive imaginary part and two of them with negative one.
In order to prove this statement, let us first assume that there exists a real solution $`k`$ for a complex $`\omega =\varpi +i\gamma `$. Then the imaginary part of Eq. (28) implies that $`\varpi =k|v|`$. Introducing this result into its real part leads to $`\gamma ^2=(k^4/4+c^2k^2)`$ which is impossible to fulfill because $`k`$ is real. So the four solutions are complex.
Because of continuity, all $`\omega `$ in the upper half complex $`\omega `$ plane have the same number of solutions with positive imaginary part. Otherwise, for some $`\omega `$ there should exist a real solution that interpolates between positive and negative imaginary part solutions, but this is not possible, as we have seen.
Now let us concentrate on small frequencies, i.e., on frequencies around $`\omega =0`$. For $`\omega =0`$, we have a double root at $`k=0`$. The other two solutions are $`k=\pm 2\sqrt{v^2c^2}`$, which are real for $`c^2<v^2`$ (i.e., inside) and pure imaginary for $`c^2>v^2`$ (i.e., outside). Let us follow these four solutions when $`\omega =i\epsilon `$. The solutions coming from the double root $`k=0`$ will now be of the form $`k=k_r+i\epsilon k_i`$. It is easy to see that $`k_r=0`$ and $`k_i=1/(|v|\pm c)`$. If $`c>|v|`$, one is positive and one is negative. If $`c<|v|`$, both of them are negative. On the other hand, the solutions $`k=\pm 2\sqrt{v^2c^2}`$, for $`c^2>v^2`$, are already complex conjugate. For $`c^2<v^2`$, we write $`k=k_r+i\epsilon k_i`$ and introduce it into Eq. (28). We then see that at first order in $`\epsilon `$, $`k_i=|v|/(v^2c^2)>0`$. Thus, we have seen that for $`\omega =i\epsilon `$, we have two solutions with positive imaginary part and two with negative one in any case (inside and outside). But if this is so for $`\omega =i\epsilon `$ it must be true in the whole upper $`\omega `$ plane and consequently in the whole complex $`\omega `$ plane. $`\mathrm{}`$
In the inside region all possible solutions are in principle allowed but outside we are only left with the two that have $`\mathrm{Im}(k)>0`$, because the other two grow exponentially.
#### Real frequencies.—
Outside ($`c^2>v^2`$), there are two real and two complex conjugate solutions. Of these two complex solutions, only one is allowed \[the one with $`\mathrm{Im}(k)>0`$\] because the other grows exponentially. Inside ($`c^2<v^2`$), for $`\omega >\omega _{\mathrm{max}}`$ there are two real and two complex conjugate solutions; for $`\omega <\omega _{\mathrm{max}}`$ there are four real solutions; the value $`\omega =\omega _{\mathrm{max}}`$ is a bifurcating point.
### D Connection formulas for complex frequencies
Since we are interested in the existence of dynamical instabilities, we will concentrate in the case in which $`\omega `$ is complex. Then, as we have seen, the dispersion equation (28) has four complex solutions for $`k`$ in each region. Inside, all four solutions $`k_{\mathrm{in},i},i=1\mathrm{}4`$ are in principle possible but outside those with $`\mathrm{Im}(k_{\mathrm{out}})<0`$ will increase exponentially. Therefore, up to corrections coming from the finite size of the condensate, which we ignore here, only modes associated with $`k_{\mathrm{out},\alpha },\alpha =1,2`$ such that $`\mathrm{Im}(k_{\mathrm{out},\alpha })>0`$ are allowed. Each mode $`u_{\mathrm{out},\alpha }(x)=e^{i(k_{\mathrm{out},\alpha }v_0/\sigma ^2)(xL)}`$ will match a linear combination $`u_{\mathrm{in},\alpha }(x)=_iF_{\alpha i}u_{\mathrm{in},i}(x)`$ of modes $`u_{\mathrm{in},i}(x)`$ inside, i.e.,
$$u_{\mathrm{in},\alpha }(x)=\underset{i}{}F_{\alpha i}e^{i(k_{\mathrm{in},i}v_0)(xL)}$$
and similarly for $`v_{\mathrm{out},\alpha }`$ and $`v_{\mathrm{in},\alpha }`$:
$$v_{\mathrm{in},\alpha }(x)=\underset{i}{}F_{\alpha i}h_{\mathrm{in},i}e^{i(k_{\mathrm{in},i}+v_0)(xL)}.$$
After some straightforward calculations, it can be seen that these connecting coefficients $`F_{\alpha i}`$ are given by $`F_{\alpha i}=_j(M^1)_{ij}C_{\alpha j}`$, where
$$M=\left(\begin{array}{cccc}1& 1& 1& 1\\ k_{\mathrm{in},1}^{}& k_{\mathrm{in},2}^{}& k_{\mathrm{in},3}^{}& k_{\mathrm{in},4}^{}\\ h_{\mathrm{in},1}& h_{\mathrm{in},2}& h_{\mathrm{in},3}& h_{\mathrm{in},4}\\ h_{\mathrm{in},1}k_{\mathrm{in},1}^+& h_{\mathrm{in},2}k_{\mathrm{in},2}^+& h_{\mathrm{in},3}k_{\mathrm{in},3}^+& h_{\mathrm{in},4}k_{\mathrm{in},4}^+\end{array}\right),$$
$$C_\alpha =\left(\begin{array}{c}1/\sigma iϵk_{\mathrm{out},\alpha }^{}\\ \sigma k_{\mathrm{out},\alpha }^{}\\ (1/\sigma iϵk_{\mathrm{out},\alpha }^+)h_{\mathrm{out},\alpha }\\ \sigma k_{\mathrm{out},\alpha }^+h_{\mathrm{out},\alpha }\end{array}\right).$$
In these equations,
$$k_{\mathrm{in},i}^\pm =k_{\mathrm{in},i}\pm v_0\text{and}k_{\mathrm{out},\alpha }^\pm =k_{\mathrm{out},\alpha }\pm v_0/\sigma ^2.$$
### E Boundary conditions at the sink. Complex eigenvalues
We have already found the modes in the inner and outer region as well as their relation. To determine which (complex-frequency) modes will be present, there only remains to impose the boundary conditions dictated by the presence of the sink at $`x=0`$.
As we have already mentioned, the symmetry of the system under reflection ($`xx`$) allows us to study only the region $`x>0`$, provided that we study the even and odd perturbations separately. For odd fluctuations \[$`\varphi _\mathrm{o}(x,\tau )=\varphi _\mathrm{o}(x,\tau )`$\], the boundary conditions (17) become
$$\varphi _\mathrm{o}(0,\tau )=0$$
at all times $`\tau `$. This implies that the $`u`$ and $`v`$ components of $`\varphi `$ must separately satisfy the boundary condition. Since we can have any linear combination of the two solutions that decay outside the horizon, we therefore have a two-by-two matrix constraint. The condition that a non-zero solution exists is that the determinant
$$\mathrm{det}\left(\begin{array}{cc}u_{\mathrm{in},1}(0)& u_{\mathrm{in},2}(0)\\ v_{\mathrm{in},1}(0)& v_{\mathrm{in},2}(0)\end{array}\right)=0$$
and therefore
$$\underset{ij}{}F_{1i}F_{2j}(h_{\mathrm{in},i}h_{\mathrm{in},j})e^{i(k_{\mathrm{in},i}+k_{\mathrm{in},j})L}=0.$$
(29)
For even fluctuations \[$`\varphi _\mathrm{e}(x,\tau )=\varphi _\mathrm{e}(x,\tau )`$\], the boundary conditions (17) become
$$\varphi _\mathrm{e}^{}(0,\tau )+iv_0\varphi _\mathrm{e}(0,\tau )=0,$$
which implies that
$$\underset{ij}{}F_{1i}F_{2j}(h_{\mathrm{in},i}h_{\mathrm{in},j})k_{\mathrm{in},i}k_{\mathrm{in},j}e^{i(k_{\mathrm{in},i}+k_{\mathrm{in},j})L}=0.$$
(30)
For fixed $`L`$, $`𝒰`$, $`v_0`$, and $`\sigma `$, the quantities $`F`$, $`h_{\mathrm{in}}`$, and $`k_{\mathrm{in}}`$ that appear in Eqs. (29) and (30) are only functions of $`\omega `$. Therefore the solutions to these equations are all the possible complex eigenfrequencies, which depend on the free parameters that determine the model, namely, the size $`2L`$ of the inner region, the speed of sound inside $`c_0`$, the relative change of the speed of sound between the inner and the outer regions $`\sigma `$, and the flow velocity inside $`v_0`$ (related to the characteristics of the outcoupler laser). In practice, there are also other parameters of the condensate such as its size $`2D`$ (which has been made arbitrarily large) and the size of the intermediate regions $`ϵ`$ (which has been made arbitrarily small).
Equations (29) and (30) can be solved numerically for different values of the parameters $`\sigma `$, $`𝒰`$, $`v_0`$, and $`L`$. The numerical method employed is the following. The equations above have the form
$$f(\omega ;\sigma ,𝒰,v_0,L)=0,$$
where $`f`$ and $`\omega `$ are both in general complex. We plot contours of constant absolute value of $`f`$ in the complex $`\omega `$ plane; where $`|f|`$ approaches zero, we have an eigenfrequency.
The distribution of complex solutions in the complex $`\omega `$ plane depends on the size of the inner region $`L`$, for given $`\sigma `$, $`𝒰`$, and $`v_0`$. Direct inspection of the numerical results shows that the number of instabilities increases by one when the black-hole size $`L`$ is increased by $`\pi /k_0`$ where $`k_0=\sqrt{v_0^2c_0^2}`$. More explicitly, for $`L`$ smaller than $`\pi /k_0\delta `$ ($`\delta `$ being much smaller than $`\pi /k_0`$) there are no complex eigenfrequencies; for $`(L+\delta )k_0/\pi [n,n+1]`$ with $`n=1,2\mathrm{}`$, we have $`n`$ complex solutions except for $`L=(n+1/2)\pi /k_0`$, where we find $`n1`$ complex solutions instead of $`n`$ \[i.e, there is one mode for which $`\mathrm{Im}(\omega )=0`$ within numerical resolution\]. This can be easily interpreted qualitatively since the unstable modes are basically the bound states in the black hole, and the highest wave number $`k`$ on the positive norm upper branch, for the barely bound state with $`\omega 0^{}`$, is exactly $`k_0`$. So the threshold is simply when the well becomes big enough to have a bound state; the small $`\delta `$ displacement comes in because the horizon is not exactly a hard wall; and similarly for the extra bound state every $`\pi /k_0`$. Thus stability can only be achieved for small sizes of the inner region, $`L\pi /k_0`$. As we discussed in Sec. II, the wave length $`2\pi /k`$ of the perturbations must be smaller than this size, which implies $`k>2\pi /L2k_0`$. However, for these perturbations the hydrodynamic approximation, which requires $`k2k_0`$, is not valid. Therefore there are no stable black hole configurations in a strict sense. The sizes of the imaginary parts of the complex solutions decrease as the size $`L`$ of the interior of the black hole increases. Thus although a larger hole has more unstable modes, it is actually less unstable (and might even became quasistable in the sense that its instability time scale would be longer than the experimental duration).
## V Conclusions
We have seen that dilute Bose-Einstein condensates admit, under appropriate conditions, configurations that closely ressemble gravitational black holes. We have analyzed in detail the case of a condensate in a ring trap, and proposed a realistic scheme for adiabatically creating stable sonic black/white holes and we have seen that there exist stable and unstable black hole configurations. We have also studied a model for a sink-generated sonic black hole in an infinite one-dimensional condensate. The dynamical instabilities can be interpreted as coming from quasiparticle pair creation, as in the well-known suggested mechanism for black hole evaporation. Generalizations to spherical or quasi-two-dimensional traps, with flows generated by laser-driven atom sinks, should also be possible, and should behave similarly. While our analysis has been limited to Bogoliubov theory, the further theoretical problems of back reaction and other corrections to simple mean field theory should be more tractable for condensates than for other systems analogous to black holes. And we expect that experiments along the lines we have proposed, including both creation and evaporation of sonic black holes, can be performed with state-of-the-art or planned technology.
###### Acknowledgements.
We thank the Austrian Science Foundation and the European Union TMR networks ERBFMRX–CT96–0002 and ERB–FMRX–CT96–0087. J.R.A. is grateful to Ted Jabobson for useful discussions.
## Complex frequencies. Redundancy and normalization
In this Appendix, we will analyze the issue of the redundancy and normalization of the Bogoliubov modes in the presence of complex frequencies from a general point of view. Dynamical instabilities in quantum field theory, and the quantization of dynamically unstable modes, do not seem to be widely understood: it is for instance common to read axiomatic statements that one must only quantize positive norm modes, even though this implicitly neglects dynamical instabilities, and does not follow in general from the fundamental commutation relations. But some explicit treatments of quantum instabilities have been available in the literature for some time; here we review this subject in the specific Bogoliubov context.
We will begin by writing the Bogoliubov equations in their most usual form
$$\mathrm{}\omega _j\left(\begin{array}{c}u_j\\ v_j\end{array}\right)=\left(\begin{array}{cc}h_0(𝐱)& c(𝐱)^2e^{2i\vartheta (𝐱)}\\ c(𝐱)^2e^{2i\vartheta (𝐱)}& h_0(𝐱)\end{array}\right)\left(\begin{array}{c}u_j\\ v_j\end{array}\right),$$
(31)
where $`h_0(𝐱)=\mathrm{}^2/2m^2+V_{\mathrm{ext}}(𝐱)+2mc(𝐱)^2\mu `$. In terms of these modes, the atomic second quantized field operator has the well-known form $`\widehat{\mathrm{\Psi }}(𝐱,t)=\mathrm{\Psi }_s(𝐱,t)+\widehat{\psi }(𝐱,t)`$ with
$$\widehat{\psi }(𝐱,t)=\underset{j}{}\left[\widehat{a}_ju_j(𝐱)e^{i\omega _jt}+\widehat{a}_j^{}v_j(𝐱)^{}e^{i\omega _j^{}t}\right].$$
(32)
If there is a solution $`(u_j,v_j)`$ to Eq. (31) with mode frequency $`\omega _j`$, then straightforward substitution shows that $`(u_j^{},v_j^{})=(v_j^{},u_j^{})`$ must be a solution with frequency $`\omega _j^{}=\omega _j^{}`$. If we examine the contributions of these two solutions, however, we find that together they yield but a single term, of the form $`(\widehat{a}_j+\widehat{a}_j^{}^{})e^{i\omega t}u_j(𝐱)+(\widehat{a}_j^{}+\widehat{a}_j^{})e^{i\omega ^{}t}v_j^{}(𝐱)`$. It is thus a quite trivial fact that the two modes $`j`$ and $`j^{}`$ are redundant. We are free to simplify our notation by redefining $`\widehat{a}_j+\widehat{a}_j^{}^{}\widehat{a}_j`$, and eliminating mode $`j^{}`$ (leaving it out of the sum over frequencies). Alternatively we could of course eliminate $`j`$ and keep $`j^{}`$. Which of these two notational conventions we should take is best determined by the commutator $`[\widehat{a}_j+\widehat{a}_j^{}^{},\widehat{a}_j^{}+\widehat{a}_j^{}]`$, which will tell us whether the coefficient of $`u_j=v_j^{}^{}`$ is properly an annihilation operator or a creation operator.
Since the only commutation relations that we are given are those of $`\widehat{\psi }`$ and $`\widehat{\psi }^{}`$, we must derive the orthogonality relation for solutions of Eq. (31), and use it to invert Eq. (32). We can use Eq. (31) to show that
$`(\omega _j+\omega _k)M_{jk}`$ $``$ $`{\displaystyle d^3x(u_jv_kv_ju_k)}=0,`$ (33)
$`(\omega _j\omega _k^{})N_{jk}`$ $``$ $`(\omega _j\omega _k^{}){\displaystyle d^3x(u_ju_k^{}v_jv_k^{})}=0,`$ (34)
where in the case of infinite volume the right hand sides are zero in the distributional sense, being infinitely rapidly oscillating boundary terms. This obviously implies that $`M_{jk}`$ vanishes unless $`\omega _k=\omega _j`$, and $`N_{jk}`$ vanishes unless $`\omega _k=\omega _j^{}`$. One can then show that it is always possible to take linear combinations among degenerate modes, and to eliminate redundant modes as just discussed, in such a way as to make $`M_{jk}`$ always vanish, and $`N_{jk}=\delta _{k\overline{ȷ}}`$, where for every $`j`$ there is a single dual mode $`\overline{ȷ}`$, with $`\omega _{\overline{ȷ}}=\omega _j^{}`$. In the case of real $`\omega _j`$, but only then, we have $`\overline{ȷ}=j`$. In general, though, duality is reciprocal (the dual mode of $`\overline{ȷ}`$ is always $`j`$).
The result is that we can now insert Eq. (32) into the second quantized Hamiltonian, with the $`T`$-matrix approximation for the interparticle interaction, to obtain the linearized Bogoliubov Hamiltonian for the perturbations
$$\widehat{H}=\mathrm{}\underset{j}{}\omega _j\widehat{a}_{\overline{ȷ}}^{}\widehat{a}_j.$$
(35)
Since the sum over all modes $`j`$ also includes the dual to every mode with complex $`\omega _j`$, $`\widehat{H}`$ is manifestly Hermitian even though $`\omega _j`$ need not be real. We can also invert Eq. (32) to learn that
$$\widehat{a}_j=d^3x[u_{\overline{ȷ}}^{}\widehat{\delta \psi }+v_{\overline{ȷ}}\widehat{\delta \psi }^{}],$$
(36)
which with Eq. (33) implies the commutation relations
$$[\widehat{a}_j,\widehat{a}_k^{}]=\delta _{k\overline{ȷ}},[\widehat{a}_j,\widehat{a}_k]=0.$$
(37)
For all $`j`$ with real $`\omega _j`$, Eq. (37) are merely the standard canonical commutation relations; and our normalization conventions $`M_{jk}=0`$ and $`N_{jk}=\delta _{jk}`$ are likewise the ones most often presented. In the case of complex $`\omega _j`$ where $`\overline{ȷ}j`$, however, Eq. (37) implies that the canonical conjugate of $`\widehat{a}_j`$ is $`\widehat{a}_{\overline{ȷ}}^{}`$, and this is no longer the same as the Hermitian conjugate $`\widehat{a}_j^{}`$. In fact for complex $`\omega _j`$ we have $`[\widehat{a}_j^{},\widehat{a}_j]=0`$; this already follows from the second line of Eq. (33), which implies that the norm $`N_{jj}`$ of any mode with complex $`\omega _j`$ is zero. But if $`\widehat{a}_j`$ and $`\widehat{a}_j^{}`$ commute, then it is clear that neither $`\widehat{a}_j`$ nor $`\widehat{a}_{\overline{ȷ}}`$ is really a harmonic oscillator annihilation operator in the usual sense, nor are $`\widehat{a}_j^{}`$ or $`\widehat{a}_{\overline{ȷ}}^{}`$ proper creation operators. The commutation relations (37) are validly derived from the fundamental commutation relations for $`\widehat{\psi }`$ and $`\widehat{\psi }^{}`$; but they do not imply, for instance, that either $`\widehat{a}_j^{}\widehat{a}_j`$ or $`\widehat{a}_{\overline{ȷ}}^{}\widehat{a}_j`$ has the discrete, equally spaced spectrum, bounded from below, that one expects of a quasiparticle number operator.
To understand the dual pairs of modes with complex frequencies, we can define the ordinary annihilation operators
$$\widehat{b}_j=\frac{1}{\sqrt{2}}(\widehat{a}_j+\widehat{a}_{\overline{ȷ}}),\widehat{b}_{\overline{ȷ}}=\frac{i}{\sqrt{2}}(\widehat{a}_j^{}\widehat{a}_{\overline{ȷ}}^{})$$
(38)
and their Hermitian conjugates, among which the only nonvanishing commutators are the ordinary
$$[\widehat{b}_j,\widehat{b}_j^{}]=[\widehat{b}_{\overline{ȷ}},\widehat{b}_{\overline{ȷ}}^{}]=1.$$
(39)
In terms of these operators, which are harmonic oscillator annihilation and creation operators with all the familiar properties of such, the $`j,\overline{ȷ}`$ subsector of the Bogoliubov Hamiltonian $`\widehat{H}`$ appears as
$$\widehat{H}_{j,\overline{ȷ}}=\mathrm{Re}(\omega _j)[\widehat{b}_j^{}\widehat{b}_j\widehat{b}_{\overline{ȷ}}^{}\widehat{b}_{\overline{ȷ}}]\mathrm{Im}(\omega _j)[\widehat{b}^{}\widehat{b}_{\overline{ȷ}}^{}+\widehat{b}_j\widehat{b}_{\overline{ȷ}}].$$
(40)
Note that Eq. (40) is only the simplest form in which one may write the $`j,\overline{ȷ}`$ sector of the Hamiltonian: by introducing appropriate factors of $`e^{\pm i\alpha /2}/\mathrm{cos}\alpha `$ into Eq. (38), for any $`\alpha `$, we can make $`\mathrm{Im}(\omega _j)\mathrm{Im}(\omega _j)/\mathrm{cos}\alpha `$ and add a term $`\mathrm{Im}(\omega _j)\mathrm{tan}\alpha (\widehat{b}_j^{}\widehat{b}_j+\widehat{b}_{\overline{ȷ}}^{}\widehat{b}_{\overline{ȷ}})`$.
We can now examine the spectrum of $`\widehat{H}_{j,\overline{ȷ}}`$ by considering it in the basis of eigenstates of $`\widehat{n}=\widehat{b}_j^{}\widehat{b}_j+\widehat{b}_{\overline{ȷ}}^{}\widehat{b}_{\overline{ȷ}}`$ and $`\widehat{\mathrm{\Delta }}=\widehat{b}_j^{}\widehat{b}_j\widehat{b}_{\overline{ȷ}}^{}\widehat{b}_{\overline{ȷ}}`$. In fact $`\widehat{\mathrm{\Delta }}`$ commutes with $`\widehat{H}_{j,\overline{ȷ}}`$, so defining
$`|E_\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}c_n|n+\mathrm{\Delta }|n\mathrm{\Delta }0,`$ (41)
$`|E_\mathrm{\Delta }`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}c_n|n|n+\mathrm{\Delta }\mathrm{\Delta }0,`$ (42)
where $`\widehat{b}_j^{}\widehat{b}_j|m|n=m|m|n`$ and $`\widehat{b}_{\overline{ȷ}}^{}\widehat{b}_{\overline{ȷ}}|m|n=n|m|n`$, we find $`\widehat{H}_{j,\overline{ȷ}}|E_\mathrm{\Delta }=[\mathrm{}\mathrm{\Delta }\mathrm{Re}(\omega _j)+E_\mathrm{\Delta }]|E_\mathrm{\Delta }`$. And we have the recursion relation
$`E_\mathrm{\Delta }c_n=\mathrm{Im}(\omega _j)[`$ $`\sqrt{n(n+\mathrm{\Delta })}c_{n1}`$ (44)
$`+\sqrt{(n+1)(n+\mathrm{\Delta }+1)}c_{n+1}].`$
As $`n\mathrm{}`$, we have $`c_{n+1}c_{n1}`$, and so $`_n|c_n|^2`$ does not converge: none of the eigenstates of $`\widehat{H}_{j,\overline{ȷ}}`$ is normalizable. One can however obtain delta-function normalization for a continuous spectrum of real $`E_\mathrm{\Delta }`$, bounded neither above nor below.
That the Hamiltonian $`\widehat{H}`$ is unbounded from below does not indicate anything unphysical about our model: we have simply linearized about an unstable excited state of the nonlinear full Hamiltonian, which is bounded from below. Real negative frequencies $`\omega _j`$, where our convention $`N_{jj}=1`$ has been imposed, indicate energetic instabilities, whereby the system will decay in the presence of dissipation. Complex $`\omega _j`$, on the other hand, indicate dynamical instabilities. Classically, a dynamically unstable system will exponentially diverge from the initial stationary state if is perturbed, even without dissipation. Quantum mechanically, we have just seen that a dynamically unstable system has no normalizable stationary states. If an initially stable system is driven into a state which is stationary but dynamically unstable at the classical (mean field) level, the initial state will have had finite Hilbert space norm, and hence under unitary evolution the final state will have the same norm. Thus it will not be a stationary state; one may say that quantum fluctuations will always trigger the dynamical instability. For a logarithmically long period of time, however, the linearized theory will still remain valid. In this sense, our linearized description of quantum dynamical instabilities is sound.
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# Cross sections for Coulomb and nuclear breakup of three-body halo nuclei
## Introduction.
Halo nuclei are spatially extended bound systems . Their existence is revealed by large total interaction cross sections . The often small one or two-nucleon separation energies led to a successful description in terms of two or three-body structures . This in turn emphasizes the importance of reactions breaking these halo structures into their constituent particles. We shall here concentrate on two-neutron halos where the three-body problem is then inherent. Absolute values for dissociation cross sections of these systems are available in a number of cases . They may be divided into processes like one or two neutron knockout arising from Coulomb or nuclear dissociation. Also core destruction processes were studied and the total interaction cross section obtained .
Theoretical investigations of two-neutron dissociation cross sections are available on light as well as heavy targets, but core destruction processes are usually not computed at all. Total interaction cross sections for light targets are also calculated for three-body halos . However, systematic studies of all these processes within one consistent model are not available. Such results even for relative quantities are rather scarce . Reliable nuclear model estimates of absolute values are in general very difficult. In particular the necessary simultaneous treatment of the Coulomb and nuclear interactions has not been available for these three-body halo reactions.
The lack of systematic experimental and theoretical information about the many different absolute dissociation cross sections is perhaps surprising, but in any case unfortunate, since characteristic features of the reaction mechanism probably can be uncovered by such investigations. The purpose of this letter is then to study all possible three-body dissociation cross sections of two-neutron halo nuclei as functions of target and beam energy. The predictions employ a consistent three-body model and an appropriate reaction framework including simultaneous treatment of Coulomb and nuclear interactions.
## Theoretical formulation.
We shall here only give a brief sketch of the model leaving the details for a more comprehensive publication. We want to describe collisions of a target and a weakly bound two-neutron halo nucleus and classify the cross sections according to the particles left in the final state, i.e. $`\sigma _{nnc}`$, $`\sigma _{nc}`$, $`\sigma _{nn}`$, $`\sigma _c`$, $`\sigma _n`$, $`\sigma _0`$, where $`c`$ and $`n`$ denote core and neutron, respectively and $`0`$ indicates that all particles are absorbed.
The collisions can approximately be described by three independent collisions of one particle (participant) at a time and undisturbed motion of the remaining two projectiles particles (spectators) . The interconnected main assumptions are that the halo is weakly bound, spatially extended and the intrinsic halo motion is slow compared to the relative projectile-target motion.
In this simple picture the spectators are always found in the final state while the participant can be either absorbed or scattered by the target. However the finite extension of the projectile constituents and the target is destroying this simple picture, since simultaneous collisions of more than one projectile constituent with the target are possible.
The simplest description of the constituent-target interaction can be made through the black sphere model. This model assumes that a particle is absorbed when passing inside a cylinder with the axis along the beam direction and left untouched otherwise. The radius of this cylinder is approximately equal to the sum of the target and constituent radii. Different contributions to the cross section have to be considered, $`i)`$ three contributions where only one of the constituents is inside the cylinder (one participant and two spectators). In this case the interaction between the participant and the target is described by the phenomenological optical model, and only the part of the three-body wave function where the two spectators are far enough from the participant is included . $`ii)`$ Three contributions where two constituents are inside the cylinders. In this case, the interaction between one of them (participant) and the target is described by the optical model, while for the second constituent-target interaction we use the black sphere model. To implement this we include only the part of the wave function where the second constituent is close to the participant and the third one is far enough away. $`iii)`$ One contribution where the three constituents are inside the cylinders. Again the interaction of one of the constituents (participant) with the target is described by the optical model, and for the other two we use the black sphere model (only the part of the wave function where the three particles are close enough is included). It is then clear that the distances between the participant and the other two constituents appear as the decisive quantities determining if more than one particle interacts simultaneously with the target. These distances depend on the sizes of the constituents and the target and they are determined as in .
In the contributions to $`ii)`$ and $`iii)`$ the participant has to be chosen, and the participant-target interaction is described by the optical model. Since the core-target interaction contains nuclear and Coulomb forces a more careful treatment is required in this case. Therefore the contributions arising from simultaneous interactions of the core and one or two neutrons with the target are computed with the core as participant and the core-target interaction described by the optical model.
The division into different contributions implicitly assumes that for large impact parameters the constituent and the target are not interacting. Obviously this is not true for the long-range Coulomb interaction. To solve this problem, when the core is the participant, we include into the contributions $`i)`$, $`ii)`$ and $`iii)`$ only the large momentum transfer part (low impact parameter) of the Coulomb interaction. The low momentum transfer contribution (large impact parameter) is considered as a process where the core is elastically scattered by the target (via Coulomb interaction) and the two neutrons survive untouched in the final state. The value of the momentum transfer dividing into low and large impact parameters is given by $`q_g=Z_0Z_ie^2(\gamma +1)/(c\gamma \beta (R_0+R_i+\pi a/2))`$. Here $`R_0`$ and $`R_i`$ are charge root mean square radii of the target and the core participant (<sup>9</sup>Li or <sup>4</sup>He) and $`a`$ is half the distance of closest core-target approach, $`eZ_0`$ and $`eZ_i`$ are the charges of the target and participant, $`\beta =v/c`$ and $`\gamma =1/\sqrt{1\beta ^2}`$. The momentum cutoff parameter $`q_g`$ separates between impact parameters smaller and larger than the sum of participant and target radii. In the present context this means distinction between absorption (destruction) and survival of the participant.
The differential cross sections arising from participant $`i`$ with mass $`m_i`$ and charge $`eZ_i`$ has contributions from elastic scattering $`\sigma _{el}^{(0i)}`$ (diffraction) and absorption $`\sigma _{abs}^{(0i)}`$ (stripping) on the target. For a spinless target of mass $`m_0`$ and charge $`eZ_0`$ we get in the rest system of the halo
$`{\displaystyle \frac{d^6\sigma _{abs}^{(i)}(𝒑_{0i,jk}^{},𝒑_{jk}^{})}{d𝒑_{0i,jk}^{}d𝒑_{jk}^{}}}=\sigma _{abs}^{(0i)}(p_{0i})|M_s(𝒑_{i,jk},𝒑_{jk}^{})|^2,`$ (1)
where $`M_s`$ is the normalized overlap matrix element between initial and final state spectator wave functions, $`𝒑_{0i,jk}^{}`$ is the relative momentum in the final state between center of mass of target-participant and the spectators $`j`$ and $`k`$, while $`𝒑_{jk}^{}`$, $`𝒑_{0i}`$ and $`𝒑_{i,jk}`$ correspondingly are relative momenta between particles $`j`$ and $`k`$, $`0`$ and $`i`$, $`i`$ and center of mass of $`j`$ and $`k`$. Primes denote final states. Momentum conservation in the rest frame of the projectile gives the relation $`𝒑_{0i,jk}^{}=𝒑_{i,jk}+𝒑_0(m_j+m_k)/(m_0+m_i+m_j+m_k)`$, where $`𝒑_0`$ is the momentum of the target.
The differential elastic cross section is
$`{\displaystyle \frac{d^9\sigma _{el}^{(i)}(𝒑_{0i,jk}^{},𝒑_{jk}^{},𝒑_{0i}^{})}{d𝒑_{0i,jk}^{}d𝒑_{jk}^{}d𝒑_{0i}^{}}}={\displaystyle \frac{d^3\sigma _{el}^{(0i)}(𝒑_{0i}𝒑_{0i}^{})}{d𝒑_{0i}^{}}}`$ (2)
$`\times \left(1|\mathrm{\Psi }|\mathrm{exp}(i\delta 𝒒𝒓_{i,jk})|\mathrm{\Psi }|^2\right)|M_s(𝒑_{i,jk},𝒑_{jk}^{})|^2,`$ (3)
where $`\mathrm{\Psi }`$ is the initial three-body halo state. We have now a 9-dimensional differential cross section, since the participant explicitly is included in the final state. The momentum $`\delta 𝒒=(𝒑_{i,jk}^{}𝒑_{i,jk})(m_j+m_k)/(m_i+m_j+m_k)`$ is the transfer into the participant-spectators relative motion described by the coordinate $`𝒓_{i,jk}`$. The second factor in eq.(3) then expresses the probability for the halo not ending up in its ground state. In this way we remove elastic scattering of the halo as a whole.
When the participant is charged the Coulomb interaction produces a logarithmic divergence in the total cross section Eq.(3). The corresponding adiabatic motion related to virtual excitations at large impact parameters should be removed from the dissociation cross sections . We therefore exclude contributions from momentum transfer smaller than the adiabatic cutoff $`q_a=\mathrm{}Z_0Z_ie^2/(\pi \mathrm{}c)B_{ps}/(\mathrm{}c)(\gamma +1)\gamma ^2\beta ^2`$, where $`B_{ps}`$ is the binding energy between participant and the system consisting of the spectators, i.e., $`B_{ps}=BB_{2s}`$, where $`B`$ is the three-body binding energy and $`B_{2s}`$ is the two–body binding energy of the two spectators. Note that in a Borromean nucleus $`B_{2s}`$ is negative.
The energy transferred from target to participant, $`\delta E\sqrt{𝐩_0^2+m_0^2}\sqrt{𝐩_{\mathrm{𝟎}}^{}{}_{}{}^{2}+m_0^2}`$, must be larger than $`B`$. When $`𝐩_\mathrm{𝟎}`$ and $`𝐪𝐩_\mathrm{𝟎}𝐩_{\mathrm{𝟎}}^{}{}_{}{}^{}`$ are parallel $`\delta E`$ is maximized. For this geometry we find for small $`B`$ compared to the target rest mass that $`\delta E=B`$ implies that $`qcq_LcB\sqrt{1+m_0^2c^2/p_0^2}`$ which reduces to $`B/v`$ in the non-relativistic limit. Thus $`q`$ must be larger than $`q_L`$ to produce dissociation, but on the other hand dissociation is not the necessary outcome for all $`q>q_L`$. In the computations we exclude contributions from momentum transfer $`q`$ smaller than the largest of $`q_L`$ and $`q_a`$.
## Cross sections.
The model is now completely defined with both Coulomb and nuclear interactions included for weakly bound three-body halo reactions. We shall study breakup reactions of <sup>6</sup>He and <sup>11</sup>Li on C, Cu, and Pb targets. The parameters corresponding to the <sup>6</sup>He and <sup>11</sup>Li wave functions are obtained from . The optical model parameters are from for neutrons, from for $`\alpha `$-particles and for <sup>9</sup>Li also from but using range and diffuseness parameters from . We furthermore drastically reduce the energy dependence of the real part of the potential in , i.e. $`a_2=0.014`$, to allow for the required huge beam energy variation. The measured core-target interaction cross sections are reproduced within error bars . The binding energy $`B_{ps}`$ between the <sup>9</sup>Li and <sup>4</sup>He cores and the two neutrons must be introduced for the adiabatic cutoff. We use the scaling relation in to obtain $`B_{ps}/B3`$ for <sup>6</sup>He and 1.4 for <sup>11</sup>Li.
The corresponding contributions to the dissociation cross sections are obtained as indicated in $`i)`$, $`ii)`$ and $`iii)`$ by integration of Eqs.(1) and (3). The cross sections are then classified according to the particles left in the final state, i.e. $`\sigma _{nnc}`$, $`\sigma _{nc}`$, $`\sigma _{nn}`$, $`\sigma _c`$, $`\sigma _n`$, and $`\sigma _0`$. Specific interesting cross section combinations are those of two-neutron removal $`\sigma _{2n}\sigma _{nnc}+\sigma _{nc}+\sigma _c`$, core destruction $`\sigma _c\sigma _{nn}+\sigma _n+\sigma _0`$ and the sum of these, the total interaction cross section $`\sigma _I\sigma _{2n}+\sigma _c`$.
In fig. 1 we show the results for a lead target and <sup>11</sup>Li (upper part) and <sup>6</sup>He (lower part) projectiles. At low beam energies $`\sigma _{nnc}`$ is the dominant cross section, especially in the <sup>11</sup>Li case. This is due to the large Coulomb interaction, that highly increases the large impact parameter contribution (where the two neutrons both survive
in the final state). The $`\sigma _{nnc}`$ cross section rapidly decreases with the beam energy, and for energies larger than 200 MeV/nucleon in the <sup>11</sup>Li case and 100 MeV/nucleon in the <sup>6</sup>He case $`\sigma _0`$ dominates (no particles in the final state). The reason for this is the large radius of the Pb-target, resulting in a very high probability for finding all the three projectile constituents inside the absorption cylinders. This probability is higher for <sup>6</sup>He projectile than for <sup>11</sup>Li, because from the three-body wave function we get $`r_{cn}^2^{1/2}=4.2`$ fm in the first case and 5.9 fm in the second ($`r_{cn}=`$ neutron-core distance). As a consequence, when the core is the participant, 85% of the <sup>6</sup>He wave function corresponds to all three constituents inside the cylinders and only 65% for the <sup>11</sup>Li projectile. This is reflected in the figure by the fact that $`\sigma _0`$ takes very similar values in both cases, while as a general rule the cross sections for <sup>11</sup>Li should be larger than the ones for <sup>6</sup>He. The core destruction cross section $`\sigma _c=\sigma _0+\sigma _n+\sigma _{nn}`$ is shown in the insets by the short-dashed line. Its behavior is determined by the core–Pb absorption cross section, and changes very little with the beam energy. The two-neutron removal cross section $`\sigma _{2n}=\sigma _{nnc}+\sigma _{nc}+\sigma _c`$ is shown as the long-dashed lines in the insets, and it is given by all the processes where the core survives and contains therefore the contribution from the Coulomb interaction. Therefore this cross section decreases with beam energy. Finally the solid lines in the insets show the
interaction cross section $`\sigma _I=\sigma _{2n}+\sigma _c`$. The agreement with the experimental data is remarkably good.
In fig. 2 we show the results for a light target (carbon) where the cross sections are up to a factor of 10 smaller than for a lead target. The main difference is that the effect produced by the Coulomb interaction is now small, and for example $`\sigma _{2n}`$ is determined by Coulomb for Pb and by nuclear interactions for C targets. The dominating cross section is now $`\sigma _{nn}`$ for C instead of $`\sigma _{nnc}`$ or $`\sigma _0`$ for Pb targets. The reason is that the small radius of the target diminishes the probability of simultaneous interaction of more than one constituent with the target, and absorption of the core together with two truly undisturbed spectator neutrons becomes the most likely process. The variation with beam energy is in general rather similar to that of Pb, in particular revealed by $`\sigma _c`$, $`\sigma _{2n}`$ and $`\sigma _I`$ shown in the insets. Also in this case we obtain good agreement with the experimental data.
In fig. 3 we show the results for a copper target, i.e. an intermediate mass with comparable Coulomb and nuclear contributions. The cross sections are between those obtained for carbon and lead with similar energy dependences. However, <sup>11</sup>Li has a larger radius than <sup>6</sup>He. Thus the dominating cross section for <sup>11</sup>Li is $`\sigma _{nn}`$ as for carbon, while $`\sigma _0`$ dominates for <sup>6</sup>He as for a lead target.
Only a few of the computed cross sections are experimentally available. In experimental data for $`\sigma _{nnc}`$, $`\sigma _{nc}`$ and $`\sigma _c`$ are given for <sup>6</sup>He on C and Pb at 240 MeV/nucleon. In the experimental data for the same cross sections are given for <sup>11</sup>Li fragmentation on C and Pb at 280 MeV/nucleon. Also in $`\sigma _{nnc}`$, $`\sigma _{nc}`$ and $`\sigma _c`$ are calculated for <sup>11</sup>Li and <sup>6</sup>He on C by use of the eikonal theory. In table I we compare the results of our calculation with the ones given in and the experimental data in . For carbon target the agreement with the experimental data and previous calculations is good. The only significant deviation is in $`\sigma _{nc}`$ for the case of lead target, where our calculation gives a cross section at least a factor of two smaller than the experimental value.
The core destruction processes entering in $`\sigma _c`$ all must take place at small impact parameters. The Coulomb contributions are therefore relatively unimportant in contrast to the large impact parameter processes where purely Coulomb dissociation reactions take place. Coulomb dissociation cross sections have been previously estimated in models where the three–body system is treated as an effective two–body (dineutron plus core). For <sup>11</sup>Li on lead at 800 MeV/nucleon we get 765 mb, while in they obtain 960 mb, and in they report values ranging from 610 to 660 mb. For <sup>11</sup>Li on Cu at 800 MeV/nucleon we obtain a Coulomb dissociation cross section of 81 mb, while in they give values from 86 to 92 mb. Calculations where the final continuum three–body wave function is considered are also available . For <sup>11</sup>Li on Pb at 180, 280, and 800 MeV/nucleon they get 2128 mb, 1429 mb and 971 mb, respectively, while our calculations for the same energies give 2050 mb, 1401 mb and 765 mb. These rather few previous computations where a comparison is possible are in general in rough agreement with our systematic results.
## Conclusion.
We have computed all possible three-body dissociation cross sections of two-neutron halos as function of beam energy and target. Coulomb and nuclear interactions are treated within the same framework. The available experimental information, mostly about total interaction and two-neutron separation cross sections, compares rather well with the calculated results, especially at high energies. At low energies the three–body continuum final state might turn out to be more appropriate. The further division into a specific number of particles in the final states carries detailed information about the reaction mechanism. We predict the absolute sizes of all these cross sections to encourage new measurements.
## Acknowledgement.
We thank K. Riisager for continuous discussions and suggestions.
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# Dynamics and transport in random quantum systems governed by strong-randomness fixed points
## I Introduction
Disorder effects arising from quenched randomness are at the heart of many interesting and novel phenomena observed in condensed matter systems—examples include Griffiths singularities near phase transitions in disordered magnets (and the related phenomenon of local-moment formation in disordered electronic systems ), metal-insulator transitions in disordered electronic systems, and two-dimensional phenomena such as weak localization and the quantum Hall plateau transitions.
In particular, the interplay between disorder and quantum interference leads to unusual dynamics and transport in these systems. Such effects are well understood for disordered quantum systems in which many-body correlations are not significant (such as disordered Fermi liquids). In contrast, relatively little is reliably known about the effects of strong disorder in the presence of strong correlations (say, due to electron-electron interactions in an itinerant electronic system, or due to exchange interactions in a system with localized spin degrees of freedom).
However, there does exist one class of systems where theoretical tools are available to analyze this interplay between strong disorder, correlations and quantum fluctuations; important examples include one-dimensional random antiferromagnetic spin chains and random quantum Ising models in one and two dimensions . In these quantum systems, it is possible to systematically treat disorder and correlation using a strong-disorder renormalization group (RG) technique that is designed to be accurate when the strength of the disorder, as measured by the widths of the distributions of the various couplings, is large. Such a strong-disorder approach works in these problems because these systems, when studied at ever larger length scales (and correspondingly lower energy scales), appear more and more disordered. More precisely, the low-energy effective theory obtained from the RG has the remarkable property that the widths of the distributions of the various couplings in the theory grow rapidly as the energy cutoff is lowered; this means that the RG procedure gives reliable results for the effective Hamiltonian that governs the low-energy properties of the system. Moreover, the extremely strong disorder present at low energies in the effective theory actually allows one to straightforwardly calculate some thermodynamic properties and ground state correlators within the effective theory—this is, in essence, because strong disorder implies that some particular terms in the effective Hamiltonian dominate over all others; calculations can then be performed by treating these terms first and including the effects of the other terms perturbatively. This approach has been used successfully in the past to obtain a wealth of information about the low-temperature thermodynamics and ground state correlators in such systems.
Here, we exploit this simplicity that emerges at strong disorder to obtain the first analytical results on the low-frequency dynamics and transport in these systems at low temperature $`T`$. Most of our results are obtained for $`T=0`$; these are expected to be exact at zero temperature in the low-frequency limit, and to remain valid at non-zero temperatures for low frequencies $`\omega T`$. Moreover, in certain special cases, we can also access the regime $`\omega <T`$.
In the remainder of this section, we introduce the various systems that are studied in this paper, and describe the organization of the rest of the paper. A brief summary of some of our results has already appeared elsewhere.
Our focus is on three model systems. The first model we consider is the one-dimensional random antiferromagnetic XXZ spin-1/2 chain with the Hamiltonian
$$_{\mathrm{XXZ}}=\underset{j}{}\left[J_j^{}(s_j^xs_{j+1}^x+s_j^ys_{j+1}^y)+J_j^zs_j^zs_{j+1}^z\right],$$
(1)
where $`\stackrel{}{s}_j`$ are spin-1/2 operators at lattice sites $`j`$ separated by spacing $`a`$, and both $`J_j^{}`$ and $`J_j^z`$ are random positive exchange energies drawn from some probability distributions. Such a Hamiltonian describes the low-energy (magnetic) dynamics of insulating antiferromagnetic spin-1/2 chain compounds with chemical disorder that affects the bond strengths. We will also consider chains with slightly different probability distributions of the even and the odd bonds and study the effects of such enforced dimerization. The strength of the dimerization in the bonds is conveniently characterized by a dimensionless parameter $`\delta `$ defined as
$$\delta =\frac{\overline{\mathrm{ln}J_\mathrm{e}}\overline{\mathrm{ln}J_\mathrm{o}}}{\mathrm{var}(\mathrm{ln}J_\mathrm{e})+\mathrm{var}(\mathrm{ln}J_\mathrm{o})},$$
(2)
where $`J_\mathrm{e}`$ ($`J_\mathrm{o}`$) represents even (odd) bonds, and the overline and “var” denote correspondingly the average and variance over the distribution of bonds. Thus, we have $`\delta >0`$ ($`\delta <0`$) if even (odd) bonds are stronger on average. For future reference, we also introduce the basic length scale in this system,
$$l_v=\frac{2a}{\mathrm{var}(\mathrm{ln}J_\mathrm{e})+\mathrm{var}(\mathrm{ln}J_\mathrm{o})}.$$
(3)
Detailed information about the spin dynamics in such systems can be obtained by inelastic neutron scattering (INS) experiments that directly probe the frequency and momentum dependent dynamic structure factor $`S^{\alpha \beta }(k,\omega )`$. At $`T=0`$, $`S^{\alpha \beta }(k,\omega )`$ has the spectral representation
$$S^{\alpha \beta }(k,\omega )=\frac{1}{L}\underset{m}{}0|\widehat{s}_k^\alpha |mm|\widehat{s}_k^\beta |0\delta (\omega E_m),$$
(4)
where $`\widehat{s}_k^\alpha =_je^{ikx_j}s_j^\alpha `$, and $`\{|m\}`$ denote the exact eigenstates of the system with excitation energies $`E_m`$ relative to the ground state $`|0`$. The symmetry of $`_{\mathrm{XXZ}}`$ under rotations about the $`z`$ axis implies that we can restrict our attention to two independent components $`S^{zz}`$ and $`S^+`$. The same symmetry also implies that the total $`s_{\mathrm{tot}}^z=_js_j^z`$ is conserved—it then makes sense to talk of the spin transport in such a system. We characterize the transport of $`s^z`$ in terms of the dynamical spin conductivity $`\sigma (\omega )`$. The real part $`\sigma ^{}(\omega )`$ of $`\sigma (\omega )`$ is defined by the relation $`P(\omega )=\sigma ^{}(\omega )|H|^2(\omega )`$, where $`P(\omega )`$ is the power absorbed per unit volume by the system when magnetic field with a uniform gradient $`H(\omega )`$ (with the field $`H`$ always in the $`z`$ direction) oscillating at frequency $`\omega `$ is applied along the length of the chain. From standard linear response theory, we have the following Kubo formula for $`\sigma ^{}(\omega )`$ at $`T=0`$:
$$\sigma ^{}(\omega )=\frac{1}{\omega L}\underset{m}{}|m|\underset{j=1}{\overset{L}{}}\tau _j|0|^2\delta (\omega E_m).$$
(5)
In the above, $`\tau _j=iJ_j^{}(s_j^+s_{j+1}^{}s_{j+1}^+s_j^{})/2`$ is the current operator on link $`j`$ that transfers one unit of the $`s^z`$ from one site to the next. Here and everywhere in the following, the frequency $`\omega `$ is taken positive for notational convenience. Note that both $`S^{\alpha \beta }(k,\omega )`$ and $`\sigma ^{}(\omega )`$ as defined here are self-averaging in the thermodynamic limit.
The second model we consider is the random antiferromagnetic Heisenberg spin-1 chain with the Hamiltonian
$$_{\mathrm{S1}}=\underset{j}{}J_j\stackrel{}{S}_j\stackrel{}{S}_{j+1},$$
(6)
where $`\stackrel{}{S}_j`$ are spin-1 operators on lattice sites $`j`$, and the $`J_j`$ are random positive nearest-neighbor exchanges; randomness in the system is characterized by a width $`W`$ of the corresponding distribution of log-exchanges $`\mathrm{ln}(J_j)`$. As in the spin-1/2 case, we can characterize spin dynamics and transport in terms of the dynamic structure factor and the dynamical conductivity; the definitions remain the same except for the obvious replacement of all spin-1/2 operators with their spin-1 counterparts. Experimental realizations of pure Heisenberg spin-1 chains are known, and experimental studies of systems with randomness have also been reported in the recent literature. We caution, however, that the degree of disorder needed to destroy the gapped Haldane phase of a pure spin-1 chain appears to be quite strong, and that all our calculations are done only in this strong-disorder regime.
The third problem that we consider is the one-dimensional random transverse field Ising model
$$_{\mathrm{RTFIM}}=\underset{j}{}J_j\sigma _j^z\sigma _{j+1}^z\underset{j}{}h_j\sigma _j^x,$$
(7)
with random ferromagnetic interactions $`J_j`$ and positive random transverse fields $`h_j`$; here $`\sigma _j`$ are Pauli spin matrices. The strong-disorder RG approach, and its consequences for the low-temperature thermodynamics and static correlators, have been analyzed in greatest detail for this particular model. Also, there are extensive numerical results available for some dynamical properties. This model thus serves as a benchmark to test reliability of our approach to the calculation of dynamical properties in these strong-disorder systems—we will analyze various average autocorrelation functions in considerable detail and compare our results with the earlier numerical work.
The paper is organized as follows: We begin in Sec. II with a general discussion of the various types of states that we encounter in these models, along with an overview of our most important results for the dynamics and transport in various regimes; the last part of this section is devoted to a general outline of the basic approach that is used to obtain these results. The following Sections IIIIV, and V, present careful derivations of our results for the zero-temperature dynamical properties of the three model systems that we consider, with each section starting with a review of the basic RG approach used to study the corresponding system. In Sec. III we evaluate the dynamic structure factor and the dynamical conductivity in the various phases of the random XXZ spin-1/2 chain. This is followed, in Sec. IV, by an analysis of the spin conductivity in the strongly-random Heisenberg antiferromagnetic spin-1 chains, and, in Sec. V, by an analysis of the average local dynamical properties of the random quantum Ising model in the vicinity of its critical point. Section VI is devoted to a qualitative analysis of the dynamical and transport properties of the XXZ spin-1/2 chains at non-zero temperatures in the regime $`\omega <T`$, along with some quantitative calculations in the XX spin-1/2 chain that are possible in this case because of the mapping to free fermions. We conclude, in Sec. VII, with a discussion of the possible experimental tests of some of our predictions for the one-dimensional random-exchange antiferromagnetic spin chains. Some technical details, as well as some additional developments slightly removed from the main thrust of the paper, are relegated to the appendices.
## II Overview
Broadly speaking, our results are for two types of states. First, there are ground states governed (and therefore best described by some suitable strong-disorder RG approach) by infinite-randomness fixed points—examples include the random singlet states of the spin-1/2 antiferromagnetic chains and the critical point of the random transverse field Ising model. Then, there are the so-called “Griffiths” phases in the immediate vicinity of these critical states; in these phases, the low-energy renormalized randomness is strong, but not infinite.
In both cases the low-energy excitations are localized, but with a characteristic ‘localization length’—i.e. the ‘size of the excitation’—that diverges as a power of $`\mathrm{ln}\omega `$ for energy $`\omega 0`$. \[We emphasize that this is the statement about the (rare) low-energy excitations and is indeed valid in the Griffiths phases, even though in this case all equal-time correlators at $`T=0`$ indicate a finite localization length; for details see the main body of the paper.\] Apart from this logarithmically divergent ‘localization length’, we can also define, from the integrated density of states $`n_\omega `$ for excitations up to energy $`\omega `$, a length $`L_\omega n_\omega ^{1/d}`$ that is the typical spacing between these excitations in $`d`$ dimensions \[the results we report here are for $`d=1`$, but similar phases do occur for $`d>1`$ \].
For a ground state governed by an infinite-randomness fixed point, $`L_\omega `$ diverges at low energies with the same power of $`\mathrm{ln}\omega `$ as the typical size of the excitation. This means a strongly-divergent density of states at low energy, which allows the system to behave as a conductor if there is a conserved quantity (e.g., spin or particle number) to be transported. In a Griffiths phase, on the other hand, $`L_\omega \omega ^{1/z}`$, with $`z`$ a nonuniversal dynamical exponent that varies continuously within the phase. Here, the low-energy excitations are rare; they are typically spaced by distance $`L_\omega `$, which diverges as a power-law at low energy and thus is much larger than the excitation’s typical size, which is diverging only logarithmically. In the RG language, the Griffiths phases are governed by lines of fixed points ending in the infinite-randomness critical fixed point; along such a line, the dynamical exponent $`z`$ varies continuously and diverges near the critical point.
In terms of the original microscopic model, the low-lying excitations in the Griffiths phases come from regions where the local quenched random variables deviate strongly from their global averages. These deviations are such that the local averages would put that region in a different phase. If the system is not at a phase transition, the probability of such a rare region occuring and being of linear size $`L`$ behaves as $`e^{c_1L^d}`$ for large $`L`$, for some constant $`c_1`$. Such a rare region typically results in a low-lying mode with a sharply defined (in the sense that $`c_2`$, introduced below, is sharply defined) characteristic frequency proportional to $`e^{c_2L^d}`$. This gives rise to a power-law low-energy density of states, with the dynamical exponent $`z`$ being determined by the ratio of the constants $`c_1/c_2`$. For a disordered Griffiths phase, the rare regions are finite “islands” of either an ordered phase, or a different disordered phase. The resulting low-lying excitations localized on these rare regions produce a low-frequency conductivity $`\sigma ^{}(\omega )`$ or scaled dynamic structure factor $`\omega S(k,\omega )`$ vanishing as $`\omega ^{1/z}`$ at low frequencies (apart from possible logarithmic factors attributable to singular low-energy behavior of the relevant matrix elements that may, in some cases, be sensitive to the logarithmically divergent size of the relevant excitations).
For one-dimensional systems, there are also power-law Griffiths effects in Ising-ordered phases. These occur because of rare regions locally in the disordered phase. The low-energy excitation associated with such a region is a domain wall (or “kink”). To produce a single such low-energy domain wall requires flipping the spontaneous magnetization on one side of the the wall, which is tantamount to flipping a semi-infinite piece of the chain. Such a flip of an infinite domain cannot occur at a finite (non-zero) frequency. The leading contribution to the low-frequency dynamics is then associated with two nearby such rare low-energy domain walls which allow the ordered domain between them to flip at a low but nonzero frequency. The result of this is that the low-frequency $`\sigma ^{}(\omega )`$ and $`\omega S(k,\omega )`$ vanish as $`\omega ^{2/z}`$ at low frequency in these one-dimensional Ising-ordered Griffiths phases (we are again ignoring possible logarithmic factors which can arise for precisely the same reasons as in the disordered phase). Note however that the Griffiths singularities in Ising-ordered phases in $`d>1`$ are of a very different character—in these cases, the low-energy density of states vanishes faster than any power of $`\omega `$, as is discussed in Ref. .
In Sections III-V we will provide a detailed justification of these general observations by explicitly calculating the low-frequency dynamical properties in a variety of cases. In the rest of this Section, we review the phase diagrams of our model systems, and highlight our most important results in each case.
### A Random antiferromagnetic XXZ spin-1/2 chains
#### 1 Phase diagram
The phase diagram of the random antiferromagnetic XXZ spin-1/2 chains is best understood as a product of the competition between the transverse part of the coupling $`J^{}`$, which favors singlet formation, and the ‘classical’ interaction $`J^z`$, which favors a ground state with Ising antiferromagnetic order.
When the $`J^{}`$ dominate, the ground state can be loosely thought of as being made up of singlet pairs. In this Random Singlet (RS) state, the interplay of disorder and quantum fluctuations locks each spin into a singlet pair with another spin; the two spins in a given singlet pair can have arbitrarily large spatial separation, with the disorder determining the particular pattern of the singlet bonds in a given sample. On the other hand, when the $`J^z`$ dominate, the system has Ising antiferromagnetic (IAF) order in the ground state (with the spins all oriented parallel to the $`z`$ axis), although Griffiths effects can fill in the gap leading to an IAF ordered Griffiths phase.
These two states are separated by a quantum phase transition that occurs when the couplings $`J^{}`$ and $`J^z`$ have roughly similar distributions (have roughly equal strengths). A special feature of this system is that the ground state at any point on the critical manifold is also a Random Singlet state, though the details of the excitation spectrum are somewhat different.
If we now turn on enforced bond dimerization starting with the RS state that obtains for small $`J^z`$, or the RS state of the Heisenberg chain, the system moves into a Griffiths phase dubbed the Random Dimer (RD) phase; in this phase the singlet bonds in the ground state now preferentially start on one sublattice and end on the other.
Schematic phase diagrams summarizing the above are shown in Figs. 1 and 2.
#### 2 Spin transport
We characterize the spin transport properties of the various phases in terms of the low-frequency behavior of the dynamical conductivity: We find that the $`T=0`$ dynamical conductivity diverges at low frequencies in the RS phase as well as at the RS critical points as
$$\sigma ^{}(\omega )=𝒦_{\mathrm{RS}}l_v\mathrm{\Gamma }_\omega ,$$
(8)
where we have taken the opportunity to introduce the log-energy scale
$$\mathrm{\Gamma }_\omega \mathrm{ln}(\mathrm{\Omega }_0/\omega ).$$
(9)
Here and henceforth we use $`\mathrm{\Omega }_0`$ to denote the non-universal microscopic energy cutoff, which corresponds roughly to the energy scale in the bare Hamiltonian for our various models; also, we use $`l_v`$ to denote the non-universal microscopic length scale in the problem. For the XXZ spin-1/2 system near the RS phase and with sufficiently strong disorder, which is what we assume in the following, the microscopic length $`l_v`$ is given by Eq. (3). \[If, on the other hand, the bare disorder is weak and the system flows to strong disorder, then $`l_v`$ is the length scale at which the strength of the disorder becomes of order one.\] $`𝒦_{\mathrm{RS}}`$ in Eq. (8) is an order one numerical constant. The RS phase and the RS critical points separating it from the IAF phase are thus unusual spin conductors.
On the other hand, the IAF Griffiths phase is a spin insulator with the low-frequency $`T=0`$ dynamical conductivity
$$\sigma ^{}(\omega )=𝒦_{\mathrm{IAF}}l_v(\omega /\mathrm{\Omega }_0)^{2/z_{\mathrm{IAF}}}\mathrm{ln}^2(\mathrm{\Omega }_0/\omega ),$$
(10)
where $`z_{\mathrm{IAF}}(\delta _{\mathrm{IAF}})`$ is a (continuously varying) dynamical exponent diverging at the critical point as $`z_{\mathrm{IAF}}\delta _{\mathrm{IAF}}^{(2\psi )/\lambda }`$, and $`𝒦_{\mathrm{IAF}}`$ is a non-universal amplitude vanishing at the transition as $`𝒦_{\mathrm{IAF}}\delta _{\mathrm{IAF}}^{(2\psi )/\lambda }`$. Here we parametrized the distance from the transition to the RS phase by $`\delta _{\mathrm{IAF}}\overline{\mathrm{\Delta }}\overline{\mathrm{\Delta }}_c`$ (where $`\mathrm{\Delta }J^z/J^{}`$). The exponent $`\lambda `$ is the relevant RG eigenvalue controlling the flow away from the critical fixed point describing the generic transition between the RS phase and the IAF phase, and the exponent $`\psi `$ characterizes the low-energy spectrum above the RS ground state at this critical point (see Ref. and Sec. III A for details). The above result is expected to hold in the frequency regime $`\omega \mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$ with the cross-over scale $`\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$ given in terms of $`\delta _{\mathrm{IAF}}`$ as $`\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega }_{\delta _{\mathrm{IAF}}})\delta _{\mathrm{IAF}}^{(2\psi )/\lambda }`$.
Similarly, the RD phases are also spin insulators, with the $`T=0`$ low-frequency dynamical conductivity
$$\sigma ^{}(\omega )=𝒦_{\mathrm{RD}}l_v(\omega /\mathrm{\Omega }_0)^{1/z_{\mathrm{RD}}}\mathrm{ln}^2(\mathrm{\Omega }_0/\omega );$$
(11)
the dynamical exponent $`z_{\mathrm{RD}}(\delta )`$ in the RD phase diverges at the transition as $`z_{\mathrm{RD}}|\delta |^1`$, and the non-universal amplitude $`𝒦_{\mathrm{RD}}`$ vanishes at the transition as $`𝒦_{\mathrm{RD}}|\delta |`$. As in the IAF phase, this result is valid at frequencies well below the corresponding crossover scale $`\mathrm{\Omega }_\delta `$ (which can be also viewed as the conductivity pseudo-gap scale); in the RD phases $`\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega }_\delta )1/|\delta |`$.
Thus, in both the IAF phase and the RD phase, the conductivity has the functional form
$$\sigma ^{}(\omega )\omega ^\alpha \mathrm{ln}^2\omega ,$$
(12)
with the non-universal exponent $`\alpha `$ vanishing at the corresponding transition. Note that a similar form but with fixed $`\alpha =2`$ —the Mott formula— is obtained via the usual Mott argument for the $`T=0`$ dynamical conductivity of the one-dimensional Anderson insulator (the fixed value of $`\alpha `$ in this case simply reflects the fact that low-energy density of states in the Anderson insulator is constant, in contrast to the situation in the Griffiths phases of interest to us here).
#### 3 Spin dynamics
Turning to the spin dynamics, we find that the $`T=0`$ dynamic structure factor in the RS states in the vicinity of $`k=\pi /a`$ can be written in the following unusual scaling form
$$S^{\alpha \beta }(k=\frac{\pi }{a}+q,\omega )=\frac{𝒜}{l_v\omega \mathrm{ln}^3(\mathrm{\Omega }_0/\omega )}\mathrm{\Phi }\left(|ql_v|^{1/2}\mathrm{ln}(\mathrm{\Omega }_0/\omega )\right)$$
(13)
for $`|q|a^1`$ and $`\omega \mathrm{\Omega }_0`$; here $`\alpha \beta +`$ or $`zz`$, $`𝒜`$ is an order one numerical constant, $`l_v`$ is the microscopic length defined earlier, and the fully universal function $`\mathrm{\Phi }(x)`$, calculated in Sec. III, interpolates smoothly between the limiting forms $`\mathrm{\Phi }(x)27x^4/90`$ for $`x1`$ and $`\mathrm{\Phi }(x)1+2x(\mathrm{cos}x+\mathrm{sin}x)e^x`$ for $`x1`$. A plot of the momentum dependence of the dynamic structure factor near $`k=\pi /a`$ (at fixed low frequency) is shown in Fig. 3; an interesting aspect is the non-monotonic nature of the lineshape. We will see in Sec. III that this oscillatory behavior becomes more pronounced and leads to really striking structure in the momentum dependence of the dynamic structure factor at (fixed) low frequency $`\omega \mathrm{\Omega }_\delta `$ in the Random Dimer phases—a plot of the expected $`k`$ dependence is shown in Fig. 4. A very similar dependence is also predicted in the IAF Griffiths phase close to the transition to the RS state.
As mentioned earlier, these results are expected to remain valid at small non-zero temperatures so long as the frequency $`\omega `$ satisfies $`\omega T`$. In Sec. VI, we will see that we can partially overcome even this restriction in the vicinity of the XX point.
### B Spin-1 Heisenberg antiferromagnetic chains
#### 1 Phases
The effect of randomness on antiferromagnetic Heisenberg spin-1 chains is even more interesting. For spin-1, there are, in general three distinct phases possible in the presence of disorder. If the disorder is weak, and the support of the probability distribution $`P(J)`$ of the exchanges is confined to a narrow enough region near the mean, then the system remains in the usual gapped, topologically ordered Haldane state. For stronger disorder, or when $`P(J)`$ has tails to large or small enough $`J`$, one has the ‘Gapless Haldane’ (GH) phase in which the system still has the topological order that characterizes the Haldane state, but becomes gapless due to Griffiths effects. Finally, if the disorder is extremely strong, with the (bare) distribution of exchanges broad on a logarithmic scale, a Random Singlet state completely analogous to the one encountered in the spin-1/2 case obtains. While the GH state and the RS state are separated by a quantum critical point with universal critical properties (these properties are in fact controlled by a strong-disorder fixed point ), the corresponding transition between the gapped and gapless Haldane states is a non-universal feature of the phase diagram, depending sensitively on the nature of the initial distribution of couplings (see Fig. 5 for a summary of the universal aspects of the phase diagram).
#### 2 Overview of results
In the spin-1 RS state, we obtain the same results for the dynamic structure factor and spin conductivity as in the spin-1/2 RS state, as the low-energy behavior of the RS state does not depend on the spin magnitude except through the values of some microscopic scale factors. Unfortunately, once we move away from the Random Singlet state, it is difficult to discuss reliably the momentum dependence of the dynamic structure factor of the original spin-1 chain—this is because our actual calculations are done in an effective model (see section IV A and Refs. for details) in which much of the spatial information about the original system is missing.
However, it is still possible to calculate transport properties, such as the dynamical conductivity, that are insensitive to the details of the spatial structure (this is, in essence, a consequence of spin conservation). At the critical point separating the Gapless Haldane state from the Random Singlet state, we find for $`\omega \mathrm{\Omega }_0`$
$$\sigma ^{}(\omega )=𝒦_{\mathrm{HY}}l_v\mathrm{ln}^2(\mathrm{\Omega }_0/\omega ),$$
(14)
which is a stronger divergence than in the strong-disorder RS phase; here $`l_v`$ is the non-universal microscopic length scale beyond which the effective model applies, $`\mathrm{\Omega }_0`$ is the corresponding microscopic energy scale, and $`𝒦_{\mathrm{HY}}`$ is an order one numerical constant. Thus, the critical point separating the RS phase from the GH phase is also an unconventional ‘spin metal’. The GH phase, on the other hand, is a ‘spin insulator’, not unlike the RD phase of spin-1/2 chains. We find for the conductivity in the GH phase
$$\sigma ^{}(\omega )𝒦_{\mathrm{GH}}l_v(\omega /\mathrm{\Omega }_0)^{1/z_{\mathrm{GH}}}\mathrm{ln}^2(\mathrm{\Omega }_0/\omega ).$$
(15)
The dynamical exponent $`z_{\mathrm{GH}}`$ varies continuously in the Gapless Haldane phase, diverging at the critical point as $`z_{\mathrm{GH}}(W_cW)^{\nu /3}`$, while the the non-universal amplitude $`𝒦_{\mathrm{GH}}(W)`$ remains non-zero as one approaches the critical point. In the above, $`W_c`$ is the critical value of the bare disorder (the parameter $`W`$ has already been defined in Sec. I), and the correlation length exponent $`\nu =6/(\sqrt{13}1)`$ is known from the analysis of Refs. and .
### C Random quantum Ising spin chains
#### 1 Phases
The self-dual nature of the random transverse field Ising model in one dimension implies that the system will be in a critical state if the distributions of bonds and fields are identical. The deviation from criticality may be parametrized by
$$\delta =\frac{\overline{\mathrm{ln}h}\overline{\mathrm{ln}J}}{\mathrm{var}(\mathrm{ln}h)+\mathrm{var}(\mathrm{ln}J)},$$
(16)
with $`\delta >0`$ corresponding to the quantum disordered paramagnet, and $`\delta <0`$ corresponding to the ordered ferromagnet. (Note that we use ‘$`\delta `$’ both as a dimensionless measure of dimerization in spin-1/2 chains, and in the present context; there is however no cause for confusion and the meaning will always be clear from the context in what follows.)
This quantum critical point is flanked, for small $`|\delta |`$ on either side, by paramagnetic and ferromagnetic Griffiths phases with gapless excitations.
#### 2 Overview of results
As mentioned earlier, these Griffiths phases and the quantum critical point separating them are among the best understood examples of such strong-randomness phenomena. However, all previous analyses of the dynamical properties relied on numerical results supplemented by scaling ideas.
In contrast, our approach allows us to analytically calculate the average local autocorrelations of both the spin and the energy operators at, and in the vicinity of, the quantum critical point, as well as obtain the scaling behavior of the dynamic structure factor of the spins. The main features of the average autocorrelations (as well as distributions of autocorrelations, which we do not address here) have already been noted in the earlier numerical work (Refs. ), while our results on the dynamical structure factor are new. Here, we only highlight some of the subtleties, missed in these numerical studies, that our analytical work has uncovered regarding the autocorrelations—a complete tabulation of our predictions (and their interpretation in terms of Griffiths effects) is given in Sec. V.
Our results for the $`T=0`$ imaginary-time off-critical spin autocorrelation in the bulk have the form
$$\left[C_{\mathrm{loc}}\right]_{\mathrm{av}}(\tau )\frac{|\mathrm{ln}\tau |}{\tau ^{n/z(\delta )}},$$
(17)
where $`z(\delta )`$ is the continuously varying dynamical exponent characterizing the Griffiths phases (from the results of Ref. , $`z^12|\delta |`$, for small enough $`\delta `$). In the above, the parameter $`n`$ distinguishes between the disordered and ordered phases with $`n=1`$ in the disordered phase and $`n=2`$ in the ordered phase. Thus, the exponent controlling the power-law decay in the ordered Griffiths phase is twice $`z^1`$, while the corresponding exponent in the paramagnetic Griffiths phase is $`z^1`$. This reflects the physical distinction between the disordered and the Ising ordered Griffiths phases noted in our general discussion at the beginning of this Overview. Moreover, the autocorrelations in the Griffiths phases are not pure power-law, but have a logarithmic correction, which reflects the fact that the appropriate ‘spin’ degrees of freedom relevant at a time-scale $`\tau `$ have an effective moment of order $`|\mathrm{ln}\tau |`$. Both these subtleties have been ignored when extracting the dynamical exponent from the numerical results for the average spin autocorrelations via the ansatz $`\left[C_{\mathrm{loc}}\right]_{\mathrm{av}}(\tau )\frac{1}{\tau ^{1/z(\delta )}}`$, and this could account for some of the discrepancies observed in the numerical studies. Similar remarks apply to other average autocorrelations considered, and we refer to Sec. V for details.
### D The basic strategy
We conclude with an overview of the basic strategy introduced by us in Ref. for the calculation of dynamical and transport properties—we will be using this approach over and over again in what follows, and while the details will differ from calculation to calculation, the basic approach will remain unchanged.
Consider, for concreteness, the calculation of the dynamic structure factor $`S^{\alpha \beta }(k,\omega )`$ for the Hamiltonian $`_{\mathrm{XXZ}}`$. The basic idea is to eliminate high-energy degrees of freedom using an appropriate strong-disorder renormalization group procedure (in this case, the singlet RG reviewed in Sec. III A), and trade in the spectral sum Eq. (4) for a sum over the eigenstates of the renormalized Hamiltonian $`\stackrel{~}{}_{\mathrm{XXZ}}`$, which has fewer degrees of freedom and renormalized bond strengths. This renormalized spectral sum must use the matrix elements of the renormalized versions of the spin operators; these renormalized operators are of course defined by the requirement that their matrix elements between the eigenstates of the renormalized problem reproduce the matrix elements of the original operators between the corresponding eigenstates of the original problem. In the systems of interest to us, the low-energy renormalized randomness is very large. In the renormalized problem at the energy cutoff $`\mathrm{\Omega }\mathrm{\Omega }_0`$, the effective bonds thus have a very broad distribution characteristic of the fixed point to which the system flows in the low energy limit. This allows us to reason as follows: Focus on pairs of spins coupled by ‘strong’ bonds in the renormalized problem, with strengths equal to the cutoff $`\mathrm{\Omega }`$. The broad distribution of bonds implies that these pairs are effectively isolated from their neighbors. It is therefore possible to unambiguously identify the excited states of these pairs with excitations of the full system at the same energies and work out the matrix elements connecting these to the ground state using the renormalized operators. Thus, to calculate the spectral sum Eq. (4), the RG is run till the cutoff $`\mathrm{\Omega }`$ equals $`\mathrm{\Omega }_{\mathrm{final}}`$, and the problem is reduced to calculating the renormalized spectral sum in this new theory; $`\mathrm{\Omega }_{\mathrm{final}}`$ is chosen so that the energy of such excited states (associated with these strong bonds) measured from the ground state equals $`\omega `$. The calculation of $`S^{\alpha \beta }(k,\omega )`$ then becomes a counting problem. One uses the known statistical properties of the renormalized bonds in the theory with cutoff $`\mathrm{\Omega }_{\mathrm{final}}`$ to calculate the number of such strong bonds, and simply adds up their contributions weighted by the corresponding matrix elements to obtain the required result. This result is expected to be asymptotically accurate in the limit of small $`\omega `$, since these contributions clearly dominate in the low-frequency limit. A certain simplicity thus emerges when the low-energy effective theory has strong disorder, and we will exploit this to the fullest in what follows.
## III Dynamics and transport in the S=1/2 XXZ chains
### A Detailed characterization of the phases
#### 1 Singlet RG description of the Random Singlet states: A review
We begin by noting that the weak-randomness analysis of Doty and Fisher implies that randomness is relevant for pure antiferromagnetic XXZ spin-1/2 chains for $`0J^z/J^{}1`$; any amount of randomness is thus expected to drive the system to strong disorder in this entire regime.
In the strong-disorder regime, the singlet RG proceeds as follows: We look for the bond with the largest $`J^{}`$ coupling, say $`J_{23}^{}`$ between spins $`2`$ and $`3`$; this sets the energy cutoff $`\mathrm{\Omega }\mathrm{max}\{J_j^{}\}`$. We first solve the corresponding two-spin problem and introduce the neighboring bonds later as a perturbation. So long as the $`J^z`$ couplings are not large compared to the $`J^{}`$ couplings, the ground state of the two-spin problem will always be a singlet separated by a large gap from the triplet excited states. We can then trade our original Hamiltonian in for another Hamiltonian (determined perturbatively in the ratio of the neighboring bonds to the strongest bond) which acts on a truncated Hilbert space with the two sites connected by the ‘strong’ bond removed. To leading order, this procedure renormalizes the Hamiltonian $`_{4\mathrm{s}\mathrm{i}\mathrm{t}\mathrm{e}\mathrm{s}}=_{j=1}^3[J_j^{}(s_j^xs_{j+1}^x+s_j^ys_{j+1}^y)+J_j^zs_j^zs_{j+1}^z]`$ to $`\stackrel{~}{}_{14}=\stackrel{~}{J}_1^{}(s_1^xs_4^x+s_1^ys_4^y)+\stackrel{~}{J}_1^zs_1^zs_4^z`$ with $`\stackrel{~}{J}_1^{}=J_1^{}J_3^{}/(J_2^{}+J_2^z)`$ and $`\stackrel{~}{J}_1^z=J_1^zJ_3^z/2J_2^{}`$; note that the length of this new bond is $`\stackrel{~}{l}_1=l_1+l_2+l_3`$. This procedure, if it remains valid upon iteration, thus ultimately leads to the Random Singlet state described in the overview.
A complete understanding of the possible phases then requires an analysis of the effects of iterating the basic RG procedure. Such an analysis was performed in Ref. leading to the following conclusions (see Fig. 1): So long as the $`J^z`$ couplings do not dominate over the $`J^{}`$ couplings and therefore do not produce a state with Ising antiferromagnetic order, the ground state is a Random Singlet state. In this case, a detailed characterization of the low-energy effective Hamiltonian is best couched in terms of logarithmic variables as follows: Let $`\mathrm{\Omega }\mathrm{max}\{J_j^{}\}`$ at any given stage of the RG, and define the log-cutoff $`\mathrm{\Gamma }\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })`$. Also define log-couplings $`\zeta _j\mathrm{ln}(\mathrm{\Omega }/J_j^{})`$ and log-anisotropy parameters $`D_j\mathrm{ln}(\mathrm{\Delta }_j)`$, where $`\mathrm{\Delta }_jJ_j^z/J_j^{}`$. As $`\mathrm{\Gamma }`$ increases, the fraction of remaining sites $`n_\mathrm{\Gamma }`$ at log-cutoff scale $`\mathrm{\Gamma }`$ is given as $`n_\mathrm{\Gamma }1/\mathrm{\Gamma }^2`$. When the $`J^{}`$ couplings dominate, the system rapidly flows to the ‘XX-RS’ fixed point and the probability distribution $`P(\zeta ,\mathrm{\Delta },l|\mathrm{\Gamma })`$ that determines the strengths and lengths of the bonds connecting the remaining sites in the effective Hamiltonian quickly converges to the following scaling form characteristic of the XX-RS fixed point: $`P(\zeta ,\mathrm{\Delta },l|\mathrm{\Gamma })=\frac{1}{\mathrm{\Gamma }^3}𝒫_1(\frac{\zeta }{\mathrm{\Gamma }},\frac{l}{\mathrm{\Gamma }^2})\times \delta (\mathrm{\Delta })`$. The function $`𝒫_1`$ has been characterized in detail in Ref. ; here we only note that $`𝑑y𝒫_1(x,y)=e^x`$. Between the IAF phase and this XX-RS phase lie two kinds of critical points. If the initial problem has full Heisenberg symmetry ($`J^z=J^{}`$ for each bond), the low-energy effective Hamiltonian preserves this symmetry and has bond strengths and lengths drawn from the same probability distribution: $`P(\zeta ,l|\mathrm{\Gamma })=\frac{1}{\mathrm{\Gamma }^3}𝒫_1(\frac{\zeta }{\mathrm{\Gamma }},\frac{l}{\mathrm{\Gamma }^2})`$. In the RG language, the Heisenberg system is critical and is controlled by the ‘XXX-RS’ critical fixed point. Finally, in this language, the generic critical point between the IAF phase and the XX-RS phase is controlled by the ‘XXZC-RS’ fixed point—the low-energy effective theory has bond strengths and lengths drawn from a distribution $`P(\zeta ,D,l|\mathrm{\Gamma })=\frac{1}{\mathrm{\Gamma }^{3+\psi }}𝒫_2(\frac{\zeta }{\mathrm{\Gamma }},\frac{D}{\mathrm{\Gamma }^\psi },\frac{l}{\mathrm{\Gamma }^2})`$ with $`\psi <1`$ and $`𝑑y𝒫_2(x,y,z)=𝒫_1(x,z)`$. Notice that these scaling forms imply that the distributions of the couplings become infinitely broad as $`\mathrm{\Omega }0`$; thus, the RG becomes asymptotically exact at low-energies and, in particular, predicts the ground state properties and low-temperature thermodynamics correctly.
#### 2 Scaling Description of the Ising Antiferromagnet
On the Ising Antiferromagnet side, the singlet RG becomes invalid at low energies, and the system has a ground state with IAF order. The proper characterization of the system at these low energies is in terms of IAF-ordered spin clusters, as well as the domain-wall excitations that act to disrupt this order. This section is devoted to providing such a description. In what follows, we will be considering mainly the IAF phase close to the transition to the RS state. In this regime, the system will ‘look’ IAF ordered only well below a crossover energy $`\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$, while resembling a critical system controlled by the XXZC critical point above the crossover scale. $`\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$ is the scale at which the singlet RG breaks down and is thus determined by the properties of the RG flows in the vicinity of the XXZC critical point. The corresponding log-energy scale $`\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega }_{\delta _{\mathrm{IAF}}})`$ is given as $`\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}\delta _{\mathrm{IAF}}^\theta `$, with $`\theta =(2\psi )/\lambda `$, where $`\lambda `$ is the leading relevant RG eigenvalue at the XXZC fixed point and $`\psi `$ has already been defined in the previous section. Below, we construct a scaling description of the IAF phase near criticality by combining information obtained from the singlet RG about the nature of the system at this crossover scale, with a ‘cluster RG’ approach that is designed to work in the limit of low-energies (well below $`\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$) above the IAF ordered ground state.
We begin with a sketch of our cluster RG approach. Consider the Hamiltonian $`_{\mathrm{XXZ}}`$ with $`J^z`$ couplings completely dominating the $`J^{}`$ couplings. Now, spins tend to order antiferromagnetically, and we can try formulating a cluster RG similar to that for the ordered phase of the random transverse field Ising model. Consider combining two such spins, say $`s_2`$ and $`s_3`$, coupled by a strong bond $`J_2^z`$ into a new ‘superspin’ $`\stackrel{~}{s}_{(23)}`$. If we identify the two states $`|_{(23)}`$ and $`|_{(23)}`$ of this superspin with the states $`|_2_3`$ and $`|_2_3`$ (which is not a unique choice), and treat the $`J^{}`$ couplings to second order in perturbation theory, the effective Hamiltonian that we obtain is
$`\stackrel{~}{}_{1(23)4}`$ $`=`$ $`\stackrel{~}{J}_1^zs_1^z\stackrel{~}{s}_{(23)}^z\stackrel{~}{J}_3^zs_4^z\stackrel{~}{s}_{(23)}^z+\stackrel{~}{h}_{(23)}\stackrel{~}{s}_{(23)}^x`$
$`\stackrel{~}{J}_{1(23)4}^{}(s_1^+\stackrel{~}{s}_{(23)}^{}s_4^{}+s_1^{}\stackrel{~}{s}_{(23)}^+s_4^+),`$
where $`\stackrel{~}{h}_{(23)}=J_2^{}`$, $`\stackrel{~}{J}_{1(23)4}^{}=J_1^{}J_3^{}/J_2^z`$, $`\stackrel{~}{J}_1^z=J_1^z+(J_1^{})^2/J_2^z`$, and $`\stackrel{~}{J}_3^z=J_3^z+(J_3^{})^2/J_2^z`$. Thus, we see that new terms, not present in the original Hamiltonian, are generated: an effective transverse field, which acts to flip the new spin, and also a three-spin exchange interaction. Before we proceed, a couple of comments regarding the new terms: The effective transverse field appears because the ground state of $`_{23}`$ is not exactly a degenerate doublet (the two lowest eigenstates, which are the symmetric and antisymmetric combinations of $`|_2_3`$ and $`|_2_3`$, are actually split by a small energy $`J_2^{}`$). Note also that the three-spin term does not violate spin conservation—for example, if we consider coupling the conserved total $`s_{\mathrm{tot}}^z`$ to a magnetic field, we immediately realize that the superspin $`\stackrel{~}{s}_{(23)}`$ does not couple to this field.
In principle, we may proceed with such a clustering process, keeping track of all additional one- or two- or multi-spin–flip terms that are generated. While this RG is not analytically tractable, we do not expect the generated terms to have any drastic consequences, since they generally become weaker and weaker, while the $`J^z`$ couplings remain almost unchanged. Alternatively, we can remedy this proliferation of new couplings by combining an odd number of spins at a time—because of the symmetries of the Hamiltonian, any odd length chain will have a degenerate pair of ground states with the total $`s_{\mathrm{tot}}^z=\pm \frac{1}{2}`$. In addition, three-spin terms of the form encountered previously will now be forbidden by spin conservation. More explicitly, if we combine three spins, say $`s_2`$, $`s_3`$, and $`s_4`$, with relatively strong couplings $`J_2^z`$ and $`J_3^z`$, into a new superspin $`\stackrel{~}{s}_2\stackrel{~}{s}_{(234)}`$, and treat the $`J^{}`$ couplings perturbatively, the XXZ form of the effective Hamiltonian is preserved, with the new couplings $`\stackrel{~}{J}_{12}^z=J_1^z+(J_1^{})^2/(2J_2^z)`$, $`\stackrel{~}{J}_{25}^z=J_4^z+(J_4^{})^2/(2J_3^z)`$, $`\stackrel{~}{J}_{12}^{}=2J_1^{}J_3^{}/J_2^z`$, and $`\stackrel{~}{J}_{25}^{}=2J_2^{}J_4^{}/J_3^z`$.
Either way, we will have effective spin-half objects with dominant Ising AF interactions. Almost always, we will be decimating strong $`J^z`$ couplings making larger and larger clusters, with the other $`J^z`$ couplings remaining essentially unchanged, and the remaining $`J^{}`$ couplings growing weaker and weaker. Only rarely will there be a bond with the $`J^{}`$ coupling large compared to the neighboring couplings, and this will then produce a singlet. Thus, the picture that emerges is very reminiscent of the ordered phase in the RTFIM.
We may now combine this schematic cluster RG description valid at low energies, with information about the crossover region obtainable from the singlet RG. At the crossover scale, the distribution of $`\zeta ^z\mathrm{ln}(\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}/J^z)`$ is given as $`P(\zeta ^z|\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}})\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^1\mathrm{exp}(\zeta ^z/\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}})`$. Roughly speaking, beyond the crossover scale, the cluster RG merely eliminates the strongest bonds from this distribution, but keeps the low-energy tail of the distribution unchanged. We thus expect a line of (classical) IAF fixed points, with properties varying smoothly with the distance from the criticality. The density of spin degrees of freedom $`n_\mathrm{\Gamma }`$ in the renormalized theory is expected to decrease as $`n_\mathrm{\Gamma }\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^2e^{c\mathrm{\Gamma }/\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}}`$ below the crossover scale $`\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$, with $`c`$ some order one constant. This immediately gives us the density of states $`\rho (\omega )\omega ^1n_{\mathrm{\Gamma }_\omega }\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^1\delta _{\mathrm{IAF}}^{3\theta }\omega ^{1+1/z_{\mathrm{IAF}}}`$, with the continuously varying dynamical exponent $`z_{\mathrm{IAF}}\delta _{\mathrm{IAF}}^\theta `$. The typical size of the excitations dominating the density of states scales as $`l_{\mathrm{dw}}(\omega )l_v\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}\mathrm{\Gamma }_\omega `$ and is much smaller than their typical separation $`\omega ^{1/z_{\mathrm{IAF}}}`$. This can be readily seen from the qualitative picture of ‘preformed tails’: the length $`l_{\mathrm{dw}}(\omega )`$ of a renormalized bond with $`\zeta ^z=0`$ in the theory with cutoff $`\omega \mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$ scales in the same way as the length of a bond with $`\zeta ^z\mathrm{ln}(\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}/\omega )`$ in the theory at the crossover scale $`\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$. On the other hand, the distribution of the log-couplings $`\zeta ^{}\mathrm{ln}(\mathrm{\Omega }/J^{})`$ is expected to broaden exponentially as a function of $`\mathrm{\Gamma }`$: for example, when we combine $`n`$ spins that are active at the crossover scale into a cluster, the effective transverse coupling acting on this cluster will be of order $`\zeta ^{}n\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}`$. This then is our scaling picture for the IAF phase; the important conclusion that emerges from this analysis is the fact that the transition to the RS state is preceded by a Griffiths phase—the IAF Griffiths phase—with a continuously varying power-law singularity in the low-energy density of states.
Finally, it is also possible to obtain a rather direct identification of the low-energy modes in terms of the rare regions that dominate the low-energy dynamics; we conclude by sketching this briefly here. The RG picture suggests that in the IAF phase the typical excitations at low energies $`\omega `$ are classical domain wall excitations that live on the effective bonds with weak effective couplings $`\stackrel{~}{J}^z\omega `$. Such a weak effective $`\stackrel{~}{J}^z`$ can appear only across a long region that is locally in the RS phase. More quantitatively, a region of length $`L`$ locally in the RS phase effectively corresponds to a weak bond with $`\stackrel{~}{J}^z\mathrm{\Omega }_0e^{c_zL}`$. The number density of such regions is roughly $`p^L`$ for some $`p<1`$. The density of such regions with $`\stackrel{~}{J}^z\omega `$ is thus some power of $`\omega `$ that we choose to write as $`n(\omega )\omega ^{1/z}`$ with some exponent $`z`$; the most numerous such regions will have some “optimal” (for a given $`\delta _{\mathrm{IAF}}`$) microscopic structure, but whatever this structure is, the corresponding optimal exponent $`z`$ can be directly identified with the dynamical exponent $`z(\delta _{\mathrm{IAF}})`$ of this phase. This picture thus predicts that the typical separation of such regions is of order $`\omega ^{1/z}`$, while their lengths are only of order $`|\mathrm{ln}\omega |`$, in complete agreement with the schematic RG approach.
#### 3 Singlet RG description of the Random Dimer phases: A review
While the effects of dimerization are not understood in detail in all regimes, it is possible to use the singlet RG and follow the flows for a chain with full Heisenberg symmetry and for a chain in the vicinity of the XX-RS point. In these cases, a mapping to the off-critical flows of the RTFIM provides a detailed characterization of the so-called ‘Random Dimer’ (RD) phases that result. In either case, the picture that emerges can be summarized as follows: For concreteness, assume $`\delta >0`$. If disorder is strong and $`\delta 1`$, then the even and the odd bonds renormalize essentially as in the corresponding RS state till the log-energy scale $`\mathrm{\Gamma }\mathrm{\Gamma }_\delta 1/\delta `$. Beyond this scale, the remaining odd bonds rapidly become much weaker relative to the remaining even bonds; the distribution of the even log-couplings $`P^\mathrm{e}(\zeta |\mathrm{\Gamma })=𝑑lP^\mathrm{e}(\zeta ,l|\mathrm{\Gamma })`$ approaches some limiting distribution with a finite but large width, while distribution of the odd log-couplings $`P^\mathrm{o}(\zeta |\mathrm{\Gamma })=𝑑lP^\mathrm{o}(\zeta ,l|\mathrm{\Gamma })`$ grows infinitely broad. In the RG language, the system renormalizes to some point on a line of RD fixed points (from this point of view, the RS states at $`\delta =0`$ represent critical points separating RD fixed points with opposite dimerization—see Fig. 2). The corresponding joint distributions of the log-couplings and the lengths have been worked out in Ref. ; here we only note that $`P^\mathrm{e}(\zeta |\mathrm{\Gamma })=\tau _0(\mathrm{\Gamma })e^{\zeta \tau _0(\mathrm{\Gamma })}`$ with $`\tau _0(\mathrm{\Gamma })2\delta `$, while $`P^\mathrm{o}(\zeta |\mathrm{\Gamma })=u_0(\mathrm{\Gamma })e^{\zeta u_0(\mathrm{\Gamma })}`$ with $`u_0(\mathrm{\Gamma })2\delta e^{2\delta \mathrm{\Gamma }}`$. The ground state again consists of singlet pairs made up of one spin on an even site $`i`$ and a second spin on some odd site $`j`$. Note however, that while $`i>j`$ and $`i<j`$ are equally probable in the RS state, in the RD phase with $`\delta >0`$ one almost always has $`j>i`$ (with the exception of a few high-energy pairs of small spatial extent).
### B Dynamic structure factor
In this section, we summarize our calculations of the dynamic structure factor in the different regions of the phase diagram of spin-1/2 XXZ chains. Our approach has already been reviewed in general terms in Sec. II D, and our calculations here represent one of the simplest examples of this approach at work.
We begin by considering only the leading term in the perturbative expansion for the renormalized spin operators; the results obtained in this manner give the correct leading behavior at low frequencies (some justification of this is given in Sec. III E, where we discuss the role of higher-order terms).
#### 1 Random Singlet states
The leading-order ‘operator renormalizations’ needed are particularly simple: the spin operator $`\stackrel{}{s}`$ remains unchanged for each of the ‘surviving’ spins and is effectively zero for each of the ‘decimated’ spins (i.e., spins that are already locked into singlets with other spins).
Consider first $`S^{zz}(k,\omega )`$; in our formulation of the singlet RG, Sec. III A, the following analysis applies to a general XXZ singlet state (i.e., remains valid so long as the ground state does not have IAF order). Consider two spins $`L`$ and $`R`$ connected by a strong bond $`(\stackrel{~}{J}^{},\stackrel{~}{J}^z)`$ in the renormalized theory with cutoff $`\mathrm{\Omega }_{\mathrm{final}}`$. The spin operators $`s_{L/R}^z`$ connect the singlet ground state of this pair only to the triplet state $`|t_0`$ (with $`m_z=0`$), which is separated from the singlet state by a gap $`\stackrel{~}{J}^{}`$. Therefore, the energy scale $`\mathrm{\Omega }_{\mathrm{final}}`$ at which we stop the RG is $`\mathrm{\Omega }_{\mathrm{final}}=\omega `$ in this case (remember that the cutoff was defined as $`\mathrm{\Omega }=\mathrm{max}\{J^{}\}`$). We thus consider the renormalized spectral sum
$$S^{zz}(k,\omega )=\frac{1}{L}\stackrel{~}{\underset{m}{}}|m|\stackrel{~}{\underset{j}{}}e^{ikx_j}\stackrel{~}{s}_j^z|0|^2\delta (\omega \stackrel{~}{E}_m),$$
(18)
where the tildes remind us of the fact that this spectral sum now refers to the new Hamiltonian with energy cutoff $`\mathrm{\Omega }_{\mathrm{final}}=\omega `$; this renormalized Hamiltonian has $`n_{\mathrm{\Gamma }_\omega }`$ spins per unit length with the distribution of couplings and bond lengths characteristic of the fixed point to which the system flows in the low-energy limit. The sum Eq. (18) is dominated by the excitations to the triplet state $`|t_0`$ of pairs of spins connected by the (renormalized) bonds with $`\stackrel{~}{J}^{}=\omega `$; these pairs are precisely the ones that are being eliminated at this energy scale. The corresponding matrix element for each such pair is simply $`(1e^{ik\stackrel{~}{l}})/2`$, where $`\stackrel{~}{l}`$ is the length of the bond connecting the pair; this allows us to write
$$S^{zz}(k,\omega )n(\mathrm{\Gamma }_\omega )𝑑l𝑑\zeta |1e^{ikl}|^2P(\zeta ,l|\mathrm{\Gamma }_\omega )\delta (\omega \omega e^\zeta )$$
(19)
for $`\omega \mathrm{\Omega }_0`$ in any RS state.
The calculation of $`S^+(k,\omega )`$ is more involved since the gap to the relevant triplet excited state $`|t_1`$ (with $`m_z=1`$) of a pair of spins connected by a strong bond $`(\stackrel{~}{J}^{},\stackrel{~}{J}^z)`$ is now $`(\stackrel{~}{J}^{}+\stackrel{~}{J}^z)/2`$. We consider each of the three cases (XX, XXX, and XXZC) separately: (1) In the XXX case, the Heisenberg symmetry of the problem guarantees that $`S^{xx}=S^{yy}=S^{zz}`$. (2) When the system approaches the XX point at low energies, we have $`\stackrel{~}{J}^z\stackrel{~}{J}^{}`$ implying that the relevant gap is approximately $`\stackrel{~}{J}^{}/2`$. Thus, to calculate $`S^+(k,\omega )`$ we now have to stop the RG at the scale $`\mathrm{\Omega }_{\mathrm{final}}=2\omega `$. ¿From the Eq. (19), it is clear that this leaves our answer unchanged except for the values of various non-universal scale factors. (3) The XXZC critical point needs special attention. In this case $`\stackrel{~}{J}^z/\stackrel{~}{J}^{}`$ can have a range of values. As a result, the excited states that dominate the spectral sum Eq. (4) are not simply obtained by stopping the RG at any particular $`\mathrm{\Omega }_{\mathrm{final}}`$ and looking at the singlets forming only at this scale. Instead, for any $`\mathrm{\Omega }_{\mathrm{final}}(0,2\omega )`$ there will be some singlets formed at this scale that will contribute to the spectral sum—namely, the pairs coupled by strong bonds with $`\stackrel{~}{J}^{}=\mathrm{\Omega }_{\mathrm{final}}`$ and $`\stackrel{~}{J}^z=2\omega \mathrm{\Omega }_{\mathrm{final}}`$. Note that there is no double-counting here since we are considering only the pairs that are being eliminated at each energy scale. Thus, we have
$`S^+(k,\omega )`$ $``$ $`{\displaystyle }d\mathrm{\Gamma }dldDn(\mathrm{\Gamma })|1e^{ikl}|^2P(0,D,l|\mathrm{\Gamma })\times `$ (21)
$`\times \delta (\omega \mathrm{\Omega }_0e^\mathrm{\Gamma }(1+e^D)/2).`$
Rewriting this in terms of the scaling probability distribution $`𝒫_2`$ and using the delta function to do the $`\mathrm{\Gamma }`$ integral gives us
$`S^+(k,\omega )`$ $``$ $`{\displaystyle \frac{1}{\omega \mathrm{\Gamma }_\omega }}{\displaystyle }d\overline{l}d\overline{D}{\displaystyle \frac{n\left(\mathrm{\Gamma }_\omega \mathrm{{\rm Y}}_\omega (\overline{D})\right)}{\mathrm{{\rm Y}}_\omega ^{3+\psi }(\overline{D})}}|1e^{i\overline{k}\overline{l}}|^2\times `$ (23)
$`\times 𝒫_2(0,{\displaystyle \frac{\overline{D}}{\mathrm{{\rm Y}}_\omega ^\psi (\overline{D})}},{\displaystyle \frac{\overline{l}}{\mathrm{{\rm Y}}_\omega ^2(\overline{D})}}),`$
where we have defined $`\overline{D}=D/\mathrm{\Gamma }_\omega ^\psi `$, $`\overline{l}=l/\mathrm{\Gamma }_\omega ^2`$, $`\overline{k}=k\mathrm{\Gamma }_\omega ^2`$, and
$$\mathrm{{\rm Y}}_\omega (\overline{D})=1+\frac{\mathrm{ln}(1+e^{\overline{D}\mathrm{\Gamma }_\omega ^\psi })\mathrm{ln}2}{\mathrm{\Gamma }_\omega }.$$
(24)
Now, since $`\psi <1`$, it is permissible to take the $`\mathrm{\Gamma }_\omega \mathrm{}`$ limit of $`\mathrm{{\rm Y}}_\omega (\overline{D})`$ before doing the $`\overline{D}`$ integral—in other words, we can replace $`\mathrm{{\rm Y}}`$ by $`1`$ in the low-energy limit. The $`\overline{D}`$ integral can then be done trivially, and the final expression is identical in form to Eq. (19). More physically, a given bond $`(J^{},J^z)`$ is described fairly well (on a logarithmic scale) by one of these two couplings—we chose the characteristic scale to be $`J^{}`$. Now, the random anisotropy leads to an uncertainty $`|\mathrm{ln}(J^z/J^{})|\mathrm{\Gamma }^\psi `$ in the corresponding log-energy scale. This uncertainty is much smaller than the already existing typical spread in the log-energies or the typical log-energies themselves, which are both of order $`\mathrm{\Gamma }`$. The leading behavior at low frequencies is therefore not affected.
Thus, in the limit of low frequencies both $`S^{zz}(k,\omega )`$ and $`S^+(k,\omega )`$ can be expressed in terms of the scaling probability distribution $`𝒫_1`$ as
$$S(k,\omega )\frac{n(\mathrm{\Gamma }_\omega )}{\omega \mathrm{\Gamma }_\omega }𝑑\overline{l}|1e^{i\overline{k}\overline{l}}|^2𝒫_1(0,\overline{l}).$$
(25)
Let us first focus on the regime $`|q||k\pi /a|a^1`$. Note that all unscaled lengths $`l`$ are odd multiples of the unit length $`a`$, and therefore $`e^{ikl}=e^{iql}`$. The integral in Eq. (25) can be evaluated using the characterization of the function $`𝒫_1(0,y)`$ available in Ref. ; the result is the following rather unusual scaling form at all three RS fixed points:
$$S(k=\frac{\pi }{a}+q,\omega )=\frac{𝒜}{l_v\omega \mathrm{ln}^3(\mathrm{\Omega }_0/\omega )}\mathrm{\Phi }\left(|ql_v|^{1/2}\mathrm{ln}(\mathrm{\Omega }_0/\omega )\right).$$
(26)
Note that we have suppressed the component labels on $`S(k,\omega )`$ as the two independent components obey the same scaling form, but with different values in general of the numerical constant $`𝒜`$ and the microscopic length scale $`l_v`$. The universal function $`\mathrm{\Phi }(x)`$ can be written as
$$\mathrm{\Phi }(x)=1+x\frac{\mathrm{cos}(x)\mathrm{sinh}(x)+\mathrm{sin}(x)\mathrm{cosh}(x)}{\mathrm{cos}^2(x)\mathrm{sinh}^2(x)+\mathrm{sin}^2(x)\mathrm{cosh}^2(x)}.$$
(27)
The resulting $`S(k,\omega )`$ is shown on Fig. 3. There is a fairly straightforward interpretation of the main features of this lineshape: The peak at $`q=0`$ (i.e., at $`k=\pi /a`$) reflects the predominantly antiferromagnetic character of the low-energy fluctuations; in our language, this is a direct consequence of the fact that the (renormalized) bonds all have odd lengths in units of $`a`$. The strongly damped oscillations with the period and the decay scale both of order $`\mathrm{\Gamma }_\omega ^2`$ express the properties of the distribution of lengths of the strong bonds: both the average and the RMS fluctuation of this distribution of lengths are of order $`\mathrm{\Gamma }_\omega ^2`$.
While this result is interesting, one needs to analyze the effects of higher-order terms in the operator renormalizations before accepting its consequences for possible neutron scattering experiments. We will argue in Sec. III E that higher-order corrections do not modify the functional form Eq. (27) of the features in $`S(k,\omega )`$ at fixed $`\omega `$ but only add an “incoherent” background (of strength comparable to that of the features) and suppress the amplitude of the features by a non-universal multiplicative factor of order one.
A similar scaling function can be derived for the regime $`|k|a^1`$. Repeating the above analysis gives
$$S(k,\omega )=\frac{𝒜^{}}{l_v\omega \mathrm{ln}^3(\mathrm{\Omega }_0/\omega )}\stackrel{~}{\mathrm{\Phi }}\left(|kl_v|^{1/2}\mathrm{ln}(\mathrm{\Omega }_0/\omega )\right),$$
(28)
with $`\stackrel{~}{\mathrm{\Phi }}(x)=2\mathrm{\Phi }(x)`$, and $`𝒜^{}`$ an order one numerical constant. This scaling function vanishes for $`k0`$; for small $`k`$ we have $`S(k,\omega )l_vk^2\mathrm{ln}(\mathrm{\Omega }_0/\omega )/\omega `$. We must therefore consider the possibility that higher-order corrections may overwhelm this scaling result and render it irrelevant. This is indeed expected to happen for $`S^+(k,\omega )`$ away from the XXX point. However, we expect the scaling result to be valid quite generally for $`S^{zz}(k,\omega )`$—spin conservation guarantees that the higher-order corrections to $`S^{zz}(k,\omega )`$ must also vanish as $`k0`$ (see Sec. III E for a detailed discussion of this point).
#### 2 Dynamic structure factor in Random Dimer phases: ‘sharpness’ of the Griffiths regions
Next, we consider spin dynamic structure factor in the XX and XXX Random Dimer phases introduced in Sec. III A 3. The same approach as for the RS states goes over unchanged, and we write
$`S(k,\omega )`$ $``$ $`n(\mathrm{\Gamma }_\omega ){\displaystyle }dld\zeta |1e^{ikl}|^2\times `$ (30)
$`\times [P^\mathrm{o}(\zeta ,l|\mathrm{\Gamma }_\omega )+P^\mathrm{e}(\zeta ,l|\mathrm{\Gamma }_\omega )]\delta (\omega \omega e^\zeta )`$
for both $`S^{zz}(k,\omega )`$ and $`S^+(k,\omega )`$ in both the XX and XXX RD phases (we are again being sloppy about the distinction between $`\mathrm{\Gamma }_\omega `$ and $`\mathrm{\Gamma }_{\omega /2}`$, as this can be absorbed in the definition of the non-universal scale factors that enter our expressions).
Using the results of Ref. , it is a simple matter to obtain the full crossover from the RS-like behavior of the structure factor in the regime $`1\mathrm{\Gamma }_\omega \mathrm{\Gamma }_\delta `$ to the behavior characteristic of the RD phase in the regime $`\mathrm{\Gamma }_\omega \mathrm{\Gamma }_\delta `$. Here, we focus on the behavior in the regime $`\mathrm{\Gamma }_\omega \mathrm{\Gamma }_\delta `$, as this exhibits some rather unusual features. At these low energies, the even bonds dominate over the odd bonds, and the contribution of the odd bonds to the sum Eq. (30) is negligible (we are assuming $`\delta >0`$ for concreteness). For wavevectors in the vicinity of $`k=\pi /a`$ with $`|q||k\pi /a|\delta ^2/l_v`$ (i.e., probing lengths larger than the correlation length $`\xi _{\mathrm{av}}l_v/\delta ^2`$) we obtain
$$S(k,\omega )=\frac{𝒞|\delta |^3\mathrm{\Omega }_0^{1/z_{\mathrm{RD}}}}{l_v\omega ^{11/z_{\mathrm{RD}}}}\left[1+\mathrm{cos}(l_vq\mathrm{\Gamma }_\omega /|\delta |)e^{cl_v^2q^2\mathrm{\Gamma }_\omega /|\delta |^3}\right],$$
(31)
where $`𝒞`$ and $`c`$ are some order one constants, and we have chosen to write the power-law prefactor in terms of the dynamical exponent $`z_{\mathrm{RD}}`$ (as far as our RG calculations are concerned $`z_{\mathrm{RD}}^1=2|\delta |`$ for small $`|\delta |`$—however, the effective value of $`\delta `$ that enters this expression is expected to acquire a non-universal multiplicative renormalization from the high-energy physics, and the only reliable statement we can make is that $`z_{\mathrm{RD}}^1|\delta |`$ for small enough $`|\delta |`$). This result has a striking oscillatory structure (see Fig. 4) that is not suppressed significantly by the exponential factor, since $`\sqrt{\mathrm{\Gamma }_\omega /|\delta |^3}\mathrm{\Gamma }_\omega /|\delta |`$ in the regime under consideration. This is best understood as a novel signature of the sharply defined geometry of the rare Griffiths regions that contribute to the scattering at a given low energy (i.e., that are filtered out by their energy). More precisely, the average length of such regions is of order $`l_v\mathrm{\Gamma }_\delta ^2(\mathrm{\Gamma }_\omega /\mathrm{\Gamma }_\delta )=l_v\mathrm{\Gamma }_\omega /|\delta |`$, while the RMS fluctuations in the length are only of order $`l_v\mathrm{\Gamma }_\delta ^2\sqrt{\mathrm{\Gamma }_\omega /\mathrm{\Gamma }_\delta }=l_v\sqrt{\mathrm{\Gamma }_\omega /|\delta |^3}`$. Our results thus suggest that low-energy INS experiments would be able to pick up the sharply defined geometry of such Griffiths regions in the RD phases in one dimension.
This feature of the Griffiths regions in one dimension was noted in Ref. Sec. IVB in the context of the RTFIM, where it was conjectured also that other properties of such low-energy regions are likewise sharply defined: for example, in the disordered phase of the RTFIM, the magnetic moment of the Griffiths regions with a given characteristic energy is sharply defined and proportional to the (sharply defined) length of such regions. \[In fact, similar “sharpness” is expected to hold for any bond “property” that “rides” on top of the singlet RG via recursion relation $`\stackrel{~}{x}=x_1+x_3+\mathrm{{\rm Y}}x_2`$ when bond $`J_2`$ is eliminated.\] We expect to see a signature of this sharpness of the Griffiths regions also in the dynamic structure factor $`S^{zz}`$ in the IAF Griffiths phase (see below) and also in the Griffiths phases of the one-dimensional RTFIM (Sec. V C). Finally, an interesting question, which we leave unanswered for now, is whether similar sharpness in the properties of the Griffiths regions at a given energy occurs and has observable consequences in higher dimensions as well, e.g., in the disordered phase of the $`d>1`$ RTFIM.
#### 3 IAF Griffiths phase
Let us first consider $`S^{zz}(k,\omega )`$ in the IAF Griffiths phase. As discussed in Sec. III A 2, the dominant low-energy excitations in this phase are classical domain walls. However, it is clear that such excitations do not contribute at all to $`S^{zz}(k,\omega )`$, since they cannot be excited from the ground state by the action of the operators like $`\widehat{s}_k^z`$, which conserve the total $`s_{\mathrm{tot}}^z`$. The leading excitations that do contribute to $`S^{zz}`$ can clearly be identified in the RG picture with the $`m_z=0`$ excited states of pairs of super-spins, with each pair connected by a bond with $`\stackrel{~}{J}^{}\omega `$ and forming a singlet (note that this is true regardless of the value of the corresponding $`\stackrel{~}{J}^z`$). Now, it is easy to generate a weak $`\stackrel{~}{J}^{}`$ coupling of order $`\omega `$ in the IAF phase, since any typical region of length $`L`$ will have an effective $`\stackrel{~}{J}^{}`$ of order $`\mathrm{\Omega }_0e^{c_xL}`$ (and an effective $`\stackrel{~}{J}^z`$ typically much stronger). What is more difficult is to isolate such a region from becoming a part of a larger cluster—otherwise this region can not support spin fluctuations at frequency $`\omega `$. For this, we need two rare RS-like segments (domain walls) with $`\stackrel{~}{J}^z\omega `$, one on each side of our (typical) region. Thus, we need two domain walls, which are usually separated by a large distance of order $`\omega ^{1/z}`$, to occur close to each other—the “density” of such occurrences is $`\omega ^{2/z_{\mathrm{IAF}}}`$. The separation of the two domain walls—the length of the IAF-ordered cluster that they isolate—must be of order $`|\mathrm{ln}\omega |`$. More precisely, if the IAF-ordered cluster has length $`L`$, it can be thought of as consisting of the $`nL/\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^2`$ strongly Ising coupled spins that are active at the crossover scale; the effective bonds connecting these spins at the crossover scale typically satisfy $`\mathrm{ln}(\stackrel{~}{J}^{}/\stackrel{~}{J}^z)\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}`$. The requirement that the spin-flip coupling for this cluster is $`\omega `$ fixes the length of this cluster to be $`L=l_v\mathrm{\Gamma }_\omega \mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}`$, while the uncertainty in this length can only be of order $`l_v\sqrt{\mathrm{\Gamma }_\omega \mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^3}l_v\mathrm{\Gamma }_\omega \mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}`$.
We are now ready to calculate $`S^{zz}(k=\pi /a+q,\omega )`$ in the regime $`|q|^1l_v\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^2`$, in addition to $`\omega \mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$. The leading-order renormalization of the $`s_j^z`$ in the cluster RG is simple: $`s_j^z`$ is renormalized to $`(1)^js_c^z`$ for each spin $`j`$ that is active in some cluster $`c`$, and renormalizes to zero for every spin that forms a singlet. Assuming that such clusters “look” fairly uniform on the length scales larger than $`l_v\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^2`$, and adding up the contributions from all such isolated clusters with effective spin fluctuation frequency $`\omega `$, we obtain
$`S^{zz}(k`$ $`=`$ $`{\displaystyle \frac{\pi }{a}}+q,\omega )={\displaystyle \frac{𝒞^{}|\delta _{\mathrm{IAF}}|^{7\theta }\mathrm{\Omega }_0^{2/z_{\mathrm{IAF}}}}{q^2l_v^3\omega ^{12/z_{\mathrm{IAF}}}}}`$ (32)
$`\times `$ $`\left[1\mathrm{cos}(ql_v\mathrm{\Gamma }_\omega /|\delta _{\mathrm{IAF}}|^\theta )e^{cq^2l_v^2\mathrm{\Gamma }_\omega /|\delta _{\mathrm{IAF}}|^{3\theta }}\right],`$ (33)
where $`𝒞^{}`$ and $`c`$ are some order one constants and the power of the $`\delta _{\mathrm{IAF}}`$ that appears in the prefactor has been fixed by demanding consistency with the off-critical scaling form
$$S^{zz}(k=\frac{\pi }{a}+q,\omega )=\frac{𝒜}{l_v\omega \mathrm{\Gamma }_\omega ^3}\mathrm{\Psi }(\frac{\mathrm{\Gamma }_\omega }{\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}},|ql_v|^{1/2}\mathrm{\Gamma }_\omega ),$$
(34)
with $`\mathrm{\Psi }(0,y)=\mathrm{\Phi }(y)`$. Note also that the overall $`1/q^2`$ dependence is a consequence of the fact that the spins contributing to the scattering have been taken to be distributed uniformly over a sharply defined region (the cluster); we expect this to cross over to a much faster decay at large momenta (such that $`|q|^1l_v\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}^2`$) well outside the range of validity of our scaling picture.
The situation is quite different for $`S^+(k,\omega )`$. As we shall see in Section III C, the renormalization of the $`s_j^\pm `$ spin operators is quite non-trivial, and we are unable to make an equally detailed prediction for $`S^+`$. However, we expect that the matrix element for producing domain wall excitations with energies of order $`\omega `$ by the action of $`s^\pm `$ on the ground state is strongly suppressed as some power of $`\omega `$, giving rise to a correspondingly small value for $`S^+(k,\omega )`$ at small $`\omega `$.
### C Average local autocorrelations
The same approach can be used to calculate average autocorrelation functions, and this section is devoted to a brief account of our results.
We consider the local dynamical susceptibilities
$$\chi _{jj}^{\alpha \alpha }(\omega )=\underset{m}{}|m|s_j^\alpha |0|^2\delta (\omega E_m),$$
(35)
where $`\alpha =z`$ or $`\alpha =x`$. A knowledge of the low-frequency behavior of these susceptibilities can immediately be translated into information about the long-time limit of the corresponding imaginary-time autocorrelation functions
$$C_{jj}^{\alpha \alpha }(\tau )=s_j^\alpha (\tau )s_j^\alpha (0).$$
(36)
#### 1 RS states and RD phases
As long as one is only interested in averages of such local quantities (over different realizations of disorder), it again suffices to consider only the leading-order spin operator renormalizations. We thus already have all the ingredients needed to calculate these average dynamical susceptibilities: our basic approach is familiar enough by now, and the relevant results of Ref. for the renormalized bond distributions have already been reviewed in Sec. III A. Below, we will be correspondingly brief. We first give our results for the average local dynamical susceptibilities and then translate these to results for the long-time behavior of the corresponding average autocorrelation functions. The leading behavior is the same for both $`\alpha =z`$ and $`\alpha =x`$, so we drop all superscripts.
For a bulk-spin, we obtain
$$\left[\chi _{\mathrm{loc}}\right]_{\mathrm{av}}(\omega )\frac{n(\mathrm{\Gamma }_\omega )}{\omega }(P_0^\mathrm{e}(\mathrm{\Gamma }_\omega )+P_0^\mathrm{o}(\mathrm{\Gamma }_\omega )).$$
(37)
For the critical RS states ($`P^\mathrm{e}=P^\mathrm{o}`$) we find
$`\left[\chi _{\mathrm{loc}}\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{1}{\omega |\mathrm{ln}\omega |^3}},`$ (38)
$`\left[C_{\mathrm{loc}}\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{1}{|\mathrm{ln}\tau |^2}},`$ (39)
while off-critical—in the RD phases—we find
$`\left[\chi _{\mathrm{loc}}\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{|\delta |^3}{\omega ^{11/z_{\mathrm{RD}}}}},`$ (40)
$`\left[C_{\mathrm{loc}}\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{|\delta |^3}{\tau ^{1/z_{\mathrm{RD}}}}}.`$ (41)
Similarly, for an end-spin $`s_1`$ of a semi-infinite chain (with $`j1`$) we obtain
$$\left[\chi _1\right]_{\mathrm{av}}(\omega )\frac{P_0^\mathrm{e}(\mathrm{\Gamma }_\omega )P_0^\mathrm{o}(\mathrm{\Gamma }_\omega )}{\omega }.$$
(42)
For the RS states we find
$`\left[\chi _1\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{1}{\omega |\mathrm{ln}\omega |^2}},`$ (43)
$`\left[C_1\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{1}{|\mathrm{ln}\tau |}},`$ (44)
and in the RD phases
$`\left[\chi _1\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{\delta ^2}{\omega ^{11/z_{\mathrm{RD}}}}},`$ (45)
$`\left[C_1\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{\delta ^2}{\tau ^{1/z_{\mathrm{RD}}}}}.`$ (46)
#### 2 IAF Griffiths phase
In the IAF phase, unlike in the singlet states, we need to make a distinction between $`\chi ^{zz}`$ and $`\chi ^{xx}`$. Consider first $`\left[\chi _{\mathrm{loc}}^{zz}\right]_{\mathrm{av}}(\omega )`$. ¿From our previous discussion of the IAF phase, it is clear that, in the regime $`\omega \mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$, the dominant contributions come from IAF-ordered clusters of lengths $`\mathrm{\Gamma }_\omega \mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}`$ (i.e., with effective spin-flip couplings of order $`\omega `$) that are isolated from the rest of the system by domain walls with $`\stackrel{~}{J}^z\omega `$ on either side. ¿From the scaling picture of the phase, we get
$`\left[\chi _{\mathrm{loc}}^{zz}\right]_{\mathrm{av}}(\omega )`$ $``$ $`\delta _{\mathrm{IAF}}^{4\theta }{\displaystyle \frac{\mathrm{\Omega }_0^{2/z_{\mathrm{IAF}}}\mathrm{ln}(\mathrm{\Omega }_0/\omega )}{\omega ^{12/z_{\mathrm{IAF}}}}},`$ (47)
$`\left[C_{\mathrm{loc}}^{zz}\right]_{\mathrm{av}}(\tau )`$ $``$ $`\delta _{\mathrm{IAF}}^{4\theta }{\displaystyle \frac{\mathrm{ln}(\mathrm{\Omega }_0\tau )}{(\mathrm{\Omega }_0\tau )^{2/z_{\mathrm{IAF}}}}}.`$ (48)
The analysis is more complicated for $`\left[\chi _{\mathrm{loc}}^{xx}\right]_{\mathrm{av}}(\omega )`$, and we can only make a plausible estimate for this quantity. This is because the $`x`$ and $`y`$ components of the spin operators renormalize in a non-trivial way under the cluster RG. The origin of this difficulty may be seen as follows: Consider, for example, combining three spins $`s_2`$, $`s_3`$, and $`s_4`$, connected by strong $`J_2^z`$ and $`J_3^z`$, into a super-spin $`\stackrel{~}{s}_{(234)}`$. To zeroth order, all three operators $`s_2^+`$, $`s_3^+`$, and $`s_4^+`$, renormalize to zero. To first order, $`s_2^+`$ and $`s_4^+`$ renormalize to $`(s_2^+)^{\mathrm{eff}}=s_1^+J_1^{}/J_2^z\stackrel{~}{s}_{(234)}^+2J_3^{}/J_2^z`$, $`(s_4^+)^{\mathrm{eff}}=s_5^+J_4^{}/J_3^z\stackrel{~}{s}_{(234)}^+2J_2^{}/J_4^z`$, while $`s_3^+`$ renormalizes to zero to this order. Roughly speaking, the original spin-flip operators of the (active) spins have projections onto the remaining effective cluster spin-flip operators with components given by the ratio of the corresponding effective spin-flip couplings to the original spin-flip couplings.
Now, the dominant contributions to $`\left[\chi _{\mathrm{loc}}^{xx}\right]_{\mathrm{av}}(\omega )`$ come from the low-energy (of order $`\omega `$) domain wall excitations, which are represented in the RG picture by the bonds with $`\stackrel{~}{J}^z\omega `$ connecting the effective spins (clusters) in the effective theory with the renormalized cutoff $`\omega `$. The matrix element for producing such an excitation by a bare spin-flip operator of a spin active in one of these clusters will be of order the corresponding $`\stackrel{~}{J}^{}`$, while the number of such spins contributing will be of order some effective “moment” $`\mu _x`$ of this cluster. Because of the matrix element proportional to $`\stackrel{~}{J}^{}`$, there will be a significant contribution only if this $`\stackrel{~}{J}^{}`$ is also of order $`\omega `$. As we have already seen, this can happen only if such an IAF-ordered cluster has length of order $`\mathrm{\Gamma }_\omega `$ and is isolated from the rest by RS-like regions (domain walls) with $`J^z\omega `$ on either side. We already know how to estimate the number density of such Griffiths regions. As far as the effective moment $`\mu _x`$ of such an IAF cluster is concerned, we can only make a crude estimate that bounds it from above by the number of spins that are active in this cluster: $`\mu _x(\omega )\mathrm{\Gamma }_\omega `$; however, we are unable to obtain the precise power of the logarithm that enters the energy dependence of the effective moment. We therefore leave out the logarithmic correction, and only write the dominant power-law part of our estimate:
$`\left[\chi _{\mathrm{loc}}^{xx}\right]_{\mathrm{av}}(\omega )`$ $``$ $`\omega ^{1+2/z_{\mathrm{IAF}}},`$ (49)
$`\left[C_{\mathrm{loc}}^{xx}\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{1}{\tau ^{2+2/z}}}.`$ (50)
### D Spin transport
This section is devoted to a discussion of the dynamical spin conductivity $`\sigma ^{}(\omega )`$ in the spin-1/2 XXZ chains. Our task here is to evaluate the Kubo formula Eq. (5) in the low-frequency limit. For the RS and RD states, we will use information available from the scaling solutions to the singlet RG recursion relations to achieve this, while in the IAF phase, we will use the scaling picture of the Griffiths phase we have developed earlier. Our results for the dynamical conductivity are summarized in Figs. 1 and 2.
#### 1 Random singlet states
We first need to work out the rules that govern the renormalizations of the current operators. Assume once again that $`J_{23}^{}`$ is the strongest bond. We wish to work out perturbatively the renormalized operators $`\stackrel{~}{\tau }_{1/2/3}`$ that we trade in $`\tau _{1/2/3}`$ for, when we freeze spins $`2`$ and $`3`$ in their singlet ground state (the other current operators to the left and right of this segment are left unchanged to leading order by the renormalization). Now, note that these other operators have overall scale factors in them that are nothing but the corresponding $`J^{}`$ couplings. In order to be consistent, we clearly need to work out $`\stackrel{~}{\tau }_{1/2/3}`$ correct to $`O(\stackrel{~}{J}_{14}^{})`$ (where $`\stackrel{~}{J}_{14}`$ is the effective bond connecting spins $`1`$ and $`4`$ after we freeze out spins $`2`$ and $`3`$) by adding the effects of virtual fluctuations to the projections of $`\tau _{1/2/3}`$ into the singlet subspace. An explicit calculation gives the simple result that all three operators renormalize to the same operator $`\stackrel{~}{\tau }_{1/3}=\stackrel{~}{\tau }_2=i\stackrel{~}{J}_{14}^{}(s_1^+s_4^{}s_4^+s_1^{})/2`$, which we will denote henceforth by $`\stackrel{~}{\tau }_1`$ for consistency of notation.
As we carry out the RG, the above result implies that the total current operator $`_{j=1}^L\tau _j`$ entering Eq. (5) renormalizes to $`\stackrel{~}{}_j\stackrel{~}{l}_j\stackrel{~}{\tau }_j`$, where $`j`$ now labels the remaining sites of the renormalized system, and the $`\stackrel{~}{l}_j`$ are the lengths of the corresponding renormalized bonds. \[Note that this result makes sense physically and is a consequence of spin conservation: when a magnetic field with a uniform gradient is applied along the length of the chain, the effective lengths $`\stackrel{~}{l}_j`$ measure the “phase” along the chain of this “driving potential”.\] Consider two spins connected by a strong bond $`(\stackrel{~}{J}^{},\stackrel{~}{J}^z)`$ in the renormalized theory with cutoff $`\mathrm{\Omega }_{\mathrm{final}}`$. Since the current operator living on this bond connects the singlet ground state of the pair only to the triplet state $`|t_0`$ separated from the singlet by a gap $`\stackrel{~}{J}^{}`$, we choose $`\mathrm{\Omega }_{\mathrm{final}}=\omega `$ and consider the renormalized spectral sum
$$\sigma ^{}(\omega )=\frac{1}{\omega L}\stackrel{~}{\underset{m}{}}|m|\stackrel{~}{\underset{j}{}}\stackrel{~}{l}_j\stackrel{~}{\tau }_j|0|^2\delta (\omega \stackrel{~}{E}_m).$$
(51)
This spectral sum is dominated by precisely the $`|t_0`$ triplet excitations of pairs of spins that are connected by the (effective) bonds with $`\stackrel{~}{J}^{}=\omega `$ and are being eliminated at this energy scale; the corresponding matrix element is just $`\stackrel{~}{l}\omega /2`$, where $`\stackrel{~}{l}`$ is the length of the bond connecting the pair. In the thermodynamic limit, we thus have
$`\sigma ^{}(\omega )`$ $``$ $`{\displaystyle \frac{n(\mathrm{\Gamma }_\omega )}{\omega }}{\displaystyle 𝑑l𝑑\zeta \omega ^2l^2P(\zeta ,l|\mathrm{\Gamma }_\omega )\delta (\omega \omega e^\zeta )}.`$ (52)
This immediately yields our central result
$$\sigma ^{}(\omega )=𝒦_{\mathrm{RS}}l_v\mathrm{ln}(\mathrm{\Omega }_0/\omega ),$$
(53)
valid for $`\omega \mathrm{\Omega }_0`$. Here, $`𝒦_{\mathrm{RS}}`$ is an order one numerical constant, $`l_v`$ is the microscopic length scale defined earlier, and $`\mathrm{\Omega }_0`$ is the microscopic energy cutoff. Notice that this analysis holds equally well at all three RS fixed points, which differ only in the corresponding values of the non-universal scale factors.
A brief digression is in order, before we go on to discuss this result: The real part of the dynamical conductivity can be related (on general grounds) to the behavior of the dynamic structure factor $`S^{zz}(k,\omega )`$ near $`k=0`$
$$\sigma ^{}(\omega )=\omega \frac{1}{2}\frac{d^2}{dk^2}S^{zz}(k,\omega );$$
(54)
this can be checked by comparing directly the corresponding spectral sums and noticing that the action of the two operators $`𝒯=_j\tau _j`$ and $`𝒱=_jj\sigma _j^z`$ on the eigenstates of the Hamiltonian $``$ are related through $`𝒯=i[,𝒱]`$. It is easy to check, using the scaling form Eq. (28), that our result for the conductivity is consistent, as it must be, with our previously derived result for the dynamic structure factor.
Going back to Eq. (52), we see that $`\sigma ^{}(\omega )`$ diverges logarithmically for small $`\omega `$ in the unusual ‘spin-metal’ phase controlled by the XX fixed point as well as at the critical points (XXX and XXZC) separating this phase from the ‘insulating’ phase with Ising antiferromagnetic order in the ground state. Note that this ‘metal-insulator’ transition has the curious feature that the quantum critical points separating the conducting phase from the insulating phase have the same $`T=0`$ transport properties as the conducting phase.
#### 2 IAF Griffiths phase
On the insulating side, we expect $`\sigma ^{}(\omega )`$ to be suppressed below the crossover scale $`\mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$; the dominant contributions for $`\omega \mathrm{\Omega }_{\delta _{\mathrm{IAF}}}`$ come from some rare regions that contain long finite segments locally in the ‘metallic’ phase.
We begin by providing a rough estimate of these contributions to $`\sigma ^{}(\omega )`$: In our sample, consider a (large) region of length $`L`$ locally in the RS phase; the number density of such regions is roughly $`p^L`$, with some $`p<1`$ (which depends on the distance from the transition). If these regions are effectively isolated from the rest of the system, the average power absorption per spin in each such region is proportional to the finite-size conductivity calculated in Appendix A:
$$W=L\sigma _{\mathrm{RS}}^{}(\omega ,L)L^{3/2}\mathrm{exp}(c|\mathrm{ln}\omega |^2/L),$$
(55)
where we have assumed that $`L`$, although large, satisfies $`L|\mathrm{ln}\omega |^2`$ (this assumption will turn out to be self-consistent). The total power absorbed in the sample is then obtained by summing over all such regions:
$$\sigma ^{}(\omega )𝑑Lp^LL^{3/2}\mathrm{exp}(c|\mathrm{ln}\omega |^2/L).$$
(56)
Evaluating this integral by a saddle-point method, we find that the lengths that dominate are of order $`|\mathrm{ln}\omega |`$ (our assumption about the lengths is thus valid), and arrive at the following estimate
$$\sigma ^{}(\omega )\omega ^\alpha |\mathrm{ln}\omega |^2,$$
(57)
where $`\alpha =\alpha (\delta _{\mathrm{IAF}})>0`$ is a continuously varying exponent vanishing at the transition. While this argument is suggestive, we find it more convincing to take an alternative route based on the scaling picture we have developed earlier for the IAF phase—this has the added advantage that it allows us to relate the exponent $`\alpha `$ to the dynamical exponent $`z(\delta _{\mathrm{IAF}})`$. This is what we turn to next.
We have already seen that the most numerous low-energy excitations in the IAF Griffiths phase are domain walls, with the integrated density of states $`n_\omega \omega ^{1/z_{\mathrm{IAF}}}`$. Such classical Ising excitations, however, do not contribute to the dynamical conductivity. The dominant contributions come from IAF-ordered clusters of lengths $`L\mathrm{\Gamma }_\omega \mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}`$ (i.e., with effective spin-flip couplings of order $`\omega `$) that are isolated from the rest of the system by domain walls with $`J^z\omega `$. Remembering that the number density of such Griffiths regions is $`\omega ^{2/z_{\mathrm{IAF}}}`$, and noting that the corresponding “phase lengths” are of order $`L|\mathrm{ln}\omega |`$, we immediately obtain Eq. (57) with $`\alpha =2/z_{\mathrm{IAF}}`$. More formally, we sum over the possible separations of two such domain walls, with the constraint that the typical IAF-ordered region isolated by the two has significant spin fluctuations at the characteristic frequency $`\omega `$:
$$\sigma ^{}(\omega )\frac{n_\omega ^2}{\omega }𝑑L\omega ^2L^2\delta (\omega \mathrm{\Omega }_0e^{c_xL}).$$
(58)
We thus obtain for the dynamical conductivity in the IAF phase
$$\sigma ^{}(\omega )=𝒦_{\mathrm{IAF}}l_v(\omega /\mathrm{\Omega }_0)^{2/z_{\mathrm{IAF}}}\mathrm{ln}^2(\mathrm{\Omega }_0/\omega ),$$
(59)
where $`𝒦_{\mathrm{IAF}}`$ is a numerical prefactor that depends continuously on $`\delta _{\mathrm{IAF}}`$. The scaling of $`𝒦_{\mathrm{IAF}}`$ with $`\delta _{\mathrm{IAF}}`$ for small $`\delta _{\mathrm{IAF}}`$ can be obtained by demanding consistency with the off-critical scaling form for the conductivity
$$\sigma ^{}(\omega )=𝒦_{\mathrm{RS}}l_v\mathrm{ln}(\mathrm{\Omega }_0/\omega )\mathrm{\Sigma }_{\mathrm{IAF}}(\mathrm{\Gamma }_\omega /\mathrm{\Gamma }_{\delta _{\mathrm{IAF}}}),$$
(60)
which immediately implies that $`𝒦_{\mathrm{IAF}}\delta _{\mathrm{IAF}}^{(2\psi )/\lambda }z_{\mathrm{IAF}}^1`$.
#### 3 Dynamical conductivity in RD phases
We now calculate the dynamical spin conductivity in the XX and XXX Random Dimer phases. Here, the same singlet RG can be employed all the way across the crossover scale $`\mathrm{\Gamma }_\delta 1/|\delta |`$, and into the energy regime of a well-developed RD phase. The dynamical conductivity is given by the same expression Eq. (52) as for the RS states: we simply add contributions from the even ($`P^\mathrm{e}`$) and the odd ($`P^\mathrm{o}`$) bonds in complete analogy with the calculation of the dynamic structure factor. Using the scaling solutions of Ref. , it is quite simple to calculate the full scaling function for the dynamical conductivity
$$\sigma ^{}(\omega ,\delta )=𝒦_{\mathrm{RS}}l_v\mathrm{ln}(\mathrm{\Omega }_0/\omega )\mathrm{\Sigma }_{\mathrm{RD}}\left(|\delta |\mathrm{ln}(\mathrm{\Omega }_0/\omega )\right).$$
(61)
Here, we restrict ourselves to noting that $`\mathrm{\Sigma }_{\mathrm{RD}}(x)\mathrm{const}`$ for $`x1`$, while for $`x1`$, $`\mathrm{\Sigma }_{\mathrm{RD}}(x)`$ scales as $`\mathrm{\Sigma }_{\mathrm{RD}}(x)xe^{2x}`$. Thus, at frequencies $`\omega `$ well below the crossover scale $`\mathrm{\Omega }_\delta `$, we have
$`\sigma ^{}(\omega )`$ $`=`$ $`𝒦_{\mathrm{RD}}l_v(\omega /\mathrm{\Omega }_0)^{1/z_{\mathrm{RD}}}\mathrm{ln}^2(\mathrm{\Omega }_0/\omega ),`$ (62)
with the numerical prefactor $`𝒦_{\mathrm{RD}}|\delta |`$ and the dynamical exponent $`z_{\mathrm{RD}}|\delta |^1`$ for small $`|\delta |`$. We can now interpret this form directly in terms of the rare regions that dominate the conductivity: Assume, for concreteness, that $`\delta >0`$, i.e., that the even bonds are dominating; the main contribution to the dynamical conductivity at frequency $`\omega \mathrm{\Omega }_\delta `$ then comes from the even bonds with effective $`\stackrel{~}{J}_\mathrm{e}=\omega `$. Such weak even bonds are generated only across rare long regions that are locally in the opposite dimerized phase, and these are precisely the regions that dominate the low-energy density of states and thus determine the dynamical exponent $`z_{\mathrm{RD}}(\delta )`$—this explains the factor $`\omega ^{1/z_{\mathrm{RD}}}`$ in Eq. (62). Moreover, all such bonds have a well-defined length proportional to $`\mathrm{ln}(\mathrm{\Omega }_0/\omega )`$, which explains the $`\mathrm{ln}^2(\mathrm{\Omega }_0/\omega )`$ in Eq. (62).
#### 4 Perspective: spinless interacting fermions with particle-hole symmetric disorder
To put these transport properties in perspective, we recall that the spin-1/2 XXZ chain is equivalent, via the usual Jordan-Wigner transformation, to a system of spinless interacting fermions with particle-hole symmetric disorder. More specifically, we write the spin operators $`s_j^\pm s_j^x\pm is_j^y`$ in terms of fermion creation (annihilation) operators $`c_j^{}`$ ($`c_j`$) as
$`s_j^+`$ $`=`$ $`{\displaystyle \underset{j^{}<j}{}}(12n_j^{})c_j^{},`$ (63)
$`s_j^{}`$ $`=`$ $`{\displaystyle \underset{j^{}<j}{}}(12n_j^{})c_j,`$ (64)
while $`s_j^z=n_j1/2`$ (here $`n_jc_j^{}c_j`$ is the fermion number operator at site $`j`$). In this language, $`_{\mathrm{XXZ}}`$ can be written as
$$=\underset{j=1}{\overset{L1}{}}\left[t_j(c_{j+1}^{}c_j+c_j^{}c_{j+1})+V_j(n_j1/2)(n_{j+1}1/2)\right],$$
(65)
with $`t_j=J_j^{}/2`$ and $`V_j=J_j^z`$. The coupling $`J^z`$ thus controls the strength of the nearest-neighbor particle-hole–symmetric repulsive interaction between the fermions. The IAF phase that obtains for large $`J^z`$ corresponds to a charge density wave state stabilized by interactions. In the absence of interactions (XX chain) we obtain a free-fermion random-hopping problem at zero chemical potential.
This free-fermion problem has been extensively studied in the past, and is known to have rather unusual localization properties due to the additional particle-hole symmetry present. For instance, an elementary calculation immediately reveals that the zero-temperature average Landauer conductance $`[g_L]_{\mathrm{av}}`$ of a finite segment of length $`L`$ connected to perfect leads scales as $`[g_L]_{\mathrm{av}}1/\sqrt{L}`$, in sharp contrast to the usual exponentially-localized behavior in one dimension; the corresponding conductivity, of course, scales as $`\sqrt{L}`$. Now, the strong-disorder RG predicts that lengths scale as the square of the logarithm of the energy scale in the low-energy effective theory describing the XX-RS state—our result for the dynamical conductivity is thus consistent with the elementary Landauer calculation (see also our explicit finite-size scaling calculations of Appendix A). Notice, however, that our approach is not limited to the non-interacting case. It allows us to reliably treat the effects of interactions, and follow the dynamical conductivity through a ‘metal-insulator’ transition that is driven by strong interactions in the presence of strong disorder.
#### 5 Numerical study of the dynamical conductivity at the XX fixed point
At the XX point, the Hamiltonian Eq. (65) describes non-interacting fermions with random hopping amplitudes, and we are essentially faced with the problem of finding the low-energy eigenvalues and eigenstates of the corresponding single-particle Hamiltonian (an $`L\times L`$ matrix operator) $`𝐇=_{j=1}^{L1}t_j(|j+1j|+|jj+1|)`$, which defines the Schroedinger equation for this problem. Any fermionic state can then be represented as a Slater determinant of the corresponding (normalized) single-particle eigenstates $`|\varphi _\mu `$ with eigen-energies $`ϵ_\mu `$. In the single-particle language, the Kubo formula for the conductivity $`\sigma ^{}(\omega )`$ at zero chemical potential and at a finite temperature $`T`$ reads
$`\sigma ^{}(\omega )={\displaystyle \frac{1}{\omega L}}`$ $`{\displaystyle \underset{\mu _1,\mu _2}{}}|\varphi _{\mu _2}|{\displaystyle \underset{j}{}}𝐓(j)|\varphi _{\mu _1}|^2\times `$ (67)
$`\times [f(ϵ_{\mu _1})f(ϵ_{\mu _2})]\delta (\omega ϵ_{\mu _2}+ϵ_{\mu _1}),`$
where $`𝐓(j)it_j(|jj+1||j+1j|)`$ is the current operator on the link $`j`$ and $`f(ϵ)1/(e^{ϵ/T}+1)`$. \[This version of the Kubo formula will also prove useful when we analyze the full temperature dependence of the dynamical conductivity in Sec. VI.\]
Here, we test the $`T=0`$ predictions by evaluating $`\sigma ^{}(\omega )`$ using exact numerical diagonalization of finite systems. The results of such calculations for system sizes $`L=128`$, $`256`$ and $`512`$ with the hopping amplitudes $`t_j`$ drawn independently from a uniform distribution over $`[0,1]`$ are shown on Fig. 6, where we have averaged over $`100,000`$ samples for each $`L`$. In an infinite sample we expect the conductivity to diverge logarithmically, but with the system sizes studied here, we cannot quite probe this infinite-sample regime $`1\mathrm{ln}(\mathrm{\Omega }_0/\omega )\sqrt{L}`$—rather, we are in the regime $`1\mathrm{ln}(\mathrm{\Omega }_0/\omega )\sqrt{L}`$. Nevertheless, the numerical results of Fig. 6 clearly show that the dynamical conductivity is increasing as the frequency is lowered all the way to the crossover scale $`\mathrm{ln}(\mathrm{\Omega }_0/\omega )\sqrt{L}`$, thus supporting our claim that $`\sigma ^{}(\omega )`$ diverges logarithmically at low frequencies.
For a more detailed test of our theoretical results, we need to quantitatively analyze the effects of a finite system size on our predictions for the dynamical conductivity. The calculation is summarized in Appendix A. Here, we only note that this analysis allows us to write the following scaling form for the conductivity
$$\sigma ^{}(\omega ,L)=l_v\mathrm{ln}(\mathrm{\Omega }_0/\omega )\mathrm{\Theta }\left(l_v\mathrm{ln}^2(\mathrm{\Omega }_0/\omega )/L\right);$$
(68)
the scaling function $`\mathrm{\Theta }`$ is characterized in Appendix A, and the above result is expected to hold for large enough $`L`$ and $`\mathrm{ln}(\mathrm{\Omega }_0/\omega )`$ (with no restrictions on the ratio $`\mathrm{ln}^2(\mathrm{\Omega }_0/\omega )/L`$). However, the numerical results cannot be compared directly with this scaling result since it assumes that the distribution of bond lengths has reached the form characteristic of the XX fixed point, which is not the case for the sizes that we can diagonalize numerically: the “length part” of the distribution $`P(\zeta ,l|\mathrm{\Gamma })`$ is still evolving towards the corresponding scaling form from the initial condition $`P(\zeta ,l|\mathrm{\Gamma }_I)=e^\zeta \delta (l1)`$. Nevertheless, we can compare the results of the exact numerical diagonalization with the (formal) predictions of the RG for the same systems. This can be done by either running the RG on the same samples or by evaluating the analytical (within the RG) expression, given by the inverse Laplace transform Eq. (A13) for these initial conditions. In Fig. 6, we compare the RG result obtained in this manner with the numerically evaluated conductivity. Given that the initial disorder is not very strong, the agreement of the RG predictions with the $`\sigma ^{}(\omega ,L)`$ from the exact diagonalization is fairly good.
### E On validity of results
So far, our calculations have relied on the leading-order renormalizations of the spin operators; in this subsection we will try to justify validity of this approximation. We will not address the corresponding question for the RG itself—this has been analyzed with great care in Refs. and , and we have nothing to add here. Instead, we focus on issues specific to our calculation of dynamical quantities. Here, we provide a (partial) justification of our leading-order results by analyzing the effects of the first corrections to the leading-order expressions for the renormalized operators—this can be done consistently within the framework of the RG approach. Any consistent analysis of further corrections would require that we also consider higher-order corrections to the RG rules themselves, and we stop well short of doing that.
As an illustrative example, we consider the dynamic structure factor in the RS states. Our leading-order calculations used only the zeroth-order result for the renormalized spin operators. The renormalized operators can also be easily worked out to first order—these were considered in Ref. in the discussion of typical correlations. When a pair of spins $`2`$ and $`3`$ connected by a strong bond is frozen into a singlet state, the neighboring spin operators $`\stackrel{}{s}_1`$ and $`\stackrel{}{s}_4`$ do not change even to first order, while the spin operators $`\stackrel{}{s}_2`$ and $`\stackrel{}{s}_3`$ renormalize to
$`(s_2^z)^{\mathrm{eff}}=(s_3^z)^{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{J_1^z}{2J_2^{}}}s_1^z+{\displaystyle \frac{J_3^z}{2J_2^{}}}s_4^z,`$ (69)
$`(s_2^+)^{\mathrm{eff}}=(s_3^+)^{\mathrm{eff}}`$ $`=`$ $`{\displaystyle \frac{J_1^{}}{J_2^{}+J_2^z}}s_1^++{\displaystyle \frac{J_3^{}}{J_2^{}+J_2^z}}s_4^+.`$ (70)
Thus, the decimated spins $`\stackrel{}{s}_2`$ and $`\stackrel{}{s}_3`$ obtain small “components” of order $`J_1/J_2`$, $`J_3/J_2`$, onto the neighboring spins $`\stackrel{}{s}_1`$ and $`\stackrel{}{s}_4`$. As we run the RG and renormalize down to scale $`\mathrm{\Gamma }`$, the system consists of $`n_\mathrm{\Gamma }`$ active spins per unit length, separated from each other by “dead” regions (with lengths of order $`\mathrm{\Gamma }^2`$) of decimated spins. Each decimated spin $`r`$ in the dead region between two remaining active spins $`j`$ and $`k`$ (where $`j`$ and $`k`$ are nearest neighbors in the effective Hamiltonian at scale $`\mathrm{\Gamma }`$) will have some components $`C_{rj}`$ and $`C_{rk}`$ on spins $`\stackrel{}{s}_j`$ and $`\stackrel{}{s}_k`$. ¿From the point of view of calculating the spectral sum Eq. (4), each active spin $`j`$ acquires some ($`k`$-dependent) effective moment $`\stackrel{~}{\mu }_j(k)`$ coming from all decimated spins with non-zero components on $`\stackrel{}{s}_j`$:
$$\stackrel{~}{\mu }_j(k)=1+\underset{i_1<r<i_2}{}^{}C_{rj}e^{ik(rj)},$$
(71)
where the sum is over all previously decimated spins $`r`$ between the effective neighbors $`i_1`$ and $`i_2`$ of the spin $`j`$, $`i_1<j<i_2`$. The components $`C_{rj}`$ are simply the ground-state correlations $`s_rs_j`$; such typical correlators decay as a stretched exponential $`[\mathrm{ln}C_{rj}]_{\mathrm{av}}|rj|^{1/2}`$. Note that the characteristic length scale for this decay is the microscopic length scale $`l_v`$. It is thus clear that the sum over $`r`$ in Eq. (71) converges quickly, and the renormalization of the moment $`\stackrel{~}{\mu }_j`$ away from its bare value of $`1`$ comes mainly from the nearby spins that were decimated early in the RG. This renormalization is of order one, but only weakly $`k`$-dependent.
We now analyze the consequences of this renormalization of the moments for the two scaling forms of the dynamic structure factor derived earlier in the limit of low frequencies, one in the vicinity of $`k=\pi /a`$, and the other in the vicinity of $`k=0`$. First, consider $`k=\pi /a+q`$, with $`|q|l_v^1`$. For such small values of $`q`$, we can neglect the $`q`$-dependence of the moments and evaluate them at $`k=\pi /a`$. To evaluate the spectral sum Eq. (4), we need to add up the contributions coming from the strong bonds at scale $`\mathrm{\Omega }_{\mathrm{final}}`$. Each strong bond contributes $`|\stackrel{~}{\mu }_L+\stackrel{~}{\mu }_Re^{iql}|^2`$, where $`\mu _L`$ and $`\mu _R`$ are the moments (evaluated at $`k=\pi /a`$) of the two spins connected by this strong bond. We can now proceed in two steps: First, we fix $`l`$ and average over the moments of all strong bonds with a given length. This gives us a quantity $`c_1+c_2|1+e^{iql}|^2`$ which we now need to average over the length distribution of the strong bonds; here $`c_1`$ and $`c_2`$ are now some fixed numbers of order one, since we expect that the main renormalization of each moment comes from few nearby spins and is roughly independent of the lengths of the adjoining bonds. Thus, we see that Eq. (26), with $`\mathrm{\Phi }`$ given by Eq. (27), indeed describes the dynamic structure factor for $`k`$ close to $`\pi /a`$ and fixed low $`\omega `$; the higher-order corrections renormalize the overall amplitude by a factor of order one, and also produce an “incoherent” background of a comparable strength that depends only weakly on $`k`$ (i.e., that changes significantly only when $`k`$ is changed by an amount of order $`l_v^1`$).
For $`kl_v^1`$ (i.e in the scaling regime near $`k=0`$), the discussion is very similar; each strong bond contributes $`|\stackrel{~}{\mu }_L\stackrel{~}{\mu }_Re^{ikl}|^2`$, where the moments are now evaluated at $`k=0`$. This again gives us a quantity $`c_1+c_2|1+e^{iql}|^2`$ to be averaged over the length distribution of the strong bonds. In general (away from the XXX point), we now have to consider the $`S^{zz}(k,\omega )`$ component separately from the $`S^+(k,\omega )`$ component, since the total $`s_{\mathrm{tot}}^z`$ conservation constrains the constant $`c_1`$ to be identically zero for the case of $`S^{zz}(k,\omega )`$. Thus, in the case of $`S^{zz}(k,\omega )`$, higher-order corrections only produce an order one renormalization to the overall scale of our scaling result Eq. (28); of course, there will be additional corrections, but these will vanish faster than the scaling result in the low-frequency limit. In the case of $`S^+(k,\omega )`$, an inspection of the renormalization rules Eq. (70) shows that to this order $`c_1`$ will be zero for $`S^+(k,\omega )`$ as well, even in the absence of full Heisenberg symmetry; however, this is not expected to be true in general (to all orders), and we expect a small but non-zero background to be present in the general case. Thus, $`S^+(k,\omega )`$ near $`k=0`$ will in general consist of two parts: the scaling part given by Eq. (70) with an order one non-universal overall scale (this part vanishes as $`k^2`$ for small $`k`$), and a non-scaling weakly $`k`$-dependent additive background of the same order as the scaling part.
The above arguments typify the general logic behind our justification of the leading-order results for all of our calculations; in some cases such a programme can be carried out analytically (e.g., for the average boundary spin autocorrelations—see also Sec. V D), while in other cases we have to be satisfied with arguments like the ones presented above. Such arguments can also be bolstered by numerically implementing the higher-order operator renormalizations to calculate corrections within the RG to our leading order results (indeed, we have confirmed that such a numerical check for $`S(k,\omega )`$ in the Heisenberg model is in qualitative agreement with the arguments presented above).
## IV Transport in strongly random spin-1 chains
### A Singlet RG description of the phases: A review
The strong-randomness quantum critical point, which controls the transition from the Haldane state to the Random Singlet state in the spin-1 chains, and the immediate vicinity of this critical point, can be analyzed by a somewhat extended RG procedure introduced in Refs. and , or by a variant of the same used in Ref. .
The basic idea is to replace the original spin-1 chain by an effective model that is argued to describe the low-energy physics of the original system—as we shall see later, this effective model can be made plausible by thinking in terms of a bond-diluted chain (it is also possible to arrive at essentially the same model by starting with a random antiferromagnetic spin-1 chain and using the approximate RG procedure of Ref. ). This effective model is written entirely in terms of spin-1/2 degrees of freedom coupled by nearest-neighbor Heisenberg exchange couplings. All even bonds are always antiferromagnetic and are drawn from an appropriate distribution of positive bonds, while odd bonds can be of either sign and are drawn from a different distribution.
This effective model can be analyzed using the extension of the singlet RG introduced in Refs. and . One begins by looking for the largest antiferromagnetic bond in the system, say $`J_2`$ connecting spins 2 and 3; this defines our bare energy cutoff $`\mathrm{\Omega }_0`$. Further analysis can be split into three cases: (i) If the bonds adjacent to the largest AF bond are smaller in magnitude, the two spins are frozen into a singlet state and an effective coupling $`\stackrel{~}{J}_{14}`$ is generated between spins 1 and 4 exactly as in the singlet RG for the spin-1/2 chain. (ii) If both adjacent bonds are larger in magnitude than $`J_2`$, then spins 1 and 2 and spins 3 and 4 are first combined to make effective spin-1 objects (since in this case $`J_1`$ and $`J_3`$ are necessarily ferromagnetic), and these effective spin-1 degrees of freedom are then frozen into a singlet state, generating an effective coupling $`\stackrel{~}{J}_{05}=4J_0J_4/3J_2`$ between spins 0 and 5. (iii) If only one of the adjacent bonds, say $`J_3`$, is larger in magnitude than $`J_2`$, then spins 3 and 4 are first combined into an effective spin-1 object. The system is then frozen into the subspace in which spin 2 and this effective spin-1 object are coupled together to form an effective spin-1/2 object which we label $`s_2`$ for consistency of notation. The corresponding renormalized couplings are given as $`\stackrel{~}{J}_{12}=J_1/3`$ and $`\stackrel{~}{J}_{25}=2J_4/3`$. This procedure is now iterated with the energy cutoff $`\mathrm{\Omega }`$ being gradually reduced. It is important to note that there is no inconsistency in leaving untouched ferromagnetic bonds $`J<\mathrm{\Omega }`$ that are not adjacent to any antiferromagnetic bonds at the cutoff scale—we could equally well have combined all pairs of spins connected by such strong ferromagnetic bonds into effective spin-1 objects at the cost of cluttering up our notation.
A detailed analysis of this iterative procedure can be summarized as follows: Let $`\mathrm{\Gamma }\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })`$ and let $`n_\mathrm{\Gamma }`$ be the fraction of active spins at log-cutoff $`\mathrm{\Gamma }`$. For the even bonds, we introduce the distribution $`P(\zeta |\mathrm{\Gamma })`$ of the corresponding logarithmic couplings $`\zeta \mathrm{ln}(\mathrm{\Omega }/J)`$. For the odd bonds, let $`N(\mathrm{\Gamma })`$ be the fraction of odd bonds at scale $`\mathrm{\Gamma }`$ that are strongly ferromagnetic with $`J<\mathrm{\Omega }`$; for large $`\mathrm{\Gamma }`$, the remainder of the odd bonds are symmetrically distributed about zero and are therefore described by a distribution for $`|J|`$ that we characterize by the distribution $`Q(\zeta |\mathrm{\Gamma })`$ of the corresponding logarithmic couplings $`\zeta \mathrm{ln}(\mathrm{\Omega }/|J|)`$. When $`W`$, the width of the distribution of the log-exchanges in the original spin-1 Hamiltonian Eq. (6), exceeds a critical value $`W_c`$, the system is in a spin-1 Random Singlet phase. In the language of the spin-1/2 effective model, this RS phase is described by a fixed point with $`P(\zeta |\mathrm{\Gamma })=\mathrm{\Gamma }^1e^{\zeta /\mathrm{\Gamma }}`$, $`N(\mathrm{\Gamma })=1`$, $`n_\mathrm{\Gamma }1/\mathrm{\Gamma }^2`$, and $`Q(\zeta |\mathrm{\Gamma })=Q_0e^{Q_0\zeta }`$ for large $`\mathrm{\Gamma }`$ ($`Q_0`$ is some non-universal O(1) number). As $`W`$ is decreased, the system undergoes a quantum phase transition to the so-called Gapless Haldane (GH) phase; both the quantum critical point and the GH phase in the vicinity of it are still controlled by strong-disorder fixed points. At the critical fixed point (which is an infinite-disorder fixed point) we have $`P(\zeta |\mathrm{\Gamma })=Q(\zeta |\mathrm{\Gamma })=2\mathrm{\Gamma }^1e^{2\zeta /\mathrm{\Gamma }}`$, $`n_\mathrm{\Gamma }1/\mathrm{\Gamma }^3`$, and $`N(\mathrm{\Gamma })=1/2`$. The GH phase in the vicinity of the quantum critical point is controlled by a line of fixed points; each point on this line is characterized by some constant $`P_0`$ (which depends on the strength of disorder $`W`$). At a point labeled by $`P_0`$, we have $`P(\zeta |\mathrm{\Gamma })=P_0e^{P_0\zeta }`$, $`Q(\zeta |\mathrm{\Gamma })=Q_0(\mathrm{\Gamma })e^{Q_0(\mathrm{\Gamma })\zeta }`$ where $`Q_0(\mathrm{\Gamma })e^{P_0\mathrm{\Gamma }}`$, $`N(\mathrm{\Gamma })0`$, and $`n_\mathrm{\Gamma }P_0^3e^{P_0\mathrm{\Gamma }}`$. The continuously varying $`P_0(W)`$ vanishes at the transition as $`P_0(W_cW)^{\nu /3}`$, where $`\nu `$ is the correlation length exponent obtained in Refs. and — the GH phase is thus similar to the dimerized phases of the spin-1/2 chains.
### B Spin transport
#### 1 Doing calculations in the effective model
Before we calculate anything, we need to describe how we think about the spin transport in this case. This is somewhat non-trivial, for we are working in an effective model of spin-1/2 degrees of freedom, and some thought is required to decide what is the correct quantity to calculate.
For this, we go back for a moment to the original random spin-1 chain and review an intuitive construction that leads to the effective model in terms of spin-1/2 variables only. Consider the case of dilute randomness, that is, consider a uniform spin-1 chain with a small fraction of very weak bonds that effectively break the chain into pure finite segments weakly coupled with each other. The low-energy effective degrees of freedom of such a segment are two half-spins localized near the two edges of the segment—these are the spin-1/2’s of the effective model. The coupling of the edge spins on neighboring segments is given roughly by the original coupling of the two segments, and is always antiferromagnetic—these are the even bonds of the effective model. On the other hand, the coupling of the two edge spins of the same segment can be either antiferromagnetic or ferromagnetic depending on whether the length of the segment is even or odd—these coupling are represented by the odd bonds in the effective model.
We now need to express dynamical properties of the system in terms of these effective spin-1/2 degrees of freedom. In particular, we want to analyze the low-frequency power absorption when an oscillating magnetic field with a uniform gradient is applied to the system; this will give us the dynamical conductivity $`\sigma ^{}(\omega )`$. Since the magnetic field couples to the conserved ‘charge’ in the system, the corresponding current operators that we need to use when working out the Kubo formula for the effective model are uniquely determined by spin conservation: The current operator on the odd bonds connecting the edge half-spins $`\stackrel{}{s}_1`$ and $`\stackrel{}{s}_2`$ of the same segment (which represents the total spin current operator of this segment) is $`\stackrel{}{\tau }=J_{12}l_{12}\stackrel{}{s}_1\times \stackrel{}{s}_2`$; here $`J_{12}`$ is the corresponding effective coupling and $`l_{12}`$ is some effective phase length that we expect to be given roughly by the length of the segment. Naturally, the current operators on the even bonds connecting the edge half-spins of the neighboring segments have a similar form (the argument in this case is even simpler: one only needs to know that the true edge spin-1 operator of a segment “projects” onto the corresponding effective edge spin-1/2 operator with an amplitude of order one). Note that the precise values of the phase lengths in the initial effective model (for the dilute spin-1 chain) are not important, since at still lower energies we expect the distributions of couplings and the corresponding bond lengths to approach some universal distributions characteristic of the appropriate fixed point.
#### 2 Dynamical conductivity
Having identified the appropriate current operators in the effective problem, we now work out the rules that govern their renormalization in the RG scheme used to analyze this effective model. As in the spin-1/2 case, and as discussed above, we write the part of the total current operator (in the spectral sum Eq. (5)) that is associated with a given bond $`(j,j+1)`$ in the form $`l_j\stackrel{}{\tau }_j`$ where $`\stackrel{}{\tau }_j`$ is the usual bond operator $`\stackrel{}{\tau }_j=J_j\stackrel{}{s}_j\times \stackrel{}{s}_{j+1}`$ and $`l_j`$ is the appropriate phase length. We can then follow renormalizations of the needed operators by keeping track of the phase lengths, in addition to the various bond-strengths. Unlike the spin-1/2 chains, these phase lengths need not equal the physical distances between the corresponding spins; in fact, even the physical position of an effective half-spin often cannot be specified unambiguously, as, for example, when this half-spin appears as an effective doublet formed by combining (via a strong AF bond) an effective spin-1 (which is an intermediate construction in the Hyman-Yang RG rules) and a neighboring spin-1/2. In such cases, our rules can actually be used to assign some meaning to the physical position of such an effective half-spin.
The rules for the phase lengths can be easily stated: In the cases (i) and (ii), when in the final step we form a singlet from either two spin-1/2 objects or two spin-1 objects, the phase length of the new effective bond is simply the sum of the phase lengths of all the bonds that are eliminated. In the case (iii) the phase lengths associated with effective bonds $`\stackrel{~}{J}_{12}`$ and $`\stackrel{~}{J}_{25}`$ are $`\stackrel{~}{l}_{12}=l_1+(4/3)l_2+(2/3)l_3`$ and $`\stackrel{~}{l}_{25}=l_4(1/3)l_2+(1/3)l_3`$ respectively.
The rules for the phase lengths in the case (iii) are somewhat unusual—for example, negative phase lengths can be produced. Note, however, that there are many factors that prevent this from happening too often, and the phase lengths will in many instances coincide with the corresponding geometrical lengths: decimations in the cases (i) and (ii) tend to “correct” deviations of the phase lengths from the geometrical lengths, and in both the RS and GH phases there are simply no decimations of type (iii) at low enough energies. Also, the lengths $`l_2`$ and $`l_3`$ in the above rule for the case (iii) are the lengths of the strong bonds that are eliminated and are therefore usually smaller than the lengths $`l_1`$ and $`l_4`$ of the more typical bonds. Finally, one can argue generally that the phase positions of the spins as dictated by the phase lengths have to agree—at least roughly—with their geometrical positions as inferred from the order of the (remaining) spins in the chain (i.e., from the spin labels). All of this implies that the phase lengths are roughly given by the geometrical distances between the spins; in particular their scaling with $`\mathrm{\Gamma }`$ is given by the inverse of the density of the remaining spins, $`ln(\mathrm{\Gamma })^1`$.
We can now immediately deduce behavior of the dynamical conductivity in the different phases exactly as in our previous calculations for the spin-1/2 model; as the method remains the same, and the relevant details about the statistics of the fixed point Hamiltonians have already been summarized, we merely state our results.
In the RS phase the same result Eq. (52) applies, as is true for an RS state of an arbitrary-$`S`$ spin chain at strong enough randomness.
At the critical point separating the RS phase from the GH phase, we find
$$\sigma ^{}(\omega )=𝒦_{\mathrm{HY}}l_v\mathrm{ln}^2(\mathrm{\Omega }_0/\omega ),$$
(72)
which is a stronger divergence than in the RS phase (note that this difference from the result in the RS states can be traced to the fact that the density of remaining spins behaves as $`n_\mathrm{\Gamma }1/\mathrm{\Gamma }^3`$ at the critical point, in contrast to the $`1/\mathrm{\Gamma }^2`$ decay of the corresponding quantity in the RS states).
Finally, In the GH phases parametrized by $`P_0`$ we find
$$\sigma ^{}(\omega )=𝒦_{\mathrm{GH}}l_v(\omega /\mathrm{\Omega }_0)^{1/z_{\mathrm{GH}}}\mathrm{ln}^2(\mathrm{\Omega }_0/\omega ),$$
(73)
where we have introduced the continuously varying dynamical exponent $`z_{\mathrm{GH}}P_0^1`$, and $`𝒦_{\mathrm{GH}}`$ is an order one numerical prefactor which goes to a constant as $`WW_c`$ (note that the factor $`\mathrm{ln}^2(\mathrm{\Omega }_0/\omega )`$ appears for exactly the same reasons as in the RD phases of the spin-1/2 chains: the lengths of the singlets that are decimated at scale $`\omega `$ are roughly $`\mathrm{ln}(\mathrm{\Omega }_0/\omega )`$).
## V Dynamics in the random transverse field Ising chain
### A Strong-disorder RG description of the phases: A review
The strong-randomness cluster RG of Ref. , from which the low-energy long-distance behavior of a system near the critical point ($`|\delta |1`$) can be obtained, proceeds as follows: One finds the largest coupling in the system, with energy $`\mathrm{\Omega }_0\mathrm{max}\{h_j,J_j\}`$. If the largest coupling is a field, say $`h_2`$ on spin $`2`$, this spin is frozen into the $`\sigma _2^x=+1`$ ground state of the local field term and is eliminated from the system leaving an effective coupling $`\stackrel{~}{J}_{13}=J_{12}J_{23}/h_2`$ between the neighboring spins $`1`$ and $`3`$. If the largest coupling is an interaction, say $`J_{12}`$ between spins $`1`$ and $`2`$, the two spins are combined into one new spin—a cluster—with an effective spin variable $`\stackrel{~}{\sigma }_{(12)}`$ \[representing the two classical minimum energy states $`\sigma _1^z=\sigma _2^z=\pm 1`$\] and an effective transverse field $`\stackrel{~}{h}_{(12)}=h_1h_2/J_{12}`$; the couplings of this new spin to the neighbors remain unchanged to leading order. Each such cluster $`c`$ has a moment $`\stackrel{~}{\mu }_c`$ given by the number of initial spins in the cluster; when two clusters are combined to form a bigger cluster, their moments add: $`\stackrel{~}{\mu }_{(12)}=\mu _1+\mu _2`$. This procedure is now iterated with the energy cutoff $`\mathrm{\Omega }\mathrm{max}\{\stackrel{~}{h}_j,\stackrel{~}{J}_j\}`$ of the new effective Hamiltonian being gradually lowered.
A detailed analysis of this procedure was given in Ref. , of which a summary follows: Define the log-couplings $`\beta _i\mathrm{ln}(\mathrm{\Omega }/h_i)`$, $`\zeta _i\mathrm{ln}(\mathrm{\Omega }/J_i)`$, and also the log-cutoff $`\mathrm{\Gamma }\mathrm{ln}(\mathrm{\Omega }_0/\mathrm{\Omega })`$; also, let $`n_\mathrm{\Gamma }`$ be the number density per unit length of the (remaining) clusters at scale $`\mathrm{\Gamma }`$. The essential feature of the RG near the critical point is that the distributions of the log-couplings $`R(\beta |\mathrm{\Gamma })`$ and $`P(\zeta |\mathrm{\Gamma })`$ become broader and broader as the energy cutoff is lowered; the RG flows are characterized by a special family of scaling solutions with $`R(\beta |\mathrm{\Gamma })=R_0(\mathrm{\Gamma })e^{R_0(\mathrm{\Gamma })\beta }`$ and $`P(\zeta |\mathrm{\Gamma })=P_0(\mathrm{\Gamma })e^{P_0(\mathrm{\Gamma })\beta }`$. At the critical point, $`\delta =0`$, we have $`R_0(\mathrm{\Gamma })=P_0(\mathrm{\Gamma })=1/\mathrm{\Gamma }`$; thus the widths of the two distributions grow without limit, and the number density decreases as $`n_\mathrm{\Gamma }1/\mathrm{\Gamma }^2`$. Also, magnetic moments of the clusters scale as $`\mu \mathrm{\Gamma }^\varphi `$, with $`\varphi =(1+\sqrt{5})/2`$. In the disordered phase, $`\delta >0`$, beyond the crossover scale $`\mathrm{\Gamma }_\delta 1/|\delta |`$, the width of the field distribution saturates, with $`R_0(\mathrm{\Gamma })2\delta `$ for $`\mathrm{\Gamma }\mathrm{\Gamma }_\delta `$, while the width of the bond distribution grows without limit, with $`P_0(\mathrm{\Gamma })2\delta e^{2\delta \mathrm{\Gamma }}`$. In the ordered phase, $`\delta <0`$, the situation is reversed: $`R_0(\mathrm{\Gamma })2|\delta |e^{2|\delta |\mathrm{\Gamma }}`$ and $`P_0(\mathrm{\Gamma })2|\delta |`$ for $`\mathrm{\Gamma }\mathrm{\Gamma }_\delta `$. In both phases, we have $`n_\mathrm{\Gamma }|\delta |^2e^{2|\delta |\mathrm{\Gamma }}`$. Note also that the clusters that are being eliminated at scale $`\mathrm{\Gamma }\mathrm{\Gamma }_\delta `$ all have a fairly well-defined length of order $`|\delta |^1\mathrm{\Gamma }`$ and magnetic moment of order $`|\delta |^{1\varphi }\mathrm{\Gamma }`$.
### B Average autocorrelations
In this section, we obtain the long-time asymptotics of average imaginary-time autocorrelations in the critical region of the RTFIM—we will be using heavily results of Ref. refering to sections in that paper by, e.g., F Sec. IVB. Our predictions can be compared with the extensive numerical results available in the literature, and this serves as a useful check on the validity of our basic approach—the results of such a comparison have already been summarized in Sec. II C 2.
We consider the local dynamical susceptibilities
$$\chi _{jj}^{\alpha \alpha }=\underset{m}{}|m|\sigma _j^\alpha |0|^2\delta (\omega E_m),$$
(74)
where the sum is over all excited states $`|m`$ with excitation energies $`E_m`$, and $`\alpha =x`$ or $`\alpha =z`$. The low-frequency behavior of these susceptibilities determines the long-time asymptotics of the corresponding imaginary-time autocorrelation functions
$$C_{jj}^{\alpha \alpha }(\tau )=\sigma _j^\alpha (\tau )\sigma _j^\alpha (0),$$
(75)
with $`\alpha =z`$ (local magnetic moment autocorrelation) or $`\alpha =x`$ (local energy autocorrelation); we are considering here only the fluctuating (time-dependent) parts of autocorrelations and will ignore any constant (time-independent) parts (such a constant part in, for example, spin autocorrelation represents a non-zero magnetization density in the system and is a static property). In the following, we simply write $`C_{\mathrm{loc}}(\tau )`$ for the local magnetic moment autocorrelations $`C_{jj}^{zz}(\tau )`$ and $`C_{\mathrm{loc}}^e(\tau )`$ for the local energy autocorrelations $`C_{jj}^{xx}(\tau )`$ (and similarly for susceptibilities). We first obtain (using our basic strategy) results for average susceptibilities, which can be conveniently defined as
$$[\chi _{\mathrm{loc}}^{\alpha \alpha }]_{\mathrm{av}}(\omega )=\frac{1}{L}\underset{j}{}\chi _{jj}^{\alpha \alpha }(\omega ),$$
(76)
where $`L`$ is the size of the system (in the thermodynamic limit of $`L\mathrm{}`$ this definition coincides with an ensemble average over disorder realizations). We also consider a semi-infinite chain, $`j1`$, and calculate average dynamical susceptibilities $`\chi _1^{\alpha \alpha }`$ of the boundary spin $`\sigma _1`$ (in this case, the average is over disorder realizations). These results are then immediately translated to the corresponding statements about the long-time behavior of average autocorrelations. As long as we are only interested in the asymptotic behavior of the average dynamical susceptibilities and autocorrelations, it suffices to use the leading-order results for the renormalization of the corresponding operators—this is discussed further in Sec. V D.
We emphasize from the outset that our calculations in this section are closely related to the discussion in Ref. of static response functions at finite temperature $`T`$: Such static response properties are calculated by assuming that all effective degrees of freedom that are present (in the sense of the RG) at energy scale $`T`$ contribute freely to the response at this temperature, while in our calculations of the dynamical properties at frequency $`\omega `$ only degrees of freedom that are being decimated at scale $`\omega `$ contribute to the dynamical response at this frequency; even this difference disappears when the dynamical susceptibilities are translated to the imaginary time autocorrelations, since average autocorrelations at time $`\tau `$ acquire contributions from all frequencies smaller than $`1/\tau `$. The intent here is merely to present a unified approach, within the RG of Ref. , to the calculation of such average dynamical properties. Also, these calculations, together with a detailed physical picture developed in Ref. of the phases of the system near the critical point, serve as a valuable guide to our intuition in identifying the relevant Griffiths regions that dominate a particular response; on some occasions in the previous sections (particularly in the IAF phase of spin-1/2 chains), such Griffiths arguments were our only source of information about the behavior of dynamical quantities, and the opportunity to compare such suggestive arguments against the results of controlled calculations is most welcome.
#### 1 Average local spin autocorrelation $`[C_{\mathrm{loc}}]_{\mathrm{av}}(\tau )`$
The leading-order renormalizations of the $`\sigma ^z`$ spin operators are particularly simple: As long as a given spin $`j`$ is active, the operator $`\sigma _j^z`$ is renormalized to the “spin” operator $`\stackrel{~}{\sigma }_c^z`$ of the cluster $`c`$ that the spin $`j`$ belongs to; when this cluster is decimated, the corresponding operator renormalizes to zero.
To calculate $`[\chi _{\mathrm{loc}}]_{\mathrm{av}}(\omega )`$, we run the RG down to energy scale $`\mathrm{\Omega }_{\mathrm{final}}=\omega /2`$, and rewrite the spectral sum in terms of the degrees of freedom of the renormalized problem; excitations that contribute to this new sum are clearly the $`\stackrel{~}{\sigma }^x=+1`$ excitations of the spin clusters that are being frozen into their $`\stackrel{~}{\sigma }^x=1`$ states by the transverse fields at this scale, and the spectral sum is now easily evaluated:
$`[\chi _{\mathrm{loc}}]_{\mathrm{av}}(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{L}}\stackrel{~}{{\displaystyle \underset{m}{}}}\stackrel{~}{{\displaystyle \underset{c}{}}}\stackrel{~}{\mu }_c|m|\stackrel{~}{\sigma }_c^z|0|^2\delta (\omega \stackrel{~}{E}_m)`$ (77)
$``$ $`{\displaystyle \frac{n(\mathrm{\Gamma }_\omega )}{\omega }}R_0(\mathrm{\Gamma }_\omega )\overline{\mu }_0(\mathrm{\Gamma }_\omega ),`$ (78)
where we used the fact that all $`\stackrel{~}{\mu }_c`$ spins that are active in an effective cluster $`c`$ contribute identically, and $`\overline{\mu }_0(\mathrm{\Gamma })`$ is the average magnetic moment of the clusters that are being eliminated at scale $`\mathrm{\Gamma }`$. \[Note that here, and in the following, we simply write $`\mathrm{\Gamma }_\omega `$ instead of $`\mathrm{\Gamma }_{\omega /2}`$ to avoid clutter in our notation—since we are interested only in the leading behavior, the difference is not important for our purposes\].
At criticality ($`\delta =0`$), we obtain
$$[\chi _{\mathrm{loc}}]_{\mathrm{av}}(\omega )\frac{1}{\omega |\mathrm{ln}\omega |^{3\varphi }},$$
(79)
for $`\omega \mathrm{\Omega }_0`$. For the average spin autocorrelation in imaginary time $`\tau \mathrm{\Omega }_0^1`$ we then find
$$[C_{\mathrm{loc}}]_{\mathrm{av}}(\tau )\frac{1}{|\mathrm{ln}\tau |^{2\varphi }}.$$
(80)
In the disordered phase ($`\delta >0`$), we obtain
$`\left[\chi _{\mathrm{loc}}\right]_{\mathrm{av}}(\omega )`$ $``$ $`\delta ^{4\varphi }{\displaystyle \frac{|\mathrm{ln}\omega |}{\omega ^{11/z(\delta )}}},`$ (81)
$`\left[C_{\mathrm{loc}}\right]_{\mathrm{av}}(\tau )`$ $``$ $`\delta ^{4\varphi }{\displaystyle \frac{|\mathrm{ln}\tau |}{\tau ^{1/z(\delta )}}},`$ (82)
for $`\omega \mathrm{\Omega }_\delta `$ and $`\tau \mathrm{\Omega }_\delta ^1`$. Here, we have used scaling solutions for the off-critical flows to write the answer for the local susceptibility and have chosen to express the power-law in terms of the dynamical exponent $`z(\delta )`$. From the scaling solution to the RG flow equations, we have $`z^1=2|\delta |`$ — this is to be thought of as the leading term in a small-$`\delta `$ expansion for $`z^1`$. Written in terms of $`z(\delta )`$, our result Eq. (82) is valid more generally, and can be understood directly in the simple picture of the disordered phase given in F Sec. IVB: The average spin autocorrelation at large time $`\tau `$ is dominated by the (rare) spins that belong to the rare strongly coupled clusters (Griffiths regions) with low effective “flipping rates” (i.e., effective transverse fields) smaller than $`\omega 1/\tau `$. The density of such clusters, which is also the density of the most numerous excitations at these low energies, is $`n(\omega )\omega ^{1/z}`$ (this is fixed by the relationship $`\tau l^z`$ between length and time scales, that serves as the definition of $`z`$). Most of these clusters have their effective flipping rates between $`\omega `$ and some fraction of $`\omega `$, and therefore effective moments of order $`|\mathrm{ln}\omega |`$ \[since all clusters that are being eliminated at a fixed energy scale $`\mathrm{\Omega }`$ have roughly the same magnetic moment proportional to $`|\mathrm{ln}\mathrm{\Omega }|`$\]. Estimating the contribution of such Griffiths regions clearly gives us Eq. (82) including the factor of $`|\mathrm{ln}\tau |`$.
Finally, in the ordered phase ($`\delta <0`$), we obtain
$`\left[\chi _{\mathrm{loc}}\right]_{\mathrm{av}}(\omega )`$ $``$ $`|\delta |^{4\varphi }{\displaystyle \frac{|\mathrm{ln}\omega |}{\omega ^{12/z(\delta )}}},`$ (83)
$`\left[C_{\mathrm{loc}}\right]_{\mathrm{av}}(\tau )`$ $``$ $`|\delta |^{4\varphi }{\displaystyle \frac{|\mathrm{ln}\tau |}{\tau ^{2/z(\delta )}}},`$ (84)
for $`\omega \mathrm{\Omega }_\delta `$ and $`\tau \mathrm{\Omega }_\delta ^1`$. In contrast to the case of the disordered phase, the interpretation of Eq. (84) in terms of the picture of the ordered phase presented in F Sec. IVA is more subtle. In the ordered phase, the typical excitations at low energies $`\omega \mathrm{\Omega }_\delta `$ are classical—they are “domain walls” that “break” large clusters apart in the places where the clusters are held together by weak (effective) bonds of strength of order $`\omega `$. Such weak effective bonds represent the rare large regions (Griffiths regions) that are locally in the disordered phase. These domain wall excitations are the most numerous excitations that define the relationship between the energy and the length scales and determine the dynamical exponent $`z(\delta )`$. Such excitations, however, even if they are localized in the neighborhood of site $`j`$, do not contribute to $`\chi _{jj}^{zz}(\omega )`$ since they cannot be “excited” from the (classical) ground state by the action of $`\sigma _j^z`$. Excitations which do contribute involve much more rare ferromagnetic clusters that are are flipping back and forth in isolation, with flipping rates of order $`\omega 1/\tau `$ or slower \[of course, we exclude the macroscopic cluster flipping at a rate of order $`e^{cL}`$ as we are subtracting out the time-independent part of the autocorrelation\]. In the RG language, these are precisely the clusters that are decimated at energy scales of order $`\omega `$, i.e., that happen to have (at these scales) an anomalously strong transverse fields of order $`\omega `$ \[remember that we are in the ordered phase\]. A simple construction, however, clearly shows that the density of such regions is indeed $`\omega ^{2/z}`$, as predicted by the scaling solution: For such a cluster to occur we need a ferromagnetic segment of length $`|\mathrm{ln}\omega |`$ (which is not rare in the ordered phase) that is isolated (from eventually becoming a part of the macroscopic cluster) on each side by a disordered region of comparable length. Each of the two disordered regions is actually a “typical” Griffiths region at these energy scales, and the two are required to occur much closer together than their typical separation $`\omega ^{1/z}`$; this explains the appearance of the power $`2/z`$ in Eq. (84). The factor $`|\mathrm{ln}\tau |`$ again comes from the typical magnetic moments of such ferromagnetic droplets.
#### 2 Average local energy autocorrelation $`[C_{\mathrm{loc}}^e]_{\mathrm{av}}(\tau )`$
We begin by working out the leading-order renormalizations of the $`\sigma ^x`$ operators: When a given spin $`j`$ is combined with another spin $`k`$ into a new cluster (i.e., when the strong bond $`J_{jk}`$ is being eliminated) the operator $`\sigma _j^x`$ renormalizes to $`(\stackrel{~}{h}_{(jk)}/h_j)\stackrel{~}{\sigma }_{(jk)}^x`$, where $`h_j`$ is the transverse field on the spin $`j`$ before the decimation, $`\stackrel{~}{h}_{(jk)}`$ is the effective transverse field on the new cluster $`(jk)`$, and $`\stackrel{~}{\sigma }_{(jk)}^x`$ is the effective “spin-flip” operator of this cluster (this rule ignores a constant term proportional to the identity operator, which is unimportant for our purposes as we are not interested in the time-independent constant piece of the energy autocorrelation function). On the other hand, when the spin $`j`$ is eliminated, the operator $`\sigma _j^x`$ becomes effectively zero to first order in the nearby interactions (we again ignore any constants). Iterating this, the operator $`\sigma _j^x`$ is renormalized to $`(\stackrel{~}{h}_c/h_j^{(0)})\stackrel{~}{\sigma }_c^x`$ if the spin $`j`$ is active in some cluster $`c`$ with the effective field $`\stackrel{~}{h}_c`$, and is renormalized to zero if the spin is not active; here $`h_j^{(0)}`$ is the original (bare) transverse field on the spin $`j`$.
We now run the RG down to energy scale $`\mathrm{\Omega }_{\mathrm{final}}=\omega /2`$ and rewrite the spectral sum as
$$[\chi _{\mathrm{loc}}^e]_{\mathrm{av}}(\omega )=\frac{1}{L}\stackrel{~}{\underset{m}{}}\stackrel{~}{\underset{c}{}}\stackrel{~}{g}_c|m|\stackrel{~}{\sigma }_c^x|0|^2\delta (\omega \stackrel{~}{E}_m),$$
(85)
where
$$\stackrel{~}{g}_c=\underset{jc}{}\left(\frac{\stackrel{~}{h}_c}{h_j^{(0)}}\right)^2.$$
(86)
In the last equation $`_{jc}`$ is over all spins that are active in a given cluster $`c`$. Note that at low energies the effective cluster field $`\stackrel{~}{h}_c`$ is only weakly correlated with each of the bare fields $`h_j^{(0)}`$; if the bare field distribution is not too broad, we can approximate $`h_j^{(0)}\mathrm{\Omega }_0`$, and write $`\stackrel{~}{g}_c\stackrel{~}{\mu }_c(\stackrel{~}{h}_c/\mathrm{\Omega }_0)^2`$, where $`\stackrel{~}{\mu }_c`$ is the moment of the cluster $`c`$. Doing this clearly misses some nonuniversal numerical factor of order one that depends on the bare (high-energy) physics. This factor is, in principle, a random quantity that differs from one cluster to another; however, this number is expected to be roughly the same for all clusters that contribute to the spectral sum at low frequencies due to averaging, since such clusters are all large, and in some sense similar. Thus, we expect that the low-frequency behavior is not affected. \[Note that we would have been spared this discussion if we were to analyze the spectral sum with matrix elements of $`h_j^{(0)}\sigma _j^x`$, which is anyway a more natural operator to consider when thinking of the local energy fluctuations\]. The excitations that contribute to the spectral sum Eq. (85) correspond to transitions from the $`\stackrel{~}{\sigma }_j^z=\stackrel{~}{\sigma }_k^z`$ states to the $`\stackrel{~}{\sigma }_j^z=\stackrel{~}{\sigma }_k^z`$ states of two (effective) spins $`j`$ and $`k`$ that are being combined into one cluster $`\stackrel{~}{\sigma }_{(jk)}`$ by a strong bond at the energy scale $`\mathrm{\Omega }_{\mathrm{final}}`$. Since the log-field distribution is broad \[we are near the critical point and at low energy scales\], for such a pair of spins to contribute significantly the transverse field on at least one of the two spins involved must be of order $`\omega `$; thus, we have
$$[\chi _{\mathrm{loc}}^e]_{\mathrm{av}}(\omega )\frac{\omega }{\mathrm{\Omega }_0^2}n(\mathrm{\Gamma }_\omega )P_0(\mathrm{\Gamma }_\omega )R_0(\mathrm{\Gamma }_\omega )\overline{\mu }_0(\mathrm{\Gamma }_\omega ),$$
(87)
from which we immediately read-off our results:
At criticality, for $`\omega \mathrm{\Omega }_0`$ and $`\tau \mathrm{\Omega }_0^1`$, we obtain
$`\left[\chi _{\mathrm{loc}}^e\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{\omega }{|\mathrm{ln}\omega |^{4\varphi }}},`$ (88)
$`\left[C_{\mathrm{loc}}^e\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{1}{\tau ^2|\mathrm{ln}\tau |^{4\varphi }}}.`$ (89)
Away from the critical point, both in the disordered phase and in the ordered phase, we obtain
$`\left[\chi _{\mathrm{loc}}^e\right]_{\mathrm{av}}(\omega )`$ $``$ $`|\delta |^{5\varphi }\omega ^{1+2/z(\delta )}|\mathrm{ln}\omega |,`$ (90)
$`\left[C_{\mathrm{loc}}^e\right]_{\mathrm{av}}(\tau )`$ $``$ $`|\delta |^{5\varphi }{\displaystyle \frac{|\mathrm{ln}\tau |}{\tau ^{2+2/z(\delta )}}},`$ (91)
for $`\omega \mathrm{\Omega }_\delta `$ and $`\tau \mathrm{\Omega }_\delta ^1`$; in the last formula we again used $`z(\delta )`$ as a more physical parameter characterizing the Griffiths phase at a given $`\delta `$. The off-critical energy autocorrelation function thus behaves similarly in the two phases, as expected from duality. It is again possible to interpret these results in terms of the statistical properties of appropriate rare regions that dominate the average energy autocorrelation at long times. As the results (and their interpretation) are identical in either phase, we sketch only the interpretation on the ordered side: As we have already noted, for a region to have significant energy fluctuations at the frequency scale of order $`\omega `$, it must contain two adjoining segments both having a characteristic energy of order $`\omega `$—a predominantly disordered segment (in the RG language, an effective bond with $`\stackrel{~}{J}\omega `$ across this segment) and a predominantly ordered segment (in the RG language, a cluster with $`\stackrel{~}{h}\omega `$). Clearly, the predominantly ordered segment with effective transverse field $`\omega `$ can exist only if it is also isolated on the other side from the rest of the system by another predominantly disordered segment having the same characteristic energy scale. This situation has already been analyzed in the context of the spin autocorrelations in the ordered phase, and clearly one recovers precisely Eq. (91) from such an analysis.
#### 3 Autocorrelations of the boundary spin
So far, we have calculated average autocorrelations for the spins in the bulk; calculations for the first spin $`\sigma _1`$ in a semi-infinite chain $`j1`$ proceed analogously, and we will simply state the results.
For the spin autocorrelation, we find
$$[\chi _1]_{\mathrm{av}}(\omega )\frac{E(0|\mathrm{\Gamma }_\omega )}{\omega }=\frac{P_0(\mathrm{\Gamma }_\omega )R_0(\mathrm{\Gamma }_\omega )}{\omega };$$
(92)
here $`E(\beta |\mathrm{\Gamma })d\beta `$ is the probability that $`\sigma _1`$ survives to scale $`\mathrm{\Gamma }`$ and is in a cluster with $`\mathrm{ln}(\mathrm{\Omega }/h)=\beta `$; such properties of the boundary spin are fully characterized in F Sec. V. At criticality we obtain
$`\left[\chi _1\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{1}{\omega |\mathrm{ln}\omega |^2}},`$ (93)
$`\left[C_1\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{1}{|\mathrm{ln}\tau |}},`$ (94)
while away from the critical point, we get
$`\left[\chi _1\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{\delta ^2}{\omega ^{11/z(\delta )}}},`$ (95)
$`\left[C_1\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{\delta ^2}{\tau ^{1/z(\delta )}}}.`$ (96)
As in the bulk case, we can interpret the average off-critical spin autocorrelation Eq. (96) in terms of the rare instances that dominate this average. In this case, the corresponding rare regions must start at $`\sigma _1`$—this explains the absence of a $`|\mathrm{ln}\tau |`$ factor in Eq. (96) compared to the bulk results Eqs. (82) and (84). The only other difference is that in the ordered phase we do not need to isolate the ferromagnetic droplet (containing $`\sigma _1`$) from the left.
For the energy autocorrelation we find
$$[\chi _1^e]_{\mathrm{av}}(\omega )\frac{\omega }{\mathrm{\Omega }_0^2}P_0^2(\mathrm{\Gamma }_\omega )R_0(\mathrm{\Gamma }_\omega ).$$
(97)
At criticality we obtain
$`\left[\chi _1^e\right]_{\mathrm{av}}(\omega )`$ $``$ $`{\displaystyle \frac{\omega }{|\mathrm{ln}\omega |^3}},`$ (98)
$`\left[C_1^e\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{1}{\tau ^2|\mathrm{ln}\tau |^3}},`$ (99)
while in the disordered phase
$`\left[\chi _1^e\right]_{\mathrm{av}}(\omega )`$ $``$ $`\delta ^3\omega ^{1+2/z(\delta )},`$ (100)
$`\left[C_1^e\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{\delta ^3}{\tau ^{2+2/z(\delta )}}},`$ (101)
and in the ordered phase
$`\left[\chi _1^e\right]_{\mathrm{av}}(\omega )`$ $``$ $`|\delta |^3\omega ^{1+1/z(\delta )},`$ (102)
$`\left[C_1^e\right]_{\mathrm{av}}(\tau )`$ $``$ $`{\displaystyle \frac{|\delta |^3}{\tau ^{2+1/z(\delta )}}}.`$ (103)
Thus, the average off-critical energy autocorrelation of the boundary spin differs from that of the bulk spins in exactly the same way (and for the same reasons) as in the case of the spin autocorrelation.
### C Dynamic structure factor of the spins
Let us now briefly consider the dynamic structure factor $`S^{zz}(k,\omega )`$ defined as
$$S^{zz}(k,\omega )=\frac{1}{L}\underset{m}{}|m|\underset{j}{}e^{ikx_j}\sigma _j^z|0|^2\delta (\omega E_m).$$
(104)
$`S^{zz}(k,\omega )`$ characterizes the spatial structure of the excitations at energy $`\omega `$. Proceeding as before, we find
$$S^{zz}(k,\omega )\frac{n(\mathrm{\Gamma }_\omega )}{\omega }R_0(\mathrm{\Gamma }_\omega )\overline{|\mu _0(k)|^2}(\mathrm{\Gamma }_\omega ),$$
(105)
where $`\overline{|\mu _0(k)|^2}(\mathrm{\Gamma })`$ is the average modulus squared of the effective magnetic moment at wavevector $`k`$ for the clusters that are being eliminated at scale $`\mathrm{\Gamma }`$; for a given cluster $`c`$, this effective moment is defined as $`\mu _c(k)=_{jc}e^{ikx_j}`$. The dynamic structure factor can also be written in terms of the function $`D(\beta ,x|\mathrm{\Gamma }_\omega )`$ defined in F Sec. IIIB4; we have
$$S^{zz}(k,\omega )\frac{\widehat{D}(0,k|\mathrm{\Gamma }_\omega )}{\omega },$$
(106)
where $`\widehat{D}(0,k|\mathrm{\Gamma }_\omega )`$ is the Fourier transform of $`D(0,x|\mathrm{\Gamma }_\omega )`$ at wavevector $`k`$. We have not attempted to analyze $`D(\beta ,x|\mathrm{\Gamma }_\omega )`$, even though a detailed characterization is likely to be possible (see F Sec. IIIB4). Instead, we will only analyze the behavior of the dynamic structure factor in some limiting cases using the scaling picture.
First, consider the system at criticality. Fix wavevector $`k`$. Then, for $`\mathrm{\Gamma }\mathrm{\Gamma }_k1/\sqrt{k}`$ the effective cluster moments at wavevector $`k`$ “add coherently” (more precisely, the real parts of the effective moments of the clusters that are being combined into bigger clusters are of the same sign) and therefore scale as $`\mathrm{\Gamma }^\varphi `$. At scales $`\mathrm{\Gamma }\mathrm{\Gamma }_k`$, the effective moments at $`k`$ “add incoherently” (the real parts of the moments being combined can be of any relative sign) and therefore scale as $`\mathrm{\Gamma }_k^\varphi (\mathrm{\Gamma }/\mathrm{\Gamma }_k)^{\varphi _{\mathrm{sym}}}`$, where $`\varphi _{\mathrm{sym}}=(1+\sqrt{5})/4`$ is the growth exponent for the cluster moments distributed symmetrically around zero (see Appendix of Ref. ). Thus, we arrive at the following scaling form for the dynamic structure factor at criticality
$$S^{zz}(k,\omega )\frac{\mathrm{\Gamma }_\omega ^{2\varphi }}{\omega \mathrm{\Gamma }_\omega ^3}\mathrm{\Phi }\left(k\mathrm{\Gamma }_\omega ^2\right),$$
(107)
where $`\mathrm{\Phi }(x)\mathrm{const}`$ for $`x1`$ and $`\mathrm{\Phi }(x)1/x^{\varphi \varphi _{\mathrm{sym}}}`$ for $`x1`$. We cannot, however, address the regime $`k\mathrm{\Gamma }_\omega ^21`$ by such a scaling analysis.
Now, consider the system that is not critical, either in the disordered or in the ordered phase, in the regime $`\mathrm{\Gamma }_\omega \mathrm{\Gamma }_\delta `$. The length and the magnetic moment of a cluster that is eliminated at scale $`\mathrm{\Gamma }\mathrm{\Gamma }_\delta `$ are sharply defined: $`l_0(\mathrm{\Gamma })=c_l(\mathrm{\Gamma }/\mathrm{\Gamma }_\delta )\mathrm{\Gamma }_\delta ^2+O(\sqrt{\mathrm{\Gamma }/\mathrm{\Gamma }_\delta }\mathrm{\Gamma }_\delta ^2)`$ and $`\mu _0(k=0,\mathrm{\Gamma })=c_\mu (\mathrm{\Gamma }/\mathrm{\Gamma }_\delta )\mathrm{\Gamma }_\delta ^\varphi +O(\sqrt{\mathrm{\Gamma }/\mathrm{\Gamma }_\delta }\mathrm{\Gamma }_\delta ^\varphi )`$, where $`c_l`$ and $`c_\mu `$ are numerical constants of order one. Such a cluster has some internal structure on the length scales below the correlation length $`\delta ^2`$, but “looks” fairly uniform on larger length scales. Then, for the wavevectors $`k\delta ^2`$ we have
$$\overline{|\mu _0(k)|^2}(\mathrm{\Gamma })\frac{\delta ^{42\varphi }}{k^2}\left[1+\mathrm{cos}(c_lk\mathrm{\Gamma }/|\delta |)e^{ck^2\mathrm{\Gamma }/|\delta |^3}\right];$$
(108)
note the oscillatory behavior at $`k|\delta |/\mathrm{\Gamma }`$ due to the “sharpness” of the lengths $`l_0`$ of the clusters; the gradual suppression of this oscillatory behavior at larger wavevectors comes from the uncertainty in $`l_0`$, which is much smaller than $`l_0`$ itself. To obtain the dynamic structure factor $`S^{zz}(k,\omega )`$ we simply need to multiply this “cluster structure factor” by the density $`\rho (\omega )`$ of such clusters at energy $`\omega `$: $`\rho (\omega )\delta ^3/\omega ^{11/z(\delta )}`$ in the disordered phase and $`\rho (\omega )\delta ^3/\omega ^{12/z(\delta )}`$ in the ordered phase.
### D On the validity of the results
Our analysis in Sec. III E of the effects on the (calculated) dynamic spin structure factor of higher-order corrections to the spin operator renormalization rules in the spin-1/2 antiferromagnetic chains can be carried over to the quantum Ising spin chain as well. Such higher-order corrections will only renormalize the effective magnetic moments of the remaining clusters by numerical factors of order one—this is also discussed in detail in F Sec. VIA. We do not repeat such analysis here. Instead, we consider the effect of such corrections on the average autocorrelations of the boundary spin $`\sigma _1`$. This is somewhat more tractable than the bulk case, and we can actually include these effects in an explicit calculation. Our results are not surprising: these next-order corrections only affect the values of some non-universal prefactors, and have no effect on the time dependence of the autocorrelations in the long-time limit.
The first-order renormalization of the boundary spin operator $`\sigma _1^z`$ is discussed in detail in F Sec. VB in the study of the end-point magnetization, and also in Ref. in the study of the end-to-end correlations. Following these references, we write the component of the boundary spin on the left-most remaining cluster $`\stackrel{~}{\sigma }_l`$ as $`\sigma _1^z\stackrel{~}{\sigma }_l^z=e^\mathrm{\Lambda }`$. Then, for the special family of scaling solutions, the distribution of the parameter $`\mathrm{\Lambda }`$ at scale $`\mathrm{\Gamma }`$ is described by the probability density
$$(\mathrm{\Lambda }|\mathrm{\Gamma })=\frac{P_0(\mathrm{\Gamma })}{P_0(\mathrm{\Gamma }_I)}\delta (\mathrm{\Lambda })+\frac{P_0(\mathrm{\Gamma }_I)P_0(\mathrm{\Gamma })}{P_0(\mathrm{\Gamma }_I)}P_0(\mathrm{\Gamma })e^{P_0(\mathrm{\Gamma })\mathrm{\Lambda }},$$
(109)
where the first term is simply the probability that $`\sigma _1`$ survives down to the scale $`\mathrm{\Gamma }`$ from the initial scale $`\mathrm{\Gamma }_I`$ (we already considered this probability in our zeroth-order calculation). The average autocorrelation is then given as
$$[\chi _1]_{\mathrm{av}}(\omega )=\frac{P_0(\mathrm{\Gamma }_\omega )R_0(\mathrm{\Gamma }_\omega )}{\omega P_0(\mathrm{\Gamma }_I)}\left[1+\frac{P_0(\mathrm{\Gamma }_I)P_0(\mathrm{\Gamma }_\omega )}{2+P_0(\mathrm{\Gamma }_\omega )}\right],$$
(110)
where we have kept all the numerical factors within the RG; the first term comes from the instances when the boundary spin survives down to energy scale $`\mathrm{\Gamma }_\omega `$ and is precisely our zeroth-order result Eq. (92). In all cases of interest $`P_0(\mathrm{\Gamma }_\omega )1`$, and the second term is thus roughly the first term times $`P_0(\mathrm{\Gamma }_I)`$. Now, $`1/P_0(\mathrm{\Gamma }_I)`$ is the width of the initial distribution of the logarithms of interactions, and is therefore some number of order one—the two contributions are thus of the same order. However, the wider this initial distribution is, the smaller the second term is relative to the first. The origin of the second term is easily understood: the boundary spin $`\sigma _1`$ may be decimated with some finite probability at an early stage of the RG; in such a case $`\sigma _1`$ will have a significant component on the nearest surviving spin, which may now be considered as a boundary spin in the new problem with a somewhat smaller energy cutoff; the average large-time (low-frequency) autocorrelation of this new boundary spin will be at least of order Eq. (92) from the instances when this new spin survives to the energy scale $`\omega `$. The above calculation explicitly shows that the average dynamical response of the spin $`\sigma _1`$ is dominated entirely by the instances when one of the spins within some small (of order one) distance from the boundary survives down to low energies of order $`\omega `$.
Thus, we see that our general arguments for the effect of higher-order terms in the operator renormalization rules are borne out by this detailed calculation for the boundary spin autocorrelation. This suggests that our approach gives asymptotically exact results (apart from non-universal prefactors) in the bulk case as well.
## VI A discussion of $`T0`$ properties
So far we have calculated various dynamical and transport quantities at $`T=0`$. These results are clearly valid even at $`T0`$ so long as the probe frequency $`\omega `$ satisfies $`T\omega `$. Unfortunately, it is not straightforward to generalize our calculations to the complementary low-frequency regime ($`\omega T`$) dominated by thermal effects. There is, however, one exception. As mentioned earlier, the spin-1/2 XX chain is equivalent to a model of spinless fermions with random nearest neighbor hopping and zero chemical potential. It should come as no surprise that the free-fermion nature of this problem allows us to straightforwardly calculate some dynamical and transport properties at small non-zero temperature. We begin this section by formulating a fermion analog of the RG procedure used for the spin chains. We will then use this RG approach to work out the low-frequency, low-temperature dynamical conductivity and the $`zz`$ component of the dynamic structure factor for the spin-1/2 XX chain without any restriction on the relative magnitudes of $`\omega `$ and $`T`$ (the calculation of the perpendicular component of the structure factor for $`\omega <T`$ is much more complicated, and we will only be able to discuss its qualitative behavior). Naturally, these results are not at all generic, relying as they do on the free-fermion nature of the problem. On the other hand, a weak $`J^z`$ coupling, which corresponds to the nearest-neighbor interaction between the fermions, is strongly irrelevant in the RG sense at the free-fermion point, and the system flows to the non-interacting point. Since this non-interacting limit is singular as far as finite temperature transport properties are concerned, we have here an example of a “dangerously irrelevant operator”, and the important physical question is how this weak, irrelevant interaction affects the $`T0`$ transport near the non-interacting point. This is what we turn to at the end of this section.
### A Free fermion RG
The free-fermion problem $`=_jt_j(c_j^{}c_{j+1}+c_{j+1}^{}c_j)`$ has been the subject of extensive investigation in the past using a variety of techniques (see, e.g., Ref. and references therein). For our purposes, it is most convenient to introduce a RG procedure analogous to the singlet RG used in the spin problem. We formulate this procedure directly in terms of the corresponding single-particle Shrodinger problem $`𝐇=_jt_j(|jj+1|+|j+1j|)`$; this RG is, for the case of the Hamiltonian above, essentially just an efficient way of (approximately) diagonalizing random symmetric tridiagonal matrices with zeroes on the diagonal. We begin with the observation that the particle-hole symmetry of the problem causes eigenstates to occur in pairs, with energies $`\pm ϵ`$. The strong-randomness RG proceeds by eliminating, at each step, such a pair of states with energies at the top and bottom of the band: One finds the largest (in absolute value) hopping amplitude in the system, say $`t_2`$ connecting sites $`2`$ and $`3`$; this defines the bandwidth $`\mathrm{\Omega }_0=2\times \mathrm{max}\{|t_j|\}`$ of the original problem. If the distribution of the $`t_j`$ is broad, the symmetric and antisymmetric wavefunctions living on these two sites will be good approximations to eigenstates with energies $`\pm \mathrm{\Omega }_0/2`$, as $`t_{1/3}`$ will typically be much smaller in magnitude than $`t_2`$. The couplings $`t_{1/3}`$ can then be treated perturbatively, and eliminating the high-energy states living on the sites $`2`$ and $`3`$ results in an effective hopping amplitude $`\stackrel{~}{t}_1=t_1t_3/t_2`$ between the neighboring sites $`1`$ and $`4`$. More precisely, in the effective Hamiltonian that describes the remaining $`L2`$ states, the block $`1`$-$`2`$-$`3`$-$`4`$ is represented as $`\stackrel{~}{}_{14}=\stackrel{~}{t}_1(|\stackrel{~}{1}\stackrel{~}{4}|+|\stackrel{~}{4}\stackrel{~}{1}|)`$, where the states $`|\stackrel{~}{1}`$ and $`|\stackrel{~}{4}`$ are essentially the original $`|1`$ and $`|4`$ states up to $`O(t_{1/3}/t_2)`$ corrections. This rule is essentially identical to the rule for the singlet RG at the XX point, as the additional minus sign can be ‘gauged away’ in the nearest-neighbor model in one dimension; we will, in fact, only keep track of the absolute values of the $`t_j`$. The distribution of $`|t_j|`$ in the renormalized problem with bandwidth $`\mathrm{\Omega }`$ will thus be the same as the renormalized distribution of $`J^{}`$ at cut-off scale $`\mathrm{\Omega }`$ in the singlet RG for the spin problem. The analysis of the asymptotic validity of this approach thus carries over unchanged from the singlet RG.
This procedure can therefore be iterated to reach lower and lower energies; at each stage we trade in our current problem for a new problem defined on two fewer sites. This new problem will have the same low-energy eigenvalues as our original problem. However, evaluating matrix elements of operators between two low-energy states requires some care, as the states $`|\stackrel{~}{j}`$ in terms of which the renormalized problem is written are different from the states $`|j`$ of the original problem. As in the singlet RG, this is best handled by renormalizing the operators as we go along, so that the matrix elements of the renormalized operators between the states of the new problem are the same as the matrix elements of the bare operators between the corresponding states of the original problem. This allows us to calculate various dynamical properties by evaluating the corresponding spectral sums exactly as in the spin language. At $`T=0`$, this amounts to nothing more than a restatement in terms of the fermions of our previous calculations. The new language, however, has one important advantage: thermal effects are easily incorporated into this framework, essentially because the non-interacting nature of the problem is made explicit. \[We emphasize again that the RG finds all eigenstates of the free-fermion problem. The corresponding statement can also be made in terms of the singlet RG on the XX spin chain: when eliminating a pair of spins $`2`$ and $`3`$ the effective Hamiltonians in all sectors (corresponding to the states $`|s_0`$, $`|t_0`$, and $`|t_{\pm 1}`$, of the pair) are identical up to a sign of $`\stackrel{~}{J}_{14}^{}`$ in the $`|t_{\pm 1}`$ sector.\] Moreover, it is now possible to address questions specific to the free-fermion problem that do not have a natural analog in the spin language (see Appendix B).
Finally, we note in passing that this RG procedure can be generalized to analyze other particle-hole symmetric free-fermion problems in one and two dimensions (which are not immediately equivalent to any quantum spin problem) as well as analyze the general properties of the Bogoliubov-de Gennes equation for quasiparticles in a one dimensional superconducting wire in the absence of spin-rotation symmetry (the results of such an analysis will be published separately).
### B $`T0`$ dynamics and transport at the XX point
Let us begin by working out the full $`T`$ and $`\omega `$ dependence of the dynamical conductivity, Eq. (67), at the free fermion point. Our first task is to work out the rules that govern the renormalization of the current operators $`𝐓(j)`$. Assume once again that the hopping element $`t_2`$ has maximum magnitude. We wish to work out what operators we should use in place of $`𝐓(1)`$, $`𝐓(2)`$ and $`𝐓(3)`$ when we renormalize down to lower energies by eliminating the corresponding two states at the top and bottom of the band (the other current operators to the left and right of this segment will be unchanged to leading order by this elimination). An explicit perturbative calculation immediately yields $`\stackrel{~}{𝐓}(2)=\stackrel{~}{𝐓}(1/3)=i\stackrel{~}{t}_1(|\stackrel{~}{1}\stackrel{~}{4}||\stackrel{~}{4}\stackrel{~}{1}|)`$; this is completely analogous to the rule obtained in the spin representation, and as before, we will call this operator $`\stackrel{~}{𝐓}(1)`$ for consistency of notation.
As we carry out the RG and reduce the bandwidth by removing states from the top and bottom of the band, the above result implies that $`_j𝐓(j)`$ renormalizes to $`\stackrel{~}{}_j\stackrel{~}{l}_j\stackrel{~}{𝐓}(j)`$, where $`j`$ now labels the sites of the renormalized problem with bandwidth $`\mathrm{\Omega }`$, and the $`\stackrel{~}{l}_j`$ are the lengths of the renormalized bonds in this problem. With this in hand, we run the RG until the bandwidth is reduced to $`\mathrm{\Omega }_{\mathrm{final}}=\omega `$ and rewrite the spectral sum Eq. (67) as
$`\sigma ^{}(\omega )={\displaystyle \frac{1}{\omega L}}`$ $`\stackrel{~}{{\displaystyle \underset{\mu _1,\mu _2}{}}}|\stackrel{~}{\varphi }_{\mu _2}|\stackrel{~}{{\displaystyle \underset{j}{}}}\stackrel{~}{l}_j\stackrel{~}{𝐓}(j)|\stackrel{~}{\varphi }_{\mu _1}|^2\times `$ (112)
$`\times [f(ϵ_{\mu _1})f(ϵ_{\mu _2})]\delta (\omega ϵ_{\mu _2}+ϵ_{\mu _1}).`$
Because of the extremely broad distribution of the $`\stackrel{~}{t}_j`$, the dominant contribution to the sum Eq. (112) comes from transitions between the two members, one at the bottom and the other at the top of the renormalized band, of each pair of states that is being eliminated at this energy scale. The matrix element for this transition is just $`\stackrel{~}{l}|\omega |/2`$, where $`\stackrel{~}{l}`$ is the length of the hop in question. In the thermodynamic limit, we thus have
$`\sigma ^{}(\omega )`$ $`=`$ $`[f(\omega /2)f(\omega /2)]\times `$ (114)
$`\times {\displaystyle \frac{n(\mathrm{\Gamma }_\omega )}{\omega }}{\displaystyle 𝑑l𝑑\zeta \omega ^2l^2P(\zeta ,l|\mathrm{\Gamma }_\omega )\delta (\omega \omega e^\zeta )}`$
$`=`$ $`{\displaystyle \frac{\mathrm{sinh}(\omega /2T)}{2\mathrm{cosh}^2(\omega /4T)}}\sigma _{T=0}^{}(\omega ),`$ (115)
which is the leading behavior for $`\omega ,T\mathrm{\Omega }_0`$. This result smoothly interpolates between the logarithmic frequency dependence seen earlier for $`\omega T`$ and the limiting form $`\sigma ^{}(\omega )\omega \mathrm{ln}(\mathrm{\Omega }_0/\omega )/T`$ valid for $`\omega T`$—a plot of this frequency dependence is shown in Fig. 7.
Let us now turn to the $`T0`$ spin dynamic structure factor at low frequencies in the vicinity of $`k=\pi /a`$. In the single-particle language, the spectral representation for $`S^{zz}(k,\omega )`$ reads
$`S^{zz}(k,\omega )`$ $`=`$ $`{\displaystyle \frac{1}{(1e^{\omega /T})L}}{\displaystyle \underset{\mu _1,\mu _2}{}}|\varphi _{\mu _2}|\widehat{𝐒}^z(k)|\varphi _{\mu _1}|^2\times `$ (117)
$`\times [f(ϵ_{\mu _1})f(ϵ_{\mu _2})]\delta (\omega ϵ_{\mu _2}+ϵ_{\mu _1}).`$
Here $`\widehat{𝐒}^z(k)_j𝐒^z(j)e^{ikx_j}`$ is the Fourier transform of the position-dependent matrix operator $`𝐒^z(j)`$; in the real-space basis $`𝐒^z(j)=𝐧(j)1/2`$, where $`𝐧(j)=|jj|`$. This spectral sum can also be evaluated within our RG approach. The leading-order operator renormalizations in this case are, in complete analogy with the spin problem, very simple: each $`𝐒^z(j)`$ remains unchanged unless a state living on site $`j`$ is eliminated, in which case $`𝐒^z(j)`$ renormalizes to a multiple of the identity. As before, we run the RG till the bandwidth is reduced to $`\mathrm{\Omega }_{\mathrm{final}}=\omega `$ and do the spectral sum with the renormalized operators in the new problem. This renormalized sum may be evaluated by again recognizing that it is dominated by transitions between pairs of states with energies $`\pm \omega /2`$ that live on pairs of sites connected by ‘strong’ hopping amplitudes (of magnitude $`\omega /2`$) in the renormalized problem. The corresponding matrix element is just $`(1e^{ik\stackrel{~}{l}})/2`$, where $`\stackrel{~}{l}`$ is the length of the hop in question. Counting the contributions exactly as in our zero-temperature calculations, we thus get
$`S^{zz}(k,\omega )={\displaystyle \frac{1}{(1+e^{\omega /2T})^2}}\times S_{T=0}^{zz}(k,\omega ).`$ (118)
Thus, we see that $`S^{zz}`$ is essentially unaffected by thermal fluctuations at the XX point; in particular, the low-frequency divergence is not cut-off by temperature effects even when $`\omega T`$.
A similar analysis can clearly be performed in the XX-RD phase. Again, both $`\sigma ^{}(\omega )`$ and $`S^{zz}(k,\omega )`$ at $`T>0`$ are simply given by the corresponding expressions at $`T=0`$ multiplied by simple functions of $`\omega /T`$, exactly as in Eqs. (115) and (118). Thus, though temperature effects are simple to work out at the XX point and in the XX-RD phase, the results are rather special due to the free-fermion character of the problem.
### C Going beyond the free-fermion results
What happens when we turn on the nearest-neighbor interaction? This is the question we need to address next.
Let us first consider the effects of small $`J^z`$ couplings added on to the XX model. The analysis of Ref. shows that this term is irrelevant in the RG sense; the typical value of $`\stackrel{~}{J}^z/\stackrel{~}{J}^{}`$ at log-cutoff $`\mathrm{\Gamma }`$ scales as $`(\stackrel{~}{J}^z/\stackrel{~}{J}^{})u_0\mathrm{exp}(c\mathrm{\Gamma }^\varphi )`$, where $`c`$ is an $`O(1)`$ constant, $`\varphi `$ is the golden mean $`(1+\sqrt{5})/2`$, and $`u_0`$ is the typical value of $`J^z/J^{}`$ in the microscopic model. A useful way of thinking about the low-frequency behavior of the conductivity is as follows: Imagine running the RG till the cutoff $`\mathrm{\Omega }T`$. In this renormalized problem the typical $`(\stackrel{~}{J}^z/\stackrel{~}{J}^{})u_0\mathrm{exp}(c\mathrm{\Gamma }_T^\varphi )`$, where $`\mathrm{\Gamma }_T\mathrm{ln}(\mathrm{\Omega }_0/T)`$. In the fermion language, this is the typical value of the ratio of the nearest neighbor interactions to the hopping amplitudes. In this renormalized problem, a naive Fermi’s Golden Rule estimate of the corresponding inelastic collision rate due to interactions gives $`1/\tau _{\mathrm{coll}}u_0^2T\mathrm{exp}(2c\mathrm{\Gamma }_T^\varphi )`$. This gives us a frequency scale below which our $`T0`$ free-fermion results are expected to break down as a result of the residual interaction effects.
Unfortunately, we are unable to do a controlled calculation that determines the transport properties in the frequency regime $`\omega <1/\tau _{\mathrm{coll}}`$. The best we can do is to work out what a naive scaling argument would predict for the d.c. limit of the conductivity. The basic idea is as follows: The collision rate may be converted into a corresponding dephasing length $`L_{\mathrm{coll}}`$ by appealing to the activated scaling that is a characteristic of our problem. This gives $`L_{\mathrm{coll}}\mathrm{ln}^2(\tau _{\mathrm{coll}})\mathrm{\Gamma }_T^{2\varphi }`$. This is the length scale beyond which quantum coherence is lost due to inelastic collisions. Now, we can imagine breaking up the system into blocks of length $`L_{\mathrm{coll}}`$. A d.c. current $`I`$ passing through the system will see a chain of resistors corresponding to these blocks—the resistance values of each of these blocks is simply given by the $`T=0`$ Landauer resistance of the corresponding system of length $`L_{\mathrm{coll}}`$. The voltage developed across a system of total length $`L`$ will therefore be $`V=I_{j=1}^{L/L_{\mathrm{coll}}}R_j=ILR_{\mathrm{av}}(L_{\mathrm{coll}})/L_{\mathrm{coll}}`$. Since $`R_{\mathrm{av}}(L_{\mathrm{coll}})e^{c_1L_{\mathrm{coll}}}`$ (where $`c_1`$ is an $`O(1)`$ scale factor), the d.c. conductivity works out to be $`\sigma _{\mathrm{d}.\mathrm{c}.}L_{\mathrm{coll}}e^{c_1L_{\mathrm{coll}}}\mathrm{ln}^{2\varphi }(\mathrm{\Omega }_0/T)e^{c^{}\mathrm{ln}^{2\varphi }(\mathrm{\Omega }_0/T)}`$. Note that in the absence of interactions, we had earlier found $`\sigma ^{}(\omega )0`$ as $`\omega 0`$ at $`T>0`$—our scaling argument implies that interactions render this conclusion invalid. Unfortunately, while this scaling argument is certainly plausible, the question of the true low-frequency limit can only be settled by a controlled calculation in the regime $`\omega 1/\tau _{\mathrm{coll}}`$, which is beyond our current capabilities.
The above arguments also suggest that $`S^{zz}(k,\omega )`$ will deviate from the $`T0`$ free-fermion result for $`\omega <1/\tau _{\mathrm{coll}}`$. In particular, one expects that the $`\omega 0`$ divergence of $`S^{zz}`$ would be cut-off below this frequency scale. Similar behavior is also expected of $`S^+`$, but again, what is really needed is a controlled calculation as opposed to a scaling argument. Note also that we expect something different at the XXX and XXZC quantum critical points: since the theory at these critical points already includes interactions, one expects that $`1/\tau _{\mathrm{coll}}T`$, and the relaxational behavior characteristic of an interacting system at finite temperature will set in for $`\omega T`$, in contrast to the behavior in the vicinity of the XX point.
## VII Prospects for experimental tests
Previous experimental work on one-dimensional random-exchange Heisenberg antiferromagnetic spin chains has characterized the dynamics of these systems in terms of the observed NMR $`1/T_1`$ relaxation rate, and ESR relaxation rates and linewidths.
As far as the NMR measurements are concerned, our calculations are unfortunately not directly relevant to the experimental measurements of $`1/T_1`$. This may be seen as follows: In the usual case of a pure, translationally invariant system, $`1/T_1`$ is directly related, by Fermi’s Golden Rule, to the local dynamic structure factor $`S_{\mathrm{loc}}`$ evaluated at frequency $`\omega `$ equal to the nuclear resonance frequency $`\gamma _NH`$, where $`\gamma _N`$ is the nuclear magnetic moment and $`H`$ is the external field. In a random system, with a broad variation in the value of $`S_{\mathrm{loc}}(\omega )`$, the following question immediately arises: what measure of the distribution of $`S_{\mathrm{loc}}(\omega )`$ does the experimentally measured $`1/T_1(H)`$ reflect?
Now, we have seen that the average $`S_{\mathrm{loc}}(\omega )`$ diverges strongly as $`\omega 0`$ at $`T=0`$. Naively, one might have thought that this would imply a corresponding divergence in $`1/T_1`$ at small $`H`$, at least when $`TH`$. However, the divergence in $`S_{\mathrm{loc}}(\omega )`$ comes from a few very rare sites which give a very large contribution. Clearly, the observed $`1/T_1`$ will be completely insensitive to this effect, since all that will happen is that a tiny fraction of nuclear spins (in the neighborhoods of those rare electron spins that have significant spin fluctuations at the frequency $`\omega =\gamma _NH`$) will flip almost instantaneously, while the rest of the nuclear spins will have an extremely small probability to flip, and this is what will be reflected in the spin relaxation experiments. In this sense, it is the typical value of $`S_{\mathrm{loc}}(\omega )`$ that is more relevant for comparisons with NMR $`1/T_1`$ data. A typical nuclear spin will in fact have essentially no spin fluctuations to couple to at frequency $`\omega =\gamma _NH`$—it can therefore relax only by paying an activation energy that is set by $`\gamma _eH`$ (where $`\gamma _e\gamma _N`$ is the electron magnetic moment) since the external field acts to freeze out all modes below this energy scale in most of the system (with the exception of the rare regions alluded to above). The experiments actually do see activated behavior for $`1/T_1`$ at finite temperature. However, the activation gap seems to scale as $`\mathrm{\Delta }H^{1.6}`$—the rough argument above of course cannot explain this non-trivial $`H`$ dependence of the observed activation energy.
Our second remark relates to the ESR linewidth measurements of Clark and Tippie. Here, again, our results do not address the experimentally relevant questions. This is because all our calculations for the XXX case are done within the context of the simple Heisenberg exchange Hamiltonian, while the observed linewidth in the experiments is determined by other effects such as dipolar interactions or anisotropy.
Inelastic neutron scattering experiments, on the other hand, if they can be done on these systems, provide a direct testing ground for our predictions. We conclude with some remarks on the relevance of our calculations of the dynamic structure factor to such experiments. First of all, note that we considered randomness in the exchanges only, with the spins themselves assumed positioned on regular lattice sites; thus, our results are restricted to compounds with chemical disorder in exchanges. It is clear that small randomness in the positions of spins (e.g., due to thermal fluctuations) will result only in some further suppression (by the standard Debye-Waller factor at wavevector $`k`$) of the features relative to an overall background. In the dimer phase, a possible difference in the lengths of even and odd bonds will result only in some phase factor in the cosine of Eq. (31). Also note that the non-magnetic neutron scattering from such spin chains will actually be suppressed near $`k=\pi /a`$, and this may facilitate a possible experimental observation of the predicted features. We caution, however, that while it would be extremely interesting to see the sharp oscillatory structure predicted in the Griffiths phases, this may be difficult to achieve without going to very low temperatures and energy transfers. Regarding transport, we hope that our results will motivate experiments to probe the spin conductivity in these systems.
## VIII Acknowledgements
We thank P. W. Anderson, R. N. Bhatt, D. Dhar, D. S. Fisher, F. D. M Haldane, M. Hastings, A. Millis, A. Madhav, R. Moessner, S. Sachdev, T. Senthil, S. L. Sondhi, and A. Vishwanath for useful discussions. This work was supported by NSF grants DMR-9809483 and DMR-9802468.
## A Finite-size scaling function for the conductivity
Consider a finite chain with an even number of sites $`N_I=L+1`$, where $`L`$ is the length of the chain, and with free boundary conditions (a similar analysis can be carried out for chains with an odd number of sites, and also for chains with periodic boundary conditions). We want to calculate the real part of the dynamical conductivity averaged over the distribution of bond strengths in the limit of low frequencies and large $`L`$. We work this out for the XX chain; the result in the presence of $`J^z`$ couplings will differ only in the values of some non-universal scale factors so long as the system does not develop Ising antiferromagnetic order in the thermodynamic limit. To proceed, we need to keep track of the joint distribution at scale $`\mathrm{\Gamma }`$ of the number of remaining spins $`N`$, the $`N1`$ couplings $`\zeta _i`$, and the corresponding bond lengths $`l_i`$. In a finite system, the couplings become correlated due to the constraint imposed by the finite length of the system. However, following Fisher and Young, we note that the couplings remain ‘quasi-independent’, and can be described in terms of the infinite-chain distribution $`P(\zeta ,l|\mathrm{\Gamma })`$ exactly as in Ref. . More precisely, if we also keep track of the lengths $`l_F`$ and $`l_R`$ of the ‘dead’ regions (consisting of singlet pairs formed at higher energy scales) at the left and right ends of the chain, then a distribution of the form
$`d\mathrm{Prob}`$ $`[N;\zeta _1,l_1\mathrm{}\zeta _{N1},l_{N1};l_F,l_R|L,\mathrm{\Gamma }]=`$ (A1)
$`=`$ $`a_N(L|\mathrm{\Gamma })P(\zeta _1,l_1)d\zeta _1\mathrm{}P(\zeta _{N1},l_{N1})d\zeta _{N1}`$ (A3)
$`\times (l_F)(l_R)\delta _{l_1+\mathrm{}+l_{N1}+l_F+l_R,L}`$
for even $`N2`$ has its from preserved under renormalization if $`a_N(L|\mathrm{\Gamma })`$ is independent of $`N`$ with
$$\frac{1}{a}\frac{da}{d\mathrm{\Gamma }}=2P_0(\mathrm{\Gamma })=2𝑑lP(0,l).$$
(A4)
Here, $`P(\zeta ,l|\mathrm{\Gamma })`$ satisfies the same flow equation as in the infinite chain, and $`(l|\mathrm{\Gamma })`$ satisfies
$$\frac{}{\mathrm{\Gamma }}=()_lP(0,)_l_0^{\mathrm{}}P(\zeta ,)𝑑\zeta P_0.$$
(A5)
In the above, the $`\mathrm{\Gamma }`$ dependence is left implicit, and $`f()_lg()`$ is used to denote a (discrete) convolution in the length variables. For clarity, we work explicitly with discrete lengths, with $`l_F`$ and $`l_R`$ even, and $`l_i`$ odd integers; this is clearly preserved under the RG.
We start the RG with $`\mathrm{\Omega }_0=1`$, $`\mathrm{\Gamma }_I=0`$, the initial bond distribution $`P(\zeta |\mathrm{\Gamma }_I)=e^\zeta `$ (this corresponds simply to choosing the initial $`J^{}`$ to be uniformly distributed in the interval $`[0,1]`$), $`l_i=1`$, $`l_F=l_R=0`$, and $`N_IL+1`$; with initial distributions $`P(\zeta ,l|\mathrm{\Gamma }_I)`$ and $`(l|\mathrm{\Gamma })`$ normalized to unity, the normalization factor is $`a(\mathrm{\Gamma }_I)=1`$. The dynamical conductivity is now given by
$$\sigma ^{}(\omega ,L)=\frac{1}{4}\frac{a(\mathrm{\Gamma }_\omega )}{L}A(L|\mathrm{\Gamma }_\omega ),$$
(A6)
with $`a(\mathrm{\Gamma })=(\mathrm{\Gamma }+1)^2`$ (for our specific choice of initial conditions), and
$`A(L|\mathrm{\Gamma })`$ $`=`$ $`{\displaystyle \underset{N=2}{\overset{L+1}{}}}{}_{}{}^{}(N1){\displaystyle \underset{l_1,l_2,\mathrm{}l_{N1},l_F,l_R}{}}P(0,l_1)l_1^2`$ (A9)
$`\times {\displaystyle _0^{\mathrm{}}}P(\zeta _2,l_2)d\zeta _2\mathrm{}{\displaystyle _0^{\mathrm{}}}P(\zeta _{N1},l_{N1})d\zeta _{N1}`$
$`\times (l_F)(l_R)\delta _{l_1+l_2+\mathrm{}+l_{N1}+l_F+l_R,L},`$
where the sum is over even $`N`$.
Now, multiplying $`A(L|\mathrm{\Gamma })`$ by $`e^{yL}`$ and summing over odd $`L1`$—i.e., doing a (discrete) Laplace transform in $`L`$—removes the constraint on the lengths, and we find
$$A(y|\mathrm{\Gamma })=^2(y)Q(y)\frac{1+T^2(y)}{(1T^2(y))^2},$$
(A10)
where $`Q(y)`$ and $`T(y)`$ are respectively the Laplace transforms of $`P(0,l)l^2`$ and $`_0^{\mathrm{}}P(\zeta ,l)𝑑\zeta `$. Thus, we can straightforwardly work out $`A(y)`$, given $`P(\zeta ,y)`$ and $`(y)`$. Using the results of Refs. and , we can write the following explicit expressions for these two functions:
$`P(\zeta ,y|\mathrm{\Gamma })`$ $`=`$ $`Y(y|\mathrm{\Gamma })e^{\zeta u(y|\mathrm{\Gamma })},`$ (A11)
$`(y|\mathrm{\Gamma })`$ $`=`$ $`{\displaystyle \frac{u(0|\mathrm{\Gamma })u(y|\mathrm{\Gamma }_I)}{u(y|\mathrm{\Gamma })u(0|\mathrm{\Gamma }_I)}}(y|\mathrm{\Gamma }_I),`$ (A12)
where $`u(y|\mathrm{\Gamma })=D(y)\mathrm{coth}[D(y)(\mathrm{\Gamma }+C(y))]`$ and $`Y(y|\mathrm{\Gamma })=D(y)/\mathrm{sinh}[D(y)(\mathrm{\Gamma }+C(y))]`$. The functions $`D(y)`$ and $`C(y)`$ depend on the initial distribution $`P(\zeta ,y|\mathrm{\Gamma }_I)`$, and in our case are given by $`D(y)=\sqrt{1e^{2y}}`$ and $`D(y)C(y)=y+\mathrm{ln}(1+\sqrt{1e^{2y}})`$. Also, $`(y|\mathrm{\Gamma }_I)=1`$.
With this in hand, it is a relatively simple matter to work out $`A(L|\mathrm{\Gamma })`$, $`L`$ odd, by performing the inverse Laplace transform:
$$A(L)=\frac{1}{\pi i}_{ci\pi /2}^{c+i\pi /2}A(y)e^{yL}𝑑y.$$
(A13)
In Sec. III D 5, we evaluated this integral numerically to compare the RG predictions with the results of the exact-diagonalization studies (note that in the main body of the paper we didn’t make a distinction between $`N_I`$ and $`L`$, since it is irrelevant in the thermodynamic limit; in the more detailed notation of this Appendix, our numerical results of Sec. III D 5 are for system sizes $`N_I=128`$ and $`256`$).
In the scaling limit $`\mathrm{\Gamma }1`$, $`L1`$, the integral Eq. (A13) for $`A(L|\mathrm{\Gamma })`$ is dominated by small $`y`$ and can be approximated by
$$A(L|\mathrm{\Gamma })=2\mathrm{L}T^1A(y|\mathrm{\Gamma }),$$
(A14)
where $`\mathrm{L}T^1`$ denotes the inverse of the continuous Laplace transform. Moreover, in this limit, $`A(y)`$ may be worked out using the following scaling forms for $`(y)`$ and $`P(\zeta ,y)`$:
$`(y|\mathrm{\Gamma })`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }\sqrt{2y}\mathrm{coth}[\mathrm{\Gamma }\sqrt{2y}]}},`$ (A15)
$`P(\zeta ,y|\mathrm{\Gamma })`$ $`=`$ $`{\displaystyle \frac{\sqrt{2y}}{\mathrm{sinh}[\mathrm{\Gamma }\sqrt{2y}]}}e^{\zeta \sqrt{2y}\mathrm{coth}[\mathrm{\Gamma }\sqrt{2y}]}.`$ (A16)
Putting everything together, we can now write $`A(L,\mathrm{\Gamma })=\mathrm{\Gamma }f(\mathrm{\Gamma }^2/L)`$, which immediately implies a scaling form for the conductivity: $`\sigma ^{}(\omega ,L)=\mathrm{\Gamma }_\omega \mathrm{\Theta }(\mathrm{\Gamma }_\omega ^2/L)`$. Thus, we see that the dynamical conductivity in a finite system satisfies a scaling form that reflects the activated dynamical scaling at the XX fixed point. Note that while the scaling form holds more generally, the values of the non-universal scale factors that we have used are specific to our choice of initial distribution. Analyzing the behavior starting with an arbitrary initial distribution allows one to relate these non-universal scale factors to the properties of the initial distribution under the assumption that ‘bad decimations’ early in the RG do not affect these values. Such an analysis allows us to write $`\sigma ^{}(\omega ,L)=l_v\mathrm{ln}(\mathrm{\Omega }_0/\omega )\mathrm{\Theta }\left(l_v\mathrm{ln}^2(\mathrm{\Omega }_0/\omega )/L\right)`$, with the microscopic energy scale $`\mathrm{\Omega }_0`$ and the microscopic length scale $`ł_v`$ precisely as defined in the main text. Moreover, it is clear that the same scaling function also describes the low-frequency dynamical conductivity in a large but finite system even in the presence of $`J^z`$ interactions as long as the system is in a Random Singlet state—only the values of the non-universal scale factors are expected to change.
While it is possible to calculate the full scaling function $`\mathrm{\Theta }(x)`$ by a detailed analysis of the inverse Laplace transform, we will confine ourselves here to working out $`\mathrm{\Theta }(x)`$ in two limiting cases: For $`x1`$, $`\mathrm{\Theta }(x)=7/180`$ (which correctly reproduces the infinite-size result, as it must), while in the limit $`x1`$ we have $`\mathrm{\Theta }(x)=e^{x/2}/\sqrt{2\pi x}`$. This is the result used in the Griffiths argument of Sec. III D.
## B Low-energy eigenfunctions of the one-dimensional random hopping Hamiltonian
We have already noted on several occasions that subleading terms in our evaluation of low-energy properties via the strong-randomness RG can be treated consistently within the RG picture, and are naturally related to the ground state correlations in the system. Here we discuss more explicitly one example of this connection, for the case of the XX model; our discussion is couched completely in the free-fermion language and uses the fermionic RG introduced in Sect. VI A—this will allow us to work in a more familiar context and to connect with the known results for the fermionic problem.
Our example relates to the structure of the low-energy eigenstates of the one-dimensional random hopping Hamiltonian $`𝐇=_{j=1}^{L1}t_j(|j+1j|+|jj+1|)`$; for simplicity, we work with free boundary conditions. We note at the very outset that the average Greens functions $`[G(x,x^{};E)]_{\mathrm{av}}`$ have been calculated exactly in Ref. using SUSY techniques. Moreover, a great deal is known about the statistics of the low-energy wavefunctions, via a sophisticated real space analysis. Here, we present a simple (if somewhat rough) picture of the wavefunctions which allows one to easily reproduce the previously published SUSY results, as well as provides us with the following rough picture of a typical pair of wavefunctions at energies $`\pm ϵ`$: such a pair is typically “localized” on two sites $`i`$ and $`j`$ separated by a distance of order $`\mathrm{\Gamma }_ϵ^2`$, $`\mathrm{\Psi }_{\pm ϵ}=𝒩(|i\pm |j)+\mathrm{}`$ with $`𝒩`$ of order 1, with the wavefunction amplitudes decaying as $`e^{c\sqrt{r}}`$ away from the two sites for $`r\mathrm{\Gamma }_ϵ^2`$ and as $`e^{c^{}r/\mathrm{\Gamma }_ϵ}`$ for $`r\mathrm{\Gamma }_ϵ^2`$. Moreover, as we shall see below, the above picture is intimately related to corresponding structure in the zero-energy wavefunction.
We start by considering this zero-energy wavefunction. If the number of sites $`L`$ is odd, the Hamiltonian $`𝐇`$ has a (unique) zero-energy eigenstate $`\psi `$, and this eigenstate can be written down immediately as $`\psi _{2n}0`$ and $`\psi _{2n+1}=\mathrm{\Pi }_{j=1}^n(t_{2j1}/t_{2j})\psi _1`$. Thus, the logarithm of the wavefunction $`\mathrm{ln}|\psi _{2n+1}|`$ performs a random walk (see Fig. 8). While the relative amplitudes of the wavefunction are thus understood directly, we are interested in the properties of the normalized zero-energy wavefunction $`\mathrm{\Psi }`$, that is, we also need to understand normalization of the eigenvector $`\psi `$. Such properties of the normalized wavefunction $`\mathrm{\Psi }`$ were analyzed by Balents and Fisher by reducing the calculation of correlators to the Liouville quantum mechanics. Here we suggest a more direct approach by noting that normalization of the $`\psi `$ for a given realization can be achieved by simply setting the largest amplitude $`\psi _{\mathrm{max}}=\mathrm{max}\{\psi _{2n+1}\}`$ to a number of order one. More precisely, if we set $`\psi _{\mathrm{max}}=1`$, then for any $`q>0`$ the norm $`𝒩_{2q}=_{j=1}^L|\psi _j|^{2q}`$ is a random quantity of order one with a well-behaved (“narrow”) distribution of width of order one. This can be seen directly by, for example, calculating first moments of the distribution of $`𝒩_q`$ using the properties of such a random walk near its global maximum (absorbing boundary). With this in hand, we can characterize the zero-energy wavefunction essentially quantitatively even without the knowledge of the right normalization constant, which is always of order one.
Let us now characterize the pairs of wavefunctions at energies $`\pm ϵ`$. Consider running the fermion RG on a given system; the concrete disorder realization can viewed as a random-walk pattern, while the basic RG step is a transformation on this pattern that removes the smallest-scale structure and keeps the rest of the structure unchanged—this is shown in Fig. 8 (note that this is similar to the RG approach of Ref. to the problem of random walks in random environments). In this picture, the sites that remain down to the energy scale $`\mathrm{\Gamma }_ϵ`$ are completely specified from the initial random-walk pattern as follows: an even site $`A`$ survives down to this energy if and only if the “walks” to the right and to the left from the site $`A`$ do not return the same level (of the site $`A`$) until after they crossed the level $`\mathrm{\Gamma }_ϵ`$ higher than the level of the site $`A`$—see Fig. 8; similar conditions can be formulated for odd sites. Moreover, this also gives a corresponding criterion for the pairs of sites on which our wavefunctions at energies $`\pm ϵ`$ are localized (within the zeroth order RG picture for these wavefunctions). Thus, as far as the zeroth order RG picture of these wavefunctions is concerned, one is faced with the problem of characterizing such extremal properties of random walks, and the RG itself can be viewed as a powerful tool for doing this, even though many questions like the number density of remaining sites or the distribution of the effective time steps (i.e., bond lengths) and heights (i.e., bond log-couplings) can be answered directly in the random-walk language by appealing to the properties of a random walk near its minima (absorbing walls).
So far, we have only reformulated our leading (zeroth) order result for the wavefunctions at $`\pm ϵ`$ in more visual terms. The real usefulness of this picture appears when we go beyond the leading order, and keep track of operators “measuring” the wavefunction amplitude throughout the system. Consider going beyond the leading order for the renormalization of the “participation” operators $`𝐧_j=|jj|`$. When we eliminate sites $`2`$ and $`3`$, the renormalization of the corresponding operators including the leading second-order correction to the zeroth order result is given as
$`(𝐧_2)^{\mathrm{eff}}`$ $`=`$ $`(t_3/t_2)^2𝐧_{\stackrel{~}{4}},`$
$`(𝐧_3)^{\mathrm{eff}}`$ $`=`$ $`(t_1/t_2)^2𝐧_{\stackrel{~}{1}}.`$
It is now easy to see that if at a given stage of the RG the site $`A^{}`$ is no longer active, it will have component only onto the nearest active site $`A`$ on the same sublattice, with the magnitude of the “participation” given by $`e^{2\mathrm{\Lambda }}`$, where $`\mathrm{\Lambda }`$ is the absolute “height” difference of the random walker at “times” corresponding to the two sites. The connection to the statistics of the zero energy wavefunction is now explicit. Moreover, our picture of the wavefunctions at energy $`ϵ`$ now follows: the dominant feature of such wavefunctions $`\psi _{\pm ϵ}`$ are two sites, with order one amplitudes of the $`\psi _{\pm ϵ}`$, separated by a distance of order $`\mathrm{\Gamma }_ϵ^2`$; the wavefunctions will have significant weight only near the two “peaks” and, to this order, only on the sites of the same sublattice as the nearby peak, i.e., the wavefunction can be thought of as being composed of two bumps centered on the two sites. As we move away from the two peaks by a distance $`r\mathrm{\Gamma }_ϵ^2`$, the amplitude of the wavefunction is given by the above random-walk prescription and is typically decreasing as $`e^{c\sqrt{r}}`$; for still larger distances the amplitudes are given by high-order perturbation theory, with the perturbative expression for $`\psi (r)`$ necessarily “stepping” through the “active” segments of lengths $`\mathrm{\Gamma }_ϵ^2`$ separating $`r`$ from the two bumps and acquiring a suppression of order $`e^{c\mathrm{\Gamma }_ϵ}`$ from each such segment—thus, for $`r\mathrm{\Gamma }_ϵ^2`$ we will typically have the decay $`\psi (r)e^{c^{}r/\mathrm{\Gamma }_ϵ}`$ (note that this gives us an interpretation, in terms of the typical wavefunction at energy $`ϵ`$, of the two length scales—one proportional to $`\mathrm{\Gamma }_ϵ`$ and the other to $`\mathrm{\Gamma }_ϵ^2`$—that diverge in the limit of low $`ϵ`$, and are usually identified with the typical and average localization length respectively).
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# Tau Neutrinos Underground: Signals of 𝜈_𝜇→𝜈_𝜏 Oscillations with Extragalactic Neutrinos
## I Introduction
A recent breakthrough in the study of neutrino oscillations came from the observation by the Super-Kamiokande experiment of a deficit of upward-going atmospheric muon neutrinos . The observed electron neutrino flux was found to be consistent with the theoretical expectation from models of cosmic ray production of neutrinos. Furthermore, SuperK measurements are consistent with earlier experiments which detected anomalous ratios of the $`\nu _\mu `$ to $`\nu _e`$ flux. The new high-statistics data disfavor scenarios in which $`\nu _\mu `$’s oscillate into sterile neutrinos ($`\nu _s`$) , and the data are consistent with $`\nu _\mu `$ to $`\nu _\tau `$ oscillation (99$`\%`$ CL) with a large mixing angle, $`\mathrm{sin}^2\theta >0.84`$ and a neutrino mass squared difference of $`2\times 10^3`$ eV$`{}_{}{}^{2}<\mathrm{\Delta }m^2<6\times 10^3`$ eV<sup>2</sup>.
Direct detection of $`\nu _\tau `$ appearance is extremely difficult because at low energies, the charged-current cross section for producing a tau is small and the tau has a very short lifetime. Several long-baseline experiments with accelerator sources of $`\nu _\mu `$ have been proposed with the goal of detecting tau neutrinos from oscillations, thus confirming the SuperK results. The only convincing evidence of neutrino oscillations to date involves astrophysical sources, neutrinos from the sun and atmospheric neutrinos. These observations involve indirect measurements, namely the disappearance of the expected neutrino fluxes.
We have recently discussed the possibility of using a kilometer-size neutrino telescope to detect tau neutrinos from extragalactic sources of high-energy neutrinos such as Active Galactic Nuclei (AGN) and Gamma Ray Bursts (GRB), assuming $`\nu _\mu \nu _\tau `$ with the oscillation parameters of the SuperK experiment . The probability for $`\nu _\mu \nu _\tau `$ is given by
$$P(\nu _\mu \nu _\tau ;L)=\mathrm{sin}^22\theta \mathrm{sin}{}_{}{}^{2}\left(\frac{1.27\mathrm{\Delta }\mathrm{m}^2(\mathrm{eV}^2)\mathrm{L}(\mathrm{km})}{\mathrm{E}(\mathrm{GeV})}\right).$$
(1)
Assuming two flavor oscillations, muon neutrinos produced in AGN or GRB would oscillate to tau neutrinos as they travel to the Earth. Over astronomical distances in the range of a megaparsec to thousands of megaparsecs, by measuring tau neutrino fluxes, one could, in principle, probe oscillations down to $`\mathrm{\Delta }m^210^{17}`$ eV<sup>2</sup>, nine orders of magnitude below current neutrino experiments . On the other hand, for the SuperK parameter range, with $`\mathrm{\Delta }m^2`$ on the order of $`10^3`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta 1`$, the oscillation probability is 0.5. It is this latter possibility that we explore in this paper.
We use the simplest assumption for the flavor content of extragalactic sources of neutrinos, in the absence of oscillations, for the ratio $`\nu _e:\nu _\mu :\nu _\tau `$ to be $`1:2:0`$. This is based on a counting argument applied to $`\pi \mu \nu _\mu `$ and $`\mu \nu _\mu e\nu _e`$ processes. With the two-flavor oscillations suggested by the SuperK experiment, the flavor ratio becomes $`1:1:1`$ after the neutrinos have traveled over astronomical distances. Even in the three-flavor oscillation scenario, the ratio is still $`1:1:1`$, because the path length of high energy extragalactic neutrinos is much larger than any neutrino oscillation length supported by the solar, atmospheric or accelerator data .
The ratio for $`\nu _e:\nu _\mu `$ might get modified at high energies due to muon cooling . In addition, $`\nu _e`$ from neutron decay might give significant contribution, resulting in neutrino fluxes dominated by electron neutrinos as in the case of diffuse neutrino fluxes from propagating cosmic rays . We comment qualitatively in the discussion section on how our results are altered with more realistic, flavor-dependent neutrino energy cutoffs. Regardless of the flavor content of the source, the maximal mixing suggested by the SuperK experimental results mean that there will be an appreciable tau neutrino component at the Earth, so one is interested in tau neutrino detection in high energy neutrino telescopes such as ANTARES , NESTOR , AMANDA and the next generation of large underground detectors .
Tau neutrino detection requires an understanding of the effect of propagating through the Earth on the tau neutrino flux. The propagation of ultra-high energy tau neutrinos through the Earth is quite different from muon and electron neutrinos. The Earth never becomes opaque to tau neutrinos, while muon and electron neutrinos are absorbed via charged-current interactions before reaching the opposite surface . Ultrahigh-energy tau neutrinos interact in the Earth producing taus which, due to the short lifetime, decay back into tau neutrinos with lower energy. This cascade continues until the tau neutrinos reach the detector on the opposite side of the Earth or until the energy of the neutrinos is small enough that the interaction length of the neutrino is longer than the path length through the Earth. The energy and nadir angle dependence of the extragalactic tau neutrinos fluxes have been examined quantitatively in Refs. . For certain fluxes, those that do not decrease too steeply with energy, there are significant enhancements of the tau neutrino flux relative to the muon neutrino flux at energies below $`10^6`$ GeV.
An enhancement of the tau neutrino flux does not necessarily translate to dramatic modifications of the standard model (no-oscillation) rates for upward-going muons, especially in view of uncertainties in the normalization of the extragalactic fluxes. However, by comparing rates for upward-going muons with rates for upward hadronic/electromagnetic (EM) showers, the signature of tau neutrino interactions is unambiguous for a large range of neutrino fluxes. In the next section, we briefly introduce the extragalactic and generic $`\nu _\tau `$ fluxes $`F_{\nu _\tau }^oE^n`$ for $`n=1,2`$ that are used here. After reviewing neutrino propagation through the Earth, we describe $`\nu _\tau `$ signatures. The fluxes considered here have a range of energy behaviors. Even if the normalizations of the neutrino fluxes are uncertain, and in some cases optimistic, it is useful to make quantitative comparisons of the event rates for upward muons and upward hadronic/EM showers, with and without neutrino oscillations, which we do in Section IV. The quantitative results for specific models lead to model independent conclusions, which we summarize graphically. Tau neutrino appearance would provide an independent confirmation of the SuperK results and would point towards the better understanding of physics beyond the Standard Model.
## II High Energy Neutrinos Sources
Active Galactic Nuclei are the most luminous objects in the Universe. Most of this radiation comes from their central region, indicating that the energy radiation most likely comes from accretion of matter into a superheavy black hole. Protons within the AGN may get accelerated via first order Fermi acceleration to very high energies. They interact with protons and photons in the infalling gas, or they may exist in the jets along the rotation axis and interact with photons there. Photon-proton and proton-proton interactions produce pions, which decay into charged leptons, neutrinos and photons. Energetic photons ($`E_\gamma 100`$ MeV) from about 40 AGN observed by the EGRET collaboration , and TeV photons have been detected from Mkn 421, Mkn 501 and 1ES2344+514 . Although these photons are conventionally explained by inverse Compton scattering from energetic electrons, this explanation is not without problems, and a hadronic origin of gamma-ray photons from AGN is a viable alternative . If a large fraction of the observed energy in high energy photons from AGN is produced in hadronic interactions, then AGN are also powerful sources of ultrahigh-energy (UHE) neutrinos .
In Fig. 1 we show neutrino fluxes predicted in the AGN models of Stecker and Salamon and Mannheim Model A . Both of these models predict neutrinos fluxes that represent the upper bounds for their class of the models. In particular, the Stecker-Salamon flux is an upper bound for AGN core emission, while Mannheim Model A is an upper bound for AGN jet emission models. Stecker-Salamon flux is bound by the the diffuse X-ray background, while Mannheim flux is bound by the extragalactic gamma ray background. The steep, low energy neutrino flux in Mannheim’s model is the emission from the host galaxy via $`pp`$ interactions of the AGN protons in the galactic gas disk. Since this part of the flux is derived with the assumption that all protons end up in the disk, it should be regarded as an upper bound. Stecker-Salamon flux at energies above 1 PeV may get reduced due to the cooling of pions and muons in the strong magnetic fields of AGN cores . The fluxes plotted in Fig. 1 are for the sum of muon neutrino plus antineutrino, at the source, namely, without accounting for oscillation over astronomical distances. We label the fluxes in the absence of oscillations by $`F_{\nu +\overline{\nu }}^s`$.
Another extragalactic source with powerful radiation and possibly associated high energy neutrino flux are the gamma ray bursts (GRB). Several models have been proposed in order to explain the origin of GRB’s . In the fireball model , the gamma ray bursts are produced by the dissipation of the kinetic energy of the relativistic expanding fireball with a large fraction, $`>10\%`$, of the fireball energy being converted by photopion production to high energy neutrinos . Photomeson production takes place when extremely energetic protons accelerated at high energies in the ultra-relativistic shocks interact with synchrotron photons inside the fireball. The decay of these charged pions and subsequently produced muons then produce electron and muon neutrinos. Contributions from proton-proton collision can be neglected in this model. In Fig. 1 we show the neutrino fluxes for gamma ray burst model of Waxman and Bahcall (GRB\_WB), in which they parameterize the flux by
$`F_{\nu +\overline{\nu }}^s(E)=4.0\times 10^\alpha E^n,`$ (2)
where $`\alpha =13`$ and $`n=1`$ for $`E<10^5`$ GeV and $`\alpha =8`$, $`n=2`$ for $`10^5<E<10^7`$ GeV and $`\alpha =1`$, $`n=3`$ for $`E>10^7`$ GeV.
Theoretical work has been done to set upper bounds on high energy neutrino fluxes from AGN jets and GRB . The bounds are based on the theoretical correlations between the cosmic ray flux and/or the extragalactic gamma ray flux and the neutrino flux. These bounds have some model dependence, and they tend to be weaker in the range of energies considered here than at higher energies ($`E10^710^9`$ GeV). The AGN and GRB neutrino fluxes used here satisfy these bounds.
Cosmic topological defects (TD) such as magnetic monopoles, cosmic strings and domain walls are predicted to be formed in the Early Universe as a result of symmetry breaking and phase transition in Grand Unified Theories (GUTs) of particle interactions. In the TD models, $`\gamma `$-rays, electrons (positrons), and neutrinos are produced directly at ultra-high energies via cascades initiated by the decay of a supermassive elementary “X” particle associated with some Grand Unified Theory, rather than being produced in high energy hadronic interactions. The X particle is usually thought to be released from topological monopoles left over from GUT phase transition. It decays into quarks, gluons, leptons. In this paper, we consider neutrino fluxes from topological defects models of Sigl-Lee-Schramm-Coppi (TD\_SLSC) and the model of Wichoski-MacGibbon-Brandenberger (TD\_WMB) . The main difference between these two models is the main channel for energy loss of the string network, in the former it is the gravitational radiation, while in the later it is the particle production. Both of these fluxes should be regarded as upper limits for TD models, because they have been constructed in such a way to satisfy the bound imposed by the measured cosmic ray and gamma ray fluxes . These fluxes are shown in Fig. 1, where we take representative flux of WMB model with the string mass parameter giving the largest neutrino flux that is consistent with cosmic ray data. This flux is also below the Frejus and Fly’s Eye experimental limits on the neutrino flux.
We also consider two generic fluxes that have a power law behavior. The flux
$$F_{\nu +\overline{\nu }}^s(E)=10^7(E/\mathrm{GeV})^2(\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{GeV}^1)$$
(3)
gives numerically stable results, however, our calculations with a flux with $`F^sE^1`$ is unstable at very high energies. Consequently, we use
$$F_{\nu +\overline{\nu }}^s(E)=10^{13}(E/\mathrm{GeV})^1\frac{1}{(1+E/10^8\mathrm{GeV})^2}(\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{GeV}^1)$$
(4)
as a way to cutoff the high energy behavior. We show results for attenuated fluxes for neutrino energies up to $`10^6`$ GeV. We have chosen the multiplicative factors in $`E^1`$ and $`E^2`$ fluxes in such a way that they exceed atmospheric flux at neutrino energies between 10 TeV and 100 TeV. The upper bound for strong source evolution discussed recently by Waxman and Bahcall would correspond to a limit of $`2\times 10^8E^2`$ (in the same units as Eq. (2)), a factor of 5 smaller than the choice of normalization we take in this paper.
We also show in Fig. 1 the atmospheric neutrino flux at zenith angle of $`0^{}`$ and the horizontal flux. In our evaluation of the atmospheric backgrounds, we use the atmospheric muon and electron neutrino fluxes as a function of zenith angle .
We do not consider the neutrino flux from cosmic ray interactions with the microwave background. This diffuse neutrino flux is typically present at energies higher than we consider here , and it gives low event rates . Furthermore, the cosmic ray interactions with the solar atmosphere are another source of neutrinos, however, for energies above a TeV, the flux scales as $`E^3`$ or steeper . As we see below, tau neutrino regeneration will not be a very important feature in fluxes with such large spectral indices.
## III Tau Neutrino Propagation through the Earth
For neutrino energies above 1 TeV, the oscillation probability for $`\nu _\mu \nu _\tau `$ in the Earth is less than a percent for the parameters constrained by the SuperK experiment. As a consequence, we can neglect neutrino oscillation in our evaluation of tau neutrino propagation accounting for interactions in the Earth.
The coupled transport equations for the fluxes of the tau neutrino and its charged partner are given by
$`{\displaystyle \frac{F_{\nu _\tau }(E,X)}{X}}={\displaystyle \frac{F_{\nu _\tau }(E,X)}{\lambda _{\nu _\tau }(E)}}+{\displaystyle _E^{\mathrm{}}}𝑑E_yG_{\nu _\tau \nu _\tau }(E_y,E,X)`$ (5)
$`+{\displaystyle _E^{\mathrm{}}}𝑑E_yG_{\tau \nu _\tau }(E_y,E,X)`$ (6)
and for taus as,
$`{\displaystyle \frac{F_\tau (E,X)}{X}}={\displaystyle \frac{F_\tau (E,X)}{\lambda _\tau (E)}}{\displaystyle \frac{F_\tau (E,X)}{\lambda _\tau ^{dec}(E,X,\theta )}}`$ (7)
$`+{\displaystyle _E^{\mathrm{}}}𝑑E_yG_{\tau \tau }(E_y,E,X)+{\displaystyle _E^{\mathrm{}}}𝑑E_yG_{\nu _\tau \tau }(E_y,E,X).`$ (8)
Here $`F_{\nu _\tau }(E,X)`$ and $`F_\tau (E,X)`$ are differential energy spectrum of tau neutrinos and tau respectively, for lepton energy $`E`$, at a column depth $`X`$ in the medium defined by
$$X=_0^L\rho (L^{})𝑑L^{}.$$
(9)
The density of the medium a distance $`L`$ from the Earth-atmosphere boundary, measured along the neutrino beam path, is $`\rho (L)`$. The lepton interaction length (in g/cm<sup>2</sup>) is $`\lambda (E)`$ and $`\lambda _\tau ^{dec}(E,X,\theta )`$ is the decay length of the tau.
The functions $`G_{ij}`$ schematically represent interaction or decay contributions to lepton $`j`$ from lepton $`i`$. We limit our evaluations of the tau neutrino flux to $`E_{\nu _\tau }<10^6`$ GeV. Consequently, we can ignore several terms in the coupled differential equations: the term with $`G_{\tau \tau }`$ and the term $`F_\tau /\lambda _\tau `$, both in Eq. (5). Only the decay contribution to the last term in Eq. (2) ($`G_{\tau \nu _\tau }`$) is included in our evaluation. This is justified by the fact that the tau decay is significantly more important than interactions for the energy range of interest, namely $`E<10^6`$ GeV. The neglected terms start contributing for lepton energies on the order of $`10^8`$ GeV. Detailed formulae for $`G_{ij}`$ appear in Ref. .
In our previous work, we have described an analytic method for solving these transport equations , based on the method of Naumov and Perrone . We have evaluated the upward $`\nu _\tau `$ flux for a selection of initial fluxes . We have shown that for “flat” initial neutrino fluxes ($`FE^1`$), a significant number of high energy $`\nu _\tau `$’s cascade down in energy, resulting in enhanced low energy flux relative to the attenuated $`\nu _\mu `$ flux. Here, we evaluate the tau neutrino flux for a more comprehensive selection of incident fluxes, including both neutrino and antineutrino attenuation. In all of our results below, we evaluate the sum of neutrino plus antineutrino fluxes or rates.
In Fig. 2 and 3, we show the attenuated tau neutrino plus antineutrino flux (blue line) and attenuated muon neutrino plus antineutrino flux (Red line), scaled by a factor of the neutrino energy $`E`$, assuming the equal fluxes of tau neutrinos and muon neutrinos incident on the surface of the Earth at a nadir angle of $`0^{}`$ for $`F_{\nu +\overline{\nu }}^o(E_\nu )=0.5\times 10^{13}E^1`$, $`F_{\nu +\overline{\nu }}^o(E_\nu )=0.5\times 10^7E^2`$, the Stecker-Salamon AGN model, the Mannheim AGN (Model A), the two topological defects models, the Waxman-Bahcall GRB model and the atmospheric flux. We note that the enhancement of the tau neutrino flux relative to the initial flux and also to the muon neutrino flux, is prominent for the flat fluxes, such as $`F_\nu ^o(E_\nu )E^1`$, the Stecker-Salamon AGN model and the topological defects model of Sigl et al. In case of the atmospheric flux, which represents the background, the enhancement is very small due to the steepness of the initial neutrino flux.
The angular dependence of the upward $`\nu _\tau `$ flux is also distinct. As an example, in Fig. 4, for the AGN model of Mannheim (Model A) , we show the ratio of the neutrino flux scaled by the flux at $`X=0`$ for two nadir angles, $`\theta =0^{}`$ and $`\theta =30^{}`$, as a function of neutrino energy. Because of the shape of the initial flux, steep for energies below $`10^6`$ GeV, and flat for higher energies, the enhancement of the tau neutrino flux becomes significant only for energies above $`10^6`$ GeV. At fixed energies of $`10^4`$, $`10^5`$, and $`10^6`$ GeV, we show the same flux ratios as a function of nadir angle.
## IV Detection of $`\nu _\tau `$ appearance
Detection of muon neutrinos, in general, is via their charged-current interactions. Produced muons have very large average range making the effective volume of an underground detector significantly larger than the instrumented volume. On the other hand, tau neutrino charged-current interactions produce tau, which has a very short lifetime, making its detection extremelly difficult. Only at very high energies, $`E_\nu >`$ PeV, the production and decay vertices are separated by a measurable distance providing a distinctive signature of tau neutrinos (“double-bang” events) . However, the predicted neutrino fluxes are low at these energies. For $`\nu _\tau `$’s in the energy range of $`10^310^6`$ GeV considered here, the produced tau decays after a very short pathlength back to $`\nu _\tau `$ plus leptons or hadrons. Tau neutrinos will interact via neutral currents, producing a hadronic signal as well. Therefore, the signals of tau neutrino interactions below the double-bang threshold are muons from tau decay, or hadronic/EM showers from the tau production and/or decay. In the first case the background to tau production of high energy muons is $`\nu _\mu `$ charged-current interactions. In the latter case, the backgrounds are $`\nu _\mu `$ neutral current and $`\nu _e`$ charged-current and neutral current interactions. The background rates shown below are with the assumption that the electromagnetic shower from $`\nu _ee`$ charged current interactions cannot be distinguished from hadronic showers. As a consequence, we evaluate the hadronic/electromagnetic (EM) shower rates.
We assume in the analysis presented below that the $`\nu _\mu `$ charged-current events and $`\nu _\tau \tau \mu `$ events can be rejected from the contained hadronic/electromagnetic shower signal. In both cases there is a hadronic shower which includes muons, however, the muons in the hadronic showers from pion and kaon decays are significantly less energetic than the muons from $`\nu _\mu \mu `$ and $`\nu _\tau \tau \mu `$. In the latter case, the energy of the shower is $`1/2`$ the incident neutrino energy $`E_\nu `$ and the energy of the muon is $`1/61/2E_\nu `$. On the other hand, muons coming from particle decays in the hadronic shower are considerably less energetic because of large particle multiplicities. The average charged particle multiplicities for hadronic interactions at $`\sqrt{s}>40`$ GeV are larger than $`10`$ particles , so individual muon energies from charged pion and kaon decays are less than $`5\%`$ of the incident neutrino energy. The hadronic shower and very energetic muon of the “muon signal” should stand out in comparison to the hadronic/EM shower signal in a detector with good energy resolution like the proposed kilometer-cubed detector IceCube.
We describe the evaluation of the muon and hadronic/EM shower event rates. The event rates for $`\nu _\tau \tau \mu `$ and $`\nu _\mu \mu `$ are evaluated and compared with the no oscillation rates. We evaluate the hadronic/EM shower rates for signal and background, then compare with the hadronic/EM shower rates assuming no oscillations of $`\nu _\mu `$. The relative rates of muons and hadronic/EM showers prove to be the most effective diagnostic to neutrino oscillations with the SuperK parameters.
### A Muon Event Rates
The standard evaluation of the muon event rate per solid angle for neutrino interactions with isoscalar nucleons $`N`$ ($`\nu _\mu N\mu X`$) follows from the formula
$`\mathrm{Rate}`$ $`=AN_A{\displaystyle _{E_\mu ^{\mathrm{min}}}^{\mathrm{}}}𝑑E_\nu {\displaystyle 𝑑yR_\mu (E_\nu (1y),E_\mu ^{\mathrm{min}})\frac{d\sigma _{cc}(E_\nu ,y)}{dy}}`$ (11)
$`\times F_\nu (E_\nu ,X)\mathrm{\Theta }(E_\nu (1y)E_\mu ^{\mathrm{min}}).`$
where $`y`$ is the neutrino energy loss, $`y=(E_\nu E_\mu )/E_\nu `$, and $`d\sigma _{cc}(E_\nu ,y)/dy`$ is the charged current differential cross section. $`F_\nu (E_\nu ,X)`$ is the upward neutrino or antineutrino flux which depends on angle implicitly through the pathlength $`X`$. We assume that the initial fluxes of muon neutrinos and antineutrinos that reach the Earth are equal, their sum in the oscillation scenario being half of the muon neutrino plus antineutrino flux produced at the source. The fluxes of neutrinos and antineutrinos at the detector are different because of the difference in charged and neutral current cross sections below energies of $`10^6`$ GeV , however, for these energies and fluxes, the antineutrino event rates differ for the neutrino event rates by at most $`\pm 20\%`$. The average range of a muon, $`R_\mu (E_\mu ,E_\mu ^{\mathrm{min}})`$, is the range of a muon produced in a charged-current interaction with energy $`E_\mu `$ which, as it passes through the medium, looses its energy via bremsstrahlung, ionization, pair production and photonuclear interaction and arrives in a detector with an energy above $`E_\mu ^{\mathrm{min}}`$. Avogadro’s number is $`N_A`$ and $`A`$ is the effective area of the detector. All of the event rates calculated are for the sum of neutrino plus antineutrino contributions to $`\mu ^++\mu ^{}`$ production.
The rate for muons produced by the tau neutrino charged current interactions followed by the tau leptonic decays is given by a modified equation, taking into account the branching fraction for $`\tau \nu _\tau \nu _\mu \mu `$ and the decay distribution of the muon via $`dn(E_\tau )/dz`$, where $`z=E_\mu /E_\tau `$. The decay formulae used here are listed in the Appendix A. The differential event rate is
$`\mathrm{Rate}`$ $`=AN_A{\displaystyle _{E_\mu ^{\mathrm{min}}}^{\mathrm{}}}𝑑E_\nu {\displaystyle 𝑑y𝑑zR_\mu (E_\nu (1y)z,E_\mu ^{\mathrm{min}})\frac{dn(E_\nu (1y)z)}{dz}}`$ (13)
$`\times {\displaystyle \frac{d\sigma _{cc}(E_\nu ,y)}{dy}}F_\nu (E_\nu ,X)\mathrm{\Theta }(E_\nu (1y)zE_\mu ^{\mathrm{min}}).`$
In Fig. 5 we show the neutrino processes that contribute to the muon production. In our evaluation of the muon event rates we use the Earth densities of the Preliminary Earth Model (PREM) described in Ref. . We have used the PREM to determine an average density for a given nadir angle, then used that average density to evaluate the attenuated fluxes. We use the muon range evaluated by Lipari and Stanev . The neutrino and antineutrino cross sections have been evaluated using the CTEQ5 parton distribution functions . The effective area $`A`$ is taken to be 1 km<sup>2</sup>.
In Figs. 6-9, we show muon event rates for $`F_{\nu +\overline{\nu }}^oE^1`$, $`F_{\nu +\overline{\nu }}^oE^2`$, AGN\_SS, AGN\_M95, TD\_WMB, TD\_SLSC and GRB\_WB for $`E_\mu ^{\mathrm{min}}=1,10,100`$ TeV. Blue lines correspond to the upward $`\mu ^++\mu ^{}`$ events from $`\nu _\tau +\overline{\nu }_\tau +\nu _\mu +\overline{\nu }_\mu `$ charged-current interactions (including $`\tau \mu `$ decay), while the red lines are the background contribution from $`\nu _\mu +\overline{\nu }_\mu `$ charged-current interaction only. We note that the muon enhancement due to the tau neutrino contribution for $`E^1`$ flux is almost factor of 2 for small angles and 25% for large angles with $`E_\mu ^{\mathrm{min}}=1`$ TeV. The enhancement is less pronounced at small nadir angles for increasing threshold energies, for example, the blue line is about 60% enhanced relative to the red line for $`E_\mu ^{\mathrm{min}}=10`$ TeV for the $`E^1`$ flux in Fig. 6. A similar enhancement is present for the TD\_SLSC. For steeper fluxes, such as AGN\_SS, AGN\_M95 and $`E^2`$, the enhancement due to tau neutrino contribution is much smaller, of the order of 20-25%.
The muon event rates from the atmospheric neutrino background are shown in Fig. 10a). The input flux is the angle dependent muon neutrino flux of Agrawal et al. . The atmospheric tau neutrino flux is very low, as the tau neutrinos are produced in the atmosphere by comic ray interactions with nuclei in the atmosphere, which produce $`D_s`$ whose leptonic decay, $`D_s\tau \nu _\tau `$, gives $`\nu _\tau `$ . The rates for the atmospheric tau neutrinos are shown in Fig. 10b). In the evaluation of the event rates, we neglect oscillations of atmospheric neutrinos as they travel to the Earth and the oscillations through the Earth since the oscillation probabilities are small above our minimum energy of 1 TeV.
The atmospheric neutrino events represent a background for detection of extragalactic neutrinos. For a muon energy threshold of 1 TeV, the background is large, $`4002000`$ events per year per steradian for 1 km<sup>2</sup> effective area detector. For a muon threshold of $`E_\mu ^{\mathrm{min}}=10`$ TeV, the event rates range between $`680`$ events per year per steradian. A comparison of the event rates from Figs. 6-9 with the atmospheric muon neutrino background indicates that detection of neutrinos from AGN might be possible with $`E_\mu ^{\mathrm{min}}=10`$ TeV or $`100`$ TeV.
The rates for muon events shown in Figs. 6-9 come from assuming that the tau neutrino and muon neutrino fluxes are equal and are half the flux of muon neutrinos produced at the source. Testing the oscillation hypothesis with muon neutrinos alone will be difficult. We see that with the exception of the $`E^1`$ and TD\_SLSC fluxes, the observed muon rate is about half of what one would expect in the absence of oscillations. Given the uncertainties in the normalizations of the predicted fluxes, this factor would not unambiguously signal the presence of tau neutrinos from oscillations. The situation with the $`E^1`$ and TD\_SLSC fluxes is only slightly better. There, in the oscillation scenario, the measured muon event rate is about 80% of the no oscillation prediction at $`\theta =0^{}`$, but less than 70% of the prediction for horizontal events. Testing the oscillation hypothesis by measuring upward muons only will be very difficult.
The relatively small contribution to the muon rate from $`\nu _\tau `$’s, despite the fact that the attenuated flux of tau neutrinos is larger than that of the muon neutrinos, is due to the fact that the muon carries a small fraction of the initial tau neutrino energy. Consequently, for a muon of a given energy, if it comes from a tau neutrino (which interacted producing a tau that subsequently decayed to a muon), the initial tau neutrino has a much higher energy than a muon neutrino which produces a muon directly via the charged current $`\nu _\mu N\mu X`$ process. All predicted neutrino fluxes decrease with energy. Even with some “pileup,” the tau neutrino fluxes are decreasing fast enough that the muon energy fraction results in sampling a much smaller tau neutrino flux than the corresponding muon neutrino flux. It is this observation that leads one to consider signals that carry a much larger fraction of the incident tau neutrino energy.
### B Upward Hadronic/Electromagnetic Showers and Their Detection
The hadronic/EM shower signal of $`\nu _\tau `$ interactions is a much more promising final state from the theoretical point of view than the muon signal. The benefit is that the hadronic showers include both production hadrons and tau decay hadrons, so there is a much higher fraction of the incident tau neutrino energy visible in the detector . The next generation of neutrino telescopes may not be able to distinguish between hadronic and electromagnetic showers, so we include in the signal and in the background, processes that include hadrons and electron. As mentioned above, we assume that the high energy muon associated with the target jet in $`\nu _\mu `$ charged current interactions will be used to veto the process $`\nu _\mu N\mu X`$. Distinguishing electromagnetic from hadronic showers might be possible by looking at the difference between the front to back ratio of the cascade Cherenkov light, and perhaps by the number of residual $`\pi \mu e`$ decay, although this is considered to be experimentally difficult .
The processes that go into our evaluation of $`\nu _\tau `$ hadrons are
$`\nu _\tau N\tau +\mathrm{hadrons},\tau \nu _\tau +\mathrm{hadrons},`$ (14)
$`\nu _\tau N\tau +\mathrm{hadrons},\tau \nu _\tau +e+\nu _e,`$ (15)
$`\nu _\tau N\nu _\tau +\mathrm{hadrons}.`$ (16)
For the charged-current interactions, the hadronic/electromagnetic energy is the sum of the energy carried by the hadrons in tau production, as well as the tau decay hadronic energy or tau decay electron energy.
The background for the hadronic/electromagnetic showers is due to the $`\nu _\mu `$ and $`\nu _e`$ neutral current interactions, and $`\nu _e`$ charged-current interactions are
$`\nu _{\mu ,e}+N\nu _{\mu ,e}+\mathrm{hadrons},`$ (17)
$`\nu _e+Ne+\mathrm{hadrons}.`$ (18)
For the $`\nu _e`$ flux, we assume it is equal to the $`\nu _\mu `$ flux in the SuperK oscillation scenario. All of the processes that contribute to the hadronic/EM showers are shown in Fig. 11.
The tau neutrino shower event rate per unit solid angle from charged-current interactions followed by the tau hadronic decay is given by
$`\mathrm{Rate}`$ $`=VN_A{\displaystyle _{E_{\mathrm{shr}}^{\mathrm{min}}}^{\mathrm{}}}𝑑E_\nu {\displaystyle 𝑑y𝑑z\frac{dn(E_\tau )}{dz}\frac{d\sigma _{cc}(E_{\nu _\tau },y)}{dy}F_{\nu _\tau }(E_{\nu _\tau },X)}`$ (20)
$`\times \mathrm{\Theta }(E_{\nu _\tau }(y+(1y)(1z))E_{\mathrm{shr}}^{\mathrm{min}}).`$
The hadronic energy from the broken nucleon $`E_{\mathrm{shr}}^{int}=E_\nu y`$ and the hadronic energy from the decay $`E_{\mathrm{shr}}^{\mathrm{decay}}=E_\nu (1y)(1z)`$ are added to get the total shower energy. Again, $`y=(E_\nu E_\tau )/E_\nu `$ for incident neutrino energy $`E_\nu `$, while $`z=E_\nu ^{}/E_\tau `$, where $`E_\nu ^{}`$ is the energy of the neutrino from the tau decay. The differential distributions for the hadronic decay modes are shown in the Appendix A. For the electronic decay of the tau, the differential distribution $`dn/dz`$ is replaced by the purely leptonic distribution in terms of $`z^{}E_e/E_\tau `$. The theta function is replaced by
$$\mathrm{\Theta }(E_{\nu _\tau }(y+(1y)(1z))E_{\mathrm{shr}}^{\mathrm{min}})\mathrm{\Theta }(E_{\nu _\tau }(y+(1y)z^{})E_{\mathrm{shr}}^{\mathrm{min}}).$$
(21)
The neutral current background event rate is given by
$$\mathrm{Rate}=VN_A_{E_{\mathrm{shr}}^{\mathrm{min}}}^{\mathrm{}}𝑑E_\nu 𝑑y\frac{d\sigma _{nc}(E_\nu ,y)}{dy}F_\nu (E_\nu ,X)\mathrm{\Theta }(E_\nu yE_{\mathrm{shr}}^{\mathrm{min}}),$$
(22)
while the electron neutrino charged current background rate is given by
$$\mathrm{Rate}=VN_A_{E_{\mathrm{shr}}^{\mathrm{min}}}^{\mathrm{}}𝑑E_\nu 𝑑y\frac{d\sigma _{cc}(E_\nu ,y)}{dy}F_\nu (E_\nu ,X)\mathrm{\Theta }(E_\nu E_{\mathrm{shr}}^{\mathrm{min}}).$$
(23)
In Figs. 12-15 we show the upward hadronic/EM shower event rates as a function of the nadir angle for $`E_{\mathrm{shr}}>E_{\mathrm{shr}}^{\mathrm{min}}`$ where $`E_{\mathrm{shr}}^{\mathrm{min}}=1`$ TeV, 10 TeV and 100 TeV for input fluxes: $`F_{\nu +\overline{\nu }}^0E^1`$, $`F_{\nu +\overline{\nu }}^0E^2`$ AGN\_SS, AGN\_M95, TD\_WMB, TD\_SLSC and GRB\_WB, all assuming that $`V=1`$ km<sup>3</sup>. The blue lines correspond to the event rates from $`\nu _\tau +\overline{\nu }_\tau +\nu _e+\overline{\nu }_e`$ charged-current interactions (and $`\tau \nu _\tau +`$ hadrons decay) and from $`\nu _\tau +\overline{\nu }_\tau +\nu _\mu +\overline{\nu }_\mu +\nu _e+\overline{\nu }_e`$ neutral current interactions. The red lines are the contributions from $`\nu _\mu +\overline{\nu }_\mu `$ neutral current interaction and $`\nu _e+\overline{\nu }_e`$ charged and neutral current interactions. We do not include $`\nu _\mu +\overline{\nu }_\mu `$ charged-current interactions in our calculation because these events can be vetoed by the high energy muons produced in the interactions. All of the rates shown in these figures assume equal neutrino and antineutrino fluxes. They are performed in the oscillation scenario where the ratios of the fluxes $`\nu _e:\nu _\mu :\nu _\tau `$ are $`1:1:1`$.
From Fig. 12a) we note that in the case of the $`E^1`$ flux, the contributions from tau neutrinos are large, a factor of 4 times larger than the muon neutrino plus electron neutrino contribution at zero nadir angle. For horizontal showers, the enhancement factor is smaller, about 2 for all the energy thresholds that we consider. Similarly, for $`E^2`$ flux, the tau neutrino contribution is a factor of 1.7 times larger than the muon neutrino plus electron neutrino contributions for upward showers.
Similar conclusions can be drawn from the plots of the other fluxes. The shower event rates including $`\nu _\tau +\overline{\nu }_\tau +\nu _\mu +\overline{\nu }_\mu +\nu _e+\overline{\nu }_e`$ are significantly enhanced relative to the rates from $`\nu _\mu +\overline{\nu }_\mu +\nu _e+\overline{\nu }_e`$ in the oscillation scenario. The AGN\_SS rates at zero nadir angle are comprised of 60% tau neutrino induced, decreasing to about 40% tau neutrino induced for horizontal showers, as shown in Fig. 13a). AGN\_SS flux gives 25-80 shower events for $`E_{\mathrm{shr}}^{\mathrm{min}}=10`$ TeV and 6-45 events for $`E_{\mathrm{shr}}^{\mathrm{min}}=100`$ TeV with negligible atmospheric background. In Fig. 13b) we show event rates for AGN\_M95 model. We find 3-6 shower events per year per steradian for $`E_{\mathrm{shr}}^{\mathrm{min}}=10`$ TeV, with atmospheric background of 2-16 events. Detection of events with higher energy threshold would require looking at almost horizontal showers, where the background is small.
The TD\_WMB model in Fig. 14a) shows an enhancement of between 2.1-2.3 for zero nadir angle, and a factor of 1.7 for almost horizontal showers. Fig. 14b) shows the more striking enhancement in the TD\_SLSC model, where the enhancement is a factor of between 3.7 to 6.2 at zero nadir angle, to a factor of 2 for large nadir angles. However, due to the particularly low normalization of the TD\_SLSC flux, the kilometer-size detector would not be sufficient for its detection. Fig. 15 shows the GRB\_WB model in which the enhancement factor is between 1.5 to 2, depending on energy threshold and angle. The event rates for showers with energies above 10 TeV are comparable with the background, but higher energy threshold of 100 TeV would still give a few events per year for large nadir angle with negligible background.
Fig. 16a) and 16b) show the shower event rates for the atmospheric $`\nu _\mu +\overline{\nu }_\mu +\nu _e+\overline{\nu }_e`$ and $`\nu _\tau +\overline{\nu }_\tau `$ fluxes, respectively. For showers with energies above 10 TeV, the event rates are twice as large as for the $`E^1`$ flux at small nadir angle. For $`E_{\mathrm{shr}}^{\mathrm{min}}=10`$ TeV, we find the event rates for the showers to be about 8-18 per km<sup>3</sup> per year per steradian for the $`E^2`$ flux, compared with the atmospheric background of 2-16.
The AGN rates will stand out above the atmospheric background for $`E_{\mathrm{shr}}^{\mathrm{min}}10`$ TeV. The GRB\_WB rates are more than half of the atmospheric neutrino rates at the 10 TeV shower threshold at small nadir angles. The TD rates are all quite low overall, and in comparison to the atmospheric background rates.
Since one does not measure separately the tau neutrino induced shower rates and the muon and electron induced shower rates, given a particular model, one can compare the rates with the oscillation hypothesis to the predicted rates without oscillations. We discuss the rates for specific models here, then discuss a more model independent analysis in the next section. To illustrate the effect of oscillations we plot in Figs. 17-20 the ratio of the shower event rates from $`\nu _\tau `$, $`\nu _\mu `$ and $`\nu _e`$ in the oscillation scenario to the shower rates from $`\nu _\mu `$ and $`\nu _e`$ in the standard model, with no oscillations.
From Fig. 17a) we note that for the $`E^1`$ flux, the shower event rates in the oscillation scenario are a factor of 3.3-3.7 larger than in the no oscillation case for $`E_{\mathrm{shr}}^{\mathrm{min}}=1100`$ TeV for $`\theta =0^{}`$. They are a factor of 1.6 enhanced for the horizontal shower rate. For the $`E^2`$ flux, the enhancement is a factor of 1.4-1.6 relative to the no oscillation case for $`E_{\mathrm{shr}}^{\mathrm{min}}=1100`$ TeV, shown in Fig. 17b). In the case of AGN models, if one assumes oscillations, the shower event rates are factor of 1.8-2.1 larger for AGN\_SS at zero nadir angle, decreasing to 1.5 for nearly horizontal showers, as shown in Fig. 18a). Fig. 18b) shows a ratio ranging between 1.4-1.9 for AGN\_M95 for small nadir angles.
From Fig. 19a), we note that the shower event rates for TD\_WMB are factor of 1.8-2.1 enhanced for energy thresholds of 1-100 TeV for the upward neutrinos, while the enhancement is a factor of 1.5 for almost horizontal showers. In the TD\_SLSC model, Fig. 19b), the shower event rate is a factor of 3-3.6 enhanced at small nadir angles, and a factor of 1.6 enhanced for horizontal showers. In the case of GRB\_WB model, Fig. 20 shows an enhancement of 1.4-1.7 for $`E_{\mathrm{shr}}^{\mathrm{min}}=1100`$ TeV.
### C Relative Rates
We have shown that comparison of the muon and shower rates serves as a diagnostic for $`\nu _\mu \nu _\tau `$ oscillations over astronomical distances. For example, for the $`E^2`$ flux,
$$\mathrm{Ratio}\left[(\nu _\tau +\nu _\mu +\nu _e\mathrm{shower})_{\mathrm{osc}}/(\nu _\mu +\nu _e\mathrm{shower})_{\mathrm{no}\mathrm{osc}}\right]1.5$$
(24)
while
$$\mathrm{Ratio}\left[(\nu _\tau +\nu _\mu \mu )_{\mathrm{osc}}/(\nu _\mu \mu )_{\mathrm{no}\mathrm{osc}}\right]0.5.$$
(25)
The ratios include contributions to showers and muons from antineutrinos. This feature of a deficit of muon rates and an excess of shower rates in the oscillation scenario compared to the no-oscillation scenario is a generic feature of all the neutrino spectra in Fig. 1. To demonstrate this point quantitatively, we define a ratio of ratios,
$$R\frac{\mathrm{Ratio}\left[(\nu _\tau +\nu _\mu +\nu _e\mathrm{shower})_{\mathrm{osc}}/(\nu _\mu +\nu _e\mathrm{shower})_{\mathrm{no}\mathrm{osc}}\right]}{\mathrm{Ratio}\left[(\nu _\tau +\nu _\mu \mu )_{\mathrm{osc}}/(\nu _\mu \mu )_{\mathrm{no}\mathrm{osc}}\right]}.$$
(26)
As Figs. 6-9 and 17-20 illustrate for individual fluxes, $`R`$ depends on energy threshold and angle. We show in Fig. 21 the band of $`R`$ spanned by the representative models of Fig. 1 for three thresholds in muon or shower energy: a) 1 TeV, b) 10 TeV and c) 100 TeV. We note that $`R`$ depends on nadir angle and threshold energy, however, $`R\stackrel{>}{}2.4`$ independent of the initial flux. For a given model, measured rates will be very distinct from predicted rates if the SuperK results for oscillation parameters are correct.
Determining
$$R_{exp}\frac{\mathrm{Ratio}\left[(\mathrm{shower}\mathrm{rate})_{\mathrm{measured}}/(\mu \mathrm{rate})_{\mathrm{measured}}\right]}{\mathrm{Ratio}\left[(\nu _\mu +\nu _e\mathrm{shower})_{\mathrm{no}\mathrm{osc}}/(\nu _\mu \mu )_{\mathrm{no}\mathrm{osc}}\right]}$$
(27)
nevertheless relies on theoretical input for the “no-oscillation” flux. Different energy behaviors of incident fluxes will have implications for the angular and energy dependence of the event rates of upward muons and upward hadronic/EM showers, allowing for an indirect characterization of the energy dependence of the source.
A more model independent test of the oscillation scenario would be to compare the measured ratio of showers to muons with the no-oscillation predictions, on an absolute scale. This requires a crude separation of the different energy behaviors of the fluxes of Fig. 1. By only looking at the ratio of shower to muon events, one could confuse the AGN\_M95 no-oscillation ratio with the similar GRB\_WB oscillation ratio. However, experimentally, the energy and angular dependence of the muon event rates for the two fluxes are quite different, and GRB neutrinos would reveal themselves by time correlations to observed GRB events. One category of fluxes, with not too steep energy behavior ($`E^1`$, AGN\_SS, TD\_WMB, TD\_SLSC and GRB\_WB), have a reduction in the event rates from a muon threshold of 1 TeV to 10 TeV at nadir angle of $`0^{}`$ by a factor of less than 4, whereas for the other two steeper fluxes ($`E^2`$ and AGN\_M95), the reduction is by more than a factor of 6. These reduction factors are independent of whether or not oscillations occur. The steeper fluxes are also distinguished by muon event rates with a less marked dependence on nadir angle. If one separates the steep from the less steep examples used here, then at all three threshold energies, the band of oscillation shower to muon ratios does not overlap with the band of no-oscillation ratios. This is shown graphically in Figs. 22 (a-f) for 1 TeV, 10 TeV and 100 TeV thresholds, respectively. In fact, for the 100 TeV threshold, one does not need any information about the energy dependence of the initial flux, but in this case, the event rates are expected to be low.
## V Discussion
We have studied signals for $`\nu _\mu \nu _\tau `$ oscillations with extragalactic high energy muon neutrinos. Assuming SuperK oscillation parameters, muon neutrinos convert into tau neutrinos as they travel megaparsec distances, with both fluxes being equal at the surface of the Earth. High energy muon neutrinos get absorbed as they pass through the Earth, while tau neutrinos cascade down to lower energies. We find this enhancement of the $`\nu _\tau `$ flux in the low energy region to be prominent for flat initial spectrum, such as $`E^1`$, the AGN model of Stecker and Salamon, and the topological model of Sigl et al. For steeper spectra, the enhancement is small because the number of higher energy neutrinos that contributes to the lower energy flux via tau decay is relatively small compared to the low energy flux of neutrinos.
Upward tau neutrinos, once they reach the detector, interact producing tau leptons which decay with very short lifetimes. We have considered muons from tau decay as well as its hadronic decay mode. Since the planned detectors are unable to distinguish between hadronic and electromagnetic showers, we have included all the processes that give both hadronic and electromagnetic showers. We find that upward muons alone would not be sufficient to separate the tau neutrinos contribution, due to the large background from $`\nu _\mu `$ charged-current interactions, the small branching fraction for $`\tau \mu `$ decay mode and the model uncertainty for the incident neutrino flux.
In the case of upward hadronic/EM showers, we find that tau neutrinos give significant contributions, signaling the $`\nu _\tau `$ appearance. Given the uncertainties in the normalizations of the extragalactic neutrino fluxes, combining muon rates and hadronic/EM rates offer the best chance to test the $`\nu _\mu \nu _\tau `$ oscillation hypothesis.
As concluded in earlier work , in general, an energy threshold of between 10 TeV and 100 TeV for upward muons and showers is needed in order to reduce the background from atmospheric neutrinos. We find that diffuse AGN neutrino fluxes, as described by the Stecker-Salamon and Mannheim models, as well as neutrinos from GRBs can be used to detect tau appearance. By measuring upward showers with energy threshold of 10 TeV, and upward muons, the event rates exceed the atmospheric background and are about a factor of 1.5-2 larger than in the no-oscillation scenario.
Here we also comment on the effect of muon and pion cooling to the flavor ratio. Athar et al. in Ref. have shown that with a negligible electron neutrino content at the source, the electron neutrino content at the Earth (in the three-flavor model) is reduced if not negligible compared to the nearly equal muon and tau neutrino fluxes. Keeping the energy spectrum unchanged, this means that the hadronic/electromagnetic shower background, which has significant contributions from $`\nu _eNeX`$ with $`\nu _e>E_{\mathrm{shr}}^{\mathrm{min}}`$ would be reduced. Electron (anti-) neutrinos from processes in the propagation of cosmic rays may dominate at some energies . We have not considered that possibility here because of the low rates below 1 PeV.
Steepening of the energy spectra displayed in Fig. 1 due to a neutrino energy cutoff from pion and muon cooling will have implications for the tau neutrino ‘pileup’, especially for the flatter spectra where the pileup is more pronounced. As an estimate of the lower bound on the relative enhancement of the hadronic/EM signal compared to the muon signal, one can compare the rates for horizontal events, where tau neutrino pileup is small. For example, Figures 22 (a-f) show clear distinction between oscillation and no-oscillation scenarios, even in directions near horizontal, where there is no pile up. Furthermore, for $`E^2`$ flux, where the pileup is very small , the ratio of ratios $`R`$ discussed above ranges from 2.5 to 2.8. Thus, even without the tau neutrino pileup, the oscillation scenario can be distinguished from the no-oscillation scenario.
The detection of $`\nu _\mu \nu _\tau `$ oscillations with a point source might also be possible. With the resolution for the planned neutrino telescopes of $`2^{}`$, the atmospheric background is reduced by $`3.8\times 10^3`$. For upward showers, this gives less than 1 event per year for $`E_{\mathrm{shr}}^{\mathrm{min}}=1`$ TeV, and even less for higher energy thresholds. Thus, if the point source has a flat spectrum, $`F_{\nu +\overline{\nu }}=10^{16}E^1`$, then one would be able to detect tau neutrinos by measuring upward showers with $`E_{\mathrm{shr}}^{\mathrm{min}}=1`$ TeV. In the more realistic case, when the point source has a steeper spectrum ($`E^2`$), such as Sgr A\* , a normalization of $`10^7`$/cm<sup>2</sup>/s/sr/GeV would be sufficient for the detection of tau neutrinos with threshold of 1 TeV. Time correlations with variable point sources would further enhance the signal relative to the background.
We have demonstrated that extragalactic sources of neutrinos can be used as a very-long baseline experiment, providing a source of tau neutrinos and opening up a new frontier in studying neutrinos oscillations.
Acknowledgements
The work of S.I.D. and I.S. has been supported in part by the DOE under Contracts DE-FG02-95ER40906 and DE-FG03-93ER40792. The work of M.H.R. has been supported in part by National Science Foundation Grant No. PHY-9802403.
## A Tau decay distribution
The decay distribution of the tau neutrinos from tau decay has the following form, in terms of $`z=E_\nu /E_\tau `$:
$$\frac{dn}{dz}=\underset{i}{}B_i(g_0^i+Pg_1^i).$$
(A1)
The polarization of the decaying $`\tau ^{}`$ is $`P`$, which for neutrino V-A production of $`\tau ^{}`$ is $`P=1`$. The branching fraction into decay channel $`i`$ is indicated by $`B_i`$. The distribution is normalized such that
$$\frac{dn}{dz}𝑑z=\underset{i}{}B_i=1.$$
(A2)
In Table 1, we show the functions $`g_0`$ and $`g_1`$ for each decay mode, written in terms of $`z`$ and $`r_i=m_i^2/m_\tau ^2`$. Details of the calculational procedure can be found in Ref. or in Ref. . For the multiprong tau decays, we approximate the distribution by a theta function, as indicated in the table.
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# Statistical characteristics of simulated walls
## 1 Introduction
Over the past decade immense progress was achieved in the investigation of the large scale matter distribution. Now the galaxy distribution is studied up to the redshift $`z3`$ (Steidel at al. 1996). At smaller redshifts the analysis of rich galaxy surveys with an effective depth $``$ (200 – 400)$`h^1`$Mpc, such as the Durham/UKST Galaxy Redshift Survey (Ratcliffe et al. 1996), and the Las Campanas Redshift Survey (Shectman et al. 1996), have established the existence of wall-like structure elements as a typical phenomenon in the visible galaxy distribution incorporating $``$ (40 – 50)% of galaxies (Doroshkevich et al. 1996, hereafter LCRS1; Doroshkevich et al. 1999a, hereafter LCRS2; Doroshkevich et al. 2000, hereafter DURS). The wall-like structure elements with a typical diameter $``$ (30 – 50)$`h^1`$Mpc surround low-density regions with a similar typical diameter $``$ (50 – 70)$`h^1`$ Mpc. Within the wall-like structures, the observed galaxy distribution is also inhomogeneous (see, e.g., Fig. 5 of Ramella et al. 1992), and galaxies are concentrated in high density clumps and filaments.
The galaxies occupying low density regions are concentrated within a random network of filaments. Filaments incorporate $``$ 50% of galaxies and are clearly seen in many redshift surveys (see, e.g., de Lapparent, Geller & Huchra 1988). These results extend the range of investigated scales in the galaxy distribution up to $``$ 100$`h^1`$Mpc. Further progress in the study of the observed large scale galaxy distribution could be reached with the 2dF redshift survey (Colless 1998; Cannon 1998) and the Sloan Digital Sky Survey (Loveday & Pier 1998; Maddox 1998).
The formation and evolution of structure on large scales are investigated in numerous simulations (see, e.g., Cole et al., 1997, 1998; Jenkins et al. 1998; Governato et al. 1998; Müller et al. 1998; Doroshkevich et al. 1999b, hereafter DMRT). These simulations are performed in large boxes ($``$ 350 – 500$`h^1`$Mpc) and reproduce the main properties of the observed large scale matter distribution. In particular, they confirm formation of large wall-like matter condensations due to a nonlinear anisotropic matter compression on a typical scale $``$ (20 – 30)$`h^1`$Mpc that is about one half of the typical wall separation.
The statistical characteristics of wall formation are described by an approximate theoretical model (Lee & Shandarin 1998; Demiański & Doroshkevich 1999a, b, hereafter DD99) based on the Zel’dovich nonlinear theory of gravitational instability (Zel’dovich 1970, 1978; Shandarin & Zel’dovich 1989). This approach relates the structure parameters with the main parameters of the underlying cosmological scenario and the initial power spectrum of density perturbations. The impact of large scale perturbations is found to be important throughout all evolutionary stages and some statistical characteristics of structure elements – filaments and walls formed in the course of nonlinear evolution – are directly connected with the parameters of these perturbations. Another theoretical model of large scale structure formation was discussed in Bond, Kofman & Pogosyan (1996).
The simulated large scale matter distribution does not exactly reproduce the theoretical expectations due to the influence of some essential factors, the most important ones are the small scale clustering and relaxation of compressed matter, and the large scale matter flow within sheet-like structure elements. Thus, compression of matter along one of the transversal directions transforms sheet-like elements into filaments, while expansion of matter in both transversal directions results in the erosion of pancakes. The disruption of walls and the small scale clustering of compressed matter substantially accelerate the relaxation and are responsible for strong matter concentration within walls. This is apparent from the isotropy of velocity dispersion within walls noticed in DMRT.
The combined influence of these (and other) factors complicates the statistical description of the large scale matter distribution at late evolutionary stages, what is typical for the final evolutionary stages of the standard COBE-normalized CDM (SCDM) model with $`\mathrm{\Omega }_m=1`$. For low density models, such as the open CDM (OCDM) model and the $`\mathrm{\Lambda }`$CDM model with $`\mathrm{\Omega }_\mathrm{\Lambda }>\mathrm{\Omega }_m`$, the situation is not so complex, and some statistical characteristics of structure can be successfully compared with the approximate theoretical expectations.
The investigation of wall-like massive structure elements is more promising in this respect because walls represent the first step in the process of structure formation and, so, hold more information about characteristics of the initial matter flow. Such walls are observed as superclusters of galaxies similar to the Great Wall (de Lapparent et al. 1988) and the Pisces-Perseus supercluster (Giovanelli & Haynes 1993). In simulations such wall-like structure elements are also easily identified because of their relatively high overdensity. Samples of such elements were investigated in DMRT and LCRS2. The connection between properties of walls and the amplitude and the spectrum of initial perturbations was discussed in DD99, and some of these results can be compared with measured properties of simulated wall-like structure elements. Examples considered in DD99 and DMRT had rather illustrative character, but they seem to be quite promising.
Here we will compare more accurately some of the expected and measured characteristics of wall-like matter condensations. We concentrate our attention on the physical aspects of the formation and evolution of the large scale matter distribution in order to better understand these processes and the phenomenon of wall-like matter condensations. Both theoretical and numerical estimates are inevitably approximate, but nevertheless, such comparison allows us to test the theoretical conclusions, to reveal and illustrate the influence of essential factors mentioned above, and to examine the abilities of statistical methods used to describe the large scale matter distribution.
These methods allow us to reveal, in particular, some differences in characteristics of the large scale matter distribution in the real and redshift spaces. Various aspects of this problem were widely discussed during the past decade (see, e.g., Kaiser 1987; McGill 1990 a; Davis, Miller & White 1997; Hamilton 1998; Melott et al. 1998; Hui, Kofman & Shandarin 1999, Tadros et al. 1999). Here we show that the differences between characteristics of walls in the real and redshift spaces depend on the basic cosmological model and increase during the cosmic evolution. Characteristics of walls in the real and redshift spaces are almost identical for the low-density models, but they differ more strongly for the SCDM model.
We do not discuss the application of these methods to the observed galaxy catalogues, what is a much harder problem, due to the strong influence of selection effects and other factors. We will consider this problem in the future.
This paper is organized as follows: In Sec. 2 the basic notations are introduced. In Sec. 3 the statistical characteristics of wall-like structure elements in the Zel’dovich theory are presented. In Sec. 4 we consider the methods used to measure the required characteristics of matter distribution. Our results are presented in Secs. 5 & 6 where they are also compared with the theoretical expectations. Sec. 7 contains summary and a short discussion of our main results. Some technical details are given in Appendixes A.
## 2 Statistical characteristics of large scale structure
It is generally recognized that the formation of observed large scale structure is driven by the middle part of the power spectrum, $`p(k)`$, with $`0.2h\mathrm{Mpc}^1k0.01h\mathrm{Mpc}^1`$ ($`k`$ is the comoving wave number), and it is weakly sensitive to the small and large scale perturbations. In many publications authors use an artificial smoothing of the spectrum to describe this process (see, e.g., Bardeen et al. 1986, hereafter BBKS, Coles et al. 1993). However, as was shown in DD99 it is possible to avoid this artificial smoothing if the process of structure formation is described in terms of the displacement, $`S_i(𝐪)`$, and velocity rather than density field.
Indeed, in contrast with the density field, the statistical characteristics of displacements are weakly sensitive to the small and large scale perturbations and are reasonably well described by the middle part of the initial power spectrum. Even the strong nonlinear matter clustering does not significantly influence the main characteristics of displacements and, so, such (approximate) description of structure holds during long period of structure evolution. Of course, this approach cannot describe the formation of gravitationally confined walls and their disruption into a system of high density clouds.
Bearing in mind these comments we will describe the structure parameters using characteristics directly connected with the displacement. One of them is the large scale amplitude of perturbations measured by the dispersion of displacements,
$$\sigma _s^2(z)=\frac{1}{2\pi ^2}_0^{\mathrm{}}p(z,k)𝑑k,$$
$`(2.1)`$
Other convenient parameter is the coherent length of the displacement and velocity fields, $`l_v`$, expressed through the moment $`m_2`$ of the initial power spectrum, $`p(k)`$. Suitably defined coherent length $`l_v`$ provides simple expressions for the correlation functions of these fields and the basic characteristics of the large scale structure (DD99 and Sec. 3).
### 2.1 The Zel’dovich approximation
The Zel’dovich theory connects the Eulerian, $`r_i`$, and the Lagrangian, $`q_i`$, coordinates of fluid elements (particles) by the expression
$$r_i=(1+z)^1[q_iB(z)S_i(𝐪)],$$
$`(2.2)`$
$$S_i(𝐪)=\mathrm{\Phi }(𝐪)/q_i,$$
where $`z`$ denotes the redshift, $`B(z)`$ describes growth of perturbations in the linear theory, and the random vector $`S_i`$ or the random potential $`\mathrm{\Phi }`$ characterize the spatial distribution of perturbations. The Lagrangian coordinates of a particle, $`q_i`$, are its unperturbed comoving coordinates.
The velocity of a particle can be found from (2.2) as
$$u_i(𝐪,z)=\frac{dr_i}{dt}=\frac{H(z)}{1+z}[q_i(1+\beta )B(z)S_i(𝐪)],$$
$$\beta (z)=\frac{1+z}{B}\frac{dB(z)}{dz},$$
$`(2.3)`$
$$H(z)=H_0\sqrt{\mathrm{\Omega }_m(1+z)^3+(1\mathrm{\Omega }_m\mathrm{\Omega }_\mathrm{\Lambda })(1+z)^2+\mathrm{\Omega }_\mathrm{\Lambda }},$$
where H is the Hubble constant ($`H_0=100h`$ km/s/Mpc). Analytical fits for the functions $`B(z)`$ and $`\beta (z)`$ were given in DD99. Approximately, at $`z1`$, we have
$$B(0)=1,\beta (0)\frac{2.3\mathrm{\Omega }_m}{1+1.3\mathrm{\Omega }_m}.$$
$`(2.4)`$
### 2.2 Main structure characteristics for the CDM-like power spectrum
The standard CDM-like power spectrum with a Harrison – Zel’dovich large scale asymptote
$$p_{cdm}(k)=A(z)kT^2(k/k_0),k_0=\mathrm{\Gamma }h\mathrm{Mpc}^1,$$
$`(2.5)`$
$$\mathrm{\Gamma }=\sqrt{\frac{1.7\rho _\gamma }{\rho _{rel}}}\mathrm{\Omega }_mh,$$
can be taken as a reasonable approximation of the initial power spectrum used in Zel’dovich’ theory. Here $`A(z)`$ is the amplitude of perturbations, $`T(x)`$ is a transfer function and $`\rho _\gamma \&\rho _{rel}`$ are the densities of CMB photons and relativistic particles (photons, neutrinos etc.). For this spectrum the parameters $`l_v`$ and $`\sigma _s`$ are expressed through the spectral moments, $`m_j`$, as follows:
$$l_v^2=_0^{\mathrm{}}kT^2(k/k_0)𝑑k=m_2k_0^2,$$
$`(2.6)`$
$$\sigma _s^2\frac{1}{2\pi ^2}_0^{\mathrm{}}p_{cdm}(k)𝑑k=\frac{A(z)}{2\pi ^2}k_0^2m_2=\frac{A(z)}{2\pi ^2l_v^2}$$
$$m_j=_0^{\mathrm{}}x^{3+j}T^2(x)𝑑x,m_2=_0^{\mathrm{}}xT^2(x)𝑑x$$
For the CDM transfer function (BBKS) $`m_2=0.023`$, and the expressions for the scale $`l_v`$ and the characteristic masses of DM and baryonic components associated with the scale $`l_v`$ can be written more explicitly as
$$l_v\frac{6.6}{\mathrm{\Gamma }}\sqrt{\frac{0.023}{m_2}}h^1\mathrm{Mpc},$$
$`(2.7)`$
$$M_v=\frac{4\pi }{3}<\rho >l_v^3\frac{210^{14}M_{}}{\mathrm{\Gamma }^2h^2},M_b^{(0)}=\frac{\mathrm{\Omega }_b}{\mathrm{\Omega }_m}M_v.$$
Here $`\mathrm{\Omega }_b`$ is the dimensionless mean density of the baryonic component. The same characteristic scale, $`l_v`$, as given by (2.7) can be used for the structure description as long as the Zel’dovich theory can be applied.
More details can be found in DD99. The same approach can be used for other power spectra as well.
### 2.3 The amplitude of large scale perturbations
The large scale amplitude of perturbation as measured by $`A(z)`$ in (2.5) and $`\sigma _s`$ (2.1) can be successfully used to describe the structure evolution in the framework of the Zel’dovich theory. As was shown in DD99 it is convenient to use – together with $`\sigma _s`$ – an effective dimensionless ‘time’, $`\tau (z,\mathrm{\Omega }_m,h)`$,
$$\tau (z)=\frac{\sigma _s}{\sqrt{3}l_v}$$
$`(2.8)`$
which is proportional to the large scale amplitude of perturbations and describes suitably the evolutionary stage reached in the model. This ’time’ is similar to that used in the adhesion model (Shandarin & Zel’dovich 1989).
As was noticed in DD99 the structure evolution shows strong features of self-similarity and is described by universal expressions depending on the dimensionless variables $`𝐪/l_v`$ and $`\tau `$. This is a direct consequence of the Zel’dovich approximation.
The ‘time’ $`\tau `$ can be measured by different methods, some of which are discussed below. It is sensitive to the sample under investigation and to the method of measurement. It can be used to quantify bias between spatial distributions of different objects, such as, for example, large scale bias between distributions of galaxies and the DM component.
The quadrupole component of the CMB anisotropy, $`T_Q`$, the variance of density in a sphere with radius $`8h^1`$Mpc, $`\sigma _8`$, and the velocity dispersion, $`\sigma _{vel}`$ are the more often used characteristics of the large scale amplitude. All these characteristics are proportional to each other, but their dependence on $`\mathrm{\Omega }_m`$ and $`h`$ is different, and they are sensitive to matter distribution in different scales. Thus, the quadrupole component of CMB anisotropy characterizes the perturbations on scales comparable with the horizon, while the values $`\sigma _{vel}`$ and $`\sigma _8`$ are more sensitive to the matter distribution in moderate and small scales.
The connection of these characteristics with $`\sigma _s`$ and $`\tau `$ can be summarized as follows:
1. Using the fits for the CMB anisotropy proposed by Bunn & White (1997) we obtain for the flat $`\mathrm{\Lambda }`$CDM and open OCDM models
$$\tau _T2.73h^2\mathrm{\Omega }_m^{1.2}\left(\frac{m_2}{0.023}\frac{T_Q}{20\mu K}\right),\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,$$
$`(2.9)`$
$$\tau _T2.73h^2\mathrm{\Omega }_m^{1.650.19\mathrm{ln}\mathrm{\Omega }_m}\left(\frac{m_2}{0.023}\frac{T_Q}{20\mu K}\right),\mathrm{\Omega }_\mathrm{\Lambda }=0,$$
where $`\tau _T`$ denotes the amplitude of large scale perturbations, $`\tau `$, measured by the CMB anisotropy. These estimates depend on the spectral moment $`m_2`$ only which is very stable and does not change during the considered period of evolution. But the estimates should be corrected if a possible contribution of gravitational waves is taken into account.
2. The amplitude of perturbations, $`\sigma _s`$ and $`\tau `$, can be directly expressed through the two point autocorrelation function as follows:
$$\sigma _s^2=\underset{r\mathrm{}}{lim}_0^r𝑑x\left(1\frac{x}{r}\right)x\xi (x),$$
$`(2.10)`$
and for the autocorrelation function $`\xi (r)`$ approximated by the power law
$$\xi (r)=(r_0/r)^\gamma ,rr_\xi ,$$
$`(2.11)`$
we have
$$\sigma _s^2(r_\xi )\frac{r_\xi ^{2\gamma }r_0^\gamma }{(2\gamma )(3\gamma )},\tau _\xi =\frac{\sigma _s(r_\xi )}{\sqrt{3}l_v}.$$
$`(2.12)`$
Here $`r_\xi `$ is the first zero-point of the autocorrelation function and $`\tau _\xi `$ denotes the amplitude $`\tau `$ measured by this function. The parameter $`r_\xi `$ is usually found with small precision, but for $`\gamma `$1.5 – 1.7, $`1\gamma /2`$0.25 – 0.15 even some variations of $`r_\xi `$ do not change significantly the final estimate of $`\tau `$.
3. The parameter $`\sigma _8`$ can be also expressed through the two point autocorrelation function, $`\xi (r)`$, (Peebles 1993), and for $`\xi (r)`$ approximated by a power law (2.11) we have
$$\sigma _8^2=\frac{72}{(3\gamma )(4\gamma )(6\gamma )}\left(\frac{r_0}{16h^1\mathrm{Mpc}}\right)^\gamma ,$$
$`(2.13)`$
$$\sigma _s^2\sigma _8^2(8h^1\mathrm{Mpc})^2\frac{(4\gamma )(6\gamma )}{18(2\gamma )}\left(\frac{r_\xi }{16h^1\mathrm{Mpc}}\right)^{2\gamma },$$
$$\tau _8=\sigma _8\mathrm{\Gamma }\sqrt{\frac{(4\gamma )(6\gamma )}{36.75(2\gamma )}}\left(\frac{r_\xi }{16h^1\mathrm{Mpc}}\right)^{\frac{2\gamma }{2}}.$$
$`(2.14)`$
4. The dispersion of the peculiar velocity of particles at small redshifts, z=0, can be written as in the linear theory (DD99)
$$\sigma _{vel}=u_0\sqrt{3}\tau ,u_0=l_vH_0\beta \frac{\mathrm{\Omega }_m}{\mathrm{\Gamma }}\frac{1535\text{km/s}}{1+1.3\mathrm{\Omega }_m},$$
$`(2.15)`$
and for $`\tau `$ we obtain the independent estimate
$$\tau _{vel}=\frac{\sigma _{vel}}{\sqrt{3}u_0}.$$
$`(2.16)`$
Here $`\tau _{vel}`$ denotes the amplitude $`\tau `$ measured by the velocity dispersion. $`\sigma _{vel}`$ takes into account also the high velocities generated by the gravitational compression of matter (in particular, within clusters of galaxies) and, so, it gives actually an upper limit of the amplitude.
## 3 Statistical characteristics of walls in the Zel’dovich theory
In both observed and simulated catalogues, at small redshifts, the wall-like structure elements accumulate $``$ 50% of galaxies and form the skeleton of large scale structure. So, investigation of the characteristics of these elements is important in itself. It allows us also to obtain information about processes of nonlinear structure evolution. In particular, we can find two independent measures of the large scale amplitude, $`\tau `$. As walls represent the first step of the large scale nonlinear matter compression their characteristics can be compared with predictions of the Zel’dovich theory.
In this Section we will consider five characteristics of walls, namely, the surface density of walls, $`m_w`$, defined as the amount of matter per unit of wall surface, for example, per $`h^2`$Mpc <sup>2</sup>, the thickness of walls, $`h_w`$, the wall separation, $`D_{sep}`$, the velocity dispersion of matter compressed within walls, $`w_w`$, and the dispersion of wall velocities, $`\sigma _v`$. All these characteristics can be derived from the Zel’dovich theory (DD99) and can be found for simulated point distributions as well.
### 3.1 Formation of walls
Following DD99, we will consider the intersection of two fluid particles with Lagrangian coordinates $`𝐪_1`$ and $`𝐪_2`$ as the formation of a wall (Zel’dovich pancake) with the surface density $`m_w=n_p|𝐪_1𝐪_2|`$ where $`n_p`$ is the mean particles density in the sample. In Zel’dovich theory statistical characteristics of such walls are described by the initial power spectrum (2.5) and can be expressed through the characteristic scale, $`l_v`$, the surface density of wall, $`m_w`$, or dimensionless surface density, $`q_w=m_w/l_v/n_p`$, and the ‘time’, $`\tau `$, introduced in Secs. 2. To do this, the structure functions of the initial power spectrum can be used. For the standard SDM-like power spectrum (2.5) with the BBKS transfer function these functions were introduced in DD99.
Naturally, the theoretical considerations describe the idealized model of structure evolution. Thus, it uses the rigid wall boundary though in reality such boundaries are always blured. Other important factor is the compression and expansion of pancakes in transversal directions. These motions transform pancakes into filaments and/or lead to the dissipation of poor pancakes. They are not so important for rich walls but can change the wall surface density by a factor of 1.3 – 1.5. The small scale clustering and relaxation of matter distorts also the measured characteristics of walls with respect to theoretical expectations.
These factors distort the actual power spectrum with respect to the used one and introduce differences between the expected and actually measured parameters of walls which however cannot be evaluated a priori. The actual power and limitations of this approach must be tested at first with N-body simulations.
### 3.2 Wall properties in the real space
#### 3.2.1 Surface density of walls
The most fundamental characteristic of walls is the surface density, $`m_w`$. The approximate expression for the probability distribution function (PDF) of the pancakes surface density, $`m_w`$, defined as above has been obtained in DD99 in the same manner as the well known Press-Schechter mass function. It characterizes the process of one dimensional matter compression and formation of wall-like pancakes as described by the Zel’dovich theory.
For Gaussian initial perturbations and the standard CDM-like power spectrum with the BBKS transfer function, it can be written as follows:
$$N_m=\frac{1}{\sqrt{2\pi }\tau _m}\frac{1}{\sqrt{q_w}}\mathrm{exp}\left(\frac{q_w}{8\tau _m^2}\right)\text{erf}\left(\sqrt{\frac{q_w}{8\tau _m^2}}\right),$$
$`(3.1)`$
$$q_w=\frac{m_w}{l_vn_p}=\frac{|𝐪_1𝐪_2|}{l_v},_0^{\mathrm{}}N_m(q_w)𝑑q_w=1,$$
$$q_w=_0^{\mathrm{}}q_wN_m(q_w)𝑑q_w=8(0.5+1/\pi )\tau _m^26.55\tau _m^2,$$
$$q_w^2=_0^{\mathrm{}}q_w^2N_m(q_w)𝑑q_w=128(0.375+1/\pi )\tau _m^4887\tau _m^4,$$
where $`n_p`$ is the mean particle density in the sample, $`l_v`$ is defined by (2.7), $`\tau _m`$ characterizes the amplitude of perturbations and the evolution stage of structures, $`\tau `$, as measured by the surface density of walls, and $`𝐪_1\&𝐪_2`$ are Lagrangian coordinates of wall boundaries. This relation was corrected for the merging of neighboring walls and this process is described by the erf - function in (3.1).
These expressions connect $`\tau _m`$ with the mean surface density of walls and allow us to estimate $`\tau _m`$ from measurements of $`q_w`$. For other models and/or other distributions of initial perturbations the PDFs similar to (3.1) could be obtained using the technique described in DD99.
#### 3.2.2 The wall separation
We have not been able to find a simple theoretical description of the wall separation. Nonetheless, taking into account the mainly one dimensional character of wall formation, we can roughly link the mean measured wall separation, $`D_{sep}`$, to the mean surface density of walls, $`q_w`$.
Indeed, the matter conservation law along the direction of wall compression can be approximately written as follows:
$$m_wf_wn_pD_{sep},q_wf_wD_{sep}/l_v,$$
where $`f_w`$ is the matter fraction assigned to walls. It implies that on average a fraction $`f_w`$ of particles situated at the distance $`\pm 0.5D_{sep}`$ from the center of wall will be collected by the wall. For simulations when the mean wall separation is comparable to the box size, $`L_{box}`$, we will use the more accurate relation
$$q_w\frac{f_{dq}}{l_v}\frac{D_{sep}}{1+D_{sep}/L_{box}}.$$
$`(3.2)`$
The averaging can be performed analytically assuming the exponential distribution function for the wall separation.
The factor $`f_{dq}`$ defined by Eq. (3.2) characterizes the matter fraction assigned to walls as it is determined by comparison of independently measured characteristics $`q_w`$ and $`D_{sep}`$. In turn, difference between $`f_{dq}`$ and $`f_w`$ characterizes the robustness and degree of self consistency of the model and the measurements. These estimates are only approximations because the wall formation is actually a three dimensional process.
#### 3.2.3 Velocity of structure elements
For the pancakes defined in Sec. 3.1 the 1D velocity of walls, $`v_w`$, can be found from relations (2.2) and (2.3) as follows:
$$v_w=\frac{1}{|𝐪_1𝐪_2|}_{q_1}^{q_2}𝐧[𝐮H(z)𝐫]𝑑q$$
$`(3.3)`$
where $`𝐧`$ is a unit vector normal to the wall. The small scale clustering and relaxation of compressed matter does not influence the velocities of walls and, so, they are the most stable characteristics of the evolutionary stage reached. As was shown in DD99 the mean velocity of walls, $`v_w`$, is expected to be negligible as compared with its dispersions, $`\sigma _v`$, and the expected PDF of this velocity, $`N_v`$, is Gaussian for Gaussian initial perturbations.
For the standard spectrum (2.5) with the BBKS transfer function, and for $`q_w`$ 1, the velocity dispersion is related to the amplitude of initial perturbations as follows:
$$\sigma _vu_0\tau ,\tau _v=\frac{\sigma _v}{u_0},$$
$`(3.4)`$
what is similar to (2.16) and also is identical to expectations of the linear theory. Here $`\tau _v`$ denotes the amplitude $`\tau `$ as measured by the dispersion of wall velocity and $`u_0`$ was introduced by (2.15).
#### 3.2.4 Velocity dispersion of matter compressed within walls
The variance of velocity of matter accumulated by walls,
$$w_{wz}^2=\frac{1}{|𝐪_1𝐪_2|}_{q_1}^{q_2}[\mathrm{𝐧𝐮}H(z)\mathrm{𝐧𝐫}v_w]^2𝑑q,$$
$`(3.5)`$
can be found in the framework of the Zel’dovich theory using the structure functions described in DD99. As is shown in Appendix A, it can be written as follows:
$$w_{wz}^2(q_w,\tau )u_0^2\left(\frac{q_w^2}{12}+\frac{\tau ^2(1+\beta )^2}{3\beta ^2}q_w\right),q_w1,$$
$`(3.6)`$
where $`\beta ,u_0,q_w\&v_w`$ were introduced by (2.3), (2.4), (2.15), (3.1) & (3.3). In fact, this function characterizes the mean kinetic energy of particles compressed into a wall of a given size $`q_w`$. After averaging over a sample of walls with the PDF $`N_m`$ (3.1), in the Zel’dovich theory, we obtain
$$w_z^2(\tau )=w_{wz}^2(q_w,\tau )u_0^2\tau ^4\left(7.4+\frac{2.2(1+\beta )^2}{\beta ^2}\right).$$
$`(3.7)`$
The comparison of the expected mean kinetic energy of the compressed particles with the kinetic energy measured in simulations characterizes the mean degree of relaxation of compressed matter at a given $`\tau `$.
For richer walls, with $`q_wq_w`$, the relation (3.6) is transformed into
$$w_{wz}\frac{u_0}{\sqrt{12}}q_w,$$
$`(3.8)`$
and for such walls, the PDF is similar to (3.1). For a rich sample of walls, this relation can be also used for the direct measurement of the amplitude $`\tau `$ (DMRT; DD99).
#### 3.2.5 Wall thickness
The methods discussed in DD99 allow us also, in the framework of the Zel’dovich theory, to obtain the expected thickness of walls along the direction of maximal compression, $`h_w`$. It can be characterized by the thickness of a homogeneous slice with the same surface density. The corresponding expression (Appendix A) is
$$h_{wz}(q_w,\tau )2l_v\tau \sqrt{q_w}(1+z)^1.$$
$`(3.9)`$
This relation shows that the wall thickness is strongly correlated with its surface density. After averaging with the PDF (3.1) we obtain for the mean thickness of walls
$$h_{wz}8\pi ^{1/2}l_v\tau ^2(1+z)^1.$$
$`(3.10)`$
The degree of matter compression in the Zel’dovich theory, $`\delta _z(q,\tau )`$, is characterized by the ratio
$$\delta _z=\frac{q_wl_v}{h_w}=\frac{\sqrt{q_w}}{2\tau }.$$
After averaging with the PDF (3.1) we have for the mean degree of matter compression
$$\delta _z\frac{\sqrt{q_w}}{2\tau }=\frac{2}{\sqrt{\pi }}=1.13.$$
$`(3.11)`$
So, in the Zel’dovich theory the averaged degree of matter compression is small.
### 3.3 Wall properties in the redshift space
In observed catalogues only the redshift position of galaxies along the line-of-sight is known, and therefore the parameters of observed structures with respect to those found above can differ due to the influence of the velocity field. The statistical characteristics of walls in redshift space predicted by the Zel’dovich theory can be found with the methods described above. This information is not so rich as in the real space because in the redshift space, positions of particles are determined by their velocities, and, for example, such a useful characteristic as the wall velocity cannot be found.
#### 3.3.1 Surface density of walls
In the real space (Sec. 3.1) the pancake formation was defined as an intersection of particles with coordinates $`𝐪_1\&𝐪_2`$. In the redshift space the velocity (2.3) along the line-of-sight must be used instead of the coordinate. In Zel’dovich theory the velocity dispersion exceeds the dispersion of displacement by a factor of $`(1+\beta )`$. Hence, this substitution increases the wall surface density in the redshift space in respect to that in the real space, and now we must use
$$\tau _{rd}=f_{rd}\tau =\tau \sqrt{(1+\beta )^2\mathrm{cos}\varphi ^2+\mathrm{sin}\varphi ^2},$$
$`(3.12)`$
instead of $`\tau `$. Here the factor $`f_{rd}`$1 describes the more effective matter compression in the redshift space predicted by the Zel’dovich theory, $`\beta `$ was introduced in (2.3), (2.4), and $`\varphi `$ is a random angle between the direction of wall compression and the line-of-sight ($`0\varphi \pi /2`$). Evidently, $`\tau _{rd}=\tau `$ for $`\beta =0`$, so $`f_{rd}(\beta =0)=`$1.
The PDF of wall surface densities in the redshift space is identical to (3.1) with a substitution of $`\tau _{rd}`$ for $`\tau `$, and now for the mean surface density of walls we have
$$q_w=8(0.5+1/\pi )f_{rd}^2\tau ^26.55f_{rd}^2\tau ^2,$$
$`(3.13)`$
$$1f_{rd}^2=\frac{1}{3}[2+(\beta +1)^2]3.667,$$
where $`\tau `$ characterizes the evolutionary stage as before. Probably, these relations can be used for the description of poorer pancakes and earlier evolutionary stages when the influence of other factors is less important.
At small redshifts we must take into account the influence of the high velocity dispersion of compressed matter generated by the small scale matter clustering and relaxation. The influence of this factor, well known as the ‘finger of God’ effect, is opposite to that discussed above. It changes the observed particle position within walls along the line-of-sight what blurs the wall boundary and increases the thickness of observed walls. It artificially removes the high velocity particles from the selected wall and effectively decreases the surface density of walls selected in redshift space with respect to the estimates (3.13).
The impact of this factor can be approximately described by a modification of PDF of wall surface density,
$$N_m^{rd}=\frac{1}{\sqrt{2\pi }f_{rd}\tau }\frac{1}{\sqrt{q_w}}\text{erf}\left(\sqrt{\frac{q_w}{8f_{rd}^2\tau ^2}}\right)\times $$
$`(3.14)`$
$$\left[\mathrm{exp}\left(\frac{q_w}{8f_{rd}^2\tau ^2}\right)\mathrm{exp}\left(\frac{q_w}{8\tau ^2}\right)W(q_w,\tau ,\delta _{thr})\right],$$
and a new normalization of distribution $`N_m^{rd}`$.
The second term in the square brackets describes the artificial rejection of high velocity particles from the wall with a surface density $`q_w`$ bounded by a threshold density $`\delta _{thr}`$. In this term the exponent gives the fraction of matter accumulated by the wall in real space for some $`q_w\&\tau `$, whereas $`W(q_w,\tau ,\delta _{thr})`$ is the fraction of high velocity particles which are removed from the wall in the redshift space. The function $`W(q_w,\tau ,\delta _{thr})`$ cannot be found in the Zel’dovich theory as it depends on the distributions of particles positions and velocities arising due to the small scale clustering and relaxation of matter compressed within walls.
An other factor which can suppress the expected difference of wall characteristics, measured in the real and redshift spaces at small redshifts, is the strong matter condensation within structure elements with various richnesses. The strong matter rearrangement transforms the continuous matter infall on walls into a discontinuous one, increases the separation of infalling structure elements, even in the redshift space and, so, at least partly, prevents the erosion of wall boundaries.
These comments show that in the redshift space the Zel’dovich theory with the factor $`f_{rd}`$ given by (3.12) & (3.13) overestimates the matter concentration within walls. Therefore, instead of the factor $`f_{rd}`$ in (3.13) a factor $`\kappa _{rd}(\mathrm{\Gamma },\tau ,l_{thr})`$ should be used and the more realistic relation
$$\tau _m\sqrt{\frac{q_w}{6.55\kappa _{rd}^2}},1\kappa _{rd}f_{rd},$$
$`(3.15)`$
connects the amplitude $`\tau _m`$ with the wall richness $`q_w`$ in the redshift space.
The actual value of $`\kappa _{rd}`$ depends on the parameters of the cosmological model and on the method of identification of walls. The analysis performed below shows that for the walls selected in 3D space as described in Sec. 6.1, no growth of $`q_w`$ was found, and the parameters $`q_w\&\tau _m`$ are connected by the relation (3.1) as in the real space.
#### 3.3.2 Wall separation
The separations of richer walls is not sensitive to relatively small shifts of particle positions introduced by the random velocities, but these shifts can result in an artificial merging of poorer walls. The influence of this factor can be tested with the relation (3.2) as before.
#### 3.3.3 Velocity dispersion of matter compressed within walls and the wall thickness
In the redshift space the expression for the velocity dispersion of matter compressed within walls in the Zel’dovich theory is identical to (3.6) with a substitution of $`\tau _{rd}=\tau \kappa _{rd}`$ instead of $`\tau `$, but now it characterizes also the observed thickness of walls. For walls selected from the 3D sample of particles, as is described in Sec. 6.1, we have
$$h_w=\sqrt{12}w_wH_0^1,$$
$`(3.16)`$
$$w_w=u_0\sqrt{\frac{q_w^2\beta ^2}{12}+\frac{\tau ^2\kappa _{rd}^2q_w}{3}(1+\beta )^2}.$$
$`(3.17)`$
This value exceeds the corresponding real thickness of walls given by (3.9). The expected overdensity of compressed matter is given by
$$\delta _{rd}=l_vq_w/h_w.$$
$`(3.18)`$
## 4 Measured characteristics of large scale matter distribution.
### 4.1 Core-sampling approach
The core-sampling approach was proposed by Buryak et al. (1994) for the analysis of the galaxy distribution in deep pencil beam redshift surveys. In the original form it allows to obtain the mean free-path between the filaments and walls. It was improved and described in detail in LCRS1 where some characteristics of the large scale galaxy distribution were found for the Las Campanas Redshift Survey. For simulated matter distributions as considered here these characteristics were discussed in DMRT.
The potential of the core-sampling approach is not exhausted by these applications, and it could be used to measure parameters of the large scale matter distribution discussed in the previous Sections. Here we will use this approach to obtain the characteristics of the wall-like structure component.
The core-sampling method deals with a sample of points (galaxies) lying within relatively narrow cores – rectangular and/or cylindrical in simulations, and conical in observations – and it studies the point distribution along these cores. For some applications the transversal coordinates of points can be used as well. To take into account the selection effects, which are important for observed catalogues, appropriate corrections can be incorporated. The sampling core is characterized by the size, $`D_{core}`$, that is the side of a rectangular core or the angular diameter of a conical core.
### 4.2 Measured characteristics of walls
Here we will apply the core sampling technique to the sample of wall-like structure elements selected by a 3D-cluster analysis (DMRT, Sec. 6.1). This means, the sampling cores contain only the particles assigned to walls. Further on, all particles are projected onto the core axes and are collected into a set of clusters with a linking length $`l_{link}`$. Clusters with richness larger than a threshold richness, $`N_{min}`$, are identified with walls within the sampling core.
The measured wall parameters are sensitive to the influence of small scale clustering of matter within walls. For strongly disrupted walls and a narrow core, the results depend on the random position of high density clumps, what strongly increases the scatter of measured wall properties. The influence of this factor is partly suppressed for larger sizes of the sampling core, $`D_{core}`$.
However, the random intersection of the core with a wall boundary generates artificially poor clusters. The number of such intersections increases proportionally to $`D_{core}`$ what restricts the maximal $`D_{core}`$. To suppress the influence of this factor a threshold richness of cluster, $`N_{min}`$, was used. If however $`N_{min}`$ becomes too large, the statistical estimates become unreliable. For large $`D_{core}`$ the overlapping of projections of neighboring walls becomes also important what distorts the measured wall characteristics.
It is also important to choose an optimal linking length, $`l_{link}`$, because for small $`l_{link}`$, only the high density part of walls is measured, whereas for larger $`l_{link}`$, again the impact of the random overlapping of wall projections becomes important.
The influence of these factors cannot be eliminated completely, and our final estimates of properties of walls are always distorted to some degree. These distortions can be minimized for an optimal range of parameters $`D_{core},N_{min}\&l_{link}`$. Practically, these factors do not distort the velocity dispersion of walls, $`\sigma _v`$, which therefore provides the best characteristic of the actual evolutionary stage of the wall formation. On the other hand, the comparison of results obtained for different $`l_{link}`$ and $`D_{core}`$ allows to characterize the inner structure of walls.
#### 4.2.1 Measurement and correction of wall parameters
The richness of clusters in the core measures the surface density of walls,
$$m_{sim}=\frac{N_m}{D_{core}^2},$$
$`(4.1)`$
where $`N_m`$ is the number of particles in a cluster. The velocity of walls, $`v_{sim}`$, the velocity dispersion of particles accumulated within walls, $`w_{sim}`$, and the proper sizes of walls, $`h_{sim}`$, are found as follows:
$$r_w=\frac{1}{N_m}\underset{i=1}{\overset{N_m}{}}r_i,v_{sim}=\frac{1}{N_m}\underset{i=1}{\overset{N_m}{}}(u_iHr_i),$$
$$w_{sim}^2=\frac{1}{N_m1}\underset{i=1}{\overset{N_m}{}}(u_iHr_iv_{sim})^2,$$
$`(4.2)`$
$$h_{sim}^2=\frac{12}{N_m1}\underset{i=1}{\overset{N_m}{}}(r_ir_w)^2.$$
Here $`r_i`$, $`r_w`$ and $`u_i`$ are the coordinates of a particle, of a wall, and the velocity of a particle along the sampling core, respectively. The wall separation, $`D_{sim}`$, is measured by the distance between neighboring clusters.
The parameters $`m_{sim}`$, $`v_{sim}`$, $`w_{sim}`$ and $`h_{sim}`$ as given by (4.1) and (4.2) are found along the sampling core and, so, are not identical to the parameters discussed in Sec. 3. These parameters must be corrected for the random orientation of walls with respect to the sampling core. The impact of this factor increases the measured surface density, and the corrected wall surface density, $`m_c`$, is connected with the measured one by
$$m_c=m_{sim}\mathrm{cos}\varphi ,0\varphi \pi /2,$$
$$m_c=0.5m_{sim},$$
$`(4.3)`$
where $`\varphi `$ is a random polar angle between the core and the vector orthogonal to the surface of the wall, and the averaging is performed in a spherical coordinate system. Corrected values of the wall velocity and the walls thickness are as follows:
$$v_c=v_{sim}\sqrt{3},h_c=h_{sim}/\sqrt{3}.$$
$`(4.4)`$
In the redshift space the wall thickness is connected with the velocity dispersion by Eq. (3.16). The velocity dispersion within walls was found to be almost isotropic (DMRT), and, so, we will use the measured $`w_{sim}`$ as the actual velocity dispersion across walls.
The measured PDF of the wall surface density, $`N_m(m_c)`$, and the mean wall surface density, $`m_c`$, are distorted due to the small statistics of rich walls and rejection of poor walls with a richness $`N_mN_{min}`$. The correction for these distortions can be estimated by comparing the simulated PDF with the expected PDF (3.1).
To do this we will fit the measured PDF to the function
$$N_m=\frac{a_m}{\sqrt{x_m}}e^{x_m}\text{erf}(\sqrt{x_m}),x_m=\frac{b_mm_{sim}}{m_{sim}}.$$
$`(4.5)`$
The parameter $`b_m`$ describes deviations of measured and expected mean surface density of walls $`m_{sim}`$, and $`a_m`$ is a normalization factor. If the measured PDF is well fitted to the function (4.5) then the value
$$m_t=m_c/b_m$$
$`(4.6)`$
can be taken as a measure of the ‘true’ mean surface density of walls.
Finally, the mean dimensionless surface density of walls, $`q_w`$ and the amplitudes of perturbations, $`\tau _m\&\tau _v`$, measured by the surface density and velocity of wall-like structure elements, can be estimated as follows:
$$q_w=\frac{m_{sim}}{2b_ml_vn_p},\tau _m=\sqrt{\frac{q_w}{6.55}},\tau _v=\frac{\sqrt{v_{sim}^2}}{u_0}.$$
$`(4.7)`$
The small statistics of rich and poor walls distorts also the measured wall separation, $`D_{sep}`$. The expected distribution of wall separations is exponential, and therefore it is possible to correct the mean separation using the fit of the measured PDF, $`N_{sep}(D_{sim})`$, to the function
$$N_{sep}=a_{sep}\mathrm{exp}(b_{sep}D_{sim}/D_{sim}).$$
$`(4.8)`$
As before, the parameter $`b_{sep}`$ describes deviations of the measured and expected mean separation of walls, and $`a_{sep}`$ is a normalization factor. If the measured PDF is well fitted to the function (4.8) then the value
$$D_{sep}=D_{sim}/b_{sep},$$
$`(4.9)`$
can be taken as a measure of the ‘true’ mean separation of walls.
## 5 General characteristics of the simulated matter distribution
### 5.1 Basic simulations
The theoretical model discussed above describes the evolution of the DM distribution and, so, should be tested with the simulated DM distribution as well. Here we use three simulations as a basis for our analysis – the COBE normalized standard CDM model (SCDM), a $`\mathrm{\Lambda }`$CDM with $`\mathrm{\Omega }_\mathrm{\Lambda }>\mathrm{\Omega }_m`$, and an open CDM (OCDM) model. These models were described and investigated with 3D cluster analysis and Minimal Spanning Tree technique in DMRT. It was found that the $`\mathrm{\Lambda }`$CDM and OCDM models successfully reproduce the main observed characteristics of large scale matter distribution while the SCDM model demonstrates strong signatures of overevolution. Here we study these three models bearing in mind that only the $`\mathrm{\Lambda }`$CDM and OCDM models can be considered as realistic models of the observed large scale matter distribution. The SCDM model represents the matter distribution typical for a late evolutionary stage.
The simulations were performed with a PM code in a box of (500$`h^1`$Mpc)<sup>3</sup> with (300)<sup>3</sup> particles for the Harrison-Zel’dovich primordial power spectrum and the BBKS transfer function. The force and mass resolutions are $`0.9h^1`$Mpc and $`10^{11}M_{}`$, respectively. The point distribution in redshift space was produced by adding an apparent shift to one coordinate due to the peculiar velocity of particles.
Four mock catalogues were prepared on the basis of the OCDM model with various degrees of large scale bias between the spatial DM distribution and the ‘galaxies’. These mock catalogues were constructed by identifying randomly ‘galaxies’ with DM particles, but with a probability depending on the environmental density, thereby identifying more particles as ‘galaxies’ in high density regions (walls). These catalogues were investigated also both in real and redshift spaces.
The main characteristics of the simulations are listed in Table 1. A more detailed description can be found in DMRT.
### 5.2 Large scale amplitude of perturbations
The evolutionary stages reached in the models under discussion can be suitably characterized using the methods described in Sec. 2. The value $`\tau _T`$ listed in Table 1 characterizes the large scale amplitude used for the normalization of simulated perturbations. Other measures of the amplitude, such as $`\sigma _8,\tau _\xi \&\tau _{vel}`$, are sensitive to both the actually realized sample of random perturbations (cosmic variance) and to the nonlinear distortions of power spectrum produced during the evolution. For the considered mock catalogues these measures are also sensitive to the large scale bias between the spatial DM and ‘galaxies’ distributions what allows us to characterize it quantitatively.
The spatial matter distribution and the bias between spatial distributions of DM component and ’galaxies’ can be characterized by the correlation length, $`r_0`$, and the slop of the correlation function, $`\gamma `$, introduced in (2.11). These parameters are listed in Table 1 for all samples. Using relations (2.12) and (2.13), these values allows to calculate $`\sigma _8`$ and $`\tau _\xi `$, which are also listed in Table 1.
The characteristics of correlation function, $`r_0`$ and $`\gamma `$, are sensitive to the perturbations in scales $`k`$0.5 – 0.1h Mpc<sup>-1</sup>. As is seen from (2.12), estimates $`\tau _\xi `$ are very sensitive to the value of $`2\gamma `$, and, so, to small scale perturbations. The first zero–point of autocorrelation function, $`r_\xi 40h^1`$Mpc, can be usually found with a large uncertainty ($``$20 – 30 per cent) but its impact is reduced by the small exponent $`1\gamma /20.3`$ in (2.12).
For OCDM and $`\mathrm{\Lambda }`$CDM models the impact of small scale matter clustering is moderate, and differences between $`\tau _\xi `$ and $`\tau _{vel}`$ are found to be $``$10 per cent. The differences between the same parameters and $`\tau _T`$ can be considered as a reasonable measure of simulated ‘cosmic variance’. For these models differences between the parameter $`\tau _\xi `$ calculated for the real and redshift spaces also do not exceed $``$ 10%. For the SCDM model both $`\tau _\xi `$ and $`\tau _{vel}`$ are distorted by the strong small scale clustering. This divergence indicates that for the SCDM model the successful application of methods discussed in Sec. 3 is also in question.
The progressive growth of $`\tau _\xi `$ and $`\sigma _8`$ for mock catalogues characterizes the degree of the large scale bias between the spatial distribution of DM component and ‘galaxies’.
## 6 Properties of wall – like structure elements
The main basic characteristics of walls were discussed in DMRT for three DM and four mock catalogues mentioned above both in the real and redshift spaces. In this Sec. the wall characteristics discussed in Sec. 3 are found with the core-sampling technique for the same simulations and the same samples of walls.
### 6.1 Selection of wall-like structure elements
The sample of wall-like structure elements was selected with the two-parameter method described and exploited in DMRT. It identifies the wall-like structure elements with clusters found using a threshold linking length, $`l_{thr}`$, and a threshold richness, $`N_{thr}`$. As usual, the boundary of the clusters is defined by the threshold overdensity, $`\delta _{thr}`$, which is connected with the threshold linking length by
$$\delta _{thr}=\frac{n_{thr}}{n}=\frac{3}{4\pi nl_{thr}^3}.$$
$`(6.1)`$
The threshold richness, $`N_{thr}`$, restricts the matter fraction, $`f_w`$, associated with walls.
The main characteristics of these samples both in real and redshift spaces are listed in Table 2. The values of $`f_w`$ 0.4 – 0.45 are consistent with the theoretically expected and observed matter fraction accumulated by walls (DD99, LCRS1, LCRS2). The analysis performed in DMRT shows that for the low density models the main characteristics of such wall-like elements are similar to the observed characteristics of superclusters of galaxies (Oort 1983 a,b; LCRS2; DURS).
### 6.2 DM walls in the real space
The analysis of DM catalogues in the real space is most interesting as in this case we can study the clear signal from the gravitational interaction of compressed matter and can reveal and characterize statistically the matter relaxation. Five basic characteristics of DM walls discussed in Sec. 3, namely, the wall thickness, $`h_w`$, the dispersions of wall velocities, $`\sigma _v`$, the velocity dispersion of matter compressed within walls, $`w_w`$, the dimensionless surface density, $`q_w`$, and mean separation of walls, $`D_{sep}`$, can be found with the core-sampling method and can be compared with those found in DMRT. The surface density of walls is closely connected with the size of proto-walls as discussed in DMRT.
Comparison of such characteristics of matter distribution as $`\tau _{vel}`$ listed in Table 1 and $`\tau _v`$ and $`\tau _m`$ related to the wall properties allows us to test the influence of small scale matter clustering and other random factors discussed in Sec. 4.3, and to find the optimal ranges of core size, $`D_{core}`$, and of threshold richness, $`N_{min}`$, as well as the optimal linking length, $`l_{link}`$. The results listed in Table 2 are obtained with the linking length $`l_{link}=5h^1`$Mpc, and are averaged over 7 core sizes, $`6h^1`$Mpc$`D_{core}9h^1`$Mpc, and over 7 threshold richness, 10$`N_{min}`$35 .
#### 6.2.1 Basic characteristics of DM walls
For all models, the dispersion of wall velocities, $`\sigma _v`$, is found to be the best and most stable characteristic of the evolutionary stage reached. This is the direct consequence of the discrimination between the wall velocity and the velocity dispersion of particles compressed within walls. The PDFs, $`N_v`$, plotted in Fig. 1, are well fitted to Gaussian functions with the measured dispersion.
For the OCDM and $`\mathrm{\Lambda }`$CDM models the mean dimensionless surface density of walls, $`q_w`$, and the amplitudes, $`\tau _m\tau _v\tau _{vel}`$, are found with scatters $``$10 – 15% for the used $`N_{min}`$, $`D_{core}`$, and $`l_{link}`$. This scatter characterizes the moderate action of random factors discussed in Sec. 3.2 and the procedure of measurement. The values of $`l_vq_w`$ are consistent with estimates of the size of proto walls obtained in DMRT. The PDFs of the surface density plotted in Fig. 2 are consistent with that expected form (3.1). These results demonstrate that for lower density cosmological models the Zel’dovich approximation successfully describes these basic characteristics of rich walls.
For the SCDM model, the results listed in Table 2 are more sensitive to the method of measurement and the surface density of walls is underestimated, $`\tau _m<\tau _v\tau _{vel}`$. This difference can be mainly ascribed to the strong disruption of walls occurring at late evolutionary stages in this model. Other important factors are the faster compression and/or expansion of walls in transversal directions and the existence of richer halos of evaporated particles around the walls mixed with infalling particles. Such a halo becomes richer for larger $`\tau `$, i.e. for the $`\mathrm{\Lambda }`$CDM, and especially, for the SCDM models.
The distribution function of wall separation, $`N_{sep}`$, plotted in Fig. 3 is well fitted to (truncated) exponential distribution. The mean wall separation $`D_{sep}`$ is sensitive to the threshold richness $`N_{min}`$ and to the core size $`D_{core}`$. The separation $`D_{sep}40h^1`$Mpc, found for the lower threshold richness, $`N_{min}`$=5, and larger core sizes, $`D_{core}=9h^1`$Mpc, coincides with the results obtained in DMRT. It increases with $`N_{min}`$ as the number of rich walls progressively decreases. For smaller $`D_{core}`$ and larger $`N_{min}`$ some of highly disrupted walls are lost due to their small covering factor. This parameter can be found with relatively large scatter.
Using relation (3.2) we can compare our estimates of $`D_{sep}`$ and $`q_w`$. For all models we have
$$f_{dq}(0.750.9)f_w$$
and the mean wall separation is probably overestimated.
For all models under consideration, the mean wall thickness, $`h_w`$, is similar to that found in DMRT with the inertia tensor technique, where a wall is represented by a homogeneous ellipsoid. It is about of 2 – 4 times smaller than that expected in the Zel’dovich approximation (3.10) what bears a sign of the relaxation of gravitationally bounded DM particles within walls.
For the OCDM model the velocity dispersion of matter compressed within walls is found to be similar to the mean velocity of walls and of all particles, $`w_w\sigma _v\sigma _{vel}`$. In contrast, for the SCDM and $`\mathrm{\Lambda }`$CDM models the dispersion $`w_w`$ is about 30% smaller than that obtained for the complete walls in DMRT and the dispersions $`\sigma _{vel}`$ and $`\sigma _v`$ discussed above. This divergence characterizes statistically the evaporation of high energy particles in course of the relaxation of compressed matter and is reinforced by the procedures of measurement and wall selection. The relatively small value of $`w_w`$ demonstrates that in contrast to the clusters of galaxies the moderate degree of one dimensional matter compression within walls is not accompanied by an essential growth of velocity dispersion.
#### 6.2.2 Relaxation of compressed matter
For $`\mathrm{\Lambda }`$CDM and SCDM models the wall thickness, $`h_w`$ (3 – 4)$`h^1`$Mpc, is 2 – 3 times smaller than that expected in the Zel’dovich theory (3.10). So large compression of matter within walls means that the selected particles are strongly confined and, probably, relaxed. For 1D matter compression the relaxation is expected to be weak, but in reality it is reinforced due to the small scale clustering and disruption of walls.
The degree of relaxation reached can be characterized by the parameters $`\delta \&ϵ`$,
$$\delta =\frac{l_vq_w}{h_w},ϵ=\frac{w_w^2}{w_z^2(\tau )},$$
$`(6.3)`$
listed in Table 2. Here $`\delta `$ measures the mean degree of matter compression, and $`ϵ`$ is the mean kinetic energy of compressed particles with respect to the expectations of the Zel’dovich theory. The function $`w_z(\tau )`$ given by (3.7) is evaluated at $`\tau =\tau _v`$.
The divergence between the expectations of Zel’dovich theory and simulations is moderate for the OCDM model and becomes strong for the $`\mathrm{\Lambda }`$CDM model as the evolution progresses. For the SCDM model the estimate of $`\delta `$ is artificially decreased together with $`q_w`$. The small value of $`ϵ`$ 0.1 – 0.2 confirms an essential deficit of energy of compressed particles in comparison to that expected in the Zel’dovich theory. This deficit is partly enhanced by the procedure of wall selection, as the wall boundaries are blured, and particles placed far from the wall center are not included into walls.
In Zel’dovich theory the strong correlation of $`w_w`$ and $`h_w`$ with the wall richness, $`m_w`$, is described by expressions (3.6) & (3.9). In simulations the measured linear correlation coefficients of $`q_w`$, $`w_w`$ and $`h_w`$ are also $``$ 0.4 – 0.5 what indicates that the essential mass dependence of these parameters remains also after relaxation. To discriminate the regular and random variations of functions $`w_w`$ and $`h_w`$ we will consider the reduced wall thickness, $`\zeta `$, and the reduced velocity dispersion, $`\omega `$, which can be defined as follows:
$$h_w=h_w\mu ^{p_h}\zeta ,w_w=w_w\mu ^{p_w}\omega ,$$
$`(6.4)`$
$$\mu =m_w/m_w=q_w/q_w,p_hp_w0.30.4,$$
$$\zeta \omega 1,\sigma _\zeta \sigma _\omega 0.2.$$
In all considered cases the PDFs of the reduced velocity dispersion within walls, $`N_\omega `$, and of the reduced wall thickness, $`N_\zeta `$, can be roughly fitted to Gaussian functions. The PDFs $`N_\omega `$ are plotted in Fig. 4 for all three models.
These results show that due to the strong relaxation of compressed matter the correlations between the considered characteristics of walls predicted by the Zel’dovich theory in Eqs. (3.6), (3.8) & (3.10) are replaced by relations (6.4) which are also universal.
### 6.3 DM walls in the redshift space
If the analysis of wall characteristics in the real space allows to reveal the influence of gravitational interaction of the compressed matter, then a similar analysis performed in the redshift space reveals the influence of random velocities on the observed characteristics of the large scale matter distribution.
In the redshift space the analysis of wall characteristics was performed for samples of walls selected as described in Sec. 6.1 . As was shown in DMRT, in low density cosmological models the main characteristics of these walls are similar to the observed characteristics of superclusters of galaxies. The determination of wall characteristics and their corrections are discussed in Sec. 4.3. The wall parameters were found in the same ranges of $`D_{core}`$ and $`N_{min}`$ as in the real space for $`l_{link}=5h^1`$Mpc.
The main results are listed in Table 2 and are plotted in Figs. 2 – 4.
#### 6.3.1 Walls in the real and redshift spaces
The samples of walls selected in the real and redshift spaces are not identical with each other due to influence of random velocities of particles. This difference can be suitably characterized by the fraction of the same particles assigned to walls in both spaces. Here this fraction was defined as a ratio of number of the particles, $`N_{com}`$, to the number of particles assigned to the selected walls, $`N_w`$. For all models under consideration this fraction, listed in Table 2, is
$$f_{cr}=N_{com}/N_w0.80.9.$$
Small variations of number of particles, $`N_w`$, assigned to walls in the real and redshift spaces lead to these variations.
These results indicate that the influence of high random velocities generated by the small scale wall disruption and the matter relaxation moderately distorts the sample of walls selected in the redshift space. More strong deviations between such wall parameters as the wall thickness and degree of matter compression, measured in the real and redshift spaces, are caused by the redistribution of matter within walls and the procedure of measurement rather than by the incorrect wall identification. The impact of these factors rapidly increases with $`\tau _m`$ and becomes extreme for the SCDM model.
These deviations can be sensitive to the code used for simulation (see, e.g., discussion in Splinter et al. 1998). For example, in the P<sup>3</sup>M code, these variations may increase due to larger velocities of compressed matter generated there.
#### 6.3.2 Basic characteristics of DM walls
For all three models the mean surface density of selected walls listed in Table 2 is similar to that found in the real space. This fact shows that the artificial growth of matter concentration within walls discussed in Sec. 3.2 is effectively suppressed by the influence of the velocity dispersion and the procedure of wall selection, and the relation (3.1), as before, connects the mean surface density of selected walls, $`q_w`$, with the amplitude, $`\tau _m`$. Variations of $`q_w`$ and $`\tau _m`$ with $`D_{core}`$ and $`N_{min}`$ are shown in Table 2 as a scatter of these parameters. The PDFs $`N_m`$ plotted in Fig. 2 are also similar to those found in the real space.
The mean wall separation is consistent with the estimate found in real space and, as before, for all models $`f_{dq}(0.750.9)f_w`$. The PDFs $`N_{sep}`$ plotted in Fig. 3 are also similar to those found in the real space.
In the redshift space the used method of wall identification selects mainly particles with a small relative velocity what essentially restricts the measured velocity dispersion within walls and the wall thickness. Results listed in Table 2 show that only for the OCDM model, the velocity dispersion of compressed matter is consistent with the values found in the real space and in DMRT. For $`\mathrm{\Lambda }`$CDM and SCDM models they are even smaller than those found in the real space. The measured wall thickness is now linked with the velocity dispersion by the relation (3.16).
#### 6.3.3 Characteristics of matter relaxation
In the redshift space walls are less conspicuous than in the real space but, even so, for all three models the mean overdensity, $`\delta `$, listed in Table 2, differs from the estimates based on the Zel’dovich theory (3.7). As in the real space, the velocity dispersion in the redshift space is strongly correlated with the surface density of walls, what is described by the relation (6.4) with an exponent $`p_w`$ 0.5. The PDFs of the reduced velocity dispersions, $`N_\omega `$, plotted in Fig. 4, demonstrate some excess of particles with lower $`\omega `$, but it can also be roughly fitted to a Gaussian function with $`\omega `$1 and dispersion $`\sigma _\omega 0.4`$. This dispersion is about two times larger than that in the real space.
These results show that, although in the redshift space walls are not so conspicuous as in the real space, in the range of ‘time’ $`0.2\tau 0.5`$, the relaxation of compressed dark matter can be directly recognized with these methods.
### 6.4 Walls in mock catalogues
The analysis of mock catalogues characterizes how the considered simple model of large scale bias influences the measured wall properties. These catalogues were investigated also both in the real and redshift spaces. The analysis was performed for 10 values $`N_{min}`$ (15$`N_{min}`$60) and for 7 values of the core radius $`D_{core}`$ (7$`h^1`$Mpc$`D_{core}10h^1`$Mpc) using a linking length $`l_{link}=5h^1`$Mpc. The main results averaged over these $`N_{min}`$ and $`D_{core}`$ are listed in Table 2.
#### 6.4.1 ’Galaxy’ walls in the real space
In the real space for all mock catalogues the parameters $`\tau _v`$, $`h_w`$ and $`ϵ_w`$ are similar to those found for the basic OCDM model. The velocity dispersion of ‘galaxies’ within walls, $`w_w\sigma _v\sigma _{vel}`$, exceeds that found for the basic model by about of 20 – 30%. These variations can be attributed to the preferential identification of ‘galaxies’ in the central high density regions of walls, where the relative velocities of DM particles are also larger than the mean values. The wall thickness and the velocity dispersion of ‘galaxies’ can be reduced and turned into dimensionless quantities in the same manner as in Eq. (6.4), and the PDFs for the reduced wall thickness and velocity dispersion within walls, $`N_\zeta `$ and $`N_\omega `$, are also similar to Gaussian functions. The PDFs $`N_\omega `$ are shown in Fig. 5 for the mock<sub>4</sub> catalogue.
As was expected, the mean surface density of walls, $`q_w`$, exceeds that found for the basic OCDM model, and this excess progressively increases together with the biasing factor used. This excess can be considered as a suitable measure of the bias. This means that to characterize this bias the difference between $`\tau _m`$ and $`\tau _v`$ and/or between $`\tau _m`$ and other amplitudes measured for the same catalogues can be used together with the autocorrelation function. The growth of $`q_w`$ leads to a proportional growth of $`\delta `$, as the wall thickness is only weakly distorted.
The large scale bias increases the contrast between richer and poorer walls what is seen as an essential growth of the mean wall separation. In all mock catalogues $`D_{sep}`$ is about two times larger than in the basic OCDM model. The growth of both $`q_w`$ and $`D_{sep}`$ do not distort the relation between $`f_w`$ and $`f_{dq}`$.
#### 6.4.2 ’Galaxy’ walls in the redshift space
In the redshift space the fraction of the same particles assigned to walls both in the real and redshift spaces becomes $`f_{cr}`$ 80 – 85% (Table 2), what explains the similarity of parameters $`q_w`$ and $`\tau _m`$ listed in Tables 2 for both cases. The expected growth of wall richness in the redshift space according to (3.13) is not found, and the surface densities of walls, $`q_w`$, are, in the range of errors, the same both in the real and redshift spaces. This fact shows that for ‘galaxies’ the expected growth of wall richness in the redshift space is suppressed even more strongly than for DM component due to the relaxation of compressed matter. Then Eq. (3.1) describes correctly the time dependence of the mean wall surface density.
The parameters $`h_w\&w_w`$ for ‘galaxy’ walls are similar to those found for the underlying OCDM model. The difference of $`w_w`$ found for the same samples of walls in the real and redshift spaces and a slow decrease of $`w_w`$ for stronger biased models can be assigned to the loss of a small fraction of particles with large velocities, what demonstrates the sensitivity of these functions to the method of wall identification.
In the redshift space we have not so a reliable independent estimator of the amplitude as $`\tau _v`$. There the bias is seen as a relation of the amplitudes $`\tau _m\tau _\xi `$. This makes it difficult to quantitatively estimate the relatively moderate large scale bias in observed catalogues because both $`\tau _m\&\tau _\xi `$ are sensitive to the bias.
This discussion shows that the simple algorithm used in DMRT for the ‘galaxy’ identification does not essentially distort the basic characteristics of simulated walls, and a stronger bias can be seen as an excess of the surface density of ‘galaxies’ relative to that found for the DM component in the basic model. At the same time the mean velocity dispersions of both the DM component and the ‘galaxies’ assigned to walls, $`w_w`$, tends to be smaller than $`\sigma _v\&\tau _v`$ and other characteristics of the amplitude of perturbations.
## 7 Summary and discussion
In this paper we continue the investigation of large scale matter distribution and processes of large scale structure formation and evolution. Some aspects of these problems were discussed in our previous papers (LCRS-1, LCRS-2, DURS, DMRT, Müller et al. 1998) where the 3-dimensional analysis of the observed and simulated large scale structure was performed with the core-sampling and the Minimal Spanning Tree techniques. Another approach to this problem, based on the percolation technique, was discussed in Sahni et al. (1994), Shandarin & Yess (1998) and Sathyaprakash et al. (1998). The statistical description of structure formation and evolution based on the Zel’dovich theory of nonlinear gravitational instability can be found in Lee & Shandarin (1998) and DD99.
Here we direct our attention to the physical aspects of the process of wall formation, what implies a more detailed discussion of the properties of DM walls in real space. The simulations described and investigated in DMRT are used to test the theoretical expectations, to estimate the influence of small scale clustering and relaxation of compressed matter and other random factors, and to examine the power of the statistical methods used to describe the large scale matter distribution. Three cosmological models, at different evolutionary stages, were analyzed in the same manner, and the comparison of results obtained for these models allows us to estimate the properties of walls at various $`\tau `$.
In the redshift space the influence of small scale clustering and large velocity dispersion of compressed matter noticeably distorts some characteristics of the walls. These distortions appear also in the considered mock catalogues, and can even be enhanced by the possible large scale bias between the spatial distribution of DM and galaxies.
Some of these results may depend on the code used for the simulations (see, e.g., the discussion in Splinter et al. 1998), and they should be checked with simulations employing a code with higher spatial resolution.
### 7.1 Identification of walls
The core-sampling approach described in Sec. 5 allows us to characterize, in more details, the matter distribution along the sampling core and to estimate the uncertainty in measured properties of wall-like condensations introduced by the influence of velocity dispersion and small scale clustering. The influence of these random factors is demonstrated by comparing results obtained with various $`D_{core},N_{min}\&l_{link}`$.
Results presented in Sec. 6 show that some fraction of the early compressed matter has subsequently evaporated due to relaxation processes. These DM particles together with the infalling matter form an extended halo around the walls and, therefore, it is difficult to separate the walls from the background. The same problem is met by the correct definition of boundaries of galaxies and clusters of galaxies. It was also discussed in the DMRT, LCRS-2 and DURS, where the methods of wall selection, described in Sec. 6, were applied to simulated DM and observed galaxy distributions.
The central high density part of walls is reliably selected in all the cases, but various definitions of the wall boundaries can noticeably change the measured characteristics of walls. To provide more objective comparisons of wall characteristics the same dimensionless parameters $`f_w`$ and $`\delta _{thr}`$ should be used for identification of walls in different catalogues and simulations.
### 7.2 DM walls in the real space
#### 7.2.1 Measured characteristics of walls
The results presented in Sec. 6 show that the core-sampling approach can be successfully used for the investigation and description of the large scale matter distribution and the wall-like matter condensations. It allows to estimate the surface density, thickness, velocity dispersion and other basic parameters of DM walls corrected for the influence of random curvature and shape of walls. These parameters differ from those obtained in 3D space with the Minimal Spanning Tree and inertia tensor methods, and these methods suitably complement each other.
The measured wall characteristics can be compared with predictions of the Zel’dovich theory what reveals the influence of relaxation of compressed matter on the properties of walls and allows to correct the theoretical expectations. The small scale clustering of compressed matter and the wall disruption lead to noticeable variations of measured wall characteristics for different parameters of the sampling core. These variations are not so large for the low density models, but they increase rapidly with $`\tau `$.
The dimensionless surface density of walls, $`q_w`$, is closely connected with the size of proto-walls as discussed in DMRT, LCRS2 and DURS. The high surface density of walls, $`q_w0.6`$, found above even for the low density models, demonstrates that processes of strong nonlinear matter evolution occur at a typical scale of $`q_wl_v`$ (15 – 25)$`h^1`$Mpc. This evolution is correctly described by the Zel’dovich theory. This characteristic is sensitive to the basic cosmological parameters, $`\mathrm{\Omega }_m\&h`$, what allows us to select the class of more perspective models for further investigation.
#### 7.2.2 Relaxation of compressed matter
The problem of relaxation of compressed matter is now in the forefront, and the obtained results allow to begin discussion of the statistical characteristics of this relaxation. The analysis performed in the real space is more important for the discussion of the basic physical processes which occurred during the formation of wall-like matter condensations, such as the small scale matter clustering and the relaxation of the compressed matter. These processes generate the large velocity dispersion within walls and lead to the evaporation of high velocity particles. Thus, in all these cases an essential deficit of energy in DM walls as compared with the expectations of the Zel’dovich theory – $``$ (50 – 80)% and more – was found. The growth of this deficit with $`\tau `$ from the OCDM to SCDM models demonstrates that the DM relaxation becomes more and more important for later evolutionary stages, and its influence on the observed parameters of the large scale matter distribution becomes crucial for $`\tau 0.5`$.
The relaxation is seen in rich superclusters of galaxies such as the Perseus-Pisces (Saslaw & Haque– Copilah 1998). It is essentially accelerated and amplified by the small scale clustering of compressed matter. This clustering is clearly seen in observations as, for example, a strongly inhomogeneous galaxy distribution within the Great Wall (Ramella et al. 1992). The clusters of galaxies situated within wall-like superclusters similar to the Great Wall and the Perseus-Pisces can be considered as extreme examples of this process.
The merging of earlier formed structure elements is very important for the formation of large walls (DD99). This means that actually the relaxation occurs step by step during all the evolutionary history beginning with the formation of first low mass, high density pancakes which later are successively integrated and merged to larger structure elements. This means also that the finally reached degree of relaxation and the properties of compressed matter depend on the (unknown) evolutionary history of the considered walls and, therefore, can be characterized only statistically.
The relaxation of compressed matter destroys the tight correlation between the surface density and velocity dispersion predicted by the Zel’dovich theory (3.6), but it generates other correlations between the same characteristics described by the relations (6.4). This fact indicates that the properties of compressed matter are sufficiently general, and these characteristics can be used to improve the methods of wall selection and the description of wall properties.
The velocity dispersion within walls increases gradually with $`\tau `$ from the OCDM to the SCDM model. As was discussed in Sec. 7, the particles with high velocity are gravitationally confined and occupy preferentially the high density central regions of walls. This fact confirms that, probably, these particles are relaxed and have a (quasi)stationary distribution. This distribution is not so stationary as, for example, in clusters of galaxies, and it is slowly evolving due to the large scale matter flow along the walls and the persisting merging of neighboring structure elements, but presumably, this evolution does not significantly distort the formed matter distribution.
### 7.3 DM walls in the redshift space
The matter condensation seen in the redshift space can be partly artificially enhanced by the influence of streaming velocities. The possible influence of this effect was widely discussed over the past decade (see, e.g., Kaiser 1987, McGill 1990 a; Davis, Miller & White 1997; Hamilton 1998; Hui, Kofman & Shandarin 1999) and, as applied to properties of absorption lines in the spectra of high redshift quasars, by McGill (1990b) and more recently by Levshakov & Kegel (1996, 1997). These tendencies are also clearly seen from the direct application of the Zel’dovich approximation to the wall formation in the redshift space as was discussed in Sec. 3.2. Of course, it is impossible to decide which particles belong to walls, but we can estimate statistically the properties of DM walls identified in the redshift space. However, the influence of this uncertainty cannot be separated from the influence of the relaxation and other factors discussed above.
For all models the comparison of DM walls selected in the real and redshift spaces demonstrates, that they are composed mainly from the same particles – this fraction is about $`f_{cr}`$ (70 – 80)% (Table 2). This means that in both cases we find the same walls, and the fraction of randomly added or lost particles is indeed small. In spite of this, some properties of walls in the redshift space are quite sensitive to the velocity dispersion and to the methods of wall identification. Thus, the strong growth of wall thickness – about a factor of 2 – confirms results obtained by Melott et al. (1998). This effect is quite similar to the well known ’finger of God’ effect observed in clusters of galaxies.
The wall surface density, $`q_w`$, is most interesting, as it is directly connected with the basic cosmological parameters, $`\mathrm{\Omega }_m\&h`$. Our analysis shows that for low density models – $`\mathrm{\Lambda }`$CDM and OCDM – the measured value of $`q_w`$ is similar both in the real and redshift spaces. This means that the growth of the matter condensation within walls due to streaming velocities as predicted by the Zel’dovich theory is strongly suppressed by the influence of the matter relaxation and the transformation of a continuous matter infall to a discontinuous one. Actually similar relations connect the fundamental wall characteristics such as $`q_w`$ and $`\tau _m`$.
The velocity dispersion within walls selected in the 3D redshift space can be noticeably underestimated what is a direct consequence of the method of wall selection. In the redshift space, particles with large velocities are artificially shifted to the periphery of selected walls and, so, can be omitted from the analysis.
### 7.4 Walls in mock catalogues
For the considered mock catalogues the influence of velocity dispersion is enhanced by the methods used for ‘galaxy’ selection. The large scale bias increases the ‘galaxy’ concentration within walls and, so, increases the density gradient near the wall boundary. When the ‘galaxies’ are identified preferentially in the high density central parts of the walls (in the real space), than their velocity dispersion exceeds that for the DM particles, and this excess may be as large as $``$ (20 – 30)%. In the redshift space, the parameters of ‘galaxy’ walls such as $`h_w`$ and $`w_w`$ are similar to those in the underlying DM distribution.
The bias is clearly seen both in the real and redshift spaces as an excess of the mean surface density of walls. The comparison of parameters $`q_w\&\tau _m`$ found for observed wall-like galaxy condensations with possible independent estimates of the same parameters gives us a chance to obtain a reasonable observational estimates of the large scale bias.
### 7.5 The amplitude of large scale perturbations
These results demonstrate again that all characteristics of the amplitude and evolutionary stage of large scale structure considered in Secs. 2 & 3 are similar, but not identical to each other, as they are sensitive to different properties of perturbations. The best and most stable measure, $`\tau _v`$, comes from measurements of the velocity of structure elements. It is insensitive to the nonlinear evolution of perturbations, large scale bias and small scale clustering or relaxation of the compressed matter.
The comparison of other estimates for the same parameter $`\tau `$, namely, $`\tau _{vel},\tau _\xi ,\&\tau _m`$ obtained in the same simulations demonstrates their sensitivity to various natural and artificial factors. For the low density models – $`\mathrm{\Lambda }`$CDM and OCDM – the parameters $`\tau _v`$ and $`\tau _m`$ are usually sufficiently close to each other, what is a direct consequence of the close connection of the process of wall formation with the large scale perturbations. The parameter $`\tau _m`$ is sensitive to a possible large scale bias, but to reveal this factor, we need to have independent unbiased estimates of the same amplitude.
The most interesting independent estimate of the amplitude is $`\tau _\xi `$ which is however more sensitive to small scale matter clustering. Thus, for the SCDM model where this clustering is stronger it significantly overestimates the large scale amplitude. It is less sensitive to the large scale bias than $`\tau _m`$.
Independent estimates of the large scale amplitude come from measurements of the CMB anisotropy. The COBE data are consistent with other available estimates of cosmological parameters and of the large scale amplitude, and therefore, $`\tau _T`$ can be considered as the best estimate of the combination (2.9) of $`\mathrm{\Gamma }`$ and the amplitude. It can be connected with estimates of cosmological parameters $`\mathrm{\Omega }_m0.3,\mathrm{\Omega }_\mathrm{\Lambda }0.7`$ obtained from observations of high-redshift supernovae (Perlmutter et al. 1998). Nonetheless, $`\tau _T`$ should be corrected for a possible contribution of gravitational waves.
The investigation of the space density of clusters of galaxies and its redshift evolution (see, e.g., Bahcall & Fan 1998; Eke et al. 1998; Wang & Steinhardt 1998) seems also to be promising and can give the required independent measure of the large scale amplitude. The formation and evolution of galaxy clusters is caused by large scale perturbations, and their characteristics can be connected with these perturbations. But they are sensitive to the thermal evolution of clusters and, moreover, are related to only $``$ (10 – 15)% of matter accumulated by the clusters. This means that they are not free from random variations what is seen, in particular, as the well known variations of the autocorrelation function with the cluster sample.
The critical discussion of available measurements of cosmological parameters (Wang et al. 1999; Efstathiou 1999) shows that in spite of a large progress reached during last years, we do not have yet a reliable unbiased estimate of these parameters, and these data should be tested with respect to possible random large scale variations. The application of the discussed methods to large observed redshift surveys can help to achieve this goal.
### Acknowledgments
We are grateful to our anonymous referee for the useful comments and criticism. This paper was supported in part by Denmark’s Grundforskningsfond through its support for an establishment of Theoretical Astrophysics Center and Polish State Committee for Scientific Research grant Nr. 2-P03D-014-17. AGD also wishes to acknowledge support from the Center for Cosmo-Particle Physics ”Cosmion” in the framework of the project ”Cosmoparticle Physics”.
Appendix A
Dynamical characteristics of walls
in the Zel’dovich theory
The results obtained in DD99 allow us to discuss in more details dynamical characteristics of walls predicted in the Zel’dovich theory. The comparison of these expected and actually simulated characteristics reveals the influence of interaction and relaxation of compressed matter.
Following DD99 we define the wall formation as the intersection of two DM particles with different Lagrangian coordinates, $`𝐪_1`$ and $`𝐪_2`$. The difference of these coordinates measures the size of the pancake. Using the basic relations of the Zel’dovich theory (2.2) and (2.3), linking the Lagrangian and Eulerian coordinates and velocities of particles, we obtain the coordinate and velocity of a wall as a whole (DD99):
$$r_w=\frac{1}{l_vq_w}_{q_1}^{q_2}\mathrm{𝐧𝐫}𝑑q=\frac{l_v}{1+z}\left(q_c\tau (z)\frac{\mathrm{\Delta }\mathrm{\Phi }}{q_w}\right),$$
$`(A.1)`$
$$v_w=\frac{1}{l_vq_w}_{q_1}^{q_2}𝐧[𝐮H𝐫]𝑑q=$$
$$\frac{l_vH(z)}{1+z}\left[q_c\tau (z)(1+\beta )\frac{\mathrm{\Delta }\mathrm{\Phi }(q_w)}{q_w}\right],$$
$$𝐧=\frac{𝐪_1𝐪_2}{|𝐪_1𝐪_2|},q_c=\frac{|𝐪_1+𝐪_2|}{2l_v},q_w=\frac{|𝐪_1𝐪_2|}{l_v},$$
where $`\mathrm{\Delta }\mathrm{\Phi }(q_w)`$ is the random difference of the dimensionless gravitational potential over the wall. It is convenient to introduce the relative normalized Lagrangian coordinate of a particle within a wall, $`\vartheta `$:
$$q_p=q_c+0.5q_w\vartheta ,1\vartheta 1.$$
Using the coordinate $`\vartheta `$ we will describe the relative position and velocity of the infalling particle with the Lagrangian coordinate $`q_p`$ or $`\vartheta `$ by the functions:
$$r_{inf}=\mathrm{𝐧𝐫}r_w=\frac{l_v}{1+z}\left[\frac{q_w}{2}\vartheta \tau (z)\left(S(\vartheta )\frac{\mathrm{\Delta }\mathrm{\Phi }(q_w)}{q_w}\right)\right],$$
$$v_{inf}=\mathrm{𝐧𝐯}v_w=u(z)0.5q_w\vartheta +H(z)(1+\beta )r_{inf},$$
$`(A.2)`$
$$u(z)=H(z)l_v\beta (z)(1+z)^1.$$
Here $`S=\mathrm{𝐧𝐒}`$ is thee random dimensionless longitudinal displacement of a particle from its unperturbed Lagrangian position introduced by (2.2).
For Gaussian initial perturbations the PDF of the random function $`r_{inf}`$ is also Gaussian, and the mean value and dispersion of $`r_{inf}`$ should be found using the conditional characteristics of functions $`S`$ and $`\mathrm{\Delta }\mathrm{\Phi }`$ taking into account that a wall is formed in the point $`r=r_w`$ (DD99). In this case for walls with $`q_w`$1 we have:
$$r_{inf}\frac{l_v}{1+z}\frac{q_w^3}{4}\vartheta \sqrt{r_{inf}^2}\frac{l_v\tau (z)}{1+z}\sqrt{\frac{q_w}{3}},$$
$`(A.3)`$
and $`r_{inf}^2`$ is independent from $`\vartheta `$. This means that both random functions,
$$r_{inf}\&v_{inf}+u(z)0.5q_w\vartheta =H(z)(1+\beta )r_{inf}$$
are also independent from $`\vartheta `$. Hence, for the thickness, $`h_w`$, of a wall with the surface density $`q_w`$, and for the velocity dispersion within such a wall we have
$$h_w^2=12\frac{1}{2}_1^1𝑑\vartheta r_{inf}^2=4l_v^2\tau ^2q_w(1+z)^2,$$
$`(A.4)`$
$$w_w^2=\frac{1}{2}_1^1𝑑\vartheta v_{inf}^2=\frac{H^2l_v^2}{(1+z)^2}\left(\frac{\beta ^2}{12}q_w^2+\frac{\tau ^2(1+\beta )^2}{3}q_w\right).$$
Here the wall thickness is normalized by the thickness of a homogeneous slice.
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# 3-manifolds having complexity at most 9
## 1 Introduction
This paper is devoted to the theoretical description and illustration of results of an algorithm which has enabled us to give a complete list, without repetitions, of all closed oriented irreducible $`3`$-manifolds of complexity up to $`9`$. More interestingly, we have actually been able to give a “name” to each such manifold, i.e. to recognize its canonical decomposition into Seifert fibered spaces and hyperbolic manifolds already considered by other authors. The complexity we are referring to here is that introduced by Matveev (, see also ), given by the minimal number of vertices of a simple spine (this has been proved in to be equal to the minimal number of tetrahedra in a triangulation).
Our algorithm relies on a structural result on closed 3-manifolds. Namely, we show that all closed 3-manifolds can be obtained by combining, in a suitable sense, building blocks taken from a certain list which, at least up to complexity 9, is dramatically shorter than the list of all manifolds. The building blocks are called *bricks*, they are bounded by tori, and these tori carry a “marking” given by an embedded trivalent graph. Moreover, the combination of two bricks corresponds to the identification of two boundary tori. The main definitions and results of the theory of decomposition into bricks are stated in the rest of the present introduction and proved in the body of the paper.
Before turning to bricks, let us mention the most interesting experimental results about complexity $`9`$ which our algorithm has allowed us to discover. Recall first that it was already known to Matveev that up to complexity $`8`$ all manifolds are graph-manifolds; tables up to complexity $`6`$ are in , and up to $`7`$ in . Now, we can show that there are $`1156`$ manifolds of complexity $`9`$, $`272`$ of them are lens spaces, $`863`$ are more general graph-manifolds which do not contain non-separating tori, $`17`$ of them are torus bundles over $`S^1`$, $`10`$ of them are graph-manifolds with graph , and there are also $`4`$ closed hyperbolic manifolds. More importantly, these $`4`$ manifolds turn out to be precisely those of least known volume , in accordance with the ideas about complexity and volume stated in .
### 1.1 Bricks and assemblings of bricks
Throughout this paper we will work in the PL category, and by *manifold* we will always mean a compact orientable 3-manifold, possibly with boundary. We will call *triod* the graph with two vertices and three edges all joining one vertex to the other one. Note that a triod $`\theta `$ can be embedded in a torus $`T`$ so that $`T\theta `$ is an open 2-disc. A pair $`(M,X)`$ is said to be a *manifold with triods* if $`M`$ is a manifold with boundary consisting of tori $`T_1,\mathrm{},T_n`$ and $`X`$ is a set of triods $`\{\theta _1,\mathrm{},\theta _n\}`$, with $`\theta _i`$ embedded in $`T_i`$ so that $`T_i\theta _i`$ is a disc. The case where $`n=0`$ and $`X=\mathrm{}`$, so $`M`$ is closed, is admitted.
Let $`𝒳`$ be the set of all manifolds with triods (up to equivalence induced by homeomorphism of manifolds). If $`M`$ has non-empty boundary consisting of tori, then there are infinitely many inequivalent ways to embed triods in these tori, so there are infinitely many inequivalent pairs $`(M,X)`$ based on the same $`M`$. On the contrary, if $`M`$ is closed, then there is a unique element $`(M,\mathrm{})𝒳`$ based on $`M`$. Therefore the set of all closed orientable manifolds can be viewed as a subset of $`𝒳`$.
We will now describe three operations on $`𝒳`$ and state the crucial properties of a complexity function on $`𝒳`$ introduced and discussed in detail below in Section 2.
#### Connected sum.
The operation of connected sum “far from the boundary” obviously extends from manifolds to manifolds with triods. Namely, given $`(M,X)`$ and $`(M^{},X^{})`$ in $`𝒳`$, we define $`(M,X)\mathrm{\#}(M^{},X^{})`$ as $`(M\mathrm{\#}M^{},XX^{})`$, where $`M\mathrm{\#}M^{}`$ is one of the two possible connected sums of $`M`$ and $`M^{}`$ (recall that our manifolds are orientable but not oriented). Of course $`(S^3,\mathrm{})𝒳`$ is the identity element for operation $`\mathrm{\#}`$. We will call a pair $`(M,X)`$ *prime* if $`M`$ is, *i.e.* if $`(M,X)`$ cannot be expressed as a connected sum of pairs different from $`(S^3,\mathrm{})`$.
#### Assembling.
Given $`(M,X)`$ and $`(M^{},X^{})`$ in $`𝒳`$, we pick triods $`\theta _iX`$ and $`\theta _i^{}^{}X^{}`$ and choose a homeomorphism $`\psi :T_iT_i^{}^{}`$ such that $`\psi (\theta _i)=\theta _i^{}^{}`$. We can then construct the manifold with triods $`(N,Y)=(M_\psi M^{},(XX^{})\{\theta _i,\theta _i^{}^{}\})`$. We call this an *assembling* of $`(M,X)`$ and $`(M^{},X^{})`$ and we write $`(N,Y)=(M,X)(M^{},X^{})`$. Of course two given elements of $`𝒳`$ can only be assembled in a finite number of inequivalent ways.
Operation $``$ has an identity element, and in a special case it is the inverse operation of $`\mathrm{\#}`$. Below we will need to exclude these types of assembling, so we describe them in detail. First, set $`B_0=(T\times [0,1],\{\theta \times \{0\},\theta \times \{1\}\})`$, where $`T`$ is the torus and $`\theta T`$ is a triod such that $`T\theta `$ is a disc ($`B_0`$ is well-defined up to equivalence). Of course if we assemble any $`(M,X)𝒳`$ with $`B_0`$ we get $`(M,X)`$ again.
Let $`H`$ be the solid torus and let $`(H,\{\theta \})`$ and $`(H,\{\theta ^{}\})`$ be elements of $`𝒳`$ based on $`H`$. Assume that there exists a homeomorphism $`HH`$ with $`\psi (\theta )=\theta ^{}`$ such that $`(H,\{\theta \})(H,\{\theta ^{}\})`$ performed along $`\psi `$ gives $`(S^3,\mathrm{})`$ as a result. Then for any $`(M,X)𝒳`$ we have $`((M,X)\mathrm{\#}(H,\{\theta \}))(H,\{\theta ^{}\})=(M,X)`$ if we use the same $`\psi `$.
This discussion motivates the following definition. An assembling $`(M,X)(M^{},X^{})`$ is called *trivial* if, up to interchanging $`(M,X)`$ and $`(M^{},X^{})`$, one of the following holds:
* $`(M^{},X^{})=B_0`$, or
* $`(M^{},X^{})=(H,\{\theta ^{}\})`$ is a solid torus with triod, and $`(M,X)`$ can be decomposed as $`(M,X)=(N,Y)\mathrm{\#}(H,\{\theta \})`$ so that $`(N,Y)(S^3,\mathrm{})`$ and the assembling identifies $`\theta `$ to $`\theta ^{}`$ and $`(H,\{\theta \})(H,\{\theta ^{}\})=(S^3,\mathrm{})`$.
#### Self-assembling.
Given $`(M,X)𝒳`$, we pick two distinct triods $`\theta _i,\theta _i^{}X`$, we choose a homeomorphism $`\psi :T_iT_i^{}`$ such that $`\psi (\theta _i)`$ and $`\theta _i^{}`$ intersect transversely in two points, and we construct the manifold with triods $`(N,Y)=(M_\psi ,X\{\theta _i,\theta _i^{}\})`$. We call this a *self-assembling* of $`(M,X)`$ and we write $`(N,Y)=(M,X)`$. As above, only a finite number of self-assemblings of a given element of $`𝒳`$ are possible.
In the sequel it will be convenient to refer to a combination of assemblings and self-assemblings of pairs just as an *assembling*. Note that of course we can do the assemblings first and the self-assemblings in the end.
#### A complexity on $`𝒳`$.
One of the main ingredients of the present paper is the extension of Matveev’s definition of complexity from closed manifolds to manifolds with triods. We warn the reader that Matveev’s complexity $`c(M)`$ is defined also when $`M\mathrm{}`$, but our definition will be different in this case, namely we will have $`c(M,X)=c(M)`$ only when $`X=\mathrm{}`$, *i.e.* when $`M`$ is closed. The key properties of $`c`$, proved below, are additivity with respect to connected sum and subadditivity with respect to assembling. More precisely, we will construct in Subsection 2.1 a function $`c:𝒳`$ and show in Subsection 2.2 that it enjoys the following properties:
1. $`c(M,\mathrm{})=c(M)`$ for any $`(M,\mathrm{})𝒳`$;
2. $`c((M,X)\mathrm{\#}(M^{},X^{}))=c(M,X)+c(M^{},X^{})`$;
3. $`c((M,X)(M^{},X^{}))c(M,X)+c(M^{},X^{})`$. Moreover, when equality holds and the assembling is non-trivial, we have that $`(M,X)(M^{},X^{})`$ is prime if and only if both $`(M,X)`$ and $`(M^{},X^{})`$ are;
4. $`c((M,X))c(M,X)+6`$. Moreover, when equality holds, we have that $`(M,X)`$ is prime if and only if $`(M,X)`$ is;
5. for any $`n0`$ there is only a finite number of prime pairs $`(M,X)𝒳`$ with $`c(M,X)n`$.
Now let $`𝒳^{\mathrm{pr}}𝒳`$ be the set consisting of prime pairs. An assembling is called *sharp* if it is non-trivial and the inequality of point (3) above is actually an equality. Similarly, a self-assembling is *sharp* if in (4) we have an equality. We will say that a prime pair $`(M,X)𝒳^{\mathrm{pr}}`$ is a *brick* if it cannot be expressed as the result of a sharp assembling or a sharp self-assembling. The following easy result will be proved in Subsection 2.1 (one could actually also deduce it from property (5), but we will refrain from doing this):
###### Lemma 1.1.
The pair $`B_0`$ is the only $`(M,X)𝒳`$ such that $`c(M,X)=0`$ and $`X`$ contains at least two triods.
Induction on complexity now readily implies the following:
###### Proposition 1.2.
Every prime manifold with triods can be obtained as a sharp-assembling of some bricks.
We define now $`𝒳^{\mathrm{pr}}`$ as the set of all bricks, and note that $``$ naturally splits as $`^0^1`$, where $`^0`$ is the set of all $`(M,X)`$ with $`X=\mathrm{}`$ (*i.e.* $`M`$ is closed). Pairs in $`^0`$ cannot be used for an assembling or self-assembling, since they have no boundary. Let $`_n^j^j`$, for $`j=0,1`$, and $`𝒳_n𝒳`$ be the subsets consisting of pairs having complexity $`n`$. Proposition 1.2 and the properties of $`c`$ stated above now imply that
$$𝒳_n^{\mathrm{pr}}=_n^0\{^k(B^1\mathrm{}B^h):B^i_n^1,c(B^i)+6kn\}.$$
If one can give an unambiguous name to each closed $`^k(B^1\mathrm{}B^h)`$, then the set of all closed prime manifolds having complexity at most $`n`$ is easily constructed from $`_n`$ by listing the (finite number of) closed manifolds obtained in this way, and by then removing duplicates. For $`n9`$ it turns out that $`_n`$ consists of a very few atoroidal manifolds (with triods), and it is experimentally not so hard to give a name to each closed manifold of the form $`^k(B^1\mathrm{}B^h)`$. We will provide more details below on the recognition issue (after listing the bricks explicitly), but we want to emphasize here that the vast majority of computer time in the implementation of our algorithm was taken by the determination of bricks. Taking the list of bricks for granted, the reader could with some patience reproduce the list of manifolds by himself.
### 1.2 Bricks and manifolds up to complexity 9
The algorithm which will be explained in Section 3 has enabled us to explicitly find $`_9^0`$ and $`_9^1`$. The former consists of 19 closed manifolds naturally coming in two families $`C_{i,j}`$ and $`E_k`$, and the latter consists of only 11 manifolds with triods, denoted by $`B_0,\mathrm{},B_{10}`$ (where $`B_0`$ is the same as defined above). The elements of $`_9^0`$ are all Seifert fibered over $`S^2`$ with 3 exceptional fibers. In order to describe the elements of $`_9^1`$ we need a way to encode the possible ways a triod can sit in a torus.
###### Remark 1.3.
*Let $`T`$ be a torus. Let $`𝒯`$ be the set of unordered triples $`\{a,b,c\}`$ of elements of $`H_1(T)`$, such that every pair of elements in $`\{a,b,c\}`$ generates $`H_1(T)`$, and $`a+b+c=0`$. Let $`\theta T`$ be a triod such that $`T\theta `$ is a disc: inside $`\theta `$ we can find $`3`$ distinct closed curves, which can be oriented in order to form a triple $`\{a,b,c\}𝒯`$. The only two triples we can get like this are $`\{a,b,c\}`$ and $`\{a,b,c\}`$. Conversely, each triple $`\{a,b,c\}𝒯`$ determines a triod $`\theta T`$. It follows that triods (up to isotopy) are in one-to-one correspondence with elements of $`𝒯/_2`$, where the non-trivial element of $`_2`$ acts mapping $`\{a,b,c\}`$ to $`\{a,b,c\}`$.*
#### Bricks.
In Table 1 we list the elements in $`_9^1`$, as produced by our algorithm, where $`D_k`$ is the disc with $`k`$ holes, and the usual notation for Seifert manifolds and cusped hyperbolic manifolds is employed. Note that $`c(M)c(M,X)`$ is the complexity of $`M`$ in the usual sense , defined for any compact 3-manifold. Every $`M`$ turns out to be atoroidal. In order to describe triods as elements in $`𝒯`$, we must fix a basis $`(\mu _i,\lambda _i)`$ for $`H_1(T_i)`$ for each $`T_i`$ in $`M`$. When $`M`$ is Seifert, by removing fibered neighbourhoods of the exceptional fibers we get $`D_k\times S^1`$ with the product fibration. Then we choose $`\lambda _i`$ to be a fiber and $`\mu _i`$ to be a component of $`D_k\times \{\mathrm{point}\}`$, with orientations chosen so that $`(\mu _i,\lambda _i)`$ is a positively oriented basis. When $`M`$ is hyperbolic we choose $`\mu _i`$ and $`\lambda _i`$ to be respectively the first and second shortest geodesic, with orientations such that $`(\mu _i,\lambda _i)`$ is a positively oriented basis. In both cases, taking $`(\mu _i,\lambda _i)`$ instead of $`(\mu _i,\lambda _i)`$ as a basis does not make any difference, since triples in $`𝒯`$ are defined up to sign.
#### From bricks to manifolds.
As pointed out in the introduction, a list of all closed orientable prime manifolds with complexity at most $`9`$ can be compiled by listing and recognizing all closed manifolds obtained by assembling bricks $`B^1,\mathrm{},B^h`$ of $`_9`$ and then self-assembling $`k`$ times, with $`c(B^i)+6k9`$. We explain here the points which make this listing and recognition feasible. Note first that, by the bound on complexity, only a few assemblings, and no self-assembling, will involve $`B_5,\mathrm{},B_{10}`$. We also know that $`B_0`$ must not be used for assemblings. Moreover we can eliminate from the list all assemblings which we know *a priori* not to be sharp. For instance we have the following (proved in Section 2):
###### Proposition 1.4.
If $`(M,X)𝒳^{\mathrm{pr}}`$ and $`(M,X)B_1`$ is sharp, then $`(M,X)`$ is either $`B_1`$ or $`B_2`$.
Concerning recognition, we note now that the effect of assembling $`B_2`$ or $`B_3`$ is very easy to describe. Since $`B_2`$ is a solid torus, the assembling with $`B_2`$ along some boundary component $`T_i`$ corresponds to a Dehn filling of $`T_i`$. Finitely many different fillings are possible, and they are determined by the position in $`T_i`$ of the triod $`\theta _i`$. Now $`B_3(T\times [0,1],\{\theta _0\times \{0\},\theta _1\times \{1\}\})`$ with $`\theta _0\theta _1`$. (Even if in Table 1 the triples describing the triods are the same, the triods are not the same, because they lie on different boundary components, so the bases of homology are different due to orientation.) More precisely, the assembling with $`B_3`$ along $`T_i`$ corresponds to changing the position of the triod $`\theta _i`$ as in Fig. 1. Summing up, the successive assembling along $`T_i`$ of some $`B_3`$’s followed by the assembling of one $`B_2`$ still corresponds to a Dehn filling of $`T_i`$. One actually sees that all Dehn fillings can be generated like this, but of course the bound on complexity allows to consider only finitely many of them.
Turning to $`B_4`$ and $`B_5`$, we note that they naturally come with a Seifert fibered structure, so any manifold generated by $`B_2,\mathrm{},B_5`$ is a graph manifold, whose graph and gluing matrices are readily deduced from the pattern of assemblings giving the manifold. Since there are algorithms checking whether two such set of data give the same manifold, recognition is not a problem at this level.
Getting to assemblings involving $`B_6,\mathrm{},B_{10}`$, one first notes that they can only be assembled with $`B_2`$ and $`B_3`$, and not in many ways. Next, one checks by direct comparison with the tables in that 4 of the resulting manifolds are the 4 hyperbolic closed manifolds with least known volume. The following fact (proven in Section 5) concludes our investigation:
###### Proposition 1.5.
Let $`M`$ be a closed manifold with $`c(M)9`$ obtained by assembling a brick in $`\{B_6\mathrm{},B_{10}\}`$ and some $`B_2`$’s and $`B_3`$’s. Then either $`M`$ is one of the 4 hyperbolic manifolds just described, or the assembling is not sharp.
#### Manifolds.
Table 2 contains the data our algorithm has allowed us to discover about closed orientable prime manifolds having complexity $`c`$ for $`c9`$. We have divided the manifolds into three groups, given respectively by the elements which may be obtained by sharp-assembling $`B_0,\mathrm{},B_{10}`$ but without self-assembling, by those which require a self-assembling, and by those of $`^0`$. The three groups have been further split to give a more precise idea of which bricks are needed to generate a manifold: in particular, the vast majority of manifolds (with 11 exceptions out of 1156 manifolds in complexity 9) are obtained assembling $`\{B_2,B_3,B_4\}`$, and only a few manifolds actually require a self-assembling. An important convention in the table is that manifolds already considered in a certain line are not considered again in subsequent lines: some manifolds can be split into bricks in distinct ways.
It follows from the topology of the bricks that the elements of $`B_2,B_3_{\mathrm{non}\mathrm{self}}`$ are all lens spaces, those of $`B_2,B_3,B_4_{\mathrm{non}\mathrm{self}}`$ and $`B_2,B_3,B_5_{\mathrm{non}\mathrm{self}}`$ are more general graph-manifolds whose graph is a tree, those of $`B_3_{\mathrm{self}}`$ are torus bundles over $`S^1`$ and those of $`B_2,B_4_{\mathrm{self}}`$ are graph-manifolds with graph . As already mentioned, and explained in detail below in Section 5, the elements of $`^0`$ (namely the $`C_{i,j}`$’s and $`E_k`$’s) are all Seifert fibered over $`S^2`$ with 3 exceptional fibers.
## 2 The complexity function
In this section we extend Matveev’s complexity to manifolds with triods, and we state and prove its properties.
### 2.1 Definition of complexity
A compact polyhedron $`P`$ is called *simple* if the link of every point of $`P`$ can be embedded in the space given by a circle with three radii. The points having the whole of this space as a link are called *vertices*: they are isolated and therefore finite in number.
Let $`(M,X)`$ be a manifold with triods. A sub-polyhedron $`P`$ of $`M`$ is said to be a *skeleton* of the pair $`(M,X)`$ if
* $`PM`$ is simple, and $`M(PM)`$ is an open ball;
* $`PM=X`$.
Note that each open disc $`T_i\theta _i`$ is automatically adjacent to the ball $`M(PM)`$, $`P`$ is simple, and the vertices of $`P`$ cannot lie on $`M`$. Note also that when $`\mathrm{\#}X=1`$ then $`P`$ is a *spine* of $`M`$ (*i.e.* $`M`$ collapses onto $`P`$), and when $`\mathrm{\#}X=0`$ (*i.e.* when $`M`$ is closed) then $`P`$ is a spine in the usual sense , namely $`M\{\mathrm{point}\}`$ collapses onto $`P`$. When $`\mathrm{\#}X2`$, then $`M`$ does not collapse onto $`P`$.
###### Remark 2.1.
*It is easy to prove that every $`(M,X)𝒳`$ has a skeleton: take any simple spine $`Q`$ of $`M\{\mathrm{point}\}`$, so that $`MQ=M\times [0,1)B^3`$, and assume that the various $`\theta _i\times [0,1)`$’s are incident in a generic way to $`Q`$ and to each other (here of course the $`\theta _i`$’s are the triods in $`X`$). Taking the union of $`Q`$ with the $`\theta _i\times [0,1)`$’s we get a simple $`Q^{}`$ such that $`M(Q^{}M)`$ consists of $`\mathrm{\#}X+1`$ balls. Then we get a skeleton of $`(M,X)`$ by puncturing $`\mathrm{\#}X`$ suitably chosen 2-discs embedded in $`Q^{}`$, so to get one ball only in the complement.*
###### Remark 2.2.
A definition of skeleton analogous to our one was given in for any compact manifold with any trivalent graph in its boundary. The notion of complexity we will now introduce extends to any such object.
We say that a skeleton of $`(M,X)`$ is *nuclear* if it does not collapse to a subpolyhedron which is also a skeleton of $`(M,X)`$. We say that a skeleton $`P`$ of $`(M,X)𝒳`$ is *minimal* if it is nuclear and no other skeleton of $`(M,X)`$ has fewer vertices. We define now the *complexity* $`c(M,X)`$ as the number of vertices of any minimal skeleton of $`(M,X)`$.
#### Examples with complexity zero.
* It is well-known that the only closed prime manifolds having complexity zero are $`S^3,S^2\times S^1,^3`$, and $`L_{3,1}`$.
* The trivial element $`B_0=(T\times [0,1],\{\theta \times \{0\},\theta \times \{1\}\})`$ has complexity zero, since it has the simple skeleton $`\theta \times [0,1]T\times [0,1]`$, which has no vertices.
* Let $`H`$ be the solid torus, let $`D`$ be a meridinal disc properly embedded in $`H`$ and let $`\theta H`$ be a triod containing $`D`$, as in Fig. 2-left. Then $`D\theta `$ is a skeleton of $`B_1=(H,\{\theta \})`$, which has therefore complexity zero.
* Let $`H`$ be the solid torus again, and let $`P`$ be the Möbius strip with one tongue shown in Fig. 2-centre, embedded in $`H`$ as in Fig. 2-right. Since $`P`$ has no vertices and it is a skeleton for $`B_2=(H,\{PH\})`$, then $`c(B_2)=0`$.
### 2.2 Properties of complexity
Of course we have $`c(M,\mathrm{})=c(M)`$, namely property (1) of our list. We prove in this subsection the other properties of $`c`$. This will require, together with some *ad hoc* methods, the extension to our context of some techniques used in . In the course of our arguments we will give several definitions used elsewhere in the paper, and we will prove other facts stated above.
#### Finiteness.
The proof of property (5) of complexity requires a careful discussion of the topological properties of minimal spines.
A simple polyhedron $`Q`$ is called *quasi-standard* if the link of every point is either a circle, or a circle with a diameter, or a circle with three radii (neighbourhoods of points of the three types are shown in Fig. 3).
A simple polyhedron $`Q`$ is called *quasi-standard with boundary* if in addition to these three types of points we have points having as a link either a closed segment or the union of $`3`$ closed segments with one common endpoint. Assuming $`Q`$ to be quasi-standard with boundary, we denote by $`V(Q)`$ the set of points (called *vertices* above) whose link is a circle with three radii, and by $`S(Q)`$ the union of $`V(Q)`$ with the set of points whose link is a circle with a diameter. We also denote by $`Q`$ the points of the two new types declared legal when passing from ‘quasi-standard’ to ‘quasi-standard with boundary.’ Moreover, we call *$`1`$-components* of $`Q`$ the connected components of $`S(Q)V(Q)`$ and *$`2`$-components* of $`Q`$ the connected components of $`Q(S(Q)Q)`$. If the $`2`$-components of $`Q`$ are open discs (and hence are called just *faces*), and the $`1`$-components are open segments (and hence called just edges), then we call $`Q`$ a *standard polyhedron with boundary*. For short we will often just call $`Q`$ a *standard* polyhedron, and possibly specify that $`Q`$ should or not be empty.
We state now several easy facts concerning nuclear skeleta, and prove a crucial result concerning minimal skeleta.
###### Remark 2.3.
*Let $`(M,X)`$ be a manifold with triods and let $`P`$ be a nuclear skeleton of $`(M,X)`$. Then, up to rearranging the components $`T_1,\mathrm{},T_n`$ of $`M`$, we have that $`P=Qs_1\mathrm{}s_mK`$, where:*
1. $`Q`$ is a quasi-standard polyhedron with boundary $`QX`$;
2. For $`i=1,\mathrm{},m`$ we have that $`s_i\theta _i`$ is a segment and $`Qs_i`$ appears near $`T_i`$ precisely as the minimal skeleton of $`B_1`$ appears near $`B_1`$ (see Fig. 2-left); for $`i>m`$ we have $`Q\theta _i`$;
3. $`K`$ is a graph with $`K(Qs_1\mathrm{}s_m)`$ finite and $`KV(QM)`$ empty.
###### Remark 2.4.
*Every $`(M,X)𝒳`$ has a minimal skeleton $`P^{}=Qs_1\mathrm{}s_mK^{}`$ as above, where in addition $`K^{}M=\mathrm{}`$. This is because, without changing $`\mathrm{\#}V(P)`$, we can take the ends of $`K`$ lying on $`M`$ and make them slide over $`Qs_1\mathrm{}s_m`$ until they reach $`\mathrm{int}(M)`$. Note that the regular neighbourhood of $`\theta _iX`$ in $`P^{}`$ is now either a product $`\theta _i\times [0,1]`$ or the union of an annulus and a segment, as for $`B_1`$.*
###### Remark 2.5.
*If $`P`$ is a nuclear and standard skeleton of $`(M,X)`$ then it is properly embedded, namely $`P=MP=X`$, and $`PM`$ is standard without boundary. Moreover $`PM`$ is a spine of a manifold bounded by one sphere and some tori, so $`\chi (PM)=1`$. Knowing that $`S(PM)`$ is 4-valent and denoting by $`F(P)`$ the set of faces of $`P`$, we also see that $`\mathrm{\#}F(P)\mathrm{\#}V(P)=\mathrm{\#}X+1`$.*
###### Theorem 2.6.
Let $`(M,X)𝒳`$ be prime and let $`P`$ be a minimal skeleton of $`(M,X)`$. Then:
1. If $`c(M,X)>0`$ then $`P`$ is standard;
2. If $`c(M,X)=0`$ and $`X\mathrm{}`$ then $`(M,X)\{B_0,B_1,B_2\}`$, and $`P`$ is the skeleton described in Subsection 2.1 (which is standard for $`B_0`$ and $`B_2`$ only);
3. If $`c(M,X)=0`$ and $`X=\mathrm{}`$ then $`(M,X)\{S^3,S^2\times S^1,L_{3,1},^3\}`$ and $`P`$ is not standard.
###### Proof.
Our argument closely follows . We can first rule out the case $`(M,X)=(S^2\times S^1,\mathrm{})`$, because for it we only need to show that $`P`$ is not standard. But a standard polyhedron without boundary must have vertices, while $`c(S^2\times S^1,\mathrm{})=0`$. So we proceed assuming that $`M`$ is irreducible.
We will now prove that if $`P`$ is not standard then $`(M,X)\{B_1,S^3,L_{3,1},^3\}`$, and that $`P`$ is as in Subsection 2.1 when $`(M,X)=B_1`$. To conclude we will later show that if $`P`$ is standard and $`c(M,X)=0`$ then $`(M,X)\{B_0,B_2\}`$ and $`P`$ is as prescribed.
Suppose then $`P`$ is not standard. First, if $`P`$ is a point then $`(M,X)=(S^3,\mathrm{})`$. Suppose now $`P`$ has a 1-dimensional part. So, let $`eP`$ be a segment disjoint from the 2-dimensional part of $`P`$. If $`eM`$, looking at the ball $`M(PM)`$, we deduce that there is a properly embedded disc in $`M`$ intersecting $`P`$ in a point of $`e`$. By irreducibility $`M`$ is then a solid torus, so $`(M,X)=B_1`$ and $`P`$ is as in Fig. 2-left. If $`e\mathrm{int}(M)`$, looking at the ball $`M(PM)`$ again, we see that there is a sphere $`SM`$ intersecting $`P`$ in one point of $`e`$. By irreducibility $`S`$ bounds a ball $`B`$, and $`PB`$ is easily seen to be a spine of $`B`$. Nuclearity now implies that $`PB`$ contains vertices, so $`PB`$ is a skeleton of $`(M,X)`$ with fewer vertices than $`P`$. A contradiction.
We have shown so far that $`P`$ is quasi-standard unless $`(M,X)`$ is $`S^3`$ or $`B_1`$. Since $`P`$ is not standard, either a 2-component $`f`$ is not a disc, or a 1-component is a circle. In the first case, either $`f=S^2`$, or $`f=^2`$, or $`f`$ contains a simple closed curve $`\gamma `$ which is non-trivial and orientation-preserving in $`f`$. In the first two cases we have respectively $`P=S^2`$, which is impossible, and $`P=^2`$, so $`(M,X)=^3`$. The third case is impossible: looking once more at the ball $`M(PM)`$, we deduce that there is a sphere $`SM`$ intersecting $`P`$ in $`\gamma `$, and again $`S=B`$. As above, $`PB`$ is a spine of $`B`$. By minimality $`PB`$ cannot contain vertices. It follows that $`PB`$ is a disc, which contradicts the choice of $`\gamma `$. Finally, if a 1-component of $`P`$ is a circle but all 2-components are discs, then $`P`$ must be the “triple hat,” a skeleton of $`L_{3,1}`$.
We are left to analyze the case where $`P`$ is standard and $`c(M,X)=0`$, so $`X\mathrm{}`$. Now, if $`\theta X`$ and $`p`$ is a vertex of $`\theta `$, then the three faces of $`P`$ incident to $`p`$ are the same as those incident to the other vertex of $`\theta `$. Moreover, since $`V(P)=\mathrm{}`$, again the same faces are incident to the endpoint of the edge of $`P`$ which starts at $`p`$. It easily follows that $`F(P)3`$, but $`F(P)=1+\mathrm{\#}X`$ by Remark 2.5, so $`\mathrm{\#}X`$ is either $`1`$ or $`2`$. It is now a routine matter to check that $`(M,X)`$ is respectively $`B_2`$ or $`B_0`$, with $`P`$ as prescribed.∎
The next two results show respectively property (5) of complexity and Lemma 1.1.
###### Corollary 2.7.
For any $`n0`$, only finitely many pairs in $`𝒳^{\mathrm{pr}}`$ have complexity $`n`$.
###### Corollary 2.8.
$`_0^0=\mathrm{}`$ and $`_0=_0^1=\{B_0,B_1,B_2\}`$.
###### Proof.
There are no closed bricks of complexity zero, since $`(S^2\times S^1,\mathrm{})`$, $`(S^3,\mathrm{})`$, $`(^3,\mathrm{})`$, and $`(L_{3,1},\mathrm{})`$ can be obtained assembling respectively two copies of $`B_1`$, two copies of $`B_1`$, one copy of $`B_1`$ and one of $`B_2`$, and two copies of $`B_2`$. Moreover $`B_0`$, $`B_1`$, and $`B_2`$ are not non-trivial assemblings of each other, and the conclusion follows.∎
#### Subadditivity under (self-)assembling.
Let $`(M,X)`$ and $`(M^{},X^{})`$ be two given pairs, and let $`(N,Y)`$ be obtained by assembling them. Let $`P`$ and $`P^{}`$ be minimal skeleta respectively of $`(M,X)`$ and $`(M^{},X^{})`$. The assembling is defined by an identification $`\psi :T_iT_i^{}^{}`$ with $`\psi (\theta _i)=\theta _i^{}^{}`$. Using Remark 2.3 we see that $`P_\psi P^{}`$ is simple, so it is a skeleton of $`(N,Y)`$, and that no new vertices appear. It follows that $`c(N,Y)c(M,X)+c(M^{},X^{})`$.
Let $`(M,X)`$ be a pair and let $`(N,Y)`$ be obtained from $`(M,X)`$ via a self-assembling, determined by a map $`\psi :T_iT_i^{}`$ such that $`\psi (\theta _i)`$ intersects transversely $`\theta _i^{}`$ in two points. If $`P`$ is a minimal skeleton of $`(M,X)`$ as in Remark 2.4, then $`PT_iN`$ is a skeleton for $`(N,X)`$. Moreover $`PT_i`$ has at most 6 vertices more than $`P`$ (2 from the vertices of $`\theta _i`$, 2 from those of $`\theta _i^{}`$, and 2 from $`\psi (\theta _i)\theta _i^{}`$). It follows that $`c(N,Y)c(M,X)+6`$.
#### Normal surfaces.
Let $`(M,X)`$ be a manifold with triods and let $`P`$ be a nuclear skeleton of $`(M,X)`$. The simple polyhedron $`PM`$ is now a spine of $`M`$ with a ball $`BM`$ removed. Choose a triangulation of $`PM`$, and let $`\xi _P`$ be the handle decomposition of $`MB`$ obtained thickening the triangulation of $`PM`$, as in . In this paragraph we will study closed normal surfaces in $`\xi _P`$. A connected normal surface $`S`$ is *parallel to the boundary* when it is obtained by taking one boundary component and pushing it a bit inside $`\xi _P`$. In our case, we have one such surface for each $`T_i`$, and one for $`B`$.
Two preliminary results are needed to prove our main statement on normal surfaces. The first one refers to another situation, very often considered below, where a normal surface naturally arises.
###### Proposition 2.9.
Let $`(M,X)`$ be a manifold with triods and let $`QM`$ be a quasi-standard polyhedron with $`QM=QX`$. Assume $`MQ`$ has two components $`N^{}`$ and $`N^{\prime \prime }`$. Then the faces of $`Q`$ that separate $`N^{}`$ from $`N^{\prime \prime }`$ form a closed orientable surface $`\mathrm{\Sigma }(Q)QM`$ which cuts $`M`$ into two components.
###### Proof.
Let $`e`$ be an edge of $`Q`$, and let $`\{f_1,f_2,f_3\}`$ be the triple of (possibly not distinct) faces of $`Q`$ incident to $`e`$. The number of $`f_i`$’s that separate $`N^{}`$ from $`N^{\prime \prime }`$ is even; it follows that $`\mathrm{\Sigma }(Q)`$ is a surface away from $`V(Q)Q`$. Let $`T_i`$ be a boundary component of $`M`$, containing the triod $`\theta _iX`$. Since $`T_i\theta _i`$ is a disc, which is adjacent either to $`N^{}`$ or to $`N^{\prime \prime }`$ (say $`N^{}`$), then each $`2`$-component of $`Q`$ incident to $`\theta _i`$ (there could be $`0`$, $`1`$ or 3 of them, with multiplicity) has $`N^{}`$ on both sides. So $`\mathrm{\Sigma }(Q)`$ is not adjacent to $`Q`$. Finally, since $`\mathrm{\Sigma }(Q)`$ intersects the link of each vertex either nowhere or in a loop, then $`\mathrm{\Sigma }(Q)`$ is a closed surface.
The surface $`\mathrm{\Sigma }(Q)`$ cuts $`M`$ in two components (and is thus orientable, since $`M`$ is) because $`N^{}`$ and $`N^{\prime \prime }`$ lie on opposite sides of $`\mathrm{\Sigma }(Q)`$. ∎
###### Lemma 2.10.
Let $`P`$ be a standard and nuclear skeleton of a pair $`(M,X)`$. If $`\mathrm{\#}V(P)>0`$ then every face of $`P`$ is incident to at least one vertex.
###### Proof.
Assume a face $`f`$ of $`P`$ contains no vertices, and let $`f`$ be incident to the triods $`\theta _{i_1},\mathrm{},\theta _{i_k}`$. Then $`f\theta _{i_1}\mathrm{}\theta _{i_k}`$ is a connected component of $`S(PM)`$, but $`PM`$ is standard without boundary by Remark 2.5, so $`S(PM)=f\theta _{i_1}\mathrm{}\theta _{i_k}`$, whence $`S(P)f`$ and $`V(P)=\mathrm{}`$. A contradiction. ∎
We go back now to the situation where $`P`$ is a nuclear skeleton of $`(M,X)`$.
###### Lemma 2.11.
Let $`F`$ be a closed normal surface in $`\xi _P`$. Assume that no component of $`F`$ is boundary-parallel. Then there exists a simple polyhedron $`P_F`$ embedded in $`M`$, with $`\mathrm{\#}V(P_F)\mathrm{\#}V(P)`$, such that $`P_FM=X`$ and $`M(P_FM)`$ is an open regular neighbourhood of $`F`$. Moreover, if $`P`$ is standard and $`\mathrm{\#}V(P)>0`$ then $`\mathrm{\#}V(P_F)<\mathrm{\#}V(P)`$.
###### Proof.
Being normal, $`F`$ is determined by an integer attached to each 2-component of $`PM`$. Now we cut $`PM`$ open along $`F`$ as explained in : if a 2-component bears an integer $`n`$ we replace the component by $`n+1`$ parallel ones. We get a polyhedron $`P^{}M`$ which contains $`M`$, such that $`MP^{}`$ is the disjoint union of an open ball $`B`$ and an open regular neighbourhood $`N`$ of $`F`$ in $`M`$. By removing from each torus $`T_iM`$ the open disc $`T_i\theta _i`$ we get a polyhedron $`P^{\prime \prime }`$ intersecting $`M`$ in $`X`$. Now we puncture a $`2`$-component which separates $`B`$ from $`N`$ and claim that the polyhedron $`P_F`$ is as desired. Only the inequalities between $`V(P)`$ and $`V(P_F)`$ are non-obvious.
By construction we have $`\mathrm{\#}V(P_FM)\mathrm{\#}V(PM)`$. Consider now a vertex $`v`$ of $`PM`$ contained in $`T_iM`$. Of the six germs of 2-component of $`PM`$ at $`v`$, three are actually the same $`T_i\theta _i`$, so their coefficient in $`F`$ is the same, say $`\alpha `$. Call $`\beta `$, $`\gamma `$, and $`\delta `$ the coefficients of the other three germs of 2-component at $`v`$. As we cut $`PM`$ along $`F`$ we see that $`v`$ disappears if and only if (up to permutation) $`\beta =\gamma >\delta `$. If $`v`$ does not disappear then $`\beta =\gamma =\delta `$ is even. Then we set $`k=\alpha \beta /2`$ and note that $`v`$ remains on $`M`$ if and only if $`k=0`$. Now let $`v^{}`$ be the other vertex of $`PM`$ on $`T_i`$. Since the coefficients $`(\alpha ,\alpha ,\alpha ,\beta ,\gamma ,\delta )`$ are the same at $`v^{}`$, we deduce that either $`v`$ and $`v^{}`$ both disappear, or they both stay on $`M`$, or they both move to $`\mathrm{int}(M)`$. In the last case, however, one sees that $`F`$ has $`k`$ components parallel to $`T_i`$, which is absurd. So both $`v`$ and $`v^{}`$ disappear in $`P^{\prime \prime }`$ (either already in $`P^{}`$ or when we remove $`T_i\theta _i`$). This shows that $`\mathrm{\#}V(P^{\prime \prime })\mathrm{\#}V(P)`$, so $`\mathrm{\#}V(P_F)\mathrm{\#}V(P)`$.
Suppose now $`P`$ is standard. Then $`P^{\prime \prime }`$ is the union of a quasi-standard polyhedron $`P^{\prime \prime \prime }`$ and some arcs in $`X`$. The 2-components of $`P^{\prime \prime }`$ which separate $`B`$ from $`N`$ are the same as those of $`P^{\prime \prime \prime }`$, so they give a closed surface $`\mathrm{\Sigma }P^{\prime \prime }`$ by Proposition 2.9. Since no component of $`F`$ is parallel to $`B`$ or to one of the $`T_i`$’s, the 2-component $`f`$ of $`P^{\prime \prime }`$ punctured to get $`P_F`$ cannot be a closed surface. Now if $`f`$ contains vertices of $`P^{\prime \prime }`$, we see that $`\mathrm{\#}V(P_F)<\mathrm{\#}V(P^{\prime \prime })\mathrm{\#}V(P)`$, whence the conclusion. Suppose on the contrary that $`f`$ contains a circle $`\gamma S(P^{\prime \prime })`$ with $`\gamma V(P^{\prime \prime })=\mathrm{}`$. Note that the process of cutting $`PM`$ along $`F`$ allows to define a local injection $`\psi :P^{}PM`$, and that $`P^{\prime \prime }P^{}`$. Now, if $`\psi (\gamma )`$ contains some vertex of $`P`$ then this vertex has disappeared in the passage from $`P`$ to $`P^{\prime \prime }`$, whence the conclusion. If $`\psi (\gamma )V(P)=\mathrm{}`$ then we consider the 2-component $`g`$ of $`P^{\prime \prime }\mathrm{\Sigma }`$ incident to $`\gamma `$ and note that $`\psi (g)`$ must be a face of $`P`$ without vertices, which is absurd by Lemma 2.10.∎
###### Theorem 2.12.
If $`(M,X)𝒳`$ has a standard minimal skeleton then it is prime.
###### Proof.
For $`c(M,X)=0`$ it was shown during the proof of Theorem 2.6 that $`(M,X)`$ is $`B_0`$ or $`B_2`$, so we suppose $`c(M,X)>0`$. By contradiction, assume $`M`$ is not prime and let $`P`$ be a standard minimal skeleton of $`(M,X)`$. Then $`\xi _P`$ contains an essential normal sphere $`S`$. Such a sphere cannot be parallel to the boundary in $`\xi _P`$. Applying Lemma 2.11 we get $`P_SM`$ with $`\mathrm{\#}V(P_S)<\mathrm{\#}V(P)`$, $`P_SM=X`$, and $`M(P_SM)S\times (0,1)`$. Since $`(S\{\mathrm{point}\})\times (0,1)`$ is an open 3-ball, adding to $`P_S`$ a generic segment isotopic to $`\{\mathrm{point}\}\times (0,1)`$ we get a skeleton for $`(M,X)`$ with as many vertices as $`P_S`$. This contradicts minimality of $`P`$. ∎
#### Additivity under connected sum.
Again, we follow quite closely. Let $`(M,X)`$ and $`(M^{},X^{})`$ be manifolds with triods, and set $`(N,Y)=(M,X)\mathrm{\#}(M^{},X^{})`$. Let $`P`$ and $`P^{}`$ be skeleta of $`(M,X)`$ and $`(M^{}X^{})`$, respectively. If we take points $`pP`$ and $`p^{}P^{}`$ which are not vertices and we join them with a segment, we get a skeleton of $`(N,Y)`$. This implies that $`c(N,Y)c(M,X)+c(M^{},X^{})`$.
Let us prove the opposite inequality. Let $`P`$ be a minimal skeleton of $`(N,Y)`$. Since $`(N,Y)`$ is not prime, there is a separating normal sphere $`S`$ in $`\xi _P`$ (maybe not the one which cuts $`N`$ into $`M`$ and $`M^{}`$, as customary in normal surface theory). Let $`(N_1,Y_1)`$ and $`(N_2,Y_2)`$ be obtained by cutting $`(N,Y)`$ along $`S`$ and gluing in balls. The polyhedron $`P_S`$ given by Lemma 2.11 is now the disjoint union of two polyhedra $`P_1`$ and $`P_2`$ such that $`P_i`$ is a skeleton of $`(N_i,Y_i)`$. Moreover $`\mathrm{\#}V(P_S)=\mathrm{\#}V(P_1)+\mathrm{\#}V(P_2)\mathrm{\#}V(P)`$. Therefore $`c(N_1,Y_1)+c(N_2,Y_2)c(N,Y)`$, whence $`c(N_1,Y_1)+c(N_2,Y_2)=c(N,Y)`$. We can now go on finding essential spheres, and additivity eventually follows from uniqueness of the decomposition into primes.
#### Sharp (self-)assemblings.
We are now in a position to prove the second half of properties (3) and (4) of complexity. The case of self-assembling is actually easier, so we start from it. Let a sharp $`(N,Y)=(M,X)`$ be performed along $`\psi :T_iT_i^{}`$. Let $`P`$ be a minimal skeleton of $`(M,X)`$ as in Remark 2.4. Then $`PT_i`$ is a minimal skeleton of $`(N,Y)`$, and it is easy to see that $`P`$ is standard if and only if $`PT_i`$ is. Moreover, by Theorem 2.6 and Theorem 2.12, $`P`$ is standard if and only if $`(M,X)`$ is prime (because $`\mathrm{\#}X2`$) and $`PT_i`$ is standard if and only if $`(N,Y)`$ is prime (because $`c(N,Y)>0`$). This shows the desired conclusion that $`(M,X)`$ is prime if and only if $`(N,Y)`$ is.
To deal with assembling, we need two preliminary results. The first one, together with Theorem 2.6, implies Proposition 1.4.
###### Lemma 2.13.
Let $`(M,X)𝒳`$ be prime and assume $`c(M,X)>0`$. Then no assembling $`(M,X)B_1`$ is sharp.
###### Proof.
Let $`P`$ be a minimal skeleton for $`(M,X)`$, which is standard by Theorem 2.6, and let $`P^{}`$ be the minimal skeleton of $`B_1`$. Then $`P_\psi P^{}`$ is a skeleton for $`(M,X)B_1`$ with minimal number of vertices, but $`P_\psi P^{}`$ is not nuclear: there is a face $`f`$ of $`P`$, glued to the free segment of $`P^{}`$, which is incident to some vertex of $`P`$ by Lemma 2.10. By collapsing $`f`$ we would get a skeleton with fewer vertices, which is absurd. ∎
###### Lemma 2.14.
Let $`P`$ be a minimal skeleton of $`(M,X)𝒳^{\mathrm{pr}}`$ with $`c(M,X)>0`$. Then, for each $`\theta _iX`$, the three faces of $`P`$ incident to $`\theta _i`$ are distinct from each other.
###### Proof.
By Theorem 2.6, $`P`$ is standard. Suppose a face $`f`$ is incident more than once to some $`\theta _i`$. Let $`\alpha `$ be an arc in $`f`$ having endpoints $`p_0`$ and $`p_1`$ in two distinct edges of $`\theta _i`$, and let $`\beta `$ be an essential closed curve in $`T_iM`$ with $`\beta \theta _i=\{p_0,p_1\}`$. Now $`\beta `$ is cut by $`\{p_0,p_1\}`$ into components $`\beta _0`$ and $`\beta _1`$. Since $`M(PM)`$ is a ball, we can glue to both curves $`\alpha \beta _i`$ a disc, and the two discs together form a disc $`DM`$ with $`D=\beta `$. Since $`\beta `$ is essential, $`M`$ is a solid torus and $`(M,X)=B_1`$. ∎
Now let $`(N,Y)=(M,X)(M^{},X^{})`$ be a sharp assembling along some map $`\psi :T_iT_i^{}^{}`$. Recall that we want to show that $`(N,Y)`$ is prime if and only if both $`(M,X)`$ and $`(M^{},X^{})`$ are. Assume first that $`c(N,Y)=0`$. If $`(M,X)`$ and $`(M^{},X^{})`$ are prime, by Theorem 2.6 $`(N,Y)`$ is a lens space, so it is prime. If $`(N,Y)`$ is prime, we consider the prime factorization of $`(M,X)`$ and $`(M^{},X^{})`$, and note that $`\psi `$ assembles one factor $`W`$ of $`(M,X)`$ to one factor $`W^{}`$ of $`(M^{},X^{})`$. If $`WW^{}(S^3,\mathrm{})`$, then, since $`(N,Y)`$ is prime, $`(M,X)=W`$ and $`(M^{},X^{})=W^{}`$, and we are done. Otherwise, up to permutation, $`(M,X)=Z\mathrm{\#}W`$ and $`(M^{},X^{})=W^{}`$. By additivity of $`c`$ under $`\mathrm{\#}`$ and Theorem 2.6, $`W`$ and $`W^{}`$ are solid tori, and the assembling is trivial.
Now let $`c(N,Y)`$ be positive. Up to permutation, $`c(M,X)>0`$. Let $`P`$ and $`P^{}`$ be minimal skeleta of $`(M,X)`$ and $`(M^{},X^{})`$ respectively, so $`P_\psi P^{}`$ is a minimal skeleton of $`(N,Y)`$. If $`(N,Y)`$ is prime, $`P_\psi P^{}`$ is standard by Theorem 2.6, so $`P`$ and $`P^{}`$ are, and Theorem 2.12 implies the conclusion. Conversely, let $`(M,X)`$ and $`(M^{},X^{})`$ be prime. If $`(M^{},X^{})=B_1`$, we get a contradiction to Lemma 2.13. Otherwise Theorem 2.6 implies that $`P`$ and $`P^{}`$ are standard. Now, it is not a priori obvious that $`P_\psi P^{}`$ is standard, because some annular component could appear, but Lemma 2.14 applied to $`P`$ shows that they actually do not, and our argument is complete.
###### Remark 2.15.
*Given $`(H,\{\theta \})𝒳`$ with $`H`$ the solid torus, it is easy to see that there are infinitely many $`(H,\{\theta ^{}\})`$’s such that $`(H,\{\theta \})(H,\{\theta ^{}\})=(S^3,\mathrm{})`$, so $`((M,X)\mathrm{\#}(H,\{\theta \}))(H,\{\theta ^{}\})=(M,X)`$ for any $`(M,X)`$. However, the only assemblings of this sort on which complexity is additive are those where $`c(H,\{\theta \})=c(H,\{\theta ^{}\})=0`$. This can only happen if $`\{(H,\{\theta \}),(H,\{\theta ^{}\})\}`$ is $`\{B_1,B_1\}`$ or $`\{B_1,B_2\}`$, so these are the only cases which our definition of *trivial* rules out from the notion of *sharp* assembling.*
## 3 The algorithm to find bricks
We will explain in this section how we have been able to determine $`_9`$.
### 3.1 Properties of minimal skeleta of bricks
We will introduce in this subsection two more bricks $`B_3`$ and $`B_4`$, besides the $`B_0`$, $`B_1`$ and $`B_2`$ already defined above. Then we will state some results giving strong restrictions on the shape of minimal skeleta of bricks different from $`B_0,\mathrm{},B_4`$. Later we will describe the operations which we actually have carried out by computer to determine $`_9`$.
#### Minimal skeleta for $`B_3`$ and $`B_4`$.
We define $`B_3`$ and $`B_4`$ as the elements of $`𝒳`$ based on $`D_1\times S^1`$ and $`D_2\times S^1`$ respectively, where $`D_i`$ is the disc with $`i`$ holes, and the boundary triods are as decribed in Table 1 (Subsection 1.2). A skeleton for $`B_3`$ is given by the union of an annulus $`D_1\times \{\mathrm{point}\}`$ and a ribbon, glued as in Fig. 4-left. Similarly, a skeleton for $`B_4`$ is given by the union of $`D_2\times \{\mathrm{point}\}`$ and a polyhedron as in Fig. 4-right, glued as shown. This implies that $`c(B_3)1`$ and $`c(B_4)3`$. Since $`B_3`$ is prime and it is not $`B_0`$, $`B_1`$, or $`B_2`$, we have $`c(B_3)=1`$ by Theorem 2.6-(2). Using Theorem 2.6-(1) and checking by hand all standard $`P`$’s with $`\mathrm{\#}V(P)=1`$ and $`P\mathrm{}`$, we see that $`B_3`$ is a brick and actually $`_1^1=\{B_3\}`$. For $`B_4`$ we need:
###### Lemma 3.1.
Let $`(M,X)`$ be prime and different from $`B_0,\mathrm{},B_3`$. Then $`c(M,X)\mathrm{\#}X`$.
###### Proof.
Of course we can assume $`X\mathrm{}`$. Since $`_1^1=\{B_0,\mathrm{},B_3\}`$ and the inequality is easy for any non-trivial assembling of $`B_0,\mathrm{},B_3`$, we also assume $`c(M,X)2`$.
Suppose now that a face $`f`$ is incident to two distinct triods $`\theta _i,\theta _i^{}X`$. Then there is an arc $`\lambda f`$, properly embedded in $`f`$, with endpoints $`p\theta _i,p^{}\theta _i^{}`$, and two essential loops $`\gamma T_i`$ and $`\gamma ^{}T_i^{}`$ such that $`\gamma \theta _i=\{p\}`$, $`\gamma ^{}\theta _i^{}=\{p^{}\}`$. Since $`M(PM)`$ is a ball, there is an annulus $`A`$ properly embedded in $`M`$, with $`A=\gamma \gamma ^{}`$ and $`AP=\lambda `$. If some face $`gf`$ is incident to the same $`\theta _i`$ and to some other $`\theta _{i^{\prime \prime }}`$, we can construct an annulus $`B`$ in the same way. Moreover $`B=\delta \delta ^{\prime \prime }`$ with $`\mathrm{\#}(\gamma \delta )=1`$. Irreducibility allows to assume that $`AB`$ is just one segment, hence $`\theta _i^{}=\theta _{i^{\prime \prime }}`$, and then to show that $`M=T\times [0,1]`$. So $`\mathrm{\#}X=2`$, but we are assuming $`c(M,X)2`$, and the conclusion holds in this case.
By Lemma 2.14, $`P`$ has distinct faces $`f_i^{(1)},f_i^{(2)},f_i^{(3)}`$ incident to $`\theta _i`$. By what already shown we can assume up to permutation that $`f_i^{(2)}`$ and $`f_i^{(3)}`$ are not incident to any other triod in $`X`$. So $`P`$ contains at least $`2\mathrm{\#}X+1`$ distinct faces. By Remark 2.5 we have $`\mathrm{\#}F(P)=\mathrm{\#}V(P)+\mathrm{\#}X+1`$, therefore $`\mathrm{\#}V(P)\mathrm{\#}X`$. ∎
###### Proposition 3.2.
$`c(B_4)=3`$ and $`B_4`$ is a brick.
###### Proof.
First, $`B_4`$ is prime and $`\mathrm{\#}X(B_4)=3`$, so $`c(B_4)=3`$ by the previous lemma. If $`B_4`$ were not a brick then it would split as $`B_iB^{(1)}\mathrm{}B^{(k)}`$ with $`i\{1,2,3\}`$. In all cases we must have $`\mathrm{\#}X(B^{(j)})>c(B^{(j)})`$ for some $`j`$, which contradicts the previous lemma.∎
#### Super-standard skeleta.
A standard polyhedron $`P`$ (with boundary) is called *super-standard* if every face of $`P`$ is incident to $`P`$ along one segment at most. For such a $`P`$, it is easy to prove that $`S(P)`$ must be connected if $`P`$ is. The minimal skeleta of $`B_0,\mathrm{},B_4`$ we have already described are not super-standard. The following theorem will be proved in Section 4.
###### Theorem 3.3.
Let $`(M,X)`$ be a brick different from $`B_0,\mathrm{},B_4`$. Then every minimal skeleton of $`(M,X)`$ is super-standard.
###### Corollary 3.4.
Let $`(M,X)`$ be a brick different from $`B_0,\mathrm{},B_4`$. Then $`c(M,X)2\mathrm{\#}X1`$.
###### Proof.
Let $`P`$ be a minimal skeleton of $`(M,X)`$. By Remark 2.5 we have $`\mathrm{\#}F(P)\mathrm{\#}V(P)=\mathrm{\#}X+1`$. Now $`P`$ is super-standard, so each edge in $`X`$ determines a different face of $`P`$. Then $`\mathrm{\#}F(P)3\mathrm{\#}X`$, and the conclusion follows. ∎
#### Enumeration of bricks.
Let $`(M,X)`$ be a brick different from $`B_0,\mathrm{},B_4`$, and let $`P`$ be one of its minimal skeleta. We will call *filling* of $`P`$ any of the (finitely many) polyhedra obtained by glueing to $`P`$ one copy of the Möbius strip with one tongue along each of the boundary triods of $`P`$ (so $`\mathrm{\#}X`$ strips in all are glued to $`P`$). Since the Möbius strip with one tongue is a skeleton of the pair $`B_2`$ based on the solid torus, we see that a filling of $`P`$ is automatically a skeleton of a (possibly non-sharp) assembling of $`(M,X)`$ with $`\mathrm{\#}X`$ copies of $`B_2`$, hence of a closed manifold $`(N,\mathrm{})𝒳`$ obtained by Dehn-filling all boundary components of $`M`$. Note that the glueing function $`\psi `$ used to define the filling of one component $`T_i`$ of $`M`$ must map the triod $`\theta _iT_i`$ to the triod of $`B_2`$, so indeed there are finitely many possibilities. Since $`P`$ is super-standard by Theorem 3.3, it is easy to see that the fillings of $`P`$ are standard.
We will call *loop* in $`P`$ a subpolyhedron $`\gamma P`$ homeomorphic to $`S^1`$ which intersects $`S(P)`$ transversely (in particular $`\gamma V(P)=\mathrm{})`$. We define the *length* $`l(\gamma )`$ of $`\gamma `$ as the number of its intersections with the edges of $`P`$. We denote by $`(\gamma )`$ a regular neighbourhood of $`\gamma `$ in $`P`$. Note that the core of the Möbius strip has length 1 in the Möbius strip with one tongue. The following result will be shown in Section 4.
###### Theorem 3.5.
Let $`(M,X)`$ be a brick different from $`B_0,\mathrm{},B_4`$. Let $`P`$ be a minimal skeleton of $`(M,X)`$ and let $`Q`$ be any filling of $`P`$. Let $`(Q)`$ be any set of representatives of the ambient isotopy classes of length-1 loops in $`Q`$. Then:
1. The elements of $`(Q)`$ are pairwise disjoint, and $`(\gamma )`$ is a Möbius strip with one tongue for all $`\gamma (Q)`$;
2. $`(Q)`$ consists of $`\mathrm{\#}X`$ loops and $`P=Q((Q))`$.
###### Remark 3.6.
*Condition (1) in Theorem 3.5 means that:*
* for every edge $`e`$ of $`Q`$, there is no face $`f`$ of $`Q`$ triply incident to $`e`$;
* if $`f`$ is doubly incident to $`e`$ and $`f`$ is oriented, then $`e`$ is induced the same orientation twice;
* for $`i=1,2`$ let $`f_i`$ be doubly incident and $`g_i`$ be incident to the same $`e_i`$, with $`e_1e_2`$; then $`f_1f_2`$.
*In addition, taking point (2) of Theorem 3.5 for granted, super-standardness of $`P`$ (stated by Theorem 3.3) means the following:*
* with the above notation, $`f_1,f_2,g_1,g_2`$ are pairwise distinct.
We state now a result on the singular set $`S(Q)`$ of a filling $`Q`$ of a minimal skeleton $`P`$, noting first that $`S(Q)`$ depends on $`P`$ only and it is a 4-valent graph (because $`Q=\mathrm{}`$). We refer again to Section 4 for the proof.
###### Theorem 3.7.
Let $`(M,X)`$ be a brick with non-zero complexity. Let $`P`$ be a minimal skeleton of $`(M,X)`$, and let $`Q`$ be a filling of $`P`$. Then $`S(Q)`$ is connected and satisfies the following:
1. No pair of edges disconnects $`S(Q)`$.
Suppose in addition either that every torus in $`M`$ is separating or that $`c(M,X)9`$. Then:
1. If a quadruple of edges disconnects $`S(Q)`$, then one of the two resulting components must be of one of the forms shown in Fig. 5.
An important tool of our search for bricks is the following non-minimality criterion, proved in Subsection 3.2. Let us say that a loop $`\gamma `$ in a skeleton $`P`$ of $`(M,X)𝒳`$ *bounds an external disc* if there exists a closed disc $`DM`$ with $`D=\gamma `$ and $`DP=\gamma `$. A loop is *fake* if it is contained in the link of some point of $`P`$.
###### Theorem 3.8.
Let $`P`$ be a standard skeleton of a manifold with triods. If $`P`$ contains a non-fake loop which bounds an external disc and has length at most $`3`$, then $`P`$ is not minimal.
#### Computer search for bricks.
To find $`_9\{B_0,\mathrm{},B_4\}`$ we have first listed by computer the 4-valent graphs satisfying the conditions of Theorem 3.7. For each such graph $`\mathrm{\Gamma }`$, using , we have then determined the standard spines $`Q`$ of closed manifolds with $`S(Q)\mathrm{\Gamma }`$ and satisfying the conditions of Remark 3.6. Then we have tested the $`Q((Q))`$’s for minimality using Theorem 3.8. The result has been a very short list of skeleta, but actually not all of them were minimal, and some pairs of them were minimal skeleta of the same element of $`𝒳`$. To eliminate non-minimal and duplicate skeleta we have therefore used certain moves on polyhedra which are known to transform a skeleton $`P`$ into another skeleton $`P^{}`$ of the same $`(M,X)`$. Namely, we have used the Matveev-Piergallini move and some disc-replacement moves involving discs of length at most 4 (see for definitions). The result has been a list of minimal skeleta of pairwise distinct elements of $`𝒳^{\mathrm{pr}}`$, but a few non-bricks were still present. To get rid of them we have used very technical extra criteria (such as Theorem 4.14 below). The fact that the list of 30 elements eventually obtained cannot be further reduced, so all its elements are indeed bricks, follows from the (easy) fact that no element of the list is obtained via a sharp-assembling from the other ones.
###### Remark 3.9.
*The bound $`c(M,X)9`$ in Theorem 3.7 is definitely not sharp, and we actually conjecture the theorem to be true for any complexity. Moreover, if an assembling of some bricks is a manifold in which each torus is separating, then the same happens in the individual bricks. Therefore, if one wants to search for closed atoroidal manifolds only, the search for bricks can be restricted to those in which each torus is separating, to which the whole of Theorem 3.7 applies.*
We explain now how Theorem 3.7 helps saving space in the search for bricks, by ruling out most graphs as possible $`S(Q)`$’s. Namely, let $`𝒦`$ be the set of all 4-valent graphs, let $`𝒦_{\mathrm{brick}}𝒦`$ consist of all $`S(Q)`$’s where $`Q`$ is a filling of some minimal skeleton of some brick, and let $`𝒦`$ consist of the graphs satisfying both the constraints of Theorem 3.7. We know that $`𝒦_{\mathrm{brick}}𝒦`$ (at least in complexity $`9`$, or for bricks in which all tori are separating). Table 3 lists, up to 10 vertices, the number of elements of each of these sets, showing that indeed $`\mathrm{\#}`$ is a lot smaller than $`\mathrm{\#}𝒦`$, and not so far from $`\mathrm{\#}𝒦_{\mathrm{brick}}`$. We still have not determined the bricks with 10 vertices.
### 3.2 The non-minimality criterion
We prove here Theorem 3.8.
###### Remark 3.10.
*Let $`(M,X)𝒳`$ be given together with a standard skeleton $`P`$. A closed surface $`FP`$ contains a graph $`H=FS(P)`$ with vertices having valency $`3`$ and $`4`$. If $`F`$ is orientable, then we can choose a transverse orientation and give each edge $`e`$ of $`H`$ a red or black color, depending on whether $`P`$ locally lies on the positive or on the negative side of $`F`$ near $`e`$. A vertex with valency $`3`$ is adjacent to edges with the same color, and a vertex with valency $`4`$ is adjacent to two opposite red edges and two opposite black ones.*
*Proof of Theorem 3.8.* Let $`D`$ be an external disc bounded by a loop as in the statement. If we add $`D`$ and remove a face in $`\mathrm{\Sigma }(PD)`$ we get another skeleton of $`P`$. We prove now that there is a face in $`\mathrm{\Sigma }(PD)`$ incident to more than $`l(D)`$ distinct vertices. This shows that $`P`$ is not minimal.
We consider the graph $`H=S(PD)\mathrm{\Sigma }(PD)`$, which contains $`D`$. By Proposition 2.9, the surface $`\mathrm{\Sigma }(PD)`$ is orientable; we can then choose a transverse orientation and color the edges as explained in Remark 3.10. Suppose by contradiction that each face $`f\mathrm{\Sigma }(PD)`$ meets at most $`l(D)`$ vertices.
A vertex in $`D`$ has valency $`4`$ if and only if it is adjacent to two distinct edges in $`D`$ with distinct colors. If $`l(D)=1`$, then the only vertex contained in $`D`$ would have valency $`3`$, as in Fig. 6-($`1`$). So $`f_1`$ would meet at least $`2`$ distinct vertices.
If $`l(D)=2`$, then the two vertices adjacent to $`D`$ have the same valency. Suppose they both have valency $`4`$, as in Fig. 6-($`2`$). Since each $`f_i`$ meets at most $`2`$ vertices, then $`H`$ is as in Fig. 6-($`3`$), but $`M(PD)`$ would have 3 components. Suppose both vertices of $`D`$ have valency $`3`$: then $`H`$ is as in Fig. 6-($`4`$), and $`D`$ is fake. Both cases are excluded.
If $`l(D)=3`$, either all vertices met by $`D`$ have valency $`3`$, or two of them have valency $`4`$. Suppose the first case holds. If a face $`f_i`$ meets $`2`$ distinct vertices only, then the other two faces adjacent to $`D`$ coincide, as in Fig. 6-($`5`$), and meet more than $`3`$ vertices. So each $`f_i`$ meets $`3`$ distinct vertices, and $`H`$ is the 1-skeleton of the tetrahedron $`\mathrm{\Sigma }(PD)`$ as in Fig. 6-(6); hence $`D`$ is fake, which is absurd.
Finally, suppose two vertices have valency $`4`$ and one has valency $`3`$ as in Fig. 6-(7); since $`f_2`$ is incident to at most $`3`$ distinct vertices, then the distinct edges $`e_1,e_2`$ have one common endpoint; for the same reason the edges $`e_1,e_3`$ have one common endpoint. Then $`H`$ is as in Fig. 6-(8); but this is absurd since $`M(PD)`$ would have at least $`3`$ components.
$`\mathrm{}`$
## 4 Minimal skeleta of bricks
In this section we will prove Theorems 3.33.5, and 3.7. This requires the introduction of several ideas not mentioned yet. The crucial point of our work will be the analysis of the intersection between a minimal skeleton and a closed orientable surface. We warn the reader that the proofs of Theorems 4.6 and 4.14 given below are long and not very much illuminating by themselves, so they can be skipped at first. We will only consider in this section bricks having positive complexity, without further notice. So all minimal skeleta will be standard by Theorem 2.6.
### 4.1 Traces
Let $`(M,X)`$ be a manifold with triods and let $`P`$ be a standard skeleton of $`(M,X)`$. A closed surface $`F\mathrm{int}(M)`$ is said to be *simply transverse* to $`P`$ if:
1. $`F`$ is transverse to $`P`$;
2. The intersection of $`F`$ with $`MP`$ consists of a finite number of discs.
The first condition implies that $`Y=PF`$ is a finite trivalent graph disjoint from $`V(P)`$, whose vertices lie precisely at the intersection of $`Y`$ with the edges of $`P`$ and appear as in Fig. 7-left. Such a graph is called the *trace* of $`F`$.
###### Remark 4.1.
*Let a trivalent graph $`YPV(P)`$ be given, in such a way that $`YS(P)`$ consists of all the vertices of $`Y`$, each appearing as in Fig. 7-centre. We show that $`Y`$ is the trace of an essentially unique simply transverse surface $`FM`$. First, we can uniquely find a surface $`𝒩(Y)`$ with boundary, transverse to $`P`$, which collapses to $`Y`$ (see $`𝒩(Y)`$ near an edge of $`P`$ in Fig. 7-right). The boundary of $`𝒩(Y)`$ consists of a finite number of circles that lie on the boundary of a sub-ball $`B^{}`$ of $`B`$. Now we can uniquely glue disjoint discs properly embedded in $`B^{}`$ to these circles, thus getting the desired closed surface $`F`$.*
### 4.2 Traces with 2 vertices
###### Lemma 4.2.
Let $`(M,X)`$ be a brick and let $`P`$ be a minimal skeleton of $`(M,X)`$. Let $`YP`$ be the trace of an orientable surface $`FM`$. Then each edge of $`Y`$ has distinct endpoints.
###### Proof.
Suppose $`s`$ is an edge of $`Y`$ with common endpoints; since $`F`$ is orientable, the regular neighbourhood of $`s`$ in $`F`$ is an annulus, so there is a component $`D_0FY`$ with $`D_0=s`$. Then $`D_0`$ is a length-1 loop; this is impossible by Theorem 3.8, since length-1 loops are never fake. ∎
Let $`P`$ be a standard skeleton of some $`(M,X)𝒳`$ and let $`\theta _iX`$ be the triod contained in $`T_iM`$. Pushing $`\theta _i`$ a bit inside $`\mathrm{int}(P)`$ we get the trace $`Y`$ of a torus parallel to $`T_i`$. Therefore we say that such a $`Y`$ is *parallel to the boundary* (of $`P`$).
###### Proposition 4.3.
Let $`(M,X)`$ be a brick, equipped with a minimal skeleton $`P`$. Let $`Y`$ be the trace of an orientable surface $`F`$. If $`Y`$ has two vertices then it is a triod and one of the following occurs:
1. $`F`$ is a non-separating torus;
2. $`Y`$ is parallel to the boundary;
3. $`F`$ is a sphere and $`Y`$ is the link of a point in $`S(P)V(P)`$.
###### Proof.
First, $`Y`$ is a triod by Lemma 4.2. There are two possibilities for the regular neighbourhood $`𝒩(Y)`$ of $`Y`$ in $`F`$, which are shown in Fig. 8 and lead to a sphere and a torus respectively. In the first case $`FY`$ contains three external discs $`D_i`$ with $`e(D_i)=2`$. By Theorem 3.8 all the loops $`D_i`$ are fake, so $`Y=\mathrm{lk}(p)`$ for some $`pS(P)V(P)`$.
In the second case, let $`F`$ be separating, and let $`N_1`$ and $`N_2`$ be the manifolds into which $`F`$ separates $`M`$. Set $`P_i=N_iP`$ for $`i=1,2`$. Then $`(M,X)`$ is obtained by assembling the manifolds with triods $`(N_1,X_1)`$ and $`(N_2,X_2)`$, where $`X_i=P_i`$ for $`i=1,2`$. Moreover $`P_i`$ is a skeleton of $`(N_i,X_i)`$, which implies that this assembling is sharp unless it is trivial. Since $`(M,X)`$ is a brick, the assembling is trivial. Now, $`P_1`$ and $`P_2`$ are standard, so $`(M_1,X_1)`$ and $`(M_2,X_2)`$ are prime by Theorem 2.12. Therefore, the assembling must be of the first trivial type, namely $`(N_1,X_1)`$ must be $`B_0`$ up to permutation. Hence $`P_1`$ is the unique minimal skeleton of $`B_0`$, homeomorphic to $`\theta \times [0,1]`$. It follows that $`Y`$ is parallel to the boundary in $`P`$. ∎
###### Corollary 4.4.
Let $`P`$ be a minimal skeleton of a brick. Then there is no embedded face in $`P`$ incident to $`3`$ or fewer vertices. Moreover, for every edge $`e`$ of $`P`$, the three faces of $`P`$ incident to $`e`$ are distinct from each other.
###### Proof.
Let $`f`$ be an embedded face with $`3`$ or fewer vertices. A loop in $`P`$ very close to $`f`$ and disjoint from $`\overline{f}`$ bounds a disc $`D`$ parallel to $`f`$. Moreover $`l(D)3`$ and $`D`$ is not fake since $`M(PM)`$ has only one component.
Let $`fP`$ be a face incident at least twice to an edge $`e`$ of $`P`$. It follows that there is a length-1 loop $`\gamma P`$ intersecting $`e`$ once. Length-1 loops are never fake, so, by Theorem 3.8, $`\gamma `$ does not bound a disc. Therefore its regular neighbourhood $`(\gamma )`$ is a Möbius strip with one tongue, and $`(\gamma )`$ is a trace with two vertices of the disconnecting torus in $`M`$ which bounds the regular neighbourhood of $`\gamma `$ in $`M`$. Proposition 4.3 implies that $`(\gamma )`$ is boundary-parallel, so $`P`$ has no vertices. ∎
#### Co-disconnecting surfaces.
Let $`P`$ be a standard skeleton of $`(M,X)𝒳`$. Let $`YP`$ be the trace of a simply transverse orientable surface $`F`$. Every component $`D`$ of $`FY`$ is an open disc; its boundary is the union of two parts $`_1D`$ and $`_2D`$, where $`_iD`$ is the closure of the union of all edges of $`Y`$ adjacent $`i`$ times to $`D`$. If we add $`D`$ to $`P`$ we do not get a simple polyhedron, unless $`_2D=\mathrm{}`$. It is nevertheless easy to see that Proposition 2.9 holds for $`PD`$ too, namely:
###### Proposition 4.5.
Let $`B^{}`$ and $`B^{\prime \prime }`$ be the balls given by $`M(PD)`$. The faces of $`PD`$ that separate $`B^{}`$ from $`B^{\prime \prime }`$ form a closed orientable surface $`\mathrm{\Sigma }(PD)PDM`$ which cuts $`M`$ into two components.
###### Proof.
The proof of Proposition 2.9 works away from $`_2D`$. We only need to show that $`\mathrm{\Sigma }(PD)`$ is a surface near $`_2D`$: let $`f^{}`$ and $`f^{\prime \prime }`$ be the faces other than $`D`$ incident to an edge $`e_2D`$. Since $`F`$ is orientable, $`f^{}`$ is adjacent to $`B^{}`$ on both sides and $`f^{\prime \prime }`$ is adjacent to $`B^{\prime \prime }`$ on both sides (or the converse). Therefore $`f^{}`$ and $`f^{\prime \prime }`$ are disjoint from $`\mathrm{\Sigma }(PD)`$, and $`\mathrm{\Sigma }(PD)`$ is a closed surface. ∎
In the above setting we define $`\mathrm{\Sigma }_DP`$ as $`\mathrm{\Sigma }(PD)D`$, and call it the *co-disconnecting surface* of $`D`$. By Proposition 4.5, $`\mathrm{\Sigma }_D`$ is a compact surface with boundary $`_1D`$. For a subpolyhedron $`KP`$ we will denote from now on by $`(K)`$ and $`_M(K)`$ the regular neighbourhoods of $`K`$ in $`P`$ and in $`M`$ respectively.
###### Theorem 4.6.
Under the assumptions of Proposition 4.3, assume that $`c(M,X)9`$. Then $`Y`$ cannot be the trace of a non-separating torus.
###### Proof.
By contradiction let $`Y=TP`$ with $`T`$ non-separating, and put $`D=TY`$. The co-disconnecting surface $`\mathrm{\Sigma }_DP`$ is by Proposition 4.5 a closed orientable surface, which is non-empty since $`\mathrm{\Sigma }_DT`$ disconnects $`M`$, whereas $`T`$ does not. We assume that $`\mathrm{\#}(V(P)\mathrm{\Sigma }_D)`$ is minimal among all mimimal skeleta of $`(M,X)`$ for which there exists a non-separating torus whose trace is a triod.
We focus now on a component $`\mathrm{\Sigma }`$ of $`\mathrm{\Sigma }_D`$. Choosing a transverse orientation for $`\mathrm{\Sigma }`$ as in Remark 3.10, we can trace on $`\mathrm{\Sigma }`$ two trivalent graphs $`Y_+`$ and $`Y_{}`$ which intersect transversely. These graphs represent the way the rest of $`P`$ glues to $`\mathrm{\Sigma }`$, and the sign $`+`$ or $``$ depends on whether $`P`$ locally lies on the positive or on the negative side of $`\mathrm{\Sigma }`$. We show now several properties of the triple $`(\mathrm{\Sigma },Y_+,Y_{})`$ which do not require the bound 9 on complexity. Only later we will use this bound.
1. *$`\mathrm{\Sigma }Y_\pm `$ consists of planar surfaces*. Given a point $`p`$ of $`\mathrm{\Sigma }(Y_+Y_{})`$ there are two points $`p_+,p_{}`$ of $`_M(P)`$ closest to $`p`$, with $`p_+`$ on the positive side of $`\mathrm{\Sigma }`$ and $`p_{}`$ on the negative side. It is not hard to show that the map $`pp_+`$ extends to a homeomorphism of $`\mathrm{\Sigma }Y_+`$ onto an open subset of $`_M(P)S^2`$, and similarly for $`Y_{}`$.
2. *The components of $`(Y_\pm )\mathrm{\Sigma }`$ bound discs in $`M`$.* This follows from the same argument just explained.
3. *$`\mathrm{\Sigma }(Y_+Y_{})`$ consists of discs.* This is because $`\mathrm{\Sigma }P`$, $`Y_+Y_{}=\mathrm{\Sigma }S(P)`$, and $`P`$ is standard.
4. *If a component of $`\mathrm{\Sigma }(Y_\pm )`$ is not a disc then its boundary loops are essential in $`\mathrm{\Sigma }`$.* We refer to $`Y_+`$. If one of them is not, it is very easy to see that there is a disc $`\mathrm{\Delta }`$ in $`\mathrm{\Sigma }`$ such that $`Y_+\mathrm{\Delta }=\mathrm{}`$ but $`Y_+\mathrm{\Delta }\mathrm{}`$, so in particular $`Y_+\mathrm{\Delta }`$ contains vertices of $`P`$. The move suggested in Fig. 9 then contradicts minimality of $`P`$.
5. *Not all the components of $`(Y_\pm )\mathrm{\Sigma }`$ are planar.* Again we refer to $`Y_+`$. By contradiction, from points 1 and 2 and the irreducibility of $`M`$, we would readily get that $`\mathrm{\Sigma }`$ bounds a handlebody, but $`\mathrm{\Sigma }`$ is non-separating.
6. *Every component of $`Y_+`$ intersects $`Y_{}`$, and conversely.* Otherwise, since $`\mathrm{\Sigma }`$ is connected, there would exist a component of $`\mathrm{\Sigma }(Y_+Y_{})`$ with disconnected boundary, contradicting point 3.
7. *$`Y_+Y_{}`$ contains at least two points.* Assume there is only one point $`v`$ (a crossing between $`Y_+`$ and $`Y_{}`$). If a face $`f`$ of $`\mathrm{\Sigma }`$ is incident to $`v`$, then it must be multiply incident, because faces contain an even number of quadrivalent vertices (with multiplicity). If two instances of $`f`$ are adjacent to each other at $`v`$, we find in the closure of $`f`$ a length-1 loop bounding an external disc, which contradicts minimality. If two instances of $`f`$ are opposite at $`v`$, then for the same reason there is another face $`g`$ doubly incident to $`v`$, and $`gf`$. Now in the closure of $`fg`$ we can easily find a length-2 loop bounding an external disc which meets edges opposite at $`v`$. By minimality the loop must be fake, so these edges must actually be the same. Orientability of $`\mathrm{\Sigma }`$ then implies that $`f=g`$: a contradiction.
8. *If a component of $`Y_+`$ is a circle then it intersects $`Y_{}`$ in at least 4 points, and conversely.* This readily follows from Corollary 4.4 and minimality, because this circle is precisely the boundary of a face of $`P`$.
9. *No squares as in Fig. 10-left occur in $`(\mathrm{\Sigma },Y_+,Y_{})`$.* If one such square exists, we can correspondingly apply to $`P`$ one move as in Fig. 10-right. The result is a new minimal skeleton $`P^{}`$ on which $`T`$ still has trace $`Y`$, but $`\mathrm{\#}(V(P^{})\mathrm{\Sigma }_D^{})<\mathrm{\#}(V(P)\mathrm{\Sigma }_D)`$. A contradiction.
We show now how to conclude, using the fact that $`\mathrm{\#}V(P)9`$. It follows from point 5 that both $`Y_+`$ and $`Y_{}`$ have vertices. Being trivalent, they have an even number of them, and the total is at most $`92=7`$ by point 7. So up to permutation we can assume that $`Y_+`$ has 2 vertices. In particular $`(Y_+)\mathrm{\Sigma }`$ has only one non-planar component, which is homeomorphic to a punctured torus (with a component of $`Y_+`$ sitting as a triod in this torus). From point 8 we deduce that $`Y_+`$ can have at most one circular component, and it is now easy to deduce from point 1 that $`\mathrm{\Sigma }`$ indeed is a torus. Point 4 then implies that $`Y_+`$ consists of the triod only. In the rest of our proof we will always depict $`\mathrm{\Sigma }`$ cut open along $`Y_+`$. So $`\mathrm{\Sigma }`$ appears as a hexagon, and we denote by $`\mathrm{\Delta }`$ its interior.
To conclude the proof we will first show that $`Y_{}`$ also has 2 vertices, and then that it appears in one of the two shapes shown in Fig. 11. This indeed yields a contradiction to the fact that $`(M,X)`$ is a brick, since $`Y_{}Y_+`$ consists of two points, so cutting $`P`$ along $`\mathrm{\Sigma }`$ we see that $`(M,X)`$ can be obtained via a sharp-assembling.
So, let $`Y_{}`$ have 4 vertices. We claim that all the components of $`Y_{}\mathrm{\Delta }`$ are trees. If one of them is not then there is a face of $`P`$ inside $`\mathrm{\Delta }`$ and bounded by $`Y_{}`$. Then either this face has $`3`$ vertices, which contradicts Corollary 4.4, or it is a square of the first forbidden type. Our claim is proved. Now note that if $`Y_{}\mathrm{\Delta }`$ has $`\nu `$ components then it has $`4+2\nu `$ free endpoints, which give $`2+\nu `$ vertices in $`P`$. Since $`Y_+`$ has 2 vertices and $`Y_{}`$ has 4, we deduce that $`\nu =1`$ and that $`Q=P_P(\mathrm{\Sigma })`$ has no vertices. Moreover $`Q`$ is connected and standard, and $`QY_+Y_{}`$. It is not hard to show that with these constraints the only possibility for $`Q`$ is as shown in Fig. 12,
so $`_M(Q)`$ has two components. In addition, also $`\mathrm{\Sigma }Y_{}`$ consists of discs (as $`\mathrm{\Sigma }Y_+`$), and we get a contradiction because $`_M(Q)`$ should then be a sphere with some holes.
We can now assume that $`Y_{}`$ has two vertices, and show that it appears as in Fig. 11. Knowing already that $`\mathrm{\Sigma }Y_{}`$ is a disc, it is enough to show that $`Y_{}\mathrm{\Delta }`$ is connected. Suppose by contradiction that $`Y_{}\mathrm{\Delta }`$ is disconnected. Then there exists an arc $`\alpha `$ properly embedded in $`\mathrm{\Delta }`$ which separates two components of $`Y_{}\mathrm{\Delta }`$. Let us consider the endpoints of $`\alpha `$. By minimality of $`P`$, they cannot belong to the same edge of $`\mathrm{\Delta }`$, nor to two adjacent ones, otherwise we could make $`Y_{}`$ slide on $`\mathrm{\Sigma }`$ and reduce the number of vertices, as in Fig. 13.
The ends of $`\alpha `$ also cannot belong to two edges adjacent to one and the same edge, as in Fig. 14-left.
To see this, consider how many vertices of $`Y_{}`$ can lie in $`\mathrm{\Delta }^{}`$. If there are no vertices at all, then either a face of $`P`$ contained in $`\mathrm{\Delta }^{}`$ has less than 4 vertices or there is a square of the second forbidden type. If $`Y_{}`$ has both vertices in $`\mathrm{\Delta }^{}`$, then again $`\mathrm{\Delta }^{}`$ contains either a small face or a forbidden square. These cases are excluded, so there is one vertex of $`Y_{}`$ in $`\mathrm{\Delta }^{}`$, and the only possible case is shown in Fig. 14-center. Now we let $`Y_+`$ slide over $`\mathrm{\Sigma }`$ as shown in Fig. 14-right. The result is a new minimal skeleton $`P^{}`$ on which $`T`$ still has trace $`Y`$, but $`\mathrm{\Sigma }_D^{}`$ now contains one of the forbidden squares of Fig. 10, which contradicts minimality of $`\mathrm{\#}(V(P)\mathrm{\Sigma }_D)`$.
We are left to show that the endpoints of $`\alpha `$ also cannot belong to opposite edges of $`\mathrm{\Delta }`$ (Fig. 15-left).
Denote by $`\nu `$ and $`\nu ^{}`$ the number of ends of $`Y_{}\mathrm{\Delta }`$ on $`e\mathrm{\Delta }^{}`$ and on $`e^{}\mathrm{\Delta }^{}`$ respectively. If $`\nu =\nu ^{}`$ then $`\alpha `$ can be isotoped so to give rise to a length-1 loop in $`P`$ bounding an external disc: a contradiction. If $`\nu =0`$ or $`\nu ^{}=0`$ then we can replace $`\alpha `$ by a curve disjoint from $`Y_{}`$ and having ends on edges of $`\mathrm{\Delta }`$ which are not opposite, so we get back to a case already ruled out. So up to permutation we can assume that $`\nu 2`$. Now the face of $`P`$ containing the portion of arc $`\alpha ^{}`$ shown in Fig. 15-right must meet another edge of $`Y_+=\mathrm{\Delta }`$, otherwise it is either small or forbidden (recall that $`Y_{}`$ has 2 vertices only). So $`\alpha ^{}`$ extends to a properly embedded arc disjoint from $`Y_{}`$. Either $`\alpha ^{}`$ belongs to a case already ruled out, or the corresponding $`\nu +\nu ^{}`$ is smaller, and a contradiction is reached anyway. This eventually shows that $`Y_{}`$ is connected, and the proof is complete. ∎
### 4.3 Moves on traces
The key step to check the properties of bricks will be Theorem 4.14 stated below. We introduce here more new notions which will be used to prove it.
Let $`P`$ be a standard skeleton of a manifold with triods $`(M,X)`$. Given the trace $`Y`$ of a surface $`F`$, there are some obvious moves that transform $`Y`$ into another trace $`Y^{}`$ of a surface $`F^{}`$ isotopic to $`F`$. Three such moves, denoted by $`J_1`$, $`J_2`$ and $`J_3`$ and collectively called $`J`$-moves, are shown in Fig. 16. Since we will be concerned with traces of (transversely) orientable surfaces only, we can ask a $`J`$-move to transform a trace $`Y`$ into a trace $`Y^{}`$ disjoint from $`Y`$. Let $`[Y,Y^{}]`$ be the sub-polyhedron which lies between $`Y`$ and $`Y^{}`$. A sequence of moves $`Y_1\mathrm{}Y_n`$ is called a *flow* if each move $`Y_iY_{i+1}`$ is a $`J`$-move and $`[Y_{i1},Y_i][Y_i,Y_{i+1}]=Y_i`$ for all $`i`$, namely, if the moves are performed towards the same normal direction to $`Y_i`$ for all $`i`$.
###### Remark 4.7.
*A move $`J_1`$ is determined by an edge $`s`$ of $`Y`$ and a vertex $`v`$ of $`P`$ such that $`s\mathrm{lk}(v)`$, or equivalently by the cone $`vs`$ from $`v`$ based on $`s`$ (a triangle). We will sometimes say that the move is performed *along* the triangle.*
###### Remark 4.8.
*If a move $`J_1`$ transforms a trace $`Y`$ of $`F`$ into a trace $`Y^{}`$ of $`F^{}`$ then there is a natural bijection between the components of $`FY`$ and those of $`F^{}Y^{}`$.*
Let $`Y`$ be the trace of a surface $`F`$. Given a component $`D`$ of $`FY`$, we denote by $`e(D)`$ the number of edges of $`Y`$ adjacent to $`D`$, counted with multiplicity (*i.e.* an edge of $`Y`$ is counted twice if it has $`D`$ on both sides).
###### Lemma 4.9.
Let $`P`$ be a minimal skeleton of a brick. Let $`YP`$ be the trace of an orientable surface $`F`$, and let $`D`$ be a component of $`FY`$. Consider a move $`J_1`$ determined by an edge $`sD`$ of $`Y`$ and a vertex $`v`$ of $`P`$, call $`Y^{}`$ the resulting trace and $`D^{}`$ the disc corresponding to $`D`$. Then $`e(D^{})<e(D)`$ if $`e(D)<6`$ and $`e(D^{})e(D)`$ if $`e(D)=6`$.
###### Proof.
The trace $`Y^{}`$ is obtained from $`Y`$ as shown in Fig. 17; it follows from the figure that if $`e(D^{})>e(D)`$
then $`D_1=D_2=DD_3`$ and if $`e(D^{})=e(D)`$ then $`D=D_1`$ or $`D=D_2`$. By Lemma 4.2 the edges of $`Y`$ have distinct ends. Using this fact one easily sees that $`e(D)>6`$ if $`D_1=D_2=DD_3`$ and $`e(D)6`$ if $`D=D_1`$ or $`D=D_2`$, and the conclusion follows. ∎
#### Good discs.
Let $`Y`$ be the trace of a surface $`F`$. We say that a disc $`DFY`$ is *good* if all discs in $`FY`$ other than $`D`$ are contained in the same component of $`M(PD)`$.
###### Remark 4.10.
*If $`F`$ has 2 discs then these discs are good.*
###### Remark 4.11.
*If $`F`$ is orientable, then $`_P(Y)Y\times [1,1]`$. Recall that $`\mathrm{\Sigma }_D=_1D`$. Now it is not hard to show that if $`D`$ is good then the identification $`_P(Y)Y\times [1,1]`$ can be chosen so that $`\mathrm{\Sigma }_D_P(Y)_1D\times [0,1]`$, and the converse holds if $`F`$ is connected. In other words, when $`F`$ is orientable and connected, we have that $`D`$ is good if and only if $`\mathrm{\Sigma }_D`$ lies on a definite side of $`Y`$ in $`P`$.*
###### Lemma 4.12.
Under the assumptions of Lemma 4.9, suppose that $`D`$ is good and that $`\mathrm{\Sigma }_D`$ and the triangle $`vs`$ lie on the same side of $`Y`$ in $`P`$. Then $`D^{}`$ is good and $`\mathrm{\Sigma }_D^{}=\mathrm{\Sigma }_D[Y,Y^{}]`$.
###### Proof.
The condition that $`\mathrm{\Sigma }_D`$ and $`sv`$ lie on the same of side of $`Y`$ means that $`Y`$, during its transformation into $`Y^{}`$, is pushed towards $`\mathrm{\Sigma }_D`$, and the conclusion is obvious. ∎
### 4.4 Traces with 4 vertices
We prove here the key result needed to establish the properties of bricks.
###### Remark 4.13.
If a polyhedron $`P`$ is super-standard (with boundary), then it can be uniquely reconstructed from the regular neighbourhood $`(S(P))`$ of $`S(P)`$ in $`P`$, by gluing discs to each circle in $`(S(P))`$, because the rest of $`(S(P))`$ can be identified to $`P`$. Therefore here and in the sequel we will describe such $`P`$’s by drawing $`(S(P))`$ in $`^3`$. Three-dimensional pictures will be needed when $`P`$ is only standard.
###### Theorem 4.14.
Let $`P`$ be a minimal skeleton of a brick $`(M,X)`$, and let $`YP`$ be the trace of an orientable connected surface $`F`$ with $`4`$ vertices. If $`F`$ is separating, then $`Y`$ is a boundary component of a polyhedron of one of the following types:
1. $`_P(v)`$ for some $`vV(P)`$ (type 1.1), or $`_P(\lambda )`$ for an arc $`\lambda `$ properly embedded in a face of $`P`$ (type 1.2);
2. $`_P(\gamma )`$ for a length-2 loop $`\gamma `$, which is fake if it bounds an external disc;
3. One of the 5 polyhedra shown in Fig. 18, whose boundary has two components: $`Y`$ and a triod $`\theta _iP`$;
4. A polyhedron as in Fig. 19, with $`1`$ (left) or more (right) vertices;
If $`F`$ is not separating, then it is a torus and there is a minimal skeleton $`P^{}`$ of $`(M,X)`$ on which $`F`$ has a triod as a trace.
Moreover, only two types $`A`$ and $`B`$ of $`Y`$ are possible, as shown in Fig. 20. The polyhedra of types 1.1, 3.3, and 3.4 have boundaries of type $`A`$, those of types 1.2, 3.1, 3.2, 3.5, and 4 have boundaries of type $`B`$; a polyhedron of type 2 has boundary of type $`A`$ if it is a Möbius strip with two tongues, of type $`B`$ otherwise.
###### Proof.
It is enough to show that one of the following must hold:
* $`F`$ is a non-separating torus, and $`F`$ has a triod as a trace on some $`P^{}`$;
* $`Y`$ bounds one of the polyhedra of type 1-4.
So we assume (I) does not hold and show (II). Our argument is long and organized in many steps. We first describe the overall scheme stating without proof 5 assertions. Later we will provide full proofs. Let $`DFY`$ be a component having lowest $`e(D)`$.
Fact 1. *If $`e(D)\{2,3\}`$ then $`Y`$ bounds a polyhedron of type 1, 2, 3.1, 3.2, or 3.3*
Suppose then that $`e(D)4`$. Since $`Y`$ is trivalent it has 6 vertices, so $`\chi (F)=d2`$, where $`d`$ is the number of components of $`FY`$. Each component is incident to at least 4 vertices, so $`264d`$, whence $`d3`$. It easily follows that $`F`$ is a torus and $`d=2`$. Then $`FY`$ consists of two discs $`D=D_1`$ and $`D_2`$, both good by Remark 4.10. Recalling from Lemma 4.2 that all edges of $`Y`$ have distinct endpoints one easily sees that only the types $`A`$ and $`B`$ for $`Y`$ are possible. The restriction that $`e(D)4`$ then implies that up to homeomorphism there is only one possible configuration $`(F,Y_A)`$ and only one $`(F,Y_B)`$, as shown in Fig. 21. If $`Y`$ is of type A we have $`e(D)=4`$, otherwise we have $`e(D)=6`$, and the two discs of $`F`$ are completely symmetric.
Figure 21 also contains notation used throughout the proof (note that $`s_1,\mathrm{},s_4`$ are the edges in $`_1D=\mathrm{\Sigma }_D`$ both in case $`A`$ and in case $`B`$). Let $`f_i`$ be the face of $`\mathrm{\Sigma }_DS(P)`$ incident to $`s_i`$. Moreover, let $`g_j`$ be the face of $`P`$ incident to $`t_j`$. Since $`D`$ is good, we have $`g_1,g_2\mathrm{\Sigma }_D`$. Finally, let $`e_i`$ be the edge of $`P`$ which contains $`p_i`$.
Fact 2. *Either the faces $`f_1,f_2,f_3,f_4,g_1,g_2`$ are all distinct or $`Y`$ bounds a polyhedron of type 2 or 4.1.*
Assuming that $`Y`$ does not bound a polyhedron of type 2 or 4.1, it follows that the segments $`e_i\mathrm{\Sigma }_D`$ for $`i=1,\mathrm{},4`$ are distinct. Then let $`v_iV(P)`$ be the endpoint of $`e_i\mathrm{\Sigma }_D`$ not lying on $`D`$.
Fact 3. *Up to symmetry we have $`v_1=v_2`$ in case $`A`$ and either $`v_1=v_2`$ or $`v_1=v_3`$ in case $`B`$.*
Let us now set $`u=s_1`$ in case $`A`$, and either $`u=s_1`$ or $`u=t_1`$ in case $`B`$, depending on whether $`v_1=v_2`$ or $`v_1=v_3`$, so there are two edges of $`PD`$ which start at the endpoints of $`u`$ and both end at $`v_1`$. These edges are $`e_1^{}=e_1\mathrm{\Sigma }_D`$ and $`e_m^{}=e_m\mathrm{\Sigma }_D`$, with $`m\{2,3\}`$ depending on the case. Recall now that if two edges end at the same vertex then one face incident to the first edge is also incident to the second one. Since we are assuming that the $`f_i`$’s and $`g_j`$’s are distinct, we deduce that $`ue_1^{}e_m^{}`$ bounds a disc of $`PD`$, which is a triangle, *i.e.* $`u\mathrm{lk}(v_1)`$. Following Remark 4.7 we can then perform a $`J_1`$-move to which Lemma 4.9 and Lemma 4.12 apply. Denoting by $`D^{}`$ the disc corresponding to $`D`$ after the move, we have $`e(D^{})e(D)`$, and equality can hold only if $`Y`$ is of type $`B`$.
Fact 4. *If $`e(D^{})<e(D)`$ then $`Y`$ bounds a polyhedron of type 3.4 or 3.5.*
Fact 5. *If $`e(D^{})=e(D)`$ then $`Y`$ bounds a polyhedron of type 4.2.*
This establishes the theorem. We now prove our assertions.
Proof of fact 1. By Theorem 3.8 the loop $`D`$ is fake, and we can perform a move $`J_{e(D)}`$ as explained in Subsection 4.3. The result is a trace $`Y^{}`$ with 2 vertices of a surface $`F^{}`$ isotopic to $`F`$. By Proposition 4.3 either $`F^{}`$ is a non-separating torus, or $`Y^{}`$ is boundary-parallel, or we have $`Y^{}=(p)`$ for some $`pV(P)S(P)`$. In the first case, up to isotoping $`F^{}`$ back to $`F`$, getting an isotopic copy $`P^{}`$ of $`P`$, we get a contradiction to our initial assumption. In the other cases we have to see which polyhedra can result from an inverse $`J_{e(D)}`$ move applied to $`\theta _i\times [0,1]`$ or to $`(p)`$. It is now rather easy to examine all possibilities and check the assertion.
Proof of fact 2. Of course no $`f_i`$ can be equal to a $`g_j`$, because $`f_i\mathrm{\Sigma }_D`$ and $`g_j\mathrm{\Sigma }_D=\mathrm{}`$. Let us first show that if two $`f_i`$’s coincide then $`Y`$ bounds a polyhedron of type 2 or 4.1. We refer to Fig. 21 for the notation.
Two adjacent $`f_i`$’s cannot coincide because of Corollary 4.4. Up to symmetry, the only cases we are left to deal with are $`A`$-$`(f_1=f_3)`$, $`B`$-$`(f_1=f_3)`$, and $`B`$-$`(f_2=f_3)`$. In all cases we will show that $`Y`$ bounds a polyhedron of type 2 or 4.1. The key point will be to exhibit two loops that must be fake because of Theorem 3.8.
Case $`A`$-$`(f_1=f_3)`$ is examined in Fig. 22-left: since $`\alpha ^{}`$ and $`\alpha ^{\prime \prime }`$ are fake, one sees quite easily that $`Y=(\gamma )`$, where $`(\gamma )`$ is a Möbius strip with two tongues (Fig. 22-right).
Case $`B`$-$`(f_1=f_3)`$ is similar (Fig. 23-left); we have $`Y=(\gamma )`$, where $`(\gamma )`$ is an annulus with two tongues on opposite sides (Fig. 23-right).
In case $`B`$-$`(f_2=f_3)`$ we consider the loops of Fig. 24-left. Since $`\alpha ^{}`$ and $`\alpha ^{\prime \prime }`$ are fake we deduce that all the edges $`e_i\mathrm{\Sigma }_D`$ end at the same vertex $`v`$, such that $`s_2,s_3\mathrm{lk}(v)`$. We can then apply a move $`J_1`$ whose effect on $`Y`$ is shown in Fig. 24-right.
The result is a trace $`Y^{}`$ which falls into case $`A`$-$`(f_1=f_3)`$. So $`Y^{}=(\gamma )`$ with $`(\gamma )`$ a Möbius strip with two tongues. Recalling that the inverse of a $`J_1`$-move is again a $`J_1`$-move, we only need to consider which such moves can be applied to $`(\gamma )`$. The move is determined by the edge of $`(\gamma )`$ which disappears during the move: of the 6 edges of $`(\gamma )`$, 4 lead to a situation in which $`e(D)=3`$, so we exclude them. The other 2 edges are actually symmetric, and the result is of type 4.1.
To conclude the proof of Fact 2 we must show that if the $`f_i`$’s are distinct then $`g_1g_2`$. If $`Y`$ is of type $`B`$ then $`g_1`$ has a certain component of $`M(PD)`$ on both sides, and $`g_2`$ has the other one, so $`g_1g_2`$. Assume in case $`A`$ that $`g_1=g_2`$. Referring to Fig. 21 let $`q_j`$ be the midpoint of $`t_j`$, and join $`q_1`$ to $`q_2`$ by an arc $`\lambda `$ in $`g_1=g_2`$. There are 4 distinct arcs $`\lambda _1,\mathrm{},\lambda _4Y`$ having endpoints $`q_1`$ and $`q_2`$ and intersecting $`S(P)`$ twice. For two of them the polyhedron $`(\lambda \lambda _i)`$ is an annulus with $`2`$ tongues on the same side. Then some $`\lambda \lambda _i`$ is fake, which is in contrast with the fact that the $`f_i`$’s are distinct.
Proof of fact 3. We start with case $`A`$. Assume that $`v_1v_2`$, and put $`P^{}=(PD)f_1`$. If $`f_1`$ is incident to $`x`$ different vertices of $`P`$ then $`\mathrm{\#}V(P^{})=\mathrm{\#}V(P)+42x`$. Since $`P`$ is minimal we have $`x2`$. On the other hand $`f_1`$ is incident to $`v_1`$ and $`v_2`$, so $`x=2`$. Now Fig. 25
shows a triod $`\tau `$ in $`P^{}`$, trace of a torus parallel to $`F`$. By Proposition 4.3, either $`F`$ is non-separating or $`\tau `$ is boundary-parallel. In the first case we get a contradiction to the initial assumption. In the second case we deduce that $`f_1`$ is incident to $`v_3`$ and $`v_4`$, so $`\{v_3,v_4\}\{v_1,v_2\}`$. So either $`v_3=v_4`$, or $`v_4=v_1`$, or $`v_3=v_1v_4=v_2`$. In all cases but the last one the conclusion is the desired one up to symmetry. Concentrating on the last case, we note that $`f_1\mathrm{}f_4`$ is a surface near $`v_1`$ and $`v_2`$, and that the $`f_i`$’s and $`g_j`$’s are all distinct. From these facts it is not hard to deduce that $`v_1`$ and $`v_2`$ appear as in Fig. 26.
The figure readily implies that $`f_2=f_4`$: a contradiction.
The proof in case $`B`$ is similar, except that $`D`$ cannot be used directly: a perturbed version $`D^{}`$ as in Fig. 27-left
must be employed. We are again supposing here that $`v_1v_2`$, so $`f_1`$ is incident to $`x2`$ vertices of $`P`$, but now $`\mathrm{\#}V(P^{})=\mathrm{\#}V(P)+63x`$. Since $`P`$ is minimal we have $`x3`$, so $`x\{2,3\}`$. We first claim that we can suppose $`x=3`$ up to symmetry. By contradiction, assume that both $`f_1`$ and $`f_2`$ are incident to exactly 2 vertices. We deduce that the situation is as in Fig. 27-right, where we also show a face $`f`$ incident twice to an edge, which is absurd by Corollary 4.4.
Our claim that $`x=3`$ up to symmetry is proved, so $`\mathrm{\#}V(P^{})=\mathrm{\#}V(P)`$ and $`P^{}`$ is minimal too. A figure very similar to Fig. 25 shows that a triod must exist in $`P^{}`$, and allows to conclude as above that either $`F`$ is separating or $`f_1`$ is incident to $`v_3`$ and $`v_4`$. So either $`v_3=v_4`$, which gives the desired conclusion up to symmetry, or $`\{v_3,v_4\}\{v_1,v_2\}\mathrm{}`$ (recall that $`f_1`$ is incident to exactly 3 vertices). If $`v_3=v_1`$ or $`v_4=v_2`$ we get the desired conclusion. Otherwise we can assume up to symmetry that $`v_1=v_4`$. So $`e_1`$ and $`e_4`$ have a common vertex in $`P`$, which implies that there is a face incident to both. But $`e_1`$ is adjacent to $`f_1,f_2,g_1`$ and $`e_4`$ is adjacent to $`f_3,f_4,g_2`$, and the $`f_i`$’s and $`g_j`$’s are distinct, so we get a contradiction.
Proof of fact 4. If $`Y`$ is of type $`A`$, then $`e(D^{})=3`$, so by Fact 1 (and its proof) $`Y^{}`$ bounds a polyhedron $`Q`$ of type 3.3 or of type 1, but the latter is impossible because $`Y^{}`$ is the trace of a torus. We only need to consider which $`J_1`$-moves can be applied to a $`Q`$ of type 3.3. By Lemma 4.12 the move actually takes place towards the exterior of $`Q`$ (*i.e.* its result contains 2 vertices of $`P`$). The move is determined by the edge of $`Q`$ which disappears during the move: of the 6 edges in $`Q`$, 3 lead to a situation in which $`e(D)=2`$, so we exclude them. The other 3 edges are actually symmetric, and the result is one of the polyhedra of type 3.4.
If $`Y`$ is of type $`B`$, then $`u`$ must be an edge in $`_2D`$ (otherwise $`e(D^{})=e(D)`$), so $`Y^{}`$ is of type $`A`$. Moreover $`Y^{}`$ is the trace of a torus. Combining Fact 2 and the part of Fact 4 already established we see that $`Y^{}=Q`$ with $`Q`$ either of type 3.4 or a Möbius strip with two tongues (type 2). However, if we denote by $`f_i^{}`$ the faces of $`\mathrm{\Sigma }_D^{}`$ incident to $`D^{}`$, by Lemma 4.12 we have $`f_i^{}f_i`$ up to permutation, so the $`f_i^{}`$’s are distinct. This shows that type 2 is impossible, and again we are left to analyze what can we get from a $`Q`$ of type 3.4 by a move $`J_1`$ which takes place towards the exterior. Of the 6 edges of $`Q`$, 4 lead to a situation in which $`e(D)=3`$, so we exclude them. The other 2 edges are actually symmetric, and the result is type 3.5.
Proof of fact 5. The first step of our proof is the extension of the move $`YY^{}`$ to a flow $`YY^{}Y^{\prime \prime }\mathrm{}Y^{(k)}`$ of $`J_1`$-moves. As mentioned in the proof of Fact 4 we must have $`u_1D`$ in this case, so we assume up to symmetry that $`u=s_1f_1`$, and we note that Remarks 4.10-4.11 and Lemma 4.12 apply. The situation is described in Fig. 28.
One easily sees that the faces of $`\mathrm{\Sigma }_D^{}`$ incident to $`_1D^{}`$ are $`f_3,f_4`$ and two new ones (one of which is contained in $`f_2`$), which we denote by $`f_1^{},f_2^{}`$. If $`\{f_1^{},f_2^{},f_3,f_4\}`$ are not distinct, the flow is reduced to $`YY^{}`$, and we move to the next step. Otherwise let $`v_1^{},v_2^{}`$ be the ends of $`e_1^{},e_2^{}`$ (see Fig. 28-left). If $`v_1^{}v_2^{}`$ then again the flow is reduced to $`YY^{}`$. Assume on the contrary that $`v_1^{}=v_2^{}`$, and consider Fig. 28-right. Then either $`s_1^{}`$ or $`s_2^{}`$ is contained in $`\mathrm{lk}(v^{})`$, but certainly $`s_1^{}`$ is not, for otherwise $`P`$ would contain an embedded face with two vertices, which is absurd by Corollary 4.4. Setting $`u^{}=s_2^{}`$, we are now in a position to apply a move $`J_1`$ along the triangle determined by $`v_1^{}`$ and $`u^{}`$, getting from $`Y^{}`$ to $`Y^{\prime \prime }`$. We proceed in a similar way and note that the process must come to an end because $`\mathrm{\Sigma }_{D^{(i)}}`$ contains one vertex less than $`\mathrm{\Sigma }_{D^{(i1)}}`$ by Lemma 4.12.
Our second step is to understand the final stage $`Y^{(k)}`$ of our flow. By construction either $`\{f_1^{(k)},f_2^{(k)},f_3,f_4\}`$ are not distinct or $`v_1^{(k)}v_2^{(k)}`$. In the first case, since at each step only 1 face not contained in the previous one is inserted (and 1 is deleted), precisely 3 of $`\{f_1^{(k)},f_2^{(k)},f_3,f_4\}`$ are distinct. We know by Fact 2 (and its proof) that $`Y^{(k)}`$ (which is of type $`B`$) bounds a polyhedron $`Q`$ which is either an annulus with 2 tongues on opposite sides, or of type 4.1. The first case is excluded by what just said about the $`f_i`$’s. By Lemma 4.12, $`Y`$ bounds $`Q_{Y^{(k)}}[Y,Y^{(k)}]`$. Since at each step of the construction of our flow the choice of move $`J_1`$ was forced, the polyhedron $`[Y,Y^{(k)}]`$ is defined unambiguously (it depends on $`k`$ only). We only need to explain which edge of $`Q`$ determines the $`J_1`$-move which glues $`Q`$ to $`[Y,Y^{(k)}]`$. Of the 6 edges, 2 lead to a trace of type $`A`$, 2 give rise to an embedded face with 2 edges (excluded by Corollary 4.4) and the other 2 are symmetric, so $`Q_{Y^{(k)}}[Y,Y^{(k)}]`$ also depends on $`k`$ only. It is now a routine matter to check that indeed $`Q_{Y^{(k)}}[Y,Y^{(k)}]`$ is the polyhedron of type 4.2 with $`k`$ vertices.
Having understood the case where $`\{f_1^{(k)},f_2^{(k)},f_3,f_4\}`$ are not distinct, we assume that they are. The rest of the proof is devoted to showing that it is actually impossible that $`v_1^{(k)}v_2^{(k)}`$. Let us first assume that $`v_3v_4`$. By Fact 3 we then have $`v_1=v_3`$ up to symmetry, and we can apply a move $`J_1`$ which reduces $`e(D)`$. Fact 4 shows that $`Y^{(k)}`$ bounds a polyhedron $`Q`$ of type 3.4 or 3.5, but $`Q`$ is of type $`B`$, so it must be of type 3.5. Once again we must analyze the possible results of a move $`J_1`$, towards the exterior of a $`Q`$ of type 3.5. Of the 6 edges of $`Q`$, 2 lead to a trace of type $`A`$, and therefore are excluded. The 4 other edges come in 2 symmetric pairs. For one type, the result of the move $`J_1`$ contains an embedded face with 3 vertices, which is absurd by Corollary 4.4.
For the other type, the result contains an embedded face with 4 vertices. We can then apply a disc-replacement move as in Fig. 10, getting a new minimal skeleton $`P^{}`$ of $`(M,X)`$. The evolution of the singular set is shown in Fig. 29, where the two white dots lie on some $`\theta _i`$, the black dots are vertices, and the gray dots lie on $`Y`$. Since the edges leaving $`\theta _i`$ end at the same vertex, a $`J_1`$-move transforms $`\theta _i`$ into a triod which is not boundary-parallel. This contradicts Proposition 4.3.
We are left to deal with the case where $`\{f_1^{(k)},f_2^{(k)},f_3,f_4\}`$ are distinct, $`v_1^{(k)}v_2^{(k)}`$, and $`v_3=v_4`$. In this case we can perform a $`J_1`$-move along either $`s_3`$ or $`s_4`$, and we can proceed just as above, constructing a flow $`Y^{(k)}Y^{(k+1)}\mathrm{}Y^{(k+h)}`$. During this process the faces $`f_1^{(k)},f_2^{(k)}`$, and the vertices $`v_1^{(k)},v_2^{(k)}`$ remain unaffected, while $`f_3,f_4,v_3=v_4`$ get transformed into $`f_3^{(h)},f_4^{(h)},v_3^{(h)},v_4^{(h)}`$. As above, we have at the end of the sequence either that $`\{f_1^{(k)},f_2^{(k)},f_3^{(h)},f_4^{(h)}\}`$ are not distinct or that $`v_3^{(h)}v_4^{(h)}`$. In the first case, Fact 2 implies that $`Y^{(k+h)}`$ bounds a polyhedron of type 2 or 4.1. Such a polyhedron has at most 1 vertex, but $`\mathrm{\Sigma }_{D^{(k+h)}}`$ contains at least $`v_1^{(k)}v_2^{(k)}`$, and we get a contradiction. In the second case we are precisely in the situation $`v_3v_4`$ previously considered, and again we get a contradiction. ∎
### 4.5 Conclusion of proofs
If $`Y`$ is a trace in $`P`$, we denote by $`P_Y`$ the polyhedron $`P(Y)`$.
*Proof of Theorem 3.3.* Let $`P`$ be a minimal skeleton of $`(M,X)`$. By Corollary 2.8 we have $`c(M,X)>0`$, so $`P`$ is standard. Suppose a face $`f`$ of $`P`$ is incident to $`P`$ in at least two distinct edges $`e\theta _i`$ and $`e^{}\theta _i^{}`$. We note that $`ii^{}`$ by Lemma 2.14, and choose an arc $`\alpha `$ in $`f`$ having one end on $`e`$ and one on $`e^{}`$. Then $`Y=(\theta _i\theta _i^{}\alpha )`$ is a trace with 4 vertices of a surface $`F`$. Moreover $`P_Y=P_1P_2`$ is disconnected, so $`F`$ separates $`M`$ and hence it is orientable. Let $`P_2`$ be the component containing $`\alpha `$.
The graph $`Y`$ is of type B (see Fig. 20) and $`P_2`$ has 3 boundary components (namely, $`\theta _i`$, $`\theta _i^{}`$, and $`Y`$). Now either $`P_1`$ or $`P_2`$ is of one of the types listed by Theorem 4.14, but no such type has 3 boundary components, so $`P_1`$ must be of one such type. The only polyhedra among those listed in Theorem 4.14 having at least one vertex and boundary of type B are those of type 3.5 (Fig. 18) and 4 (Fig. 19). If $`P_1`$ is of type 3.5 then $`P`$ is the skeleton of $`B_4`$, and if $`P_1`$ is of type 4 with 1 vertex then it is the skeleton of $`B_3`$. Otherwise $`P_1`$ is of type 4 with $`k2`$ vertices, and the two edges of $`S(P)`$ adjacent to $`\theta _i`$ have a common endpoint. It easily follows that via a $`J_1`$-move we can transform $`\theta _i`$ into a triod which is not boundary-parallel and is the trace of a separating torus. This contradicts Proposition 4.3.
$`\mathrm{}`$
*Proof of Theorem 3.5.* Set $`=\{\gamma _1,\mathrm{},\gamma _n\}`$, where $`\gamma _i`$ is the core of the Möbius strip with one tongue attached to $`\theta _iP`$. By Theorem 3.3, even if we modify each $`\gamma _i`$ within its isotopy class, the $`\gamma _i`$’s stay disjoint. Moreover, each $`(\gamma _i)`$ is a Möbius strip with one tongue. Therefore it is enough to show that $``$ is a set of representatives of length-1 loops in $`Q`$. If not, there is a length-1 loop $`\gamma `$ not isotopic to any $`\gamma _i`$.
If $`\gamma `$ is disjoint from all $`\gamma _i`$’s, then $`\gamma P`$, so a face of $`P`$ is doubly incident to some edge, and we get a contradiction to Corollary 4.4. If $`\gamma `$ meets some $`\gamma _i`$ then, by Theorem 3.3, it meets only one, and we can assume that $`\gamma \gamma _i`$ is one point away from $`S(Q)`$. Set $`R=_Q(\gamma \gamma _i)`$. We need now to distinguish two cases, depending on whether $`_Q(\gamma )`$ is a Möbius strip or an annulus with one tongue. In the first case there exists a curve $`\alpha `$ contained in $`R`$, and therefore in $`P`$, such that $`l(\alpha )=2`$ and $`\alpha `$ bounds an external disc ($`\alpha `$ is homologous to $`\gamma +\gamma _i`$ in $`R`$). By Theorem 3.8 $`\alpha `$ is fake, and it easily follows that $`\gamma `$ is isotopic to $`\gamma _i`$.
Assume now that $`_Q(\gamma )`$ is an annulus with one tongue. Note that $`RP`$ is a trace with 4 vertices of a separating, and hence orientable, surface $`F`$. Moreover $`R`$ is of type $`A`$, so, by Theorem 4.14, $`R`$ bounds in $`P`$ a polyhedron $`S`$ of type 1.1, 3.3, 3.4, or 2 based on a Möbius strip. But $`RP`$ is not of such a type, so the rest of $`P`$ is, hence $`\mathrm{\#}V(P)1`$. But $`_1^1=\{B_0,\mathrm{},B_3\}`$, and we are done.
$`\mathrm{}`$
Before proving Theorem 3.7 we establish a general fact.
###### Lemma 4.15.
Let $`Q`$ be a filling of a minimal skeleton $`P`$ of a brick. Let $`\{e_1,\mathrm{},e_{2m}\}`$ be a set of edges which disconnects $`S(Q)`$ in two components. Then there is a trace $`Y`$ contained in $`P`$ which has $`2m`$ vertices $`p_ie_i`$ for $`i=1,\mathrm{},2m`$, and $`Y`$ is the trace of an orientable separating surface.
###### Proof.
Take points $`p_ie_i`$; we have $`S(Q)\{p_i\}=K_1K_2`$. Let $`f`$ be a face of $`Q`$ incident to some $`e_i`$. The gluing path of $`f`$ to $`S(Q)`$ can be split into arcs $`s_1,\mathrm{}s_{2\nu }`$, meeting at points $`q_1,\mathrm{},q_{2\nu }`$, where $`s_{2j+1}K_1`$ and $`s_{2j}K_2`$ for all $`j`$, and each $`q_k`$ is glued to one $`p_{\beta (k)}`$. The map $`\beta `$ is not necessarily injective, since $`f`$ can be multiply incident to an edge $`e_i`$. We can give the points $`q_k`$ alternating (red and black) colors.
Since $`P=Q((Q))`$ is super-standard, $`f`$ can intersect at most one loop $`\gamma `$ among those in $`(Q)`$. Now take $`\nu `$ pairwise disjoint segments $`\lambda _1,\mathrm{},\lambda _\nu `$, properly embedded in $`f`$, such that $`_{j=1}^\nu \lambda _j=_{k=1}^{2\nu }q_k`$. We can ask the $`\lambda _j`$’s to be disjoint from $`\gamma `$, since the points on $`f`$ are separated into two even subsets by $`\gamma `$. It is easy to see that the two endpoints of each $`\lambda _j`$ automatically have distinct colours. If we do this for each face $`f`$ incident to some $`e_i`$, the union of all the chosen segments is a trace $`Y`$ disjoint from $`(Q)`$ and hence contained in $`P`$.
We claim that $`Y`$ has a product regular neighbourhood in $`P`$: take for $`i=1,\mathrm{},2m`$ a vector $`v_i`$ at $`p_i`$, tangent to $`e_i`$ and directed towards $`K_2`$. Each segment of $`Y`$ is a $`\lambda _j`$, properly embedded in a face $`f`$ such that $`\lambda _j`$ consists of points with distinct colors. It follows that the vectors at the ends of $`\lambda _j`$ extend along $`\lambda _j`$ to a non-vanishing field tangent to $`f`$. The existence of such a field on $`Y`$ easily implies that $`F`$ is orientable and that $`F`$ cuts $`M`$ into two components.∎
*Proof of Theorem 3.7.* Suppose $`S(Q)`$ contains a pair $`\{e_0,e_1\}`$ of separating edges. By Lemma 4.15 there is a trace $`Y`$ of a separating (and hence orientable) surface $`F`$ with two vertices, intersecting both $`e_0`$ and $`e_1`$. Proposition 4.3 applies, and possibility (1) is ruled out because $`F`$ separates. Both other possibilities imply that the vertices of $`Y`$ lie on the same edge of $`Q`$, but $`e_0e_1`$ by assumption.
Suppose $`S(Q)`$ contains a separating quadruple $`\{e_0,e_1,e_2,e_3\}`$ of edges. By Lemma 4.15 there is a trace $`Y`$ of a separating (and hence orientable) surface with 4 vertices intersecting them. If $`Y`$ is connected then Theorem 4.14 applies, and we are done because the singular sets of polyhedra of types 1-4 indeed are as shown in Fig. 5. If $`Y=Y_0Y_1`$, then $`Y_0`$ is a trace with two vertices to which Proposition 4.3 applies. Now possibility (1) is ruled out either because every torus in $`M`$ is separating or by Theorem 4.6, and as above the other two possibilities lead to a contradiction.
$`\mathrm{}`$
## 5 Bricks and skeleta up to complexity 9
We provide in this section a complete description of the bricks in $`_n`$ for $`n9`$ anticipated in Subsection 1.2. Recall that $`_n`$ was split as $`_n^0_n^1`$, where $`_n^0`$ consists of the elements of $`_n`$ without boundary. We describe now $`_9^1`$, postponing $`_9^0`$ for a moment, because to discuss it we will first need to introduce a new move on skeleta.
Our computations show that the set $`_9^1`$ consists of $`11`$ bricks $`B_0,\mathrm{},B_{10}`$. Moreover, for $`i9`$ there is a unique minimal skeleton of $`B_i`$, while for $`i=10`$ there are two. Minimal skeleta for $`B_0,\mathrm{},B_4`$ were shown in Figg. 2 and 4, and for $`B_5,\mathrm{},B_{10}`$ they are now shown in Fig. 30.
Using Remark 4.13, in this figure we only draw $`_P(S(P))`$, and we use a thicker line for the $`Y`$-shaped portions of $`_P(S(P))`$ lying on $`P`$. Each component of $`B_i`$ contains two such $`Y`$’s (shown close to each other when $`B_i`$ has more than one component).
Having described $`B_0,\mathrm{},B_{10}`$, we can now prove Proposition 1.5.
*Proof of Proposition 1.5.* Suppose $`(M,\mathrm{})`$ is a sharp assembling of $`B_i`$ with $`i6`$ and some $`B_2`$’s and $`B_3`$’s. Since $`c(M)9`$, one $`B_3`$ can occur if $`i=6`$ only. Let $`P_i`$ be the minimal skeleton of $`B_i`$ shown in Fig. 30. A minimal skeleton $`P`$ for $`M`$ is then a filling of $`P_i`$, possibly after glueing one copy of the minimal skeleton of $`B_3`$ if $`i=6`$. If we check all the polyhedra which can be built in this way, we see that many of them contain embedded faces with no more than 3 vertices, which contradicts Theorem 3.8. Only 16 of them do not contain such a face. Now 9 of these 16 are shown to be non-minimal by checking that small faces appear after suitable disc-replacement moves. The 7 polyhedra left out are skeleta of the 4 mentioned hyperbolic manifolds (there are some duplicates).
$`\mathrm{}`$
### 5.1 Twists
We introduce here a notion needed below to describe $`_9^0`$. Let $`P`$ be a quasi-standard skeleton of a closed manifold $`(M,\mathrm{})`$, and let $`\gamma `$ be a length-$`2`$ loop in $`P`$ such that $`(\gamma )`$ is an annulus with $`2`$ tongues. For $`k1`$ let $`W_k`$ be the polyhedron of type 4 with $`k`$ vertices (Fig. 19). The boundaries $`(\gamma )`$ and $`W_k`$ are homeomorphic (of type $`B`$). We can then choose a homeomorphism $`\psi :W_k(\gamma )`$ and form a polyhedron $`P_k=P(\gamma )_\psi W_k`$. Note now that $`W_k`$ naturally sits in a solid torus $`H`$, with $`W_k=W_kH`$.
###### Proposition 5.1.
The homeomorphism $`\psi :W_k(\gamma )`$ can be chosen so that it extends to a homeomorphism $`\mathrm{\Psi }:H_M(\gamma )`$. For these choices $`P_k`$ is a skeleton of the Dehn surgered manifold $`M_k=M_M(\gamma )_\mathrm{\Psi }H`$.
###### Proof.
The first assertion is easy and taken for granted. By construction $`P_k`$ sits in $`M_k`$ and it is simple, so we only need to show that $`M_kP_k`$ is an open $`3`$-ball. To this end we note that $`MP_M(\gamma )`$ is a ball $`B`$. Moreover $`HW`$ consists of two discs $`D^{}`$ and $`D^{\prime \prime }`$, and $`H(HW)`$ consists of two balls $`B^{}`$ and $`B^{\prime \prime }`$, with $`B^{}H=D^{}`$ and $`B^{\prime \prime }H=D^{\prime \prime }`$. So $`M_kP_k=B_{\mathrm{\Psi }|_D^{}}B^{}_{\mathrm{\Psi }|_{D^{\prime \prime }}}B^{\prime \prime }`$ is a ball. ∎
We say that $`P_k`$ is obtained from $`P`$ by a $`k`$-*twist* along $`\gamma `$, and we adopt the convention that making a $`0`$-twist means leaving $`P`$ unaffected.
### 5.2 Closed bricks up to complexity 9
Our computations show that the set $`_9^0`$ consists of $`19`$ bricks which belong to the union of two classes $`\{C_{i,j}\}`$ and $`\{E_k\}`$. We describe here these manifolds and minimal skeleta $`\stackrel{~}{C}_{i,j}`$ and $`\stackrel{~}{E}_k`$ of them. (As opposed to the case of $`_{n9}^1`$, minimal skeleta are often not unique in $`_{n9}^0`$.) The polyhedron $`\stackrel{~}{C}_{0,0}`$ of Fig. 31-left is a skeleton of $`(S^2\times S^1,\mathrm{})`$ and it contains $`2`$ length-2 loops $`\gamma `$ and $`\delta `$, shown in Fig. 31-left, such that $`S^2\times S^1_{S^2\times S^1}(\gamma \delta )(A,(2,1))`$, where $`A`$ is the annulus. Both $`_{\stackrel{~}{C}_{0,0}}(\gamma )`$ and $`_{\stackrel{~}{C}_{0,0}}(\delta )`$ are annuli with two tongues on different sides. We can therefore perform an $`i`$-twist along $`\gamma `$ and a $`j`$-twist along $`\delta `$. If we do this with appropriate gluing maps we get the skeleton shown in Fig. 31-right, which we denote by $`\stackrel{~}{C}_{i,j}`$.
Using Proposition 5.1 it is not hard to check that $`\stackrel{~}{C}_{i,j}`$ is a skeleton of the Seifert manifold $`C_{i,j}=(S^2,(2,1),(1+i,1),(1+j,1),(1,1))`$. We have $`C_{i,j}=C_{j,i}`$ for all $`i,j`$.
Poincaré’s homology sphere $`(S^2,(2,1),(3,1),(5,1),(1,1))`$ has a unique minimal skeleton $`\stackrel{~}{E}_0`$ (Fig. 32-left). For any pair of non-adjacent edges of $`S(\stackrel{~}{E}_0)`$ there is a length-2 loop $`\gamma `$ intersecting them, isotopic to the singular fiber $`(5,1)`$. Since $`(\gamma )`$ is an annulus with two tongues, we can perform a $`k`$-twist along $`\gamma `$. If we do this with an appropriate gluing map we get the skeleton shown in Fig. 32-right, which we denote by $`\stackrel{~}{E}_k`$. Each $`\stackrel{~}{E}_k`$ turns out to be a skeleton of the manifold $`E_k=(S^2,(2,1),(3,1),(5+k,1),(1,1))`$. It is worth mentioning here that the minimal skeleton of the brick $`B_5`$ may be obtained from $`\stackrel{~}{E}_0`$ by an operation similar to a $`k`$-twist along $`\gamma `$, except that the polyhedron of type 3.5 (Fig. 18) is employed instead of $`W_k`$.
The set $`_9^0`$ consists of all manifolds $`C_{i,j}`$ and $`E_k`$ with $`k0`$ and $`ij1`$ having at most $`9`$ vertices (*i.e.* with $`k4`$ and $`i+j9`$), except the cases $`k=1`$ and $`(i4,j=2)`$. The skeleton $`\stackrel{~}{E}_1`$ is indeed minimal, but the associated manifold is not a brick, since it lies in $`B_0_{\mathrm{self}}`$. This is coherent with the well-known fact that $`(S^2,(2,1),(3,1),(6,1))`$ fibers over $`S^1`$ with torus fiber. Each $`\stackrel{~}{C}_{i,0}`$ is minimal (for $`i9`$), but the corresponding manifold is contained in $`B_2,B_3_{\mathrm{non}\mathrm{self}}`$. Each $`\stackrel{~}{C}_{i,2}`$ for $`i4`$ is not minimal, since $`\stackrel{~}{E}_{i4}`$ is a skeleton of the same manifold $`(S^2,(2,1),(3,1),(i,1),(1,1))`$ with one vertex less.
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# Superactivation of Bound Entanglement
\[
## Abstract
We show that, in a multi-party setting, two non-distillable (bound-entangled) states tensored together can make a distillable state. This is an example of true superadditivity of distillable entanglement. We also show that unlockable bound-entangled states cannot be asymptotically unentangled, providing the first proof that some states are truly bound-entangled in the sense of being both non-distillable and non-separable asymptotically.
\]
The joint state of more than one quantum system cannot always be thought of as a separate state of each system, nor even as a correlated mixture of separate states of each system , a situation known as quantum entanglement. Entanglement leads to the most counterintuitive effects in quantum mechanics, including the disturbing idea due to Bell that quantum mechanics is incompatible with local hidden variable theories . Even today new quantum oddities with their basis in entanglement are being found, and the study of entanglement is at the heart of quantum information theory.
A state belonging to parties $`A`$, $`B`$, $`C`$, etc. is said to be inseparable if it cannot be written in separable form
$$\rho ^{ABC\mathrm{}}=\underset{i}{}p_i\rho _i^A\rho _i^B\rho _i^C\mathrm{}$$
(1)
for any positive probabilities $`p_i`$ summing to one and set of density matrices $`\rho _i^A,\rho _i^B,\rho _i^C\mathrm{}`$, where, for example, $`\rho _i^A`$ operates on the Hilbert space belonging to party $`A`$. Notice that the superscripts $`A`$, $`B`$, $`C`$, etc. denote the parties by whom the state is shared. We say that a state is distillable if some pure entangled state shared by some subset of the parties is obtainable (asymptotically ) from it by local operations and classical communication (LOCC) amongst the parties.
It is known that many inseparable quantum mixed states are distillable, while separable states are not . More recently it has been shown that some mixed states which are entangled in the sense of being inseparable nevertheless cannot be distilled into any pure entanglement . Such states are known as bound-entangled states.
It has been an open question whether bound-entangled states, though inseparable, are actually entangled at all in an asymptotic sense. A state $`\rho `$ is said to be asymptotically unentangled ) if for any positive $`ϵ`$ there exists a number of copies $`N`$, a number $`m`$ sublinear in $`N`$ of EPR pairs shared in some way among the parties, and an LOCC method of constructing from those EPR pairs a state $`\rho ^{}`$ such that $`F(\rho ^N,\rho ^{})>1ϵ`$ for some sensible definition of the fidelity $`F`$ between two density matrices . In this letter we show the first example of a bound entangled state that can be proved not to be asymptotically unentangled. Other examples can be found it .
In the bipartite case, bound entanglement may sometimes be useful in a kind of quasi-distillation process known as activating the bound entanglement in which a finite number of free-entangled mixed states are distilled with the help of a large number of bound-entangled states. This is not a true distillation of the bound entanglement in that no more pure entanglement is produced than the distillable entanglement of the free-entangled mixed states, the distillable entanglement being defined as the pure entanglement distillable per state from an infinite number of copies of a state.
In the case of more than two parties the bound entanglement can be more truly activated by the presence of free entanglement. In examples given by Cirac, Tarrach and Dür , and in the equivalent formulation of unlockable bound-entangled states , when several parties share certain bound entangled states, and when some subset of the parties get to share pure entanglement then some pure entanglement may be distilled between parties where it would be impossible to obtain any without having shared the bound-entangled state. This is a kind of superadditivity of distillable entanglement, though in the known cases no more entanglement is distilled than the pure entanglement that was shared, rather it is in a different place. Later in this letter we will look at unlockable states in much more detail since we will need some of the results about them.
In this letter we present an effect we call superactivation of bound entanglement. It is “super” in the sense of superadditivity of distillable entanglement, but without the restrictions of either of the earlier types of activation of bound entanglement. In superactivation two entangled mixed states $`\rho ,\rho ^{}`$ are combined to yield more pure entanglement than the sum of what a set of parties could distill from either $`\rho `$ or $`\rho ^{}`$ on their own, even if many copies of $`\rho `$ or $`\rho ^{}`$ are shared. In particular, both states in our example are bound-entangled states from which no pure entanglement can be distilled. Our result thus provides the first example of superaddivity of distillable entanglement.
We will use the usual notation for the maximally entangled states of two qubits (the Bell states):
$$|\mathrm{\Psi }^\pm =\frac{1}{\sqrt{2}}(|\pm |),|\mathrm{\Phi }^\pm =\frac{1}{\sqrt{2}}(|\pm |)$$
(2)
For convenience we adopt the following notation as well:
$`\mathrm{\Psi }=\{\mathrm{\Psi }^{},\mathrm{\Psi }^+,\mathrm{\Phi }^+,\mathrm{\Phi }^{}\}\mathrm{with}\mathrm{elements}\mathrm{\Psi }_i,\mathrm{and}`$ (3)
$`\sigma =\{\mathrm{𝟏}_\mathrm{𝟐},\left({\displaystyle \genfrac{}{}{0pt}{}{\mathbf{1\; 0}}{\mathrm{𝟎}\mathrm{𝟏}}}\right),\left({\displaystyle \genfrac{}{}{0pt}{}{\mathrm{𝟎}\mathrm{𝟏}}{\mathbf{1\; 0}}}\right),\left({\displaystyle \genfrac{}{}{0pt}{}{\mathbf{0\; 1}}{\mathbf{1\; 0}}}\right)\}\mathrm{with}\mathrm{elements}\sigma _𝐢,`$ (4)
where $`\mathrm{𝟏}_\mathrm{𝟐}`$ is the identity matrix in $`2\times 2`$. In the text, we shall refer to a Bell state as any one of the four states (3) and to an EPR state as the standard singlet state $`|\mathrm{\Psi }^{}`$. The Bell states $`|\mathrm{\Psi }_i`$ are related to the standard EPR state $`|\mathrm{\Psi }^{}`$ by the following identities, up to an overall phase which is unimportant here:
$`|\mathrm{\Psi }^{}=\mathrm{𝟏}_\mathrm{𝟐}\sigma _𝐢|𝚿_𝐢=\sigma _𝐢\mathrm{𝟏}_\mathrm{𝟐}|𝚿_𝐢`$ (5)
$`|\mathrm{\Psi }_i=\mathrm{𝟏}_\mathrm{𝟐}\sigma _𝐢|𝚿^{}=\sigma _𝐢\mathrm{𝟏}_\mathrm{𝟐}|𝚿^{}.`$ (6)
In teleportation , $`A`$ and $`B`$ share an EPR pair $`|\mathrm{\Psi }^{}`$, and $`A`$ has another qubit in a state $`|\psi `$. $`A`$ first does a joint measurement on her two qubits in the basis formed by the Bell states. There are four equally likely outcomes corresponding to the Bell states $`|\mathrm{\Psi }_i`$. $`B`$’s half of the EPR pair after this measurement is $`\sigma _i|\psi `$ up to a phase that can be ignored. Then $`A`$ communicates $`i`$ to $`B`$ who then performs a rotation $`\sigma _i`$ on his state giving $`\sigma _i^2|\psi `$. But $`\sigma _i^2=\pm \mathrm{𝟏}_\mathrm{𝟐}`$ and thus the final state $`B`$ has is $`|\psi `$ up to a phase.
An easy lemma about teleportation is that if a state $`|\psi `$ is teleported from $`A`$ to $`B`$ using an incorrect one of the Bell states $`|\mathrm{\Psi }_i`$ rather than $`|\mathrm{\Psi }^{}`$ as normally required by the protocol, then the result of the teleportation will be $`\sigma _i|\psi `$, again up to an overall phase. This is easily seen by using (6) to write the incorrect Bell state as $`|\mathrm{\Psi }^{}`$ with a $`\sigma _i`$ operating on $`B`$’s part of the $`|\mathrm{\Psi }^{}`$ i.e. $`|\psi =\mathrm{𝟏}_\mathrm{𝟐}\sigma _𝐢|𝚿^{}`$. If A’s outcome of the Bell measurement is $`j`$ then B’s corresponding state is $`\sigma _i\sigma _j|\psi `$. Thus after B applies the rotation $`\sigma _j`$ the state becomes $`\sigma _j\sigma _i\sigma _j|\psi `$. If the rotation $`\sigma _j`$ which is the final step in teleportation could be squeezed in before the $`\sigma _i`$ the proof would be complete, but instead it follows the $`\sigma _i`$. However, the rotations used in teleportation are also the $`\sigma `$ matrices, and all the $`\sigma _i,\sigma _j`$ either commute or anticommute ($`\sigma _i\sigma _j=\pm \sigma _j\sigma _i`$) and so their order can be freely interchanged up to a phase. Thus the lemma is proved. $`\mathrm{}`$
In a four-party bound entangled state was presented:
$$\rho ^{ABCD}=\frac{1}{4}\underset{i=0}{\overset{3}{}}|\mathrm{\Psi }_i^{AB}\mathrm{\Psi }_i||\mathrm{\Psi }_i^{CD}\mathrm{\Psi }_i|$$
(7)
In other words, $`A`$ and $`B`$ share one of the four Bell states, but don’t know which one, and $`C`$ and $`D`$ share the same Bell state, also not knowing which one.
This state has several properties:
$``$ Symmetry under interchange of parties: $`\rho ^{ABCD}=\rho ^{ABDC}=\rho ^{ADBC}`$, etc. This may be verified by writing out the state as a $`1616`$ matrix and interchanging indices. A more enlightening way is to use our lemma and think of the state in terms of teleportation. First, we note that some of the symmetries are obvious, for example interchanging $`A`$ and $`B`$ because Bell states are themselves symmetric under interchange. So the only symmetry we need to consider is the interchange of $`B`$ with $`C`$ and the rest can be constructed trivially.
Consider the state in its original form, with $`A`$ and $`B`$ sharing an unknown Bell state and $`C`$ and $`D`$ sharing the same one. Now consider $`A`$ and $`C`$ getting together and performing a Bell measurement and obtaining the result $`|\mathrm{\Psi }_j`$, which we can think of as $`A`$ and $`C`$ doing the first step required to teleport $`A`$’s particle to $`D`$ using the unknown Bell state shared by $`C`$ and $`D`$. The result $`|\mathrm{\Psi }_j`$ is random since $`A`$ and $`C`$ had halves of completely separate unknown Bell states. The state being teleported is half of a Bell state given by Eq. (6) $`\sigma _i\mathrm{𝟏}_\mathrm{𝟐}|𝚿^{}`$ as is the state used in the teleportation. So, by our lemma, if the teleportation were completed an extra $`\sigma _i`$ would be introduced, and the two $`\sigma _i`$’s would cancel being self-inverse (up to a phase). Thus, $`B`$ and $`D`$ would share a standard $`|\mathrm{\Psi }^{}`$. But if the $`\sigma _i`$ needed to complete teleportation is not performed, this means that $`B`$ and $`D`$ share the Bell state $`\sigma _j^1\mathrm{𝟏}_\mathrm{𝟐}|𝚿^{}=\sigma _𝐣\mathrm{𝟏}_\mathrm{𝟐}|𝚿^{}=|𝚿_𝐣`$ (ignoring phases), which is the result obtained by $`A`$ and $`C`$. So $`AC`$ and $`BD`$ share identical random Bell states, which was the original form of the density matrix, but with $`A`$ and $`C`$ interchanged.
$``$ Non-distillability: When all four parties remain separated and cannot perform joint quantum operations, then they cannot distill any pure entanglement by LOCC, even if they share many states, each having density matrix $`\rho ^{ABCD}`$. This comes from the fact every party is separated from every other across a separable cut. This is easy to see since the state (7) is separable across the $`AB:CD`$ cut by construction and the state has the symmetry property.
$``$ Unlockability: The entanglement of the state can be unlocked. If $`A`$ and $`B`$ come together and perform a joint quantum measurement, they can determine which of the four Bell states they have (the four Bell states form an orthogonal basis) and tell $`C`$ and $`D`$ the outcome. Since $`C`$ and $`D`$ then know which Bell state they share, they can convert it into the standard $`|\mathrm{\Psi }^{}`$ state using local operations by Eq. (5). Because of the symmetry property any two parties can join together to help the other two get a $`|\mathrm{\Psi }^{}`$. Note that the unlockability property implies the state must not be fully separable, or no entanglement could be distilled between separated parties, even when some of the parties come together. Because the state is both non-distillable and entangled, it is by definition a bound-entangled state .
Now we consider the mixed state of five parties $`A`$, $`B`$, $`C`$, $`D`$, and $`E`$
$$M=\rho ^{ACBD}\rho ^{ABCE}$$
(8)
where $`\rho ^{ACBD}`$ (call it state 1) and $`\rho ^{ABCE}`$ (call it state 2) are the states of Eq. (7) but with the qubits assigned to different parties. Thus parties $`A`$, $`B`$, $`C`$ and $`D`$ each have a one qubit subsystem of state $`\rho ^{ACBD}`$ and similarly parties $`A`$, $`B`$, $`C`$ and $`E`$ each have a one qubit subsystem of state $`\rho ^{ABCE}`$. Thus the parties $`A`$, $`B`$, $`C`$, $`D`$ and $`E`$ have Hilbert spaces of size $`4`$, $`4`$, $`4`$, $`2`$, and $`2`$. Technically $`\rho ^{ACBD}`$ could be written as $`\rho ^{ABCD}`$ due to the symmetry property but it will be useful to have it explicitly written in the form where it is an unknown Bell state shared between $`A`$ and $`C`$, and the same state shared by $`B`$ and $`D`$. The state $`M`$ is illustrated in Figure 1a. $`M`$ is the tensor product of two density matrices, neither of which is independently distillable. We now show how to distill a $`|\mathrm{\Psi }^{}`$ between $`D`$ and $`E`$.
In the distillation procedure $`A`$ and $`B`$ use state 1 to “teleport” state 2 to $`C`$ and $`D`$. First, party $`A`$ teleports her half of the unknown Bell state she shares with $`B`$ (which is part of state 2 and shown by the solid arrow connecting $`A`$ and $`B`$ in Figure 1a and part of state 2) to $`C`$ using the unknown Bell state she shares with $`C`$ (which is part of state 1, shown by the dashed arrow connecting $`A`$ and $`C`$ in the figure). This results in the situation of Figure 1b, where now $`C`$ shares an unknown Bell state with $`B`$, her half of which has additionally picked up the unknown rotation $`\sigma _i`$ from having been teleported with an incorrect Bell state $`|\mathrm{\Psi }_i`$. The Bell state connecting $`A`$ and $`C`$ is gone in the figure, since it has been expended performing the teleportation. Then $`B`$ teleports his half of that state to $`D`$ using the unknown Bell state (again $`|\mathrm{\Psi }_i`$ that they share, resulting in the situation of Figure 1c, where now $`C`$ and $`D`$ share the unknown Bell state originally shared by $`A`$ and $`B`$, both halves of which having been rotated by $`\sigma _i`$. It is important to note here that because of the structure of $`\rho ^{ACBD}`$ this is the same $`\sigma _i`$. Now, using Eq. (6) and the fact that $`\sigma _i^2`$ is the identity (once again except for a phase), we can see that the $`\sigma _i`$’s cancel and we are left with the state $`\rho ^{CDCE}`$. This is the same form as the four-party unlockable state (Eq. (7)) but with one party sharing two of the qubits, and it is therefore distillable into a pure EPR pair shared by $`D`$ and $`E`$.
$`M`$ cannot be distilled into EPR pairs between any of the other parties. This is because if we give the five parties the additional power of having $`D`$ and $`E`$ in the same room, then $`M`$ is just two copies of $`\rho ^{ABCD}`$ which are known not to be distillable (by definition if $`\rho `$ is not distillable, then neither is $`\rho ^N`$). To construct a state out of tensor products of bound-entangled states that is distillable into any kind of pure entanglement, it is sufficient to symmetrize $`M`$, i.e.
$`M_\mathrm{S}=`$ (9)
$`\rho ^{ABCD}\rho ^{ABCE}\rho ^{ABDE}\rho ^{ACDE}\rho ^{BCDE}.`$ (10)
Then the distillation protocol just described can be used to obtain an EPR pair between any two of the parties, and using more copies of $`M_\mathrm{S}`$ one can obtain EPR pairs between all pairs of parties. Once this is accomplished any arbitrary multi-party entangled state can be constructed by one party creating it in his lab and teleporting the pieces as needed to the others.
Because $`M`$ (Eq. (8)) is distillable, it cannot be that the original state $`\rho ^{ABCD}`$ is asymptotically unentangled. If it were, then many copies $`N`$ of $`\rho ^{ABCD}`$ and $`\rho ^{ABCE}`$ could be created arbitrarily precisely using a number of EPR pairs sublinear in $`N`$. These could be used to create $`N`$ copies of $`M`$ which could then be distilled into $`N`$ pure EPR pairs between $`D`$ and $`E`$. These $`DE`$ EPR pairs would, to arbitrarily high probability, pass any test that pure EPR pairs would pass. Thus, an amount of entanglement sublinear in $`N`$ would have been converted into $`N`$ EPR pairs by LOCC, which is impossible .
In fact, all unlockable bound-entangled states are asymptotically inseparable. This is because when some subset $`S`$ of the parties possessing such a state come together in the same lab the state becomes distillable. If the state were asymptotically unentangled then it could be made arbitrarily precisely with asymptotically no entanglement even when parties in $`S`$ are actually together in the same lab (it cannot hurt for them to be together as they can conveniently ignore this fact as they carry out whatever procedure results in the creation of the state). But then they can distill a finite amount of arbitrarily pure entanglement per state from the sublinear amount of entanglement they started with, which is impossible. It is worth noting then that the unlockable bound-entangled states are the first states shown to be true bound-entangled states in the sense of both being non-distillable and being non-separable asymptotically.
It is clearly a necessary condition for superactivation that at least one of the states involved must not be asymptotically unentangled. It is by no means a sufficient one, however, since the states $`\rho ^{ABCD}`$ and $`\rho ^{EFGH}`$ are each not asymptotically unentangled but $`\rho ^{ABCD}\rho ^{EFGH}`$ is not distillable as the two pieces are on disconnected sets of parties.
In the individual states $`\rho ^{ABCD}`$ and $`\rho ^{ABCE}`$, every party is separated from every other party by at least one separable cut. In order for the combined state $`M`$ to be distillable into a $`DE`$ EPR pair, and for $`M_\mathrm{S}`$ to be distillable into EPR pairs between any pair of parties, it is necessary that the parties who get EPR pairs no longer be separated by any separable cut, as is indeed the case by construction for these states. Using this observation, Dür has reported a whole family of superactivated states based on the unlockable bound-entangled states of . References also discuss separable cuts and their relation to distillibility in more detail.
In conclusion, we have shown that asymptotically entangled states exist from which no pure entanglement can be distilled. This has been suspected for some time but ours is the first such example for which it has been proved. Further we have shown the surprising fact that distillable entanglement is not additive by showing two undistillable asymptotically entangled states that when combined gives a distillable state.
Many future directions are suggested by this work. Here we have shown four party example of a state that is asymptotically entangled but not distillable. An interesting question is whether such a state can be found for two parties. Since the first writing of this letter, this question has been answered in the affirmative in . Another direction for further research is to find bipartite states that show the non-additivity of distillable entanglement. Such examples have been shown to exist if the NPT-boud entangled states are truly bound entangled .
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# Anomalous dimensions and phase transitions in superconductors
## I Introduction
The superconducting phase transition has received considerable attention in recent years. All this interest is due in part to the experimentally larger critical region in the high-$`T_c`$ materials . This larger critical region, however, does not correspond to the inverted 3D XY ($`IXY_3`$ for short) universality class . Instead, the observed critical behavior belongs to the ordinary 3D XY ($`XY_3`$ for short) universality class, meaning that the phase transition is governed by the neutral non-trivial Wilson-Fisher fixed point. Concerning the charged transition (that is, the $`IXY_3`$ behavior), there is some progress from the theoretical side. Unfortunatly, the corresponding critical region remains experimentally out of reach.
Concerning the $`IXY_3`$ regime, interesting precise numerical results on the anomalous scaling dimensions have been obtained recently by Sudbø and collaborators using a lattice version of the Ginzburg-Landau (GL) model. Their results give a strong support to the duality scenario which underlies the $`IXY_3`$ behavior. The aim ot this paper is to provide an analysis of the anomalous scaling dimensions from the point of view of field theory. An important issue to be understood is the sign of the order parameter field anomalous dimension, $`\eta `$. As argued in Ref. a negative $`\eta `$, though fulfilling the inequality $`\eta >2d`$, would spoil some important properties that must be verified in any legitimate continuum (scaling) limit. A fundamental property, the positivity of the spectral weight of the Källen-Lehmann (KL) spectral representation of the 2-point correlation function, is violated if $`\eta <0`$. Kiometzis and Schakel pointed out also that unitarity is violated if $`\eta <0`$. In fact, violation of the unitarity is an immediate consequence of the violation of the positivity of the spectral weight. We should note, however, that most renormalization group (RG) calculations give in general $`\eta `$ in the range $`1<\eta <0`$ in $`d=3`$ ($`ϵ=1`$ in the context of the $`ϵ`$-expansion). The only exception is the case where a mass (Proca) term is added explicitly for the gauge field, where the inequality $`\eta 0`$ is satisfied . This last situation corresponds just to the case of the continuum dual model where the gauge symmetry is global . Since in the RG calculations $`\eta `$ is only slightly negative we may wonder if such a negativeness is not just an artifact of the approximations used. The situation is, however, much more subtle. The recent numerical simulations in the lattice of Nguyen and Sudbø gives $`\eta =0.18`$ . The results in Ref. are non-perturbative, in contrast to most RG calculations. In the RG context we can cite the work of Bergerhoff et al. and the $`1/N`$ expansion , both non-perturbative and giving also $`\eta <0`$.
From a thermodynamical point of view, the anomalous dimension $`\eta _A`$ of the vector potential has more far reaching consequences. Indeed, it plays an important role in a critical regime where the magnetic fluctuations are not negligible, such as in the $`IXY_3`$ regime. Gauge invariance allows an exact determination of $`\eta _A`$ in $`2<d<4`$ dimensions. Indeed, its value is given simply by $`\eta _A=4d`$. One important consequence of this result is the scaling $`\lambda \xi `$ where $`\lambda `$ is the penetration depth and $`\xi `$ is the correlation length.
In this paper we will discuss some interesting new aspects of the superconducting transition. We will focus the issue of the anomalous dimensions of fields, for both the scalar field and the gauge field. Our analysis should be applicable to superconductors in the type II regime, where we expect a second-order (charged) phase transition . In section II the negativeness of $`\eta `$ will be shown to be a consequence of the existence of two singularities in the scalar 2-point bare correlation function at the critical point (CP). One singularity happens at $`p=0`$ while the other one happens at a nonzero momentum $`p=p^{}`$. This second singularity is related to the existence of a first-order phase transition regime. This singularity at nonzero momentum is at the origin of the negative sign of the $`\eta `$ exponent. Indeed, $`\eta `$ is negative because the order parameter wave function renormalization $`Z_\varphi `$ is greater than one and this happens only if the corresponding critical 2-point correlation function has a pole at nonzero momentum. We argue that this behavior implies the existence of a Lifshitz point induced by gauge field fluctuations.
Another point of view is to study the small fluctuations around the Halperin-Lubensky-Ma (HLM) mean field theory . This is done in section III, where the Gaussian fluctuations are calculated in order to study the positivity properties of the propagators. It turns out that the propagators are positive definite and no pole at $`p0`$ is found at the CP. Indeed, in order to find out such a pole it is necessary to compute the non-Gaussian fluctuations. The analysis of the section III shows that the functional integral has a well defined Gaussian measure.
In section IV we discuss the physical consequences of the anomalous scaling dimension of the magnetic vector potential. After reviewing some known properties like the scaling $`\lambda \xi `$ , we analyse the consequences of the magnetic fluctuations for the frequency dependent conductivity, $`\sigma (\omega )`$. As argued by Fisher et al. , in the $`XY_3`$ regime scales as $`\sigma (\omega )|t|^{\nu (d2z)}`$ (for the sake of generality we wrote the scaling relation in dimension $`2<d<4`$, that is, a $`XY_d`$ regime). However, if the magnetic fluctuations are included the anomalous dimension of the vector potential is no longer equal to zero. This implies the dimension independent scaling $`\sigma (\omega )|t|^{\nu (2z)}`$. We point out that the scaling $`\nu ^{}=\nu `$ implied by $`\lambda \xi `$ ($`\nu ^{}`$ is the penetration depth exponent) is also dimension independent, in contrast with the dimension dependent result of the $`XY_d`$ regime, $`\nu ^{}=\nu (d2)/2`$. In section V we infer from the Monte Carlo data of Lidmar et al. that $`z3.7`$ in the $`IXY_3`$ regime, which is a translation of one unity of the result obtained by these authors ($`z2.7`$). This difference is due to the fact that the scaling $`\sigma (\omega )|t|^{\nu (d2z)}`$ was assumed in their Monte Carlo simulation of the $`IXY_3`$ regime. Finally, we discuss the relevance of these ideas to the Bose-glass transition in the direction perpendicular to the columnar defects, where a transverse Meissner effect happens .
## II Phase transitions and the order parameter anomalous dimension
In order to fix the ideas, let us consider first the case of a scalar $`O(2)`$ invariant field theory with bare Lagrangian
$$L=|_\mu \varphi |^2+m^2|\varphi |^2+\frac{u}{2}|\varphi |^4.$$
(1)
Such a theory has a non-trivial infrared stable fixed point at $`d=3`$. The 2-point bare truncated correlation function is diagonal in the color indices and is defined by
$$W^{(2)}(x,y)=Z^{(2)}(x,y)\varphi (x)\varphi ^{}(y),$$
(2)
where
$$Z^{(2)}(x,y)=\varphi (x)\varphi ^{}(y).$$
(3)
The 2-point function $`Z^{(2)}`$ has the Fourier representation
$$Z^{(2)}(x,y)=\frac{d^dp}{(2\pi )^d}e^{ip(xy)}\stackrel{~}{Z}^{(2)}(p),$$
(4)
which satisfies the KL spectral representation :
$$\stackrel{~}{Z}^{(2)}(p)=c\delta ^d(p)+_0^{\mathrm{}}𝑑\mu \frac{\rho (\mu )}{p^2+\mu ^2},$$
(5)
where $`\rho `$ is the spectral density satisfying
$$_0^{\mathrm{}}𝑑\mu \rho (\mu )=1.$$
(6)
From Eq. (5) we obtain
$$Z^{(2)}(x,y)=c+_0^{\mathrm{}}𝑑\mu \rho (\mu )\frac{e^{\mu |xy|}}{4\pi |xy|}.$$
(7)
Let us put $`y=0`$ for convenience. Then, when the symmetry is broken, $`W^{(2)}(x,0)0`$ as $`|x|\mathrm{}`$. Therefore, from Eqs. (2) and (7) we obtain that $`c=|\varphi (0)|^2`$ which is different from zero if $`T<T_c`$, vanishing otherwise.
Using Eq. (6) it follows easily that the Fourier transform of the bare truncated 2-point correlation function satisfies the infrared bound , $`\stackrel{~}{W}^{(2)}(p)1/p^2`$. Moreover, Griffiths correlation inequality implies $`\stackrel{~}{W}^{(2)}(p)0`$. Therefore,
$$0\stackrel{~}{W}^{(2)}(p)\frac{1}{p^2}.$$
(8)
The inequality (8) has an important consequence for the infrared behavior. At the CP, the bare correlation function behaves as $`\stackrel{~}{W}^{(2)}(p)1/p^{2\eta }`$ as $`p0`$ and Eq. (8) implies therefore that $`\eta 0`$. Note that in the above argument no reference is made to the global character of the symmetry group. Thus, we may think that the same rule should apply to the GL model where the gauge symmetry is local. We will see that this is not the case.
The bare Lagrangian of the GL model is
$$L=\frac{1}{4}F^2+(D_\mu \varphi )^{}(D_\mu \varphi )+m^2|\varphi |^2+\frac{u}{2}|\varphi |^4,$$
(9)
where $`F^2`$ is a short for $`F^{\mu \nu }F^{\mu \nu }`$, $`F^{\mu \nu }=_\mu A_\nu _\nu A_\mu `$, and $`D_\mu =_\mu +ieA_\mu `$. At 1-loop, we obtain for $`d=3`$, $`TT_c`$ and in the Coulomb gauge $`_\mu A_\mu =0`$,
$$\stackrel{~}{W}^{(2)}(p)=\frac{1}{p^2+m^2+\mathrm{\Sigma }(p)},$$
(10)
with the self-energy
$$\mathrm{\Sigma }(p)=\frac{m}{2\pi }(u+e^2)\frac{e^2}{4\pi |p|}(p^2m^2)\left[\frac{\pi }{2}+\mathrm{arctan}\left(\frac{p^2m^2}{2m|p|}\right)\right].$$
(11)
In writing the above equations we have absorbed in the bare mass a contribution with a linear dependence on the ultraviolet cutoff $`\mathrm{\Lambda }`$. Thus, $`m^2t`$, where $`t=(TT_c)/T_c`$ is the reduced temperature. The correlation function $`\stackrel{~}{W}^{(2)}(p)`$ at $`p=0`$ gives the susceptibility $`\chi `$. The divergence of the susceptibility at $`T=T_c`$ signals a phase transition. In terms of the correlation length $`\xi =m_r^1`$, where $`m_r`$ is the renormalized mass, the susceptibility is written as $`\chi =Z_\varphi \xi ^2`$, where $`Z_\varphi `$ is the wave-function renormalization. Here is the crucial point. For the $`O(2)`$ model, the 2-point correlation function diverges at $`T_c`$ only for $`p=0`$. The same is not true for the GL model. In fact, the above 1-loop calculation shows that for $`|p|=p^{}=e^2/4`$ the 2-point correlation function also diverges at $`T_c`$. Thus, we can define a second susceptibility $`\chi ^{}=\stackrel{~}{W}^{(2)}(p^{})`$. The existence of a second pole in $`\stackrel{~}{W}^{(2)}`$ implies that $`Z_\varphi >1`$. Thus, the infrared bound Eq. (8) does not hold. If moreover we assume that the phase transition at $`p=0`$ is of second-order, we obtain that $`\eta <0`$. The same result holds at 2-loops and also in the $`1/N`$ expansion. A negative value of $`\eta `$ is also found by means of non-perturbative RG and in a recent Monte Carlo simulation . This strange behavior needs an explanation and an interpretation. Note that not only the right hand side of (8) is violated but also its left hand side. The striking feature of this behavior is that $`|\stackrel{~}{\varphi }(p)|^20`$ for all $`p`$ but $`\stackrel{~}{W}^{(2)}(p)=|\stackrel{~}{\varphi }(p)|^2<0`$ if $`0<|p|<p^{}`$. The average of the everywhere positive operator $`|\stackrel{~}{\varphi }(p)|^2`$ is not positive everywhere! Thus, it seems that the corresponding effective Gaussian measure is not positive definite and, as a consequence, the functional integral is not well-defined. Of course, the KL representation cannot hold with a positive measure.
Let us explain the meaning of the susceptibility $`\chi ^{}`$. The fact that $`\stackrel{~}{W}^{(2)}(p^{})`$ diverges at $`T_c`$ means that a phase transition happens at finite distances. This is a typical feature of a first-order phase transition. The first- and second-order phase transition can be described at a same $`T_c`$ but at different momentum scales, $`p=0`$ for the second-order phase transition and $`|p|=p^{}`$ for the first-order one. This shed new light in the RG fixed dimension approach at the CP of Refs. , where two momentum scales are considered, defining in this way two characteristic lengths (note that for $`T<T_c`$ there are two lengths in the problem, namely, the correlation length $`\xi `$ and the penetration depth $`\lambda `$). For $`T=T_c`$, the fixed point structure is such that both phase transition regimes are contained in the RG flow diagram determined by dimensionless couplings $`\widehat{u}(\mu )=u_r(\mu )/\mu `$ and $`\widehat{e}^2(\mu )=e_r^2(\mu )/\mu `$, with $`u_r`$ and $`e_r`$ being the renormalized counterparts of $`u`$ and $`e`$. The regions of first- and second-order phase transition are separated by a line connecting the Gaussian fixed point and the so called tricritical fixed point . This fixed point is infrared stable along the tricritical line and unstable in the the direction of $`\widehat{u}`$. For momentum scales such that the couplings are at the left of the tricritical line, the phase transition is of first-order. Concerning the sign of $`\eta `$, it must be observed the following crossovers. The first one corresponds to zero charge, $`\widehat{e}^2=0`$. In this case the flow is towards the $`XY_3`$ fixed point and $`\eta 0`$ ($`\eta =0`$ at 1-loop). This situation is consistent with the infrared bound (8). The other crossover corresponds to the case where the couplings are over the tricritical line. In this situation the flow is towards the tricritical point. Both crossovers give a critical behavior consistent with a second-order like phase transition. It must be stressed, however, that the true second-order phase transition is governed by the infrared stable fixed point. The described crossovers are infrared stable only along the crossover lines, the tricritical line and the line $`\widehat{e}^2=0`$. The critical regime associated to the tricritical line leads to $`\eta <0`$, in contrast to the $`XY_3`$ crossover.
The singularity at $`|p|=p^{}`$ can be interpreted in terms of the effective action. We will write the effective action in momentum space rather than in real space. Thus, if $`\phi `$ and $`a_\mu `$ are the respective Legendre transformed fields of $`\varphi `$ and $`A_\mu `$, we have $`\mathrm{\Gamma }=d^3p/(2\pi )^3\stackrel{~}{\mathrm{\Gamma }}(p)`$, with
$`\stackrel{~}{\mathrm{\Gamma }}(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{\mathrm{\Gamma }}^{(2)}(p)\stackrel{~}{\phi }_i(p)\stackrel{~}{\phi }_i(p)+{\displaystyle \frac{1}{2}}\stackrel{~}{\mathrm{\Gamma }}_{\mu \nu }^{(2)}(p)\stackrel{~}{a}_\mu (p)\stackrel{~}{a}_\nu (p)`$ (12)
$`+`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{d^3q}{(2\pi )^3}\frac{d^3k}{(2\pi )^3}\stackrel{~}{\mathrm{\Gamma }}^{(4)}(p,q,pk,q+k)(\delta _{ij}\delta _{kl}+\delta _{ik}\delta _{jl}+\delta _{il}\delta _{jk})\stackrel{~}{\phi }_i(p)\stackrel{~}{\phi }_j(q)\stackrel{~}{\phi }_k(pk)\stackrel{~}{\phi }_l(q+k)}`$ (13)
$`+`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\stackrel{~}{\mathrm{\Lambda }}_\mu (pk,p,k)\stackrel{~}{a}_\mu (pk)\stackrel{~}{\phi }_1(p)\stackrel{~}{\phi }_2(k)}`$ (14)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }{\displaystyle \frac{d^3q}{(2\pi )^3}}{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}\stackrel{~}{\mathrm{\Omega }}(p,q,pk,q+k)\stackrel{~}{\phi }_i(p)\stackrel{~}{\phi }_i(q)\stackrel{~}{a}_\mu (pk)\stackrel{~}{a}_\mu (q+k)+(h.o.t.),`$ (15)
where summation over repeated latin and greek indices is implied and we have written $`\phi =(\phi _1+i\phi _2)/\sqrt{2}`$. Since $`\stackrel{~}{\mathrm{\Gamma }}^{(2)}(p)=1/\stackrel{~}{W}^{(2)}(p)`$, we have that at the CP the first term of the RHS of Eq. (12) vanishes when $`|p|=p^{}`$ and is negative when $`0<|p|<p^{}`$. Of course, when $`\stackrel{~}{\mathrm{\Gamma }}^{(2)}(p)`$ is negative we must have a positive $`\stackrel{~}{\mathrm{\Gamma }}^{(4)}`$ to ensure the stability of the effective action. In this paper we will not enter into the details of the stability conditions with respect to the 4-point function.
The physical picture that emerges from the behavior of the 2-point function is that of a tricritical Lifshitz point . In fact, in scalar models for the Lifshitz points the 2-point function vanishes at the CP for a nonzero momentum value. In pure scalar models this can happen only if higher derivative Gaussian terms are present already at the tree level . Remarkably, in the GL model the Lifshitz point is induced by the gauge field fluctuations . The existence of a tricritical point in the GL model was established by Kleinert using a disorder field theory obtained from duality arguments. In the disorder field theory scenario, an effective local scalar Lagrangian with disorder parameter $`\psi `$ is constructed. It has been shown that the effective quartic coupling in this model changes sign at some point in the coupling space of the original model. This characterizes an ordinary tricritical behavior in the disorder field theory. In the original GL model this tricritical point is of a Lifshitz type and that is the physical interpretation of the negative sign of $`\eta `$. Note that $`\eta _{dual}`$ is positive in the disorder field theory . In scalar theories of the Lifshitz point the sign of $`\eta `$ is negative in dimensions $`d_c1`$ where $`d_c`$ is the critical dimension of the model. For instance, a fixed dimension calculation in a $`1/N`$ expansion gives for the isotropic Lifshitz point in $`d=7`$ ($`d_c=8`$ in this case) $`\eta _{l4}0.08/N`$ .
A Lifshitz point behavior implies the existence of a modulated regime for the order parameter. This modulated regime should correspond to the type II regime and is analogous to the helical phase in scalar models of the Lifshitz point. The type I regime is analogous to the ferromagnetic or uniform order parameter phase in these models, the normal regime being the analog of the paramagnetic phase. The phase diagram should be therefore quite similar to the phase diagram of the R-S model . The phase diagram of the R-S model is drawn in the $`TX`$ plane where $`X=S/R`$ is the ratio between the couplings $`S`$ and $`R`$. In this phase diagram the line separating the helical phase from the ferromagnetic phase is a first order line. In the case of superconductors we should draw the phase diagram in a $`T\kappa ^2`$ plane, where $`\kappa ^2=u/2e^2`$ is the square of the Ginzburg parameter. The phases paramagnetic, ferromagnetic and helical of the R-S model are replaced respectively by normal, type I and type II. Experimentally, the modulated nature of the order parameter in the type II regime is seen upon applying an external magnetic field and corresponds to the Abrikosov vortex lattice .
## III Wave function renormalization from fluctuations around the Halperin-Lubensky-Ma mean field theory
The HLM mean-field theory neglects the order parameter fluctuations while including the gauge field fluctuations. For an uniform order parameter, the gauge field is integrated out exactly and a term proportional to $`|\varphi |^3`$ with negative sign is generated in the free energy. The corresponding phase transition is found to be wekly first-order. RG calculations using the $`ϵ`$-expansion confirms this scenario since no charged fixed point arises. A stable flow towards the infrared happens only at zero charge and the $`XY_3`$ regime follows by taking $`ϵ=1`$. The $`XY`$ fixed point is unstable for arbitrarily small charge. This behavior remains even at 2-loop order . Charged fixed points are obtained only by considering an order parameter with $`N/2`$ complex components and in the limit of $`N`$ sufficiently large. Indeed, at 1-loop order charged fixed points are obtained if $`N>365.9`$. Interestingly, the critical value of $`N`$ decreases considerably already at 2-loops and charged fixed points are found for $`N>36`$. More recently, by using Padé-Borel resummation of the $`ϵ`$-expansion, Folk and Holovatch succeeded in obtaining charged fixed points for the physical value $`N=2`$.
In this section we will evaluate the Gaussian fluctuations around the HLM mean-field theory. These fluctuations will not suffice for changing the order of the transition, and so it will remains first-order. Our interest here is the positivity properties of the 2-point correlation function in this fluctuation-corrected Gaussian approximation. This amounts in calculating the propagators associated to the HLM mean-field solution. Once this is done, the Gaussian measure necessary to compute the non-Gaussian fluctuations is determined. If this measure is not positive definite, then the functional integral is not well defined and all the theory is inconsistent. We will see that this is not the case.
Let us write $`\varphi =(\varphi _1+i\varphi _2)/\sqrt{2}`$. By integrating out exactly the gauge field we obtain,
$$Z=\underset{a0}{lim}D\varphi _1D\varphi _2\mathrm{exp}(S_{eff}),$$
(16)
where the effective action
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Tr\mathrm{ln}[\widehat{M}_{\mu \nu }(xy;a)]`$ (17)
$``$ $`{\displaystyle \frac{e^2}{2}}{\displaystyle d^3xd^3y[\varphi _1(x)_\mu ^x\varphi _2(x)\varphi _2(x)_\mu ^x\varphi _1(x)]\widehat{D}_{\mu \nu }(xy;a)[\varphi _1(y)_\nu ^y\varphi _2(y)\varphi _2(y)_\nu ^y\varphi _1(y)]}`$ (18)
$`+`$ $`{\displaystyle d^3x\left[\frac{1}{2}\varphi _1(\mathrm{\Delta }+\delta m^2+m^2)\varphi _1+\frac{1}{2}\varphi _2(\mathrm{\Delta }+\delta m^2+m^2)\varphi _2+\frac{u}{8}(\varphi _1^2+\varphi _2^2)^2\right]},`$ (19)
where we have introduced a mass counterterm $`\delta m^2`$, necessary to cancel tadpole divergences (see below). The operator $`\widehat{D}`$ is the inverse of $`\widehat{M}`$, the latter being given by
$$\widehat{M}_{\mu \nu }(xy;a)=\delta ^3(xy)\{[\mathrm{\Delta }+e^2(\varphi _1^2+\varphi _2^2)]\delta _{\mu \nu }+(11/a)_\mu _\nu \},$$
(20)
where $`a`$ is the gauge fixing parameter. In Eq. (16), the limit $`a0`$ is taken in order to inforce the Coulomb gauge condition. Now, we consider small fluctuations around $`\varphi _i=v\delta _{i1}`$, where $`v=const`$ is the solution of
$$\frac{\delta S_{eff}}{\delta \varphi _i}=0.$$
(21)
In this case it is legitimate to truncate $`S_{eff}`$ up to quadratic order in the fluctuating fields $`\delta \varphi _1`$ and $`\delta \varphi _2`$. The result is
$`S_{eff}`$ $`=`$ $`S_{eff}^{HLM}`$ (22)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^3x{\displaystyle }d^3y\{\delta \varphi _1(x)[(\mathrm{\Delta }+\delta m^2+3\overline{m}^2+m^2+e^2D_{\mu \mu }(0))\delta ^3(xy)`$ (23)
$``$ $`2e^2M^2D_{\mu \nu }(xy)D_{\nu \mu }(yx)]\delta \varphi _1(y)`$ (24)
$`+`$ $`\delta \varphi _2(x)[(\mathrm{\Delta }+\delta m^2+\overline{m}^2++m^2+e^2D_{\mu \mu }(0))\delta ^3(xy)M^2_\mu ^x_\nu ^yD_{\mu \nu }(xy)]\delta \varphi _2(y)\},`$ (25)
where $`\overline{m}^2=uv^2/2`$, $`M^2=e^2v^2`$ and $`S_{eff}^{HLM}`$ corresponds to the HLM mean-field free energy . Also,
$$D_{\mu \nu }(xy)=\frac{d^3p}{(2\pi )^3}e^{ip(xy)}\stackrel{~}{D}_{\mu \nu }(p)$$
(26)
is the operator $`\widehat{D}`$ for $`\delta \varphi _1=\delta \varphi _2=0`$. In the Coulomb gauge,
$$\stackrel{~}{D}_{\mu \nu }(p)=\frac{1}{p^2+M^2}\left(\delta _{\mu \nu }\frac{p_\mu p_\nu }{p^2}\right),$$
(27)
which implies
$$d^3xd^3y_\mu ^x_\nu ^yD_{\mu \nu }(xy)\delta \varphi _2(x)\delta \varphi _2(y)=0.$$
(28)
Now we see that the counterterm $`\delta m^2`$ is necessary in order to cancel the linear cutoff dependence coming from the tadpole term $`e^2D_{\mu \mu }(0)`$. Therefore, the $`\delta \varphi _1`$ propagator is
$$G_{11}(p)=\frac{1}{p^2+m^2+3\overline{m}^2+\mathrm{\Sigma }_{11}(p)},$$
(29)
where the self-energy $`\mathrm{\Sigma }_{11}(p)`$ is given by
$$\mathrm{\Sigma }_{11}(p)=e^2D_{\mu \mu }(0)2e^2M^2\frac{d^3k}{(2\pi )^3}\stackrel{~}{D}_{\mu \nu }(kp)\stackrel{~}{D}_{\nu \mu }(k).$$
(30)
By evaluating the integrals in Eq. (30) we obtain
$`\mathrm{\Sigma }_{11}(p)`$ $`=`$ $`{\displaystyle \frac{e^2M}{2\pi }}e^2[{\displaystyle \frac{M}{4\pi }}{\displaystyle \frac{p}{8}}{\displaystyle \frac{M^2}{16p}}+{\displaystyle \frac{p^4+8M^4+4M^2p^2}{8\pi pM^2}}\mathrm{arctan}\left({\displaystyle \frac{p}{2M}}\right)`$ (31)
$``$ $`{\displaystyle \frac{(p^2+M^2)^2}{8\pi pM^2}}\mathrm{arctan}\left({\displaystyle \frac{p^2M^2}{2pM}}\right)].`$ (32)
The $`\delta \varphi _2`$ propagator is given simply by
$$G_{22}(p)=\frac{1}{p^2+m^2+\overline{m}^2\frac{e^2M}{2\pi }}.$$
(33)
Note that the above calculation differs from the usual 1-loop result. In an ordinary 1-loop calculation we integrate out the quadratic fluctuations around the solution $`v=(2m^2/u)^{1/2}`$ corresponding to the tree-level with $`𝐀=0`$. Above, we integrated out $`𝐀`$ first and then we computed the Gaussian fluctuations around a solution $`v`$ given by Eq. (21), which already contain magnetic fluctuations. Eq. (21) have the following solutions:
$$v=0,$$
(34)
$$|v|=\frac{e^3}{2\pi u}\pm \frac{1}{u}\sqrt{\frac{e^6}{4\pi ^2}2um^2}.$$
(35)
Thus, in the ordered phase given by Eq. (35) we have $`m^2+\overline{m}^2e^2M/2\pi =0`$ and therefore the $`\delta \varphi _2`$ propagator $`G_22(p)`$ is massless. Then, the fluctuating field $`\delta \varphi _2`$ is the would-be Goldstone boson of the theory. When $`m^2=e^6/(8\pi ^2u)`$ the square root in Eq. (35) vanishes. This value of $`m^2`$ corresponds to a point of non-analyticity of $`v`$ as a function of $`m^2`$. Indeed, the derivative of $`v`$ with respect to $`m^2`$ diverges for $`m^2=e^6/(8\pi ^2u)`$. Thus, if we expand the denominator of $`G_{11}`$ for $`p`$ small we obtain,
$$G_{11}(p)=\frac{1}{\left(1+\frac{5}{24\pi }\frac{e^2}{M}\right)p^2+m^2+3\overline{m}^2\frac{e^2M}{\pi }+O(p^4)},$$
(36)
and we see that the susceptibilities $`\chi _i=G_{ii}(0)`$ ($`i=1,2`$) diverge together if $`m^2=e^6/(8\pi ^2u)`$. This singular behavior of the susceptibilities is not associated to any phase transition. It is just an artifact of our fluctuation-corrected Gaussian approximation. The singularity of $`\chi _i`$ for $`m^2=e^6/(8\pi ^2u)`$ is inherited from the non-analytic behavior of $`v`$ for this value of $`m^2`$. Once the non-Gaussian fluctuations are taken into account and a full renormalization of mass and coupling constants is done, this artifact disappears. On the other hand, if $`m^2<e^6/(8\pi ^2u)`$, $`\chi _2`$ diverge but not $`\chi _1`$. In this fluctuation induced phase transition scenario $`v=\varphi `$ is different from zero at the CP, a typical behavior of a first-order transition, as we have already discussed in section II. Note that in this calculation the correlation functions diverge only at $`p=0`$. As a consequence, the wave function renormalizations $`Z_i1`$. By putting $`v=e^3/(\pi u)`$ which corresponds to $`m^2=0`$ in Eq. (36), we obtain
$$Z_1=\frac{1}{1+\frac{5}{12}\kappa ^2}<1.$$
(37)
Therefore, the corresponding Gaussian measure is positive definite and the non-Gaussian fluctuations can be calculated by means of this measure. The non-Gaussian fluctuations will ultimately make $`Z_1>1`$, violating again the KL representation.
## IV The vector potential anomalous dimension and its consequences for the critical dynamics in superconductors
One important feature of the $`IXY_3`$ universality class is the scaling $`\lambda \xi `$ , where $`\lambda `$ and $`\xi `$ are the penetration depth and correlation length, respectively. This scaling contrast with the $`XY_3`$ behavior, where $`\lambda \xi ^{1/2}`$ . The reason for this different behavior comes from the magnetic fluctuations, which in the $`XY_3`$ universality class play no role. In the $`XY_3`$ regime the magnetic vector potential has no anomalous dimension. Concerning the scaling of the penetration depth, it was argued in Refs. that the vector potential anomalous dimension contributes in such a way that we have in general $`\lambda \xi ^{(\eta _A+d2)/2}`$. Thus, when the magnetic fluctuations are negligeable we have $`\eta _A=0`$ and $`\lambda \xi ^{(d2)/2}`$, implying in this way a penetration depth exponent $`\nu ^{}=\nu (d2)/2`$ with $`\nu 2/3`$ when $`d=3`$. On the other hand, if we take into account the magnetic fluctuations, we have that $`\eta _A=4d`$ and $`\lambda \xi `$ implying $`\nu ^{}=\nu `$. The critical exponent $`\nu `$ is the same in both $`XY_3`$ and $`IXY_3`$ universality classes and we obtain that $`\nu ^{}1/3`$ and $`\nu ^{}2/3`$ for the $`XY_3`$ and $`IXY_3`$ regimes, respectively. Note that only the thermodynamic exponents coincide in these two $`XY`$ regimes. As we have already seen, the anomalous dimensions are not the same.
The frequency-dependent conductivity $`\sigma (\omega )`$ scales differently in a magnetic fluctuation regime. For $`T<T_c`$ we have that $`\sigma (\omega )e^2\rho _s/(i\omega )`$, where $`\rho _s`$ is the superfluid density. Near a charged fixed point we have $`e^2\xi ^{\eta _A}`$ and therefore,
$$\sigma (\omega )\xi ^{2d+z\eta _A},$$
(38)
where $`z`$ is the dynamical exponent and we have used the Josephson relation $`\rho _s\xi ^{2d}`$ . Again, by neglecting the magnetic fluctuations we recover the usual scaling . The $`XY`$ scaling proposed by Fisher et al. was verified recently by Wickham and Dorsey , who calculated $`\sigma (\omega )`$ using the Kubo formula to $`O(ϵ^2)`$ in the $`ϵ=4d`$-expansion.
Since in the magnetic fluctuation regime $`\eta _A=4d`$, we obtain
$$\sigma (\omega )\xi ^{z2},$$
(39)
a result independent of the dimension. This independence of the dimension in the scaling behavior (39) seems to be a special feature of the charged fixed point. Note that already in the case of the penetration depth we have obtained $`\nu ^{}=\nu `$ instead of the dimension dependent result $`\nu ^{}=\nu (d2)/2`$ of the $`XY`$ regime. The scaling given in Eq. (39) has been obtained before by Mou who used a completely different argument. Our argument is much more simple and follows from the exact value of the vector potential anomalous dimension. However, the dynamical exponent $`z`$ is not be the same as in the uncharged model, as was claimed in Ref. . The Monte Carlo simulations of Lidmar et al. show very clearly that this is not the case and that the value of $`z`$ is enhanced by magnetic fluctuations. However, Lidmar et al. fitted their Monte Carlo data to $`\sigma |t|^{\nu (d2z)}`$ instead of using the scaling (39). Since $`\eta _A=1`$ in $`d=3`$, we conclude that the numerical result of Ref. should be shifted to obtain $`z3.7`$ instead of $`z2.7`$. This surprisingly high value of $`z`$ could be, however, a matter of controversy. It may be a consequence of the way the authors of Ref. modelled the $`IV`$ characteristics of the $`IXY_3`$ regime. For instance, Ampère’s law is neglected in their approach.
From the experimental side, the work of Booth et al. fit reasonably the value $`z=2.7`$ but they assume also a scaling with $`\eta _A=0`$. Anyway, in the case of Ref. it is more probable that the critical region probed does not correspond to a $`IXY_3`$ universality class. In this case $`\eta _A=0`$ would be a legitimate assumption.
The scaling Eq. (39) is also relevant in other situations. For example, in the Bose-glass transition the conductivity perpendicular to the columnar defects is argued to obey a scaling exactly as in Eq. (39) . Although Eq. (39) is a zero field scaling, it should apply in the nonzero field situation of the Bose-glass transition in the direction perpendicular to the columnar defects but not in the longitudinal direction. The reason for this behavior comes from the fact that in the perpendicular direction a transverse Meissner effect happens, implying in this way a zero field like situation. The Bose-glass transition is an example which shows that at high fields the magnetic thermal fluctuations may be experimentally important and observable.
## V Conclusion
In this paper we have discussed some new features of the superconducting phase transition. The important role of the anomalous dimensions has been emphasized. However, the critical behavior disussed in this paper is relevant only near a charged fluctuation critical regime. The relevance of the ideas discussed here to high-temperature superconductors (HTSCs) may be questioned. Usually the HTSCs have very high values of $`\kappa `$, typically in the range $`70100`$. For this reason, it is generally assumed that magnetic fluctuations do not play an important role. This is in fact the case in the extreme type II limit, that is, $`\kappa \mathrm{}`$. For extreme type II superconductors, the local magnetic induction equals the applied magnetic field and the constraint $`\times 𝐀=𝐇`$ applies . In the presently accessible critical region, the HTSCs seem to be well approximated by an extreme type II limit. In this case the $`XY_3`$ regime dominate at zero or low magnetic fields. The $`XY_3`$ behavior has been probed with considerable confidence in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> (YBCO) crystal samples . For Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> (BSCCO), however, the situation is less clear due to the experimental difficulties involved. Specific heat measurements seem to indicate that the universality class is not $`XY_3`$ . The apparent failure of the $`XY_3`$ scaling in BSCCO seems also to be corroborated by the penetration depth data . However, inhomegeneities and finite size effects can play a significant role in BSCCO and it may happen that it obeys also a $`XY_3`$ scaling .
The $`IXY_3`$ behavior, on the other hand, seems to be not presently accessible. In fact, penetration depth data from YBCO fulfill very well the scaling relation $`\nu ^{}=\nu /2`$ , agreeing with the $`XY_3`$ behavior. Thus, in order to check the theoretical predictions concerning the $`IXY_3`$ regime, we have to compare these mainly to Monte Carlo simulations. For instance, the scaling relation $`\nu ^{}=\nu `$ with $`\nu 2/3`$ was well verified by Olsson and Teitel . The value $`\eta 0.18`$ was obtained by Nguyen and Sudbø. The dynamical exponent $`z`$ was studied by Lidmar et al. both in the $`XY_3`$ and in the $`IXY_3`$ regimes. However, as discussed in section IV, they assumed the same scaling for the frequency dependent conductivity in both regimes. This does not invalidate their data, which remain useful and lead to the prediction $`z3.7`$ instead of $`z2.7`$. While presently there is little hope in checking these predictions in zero field experiments, further Monte Carlo simulations can be done in order to obtain a definitive answer. As far as real experiments are concerned, we have pointed out that the scaling given in Eq. (39) holds for the condcutivity perpendicular to the columnar defects in a Bose-glass transition . Unfortunately, in this nonzero field regime we are unable to estimate the value of $`z`$ with the arguments presented in this paper. It is worth to mention, however, that an experimental value $`z5.3`$ was probed recently by Klein et al. for the Bose-Glass transition in the fully isotropic compound (K,Ba)BiO<sub>3</sub> with columnar defects.
###### Acknowledgements.
The author would like to acknowledge A. Sudbø for sending the paper Ref. prior to publication. Much of the present work originated from discussions with A. Sudbø and Z. Tes̆anović and the author is indebted to them. The author acknowledges also C. de Calan for numerous discussions.
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# Presence of energy flux in quantum spin chains: An experimental signature
## I Introduction
The problem of thermal transport in one-dimensional systems has been much investigated, the main goal being to derive Fourier’s heat law. Analytical and numerical studies of a number of classical lattice-dynamical models indicate that the condition for Fourier’s law to hold is the presence of strong nonlinearities i.e. nonintegrability (chaoticity) of the dynamics . Although quantum systems have been less studied, it appears that similar considerations apply to quantum spin chains as well.
Integrable systems, on the other hand, show anomalous thermal transport. No internal thermal gradient is formed in a harmonic crystal or in a transverse Ising chain and, as a consequence, the energy (heat) flux is not proportional to the temperature gradient inside the sample. The origin of this anomaly may be the fact that the energy current in integrable systems often emerges as an integral of motion which automatically yields anomalous thermal transport coefficients .
The flat temperature profile in the presence of energy current is an intriguing feature of integrable systems. In effect, it points to the existence of a homogeneous state carrying finite energy current. In this paper, we shall explore the experimentally measurable properties of such a state by studying the $`XXZ`$ spin chain in the presence of an energy current. The most spectacular feature of such a state is the incommensurability of magnetic excitations. Namely, in the presense of energy flow $`j_E`$, the characteristic wave vector is shifted from its antiferromagnetic value $`\pi `$ by the amount $`\delta k\sqrt{j_E}`$. Since there are quite a few well established realizations of quasi-one-dimensional Heisenberg chains (e.g. $`\mathrm{KCuF}_3`$ , $`\mathrm{Cs}_2\mathrm{CoCl}_4`$ , Copper Benzoate , $`\mathrm{Sr}_2\mathrm{CuO}_3`$, $`\mathrm{Cs}_2\mathrm{CuCl}_4`$ ), we believe that the predicted changes in the dynamical correlation functions bear direct experimental relevance.
The basic problem of constructing a state which carries an energy current is the nonequilibrium nature of that state. Even if we assume that the flat temperature profile means the existence of equilibrium, we still face a problem that the value of the established temperature is not known . We shall avoid this problem by restricting our calculation to zero temperature $`(T=0)`$ and assuming that the ground state correlations are robust enough to survive at low temperatures.
The construction of a homogeneous state with energy current at $`T=0`$ can be done by adding the energy current with a Lagrange multiplier to the $`XXZ`$ Hamiltonian and then finding the ground state. Similar calculations have been carried out already for the transverse Ising and $`XX`$ chains and, in a different context, for the $`XXZ`$ model . The new result we report is the calculation of an experimentally accessible parameter, namely the shift, $`\delta k`$, of the characteristic wavenumber in the spin-spin correlations as a function of the energy current, $`j_E`$.
Once we have $`\delta k(j_E)`$, we turn to a realistic experimental setup and estimate $`j_E`$ flowing through a single spin chain which gives an estimate of $`\delta k`$. Our result shows that $`\delta k`$ is in the accessible range of an inelastic neutron scattering experiment.
## II The model and the characteristic wave number
The model we study is the spin-1/2 XXZ chain defined by the Hamiltonian
$$\widehat{H}_{\mathrm{xxz}}=𝒥\underset{\mathrm{}}{}\left[\sigma _{\mathrm{}}^x\sigma _{\mathrm{}+1}^x+\sigma _{\mathrm{}}^y\sigma _{\mathrm{}+1}^y+\mathrm{\Delta }\sigma _{\mathrm{}}^z\sigma _{\mathrm{}+1}^z\right],$$
(1)
where the spins $`\sigma _{\mathrm{}}^\alpha `$ ($`\alpha =x,y,z`$) are Pauli spin matrices at sites $`\mathrm{}=1,2,\mathrm{},N`$ of a one-dimensional periodic chain ($`\sigma _{N+1}^\alpha =s_1^\alpha `$). We shall use the parametrization $`\mathrm{\Delta }=\mathrm{cos}\gamma `$ and consider only the ‘antiferromagnetic’ region $`0<\gamma <\pi /2`$. In order to impose a fixed energy current $`j_E`$ in the ground state, we add the current operator to the Hamiltonian with a Lagrange multiplier
$$\widehat{H}=\widehat{H}_{\mathrm{xxz}}+\lambda \widehat{j}_E.$$
(2)
where
$`\widehat{j}_E={\displaystyle \frac{𝒥^2}{\mathrm{}}}{\displaystyle \underset{\mathrm{}}{}}\sigma _{\mathrm{}}^z[\sigma _\mathrm{}1^y\sigma _{\mathrm{}+1}^x\sigma _\mathrm{}1^x\sigma _{\mathrm{}+1}^y+`$ (3)
$`\mathrm{\Delta }(\sigma _\mathrm{}2^x\sigma _\mathrm{}1^y\sigma _\mathrm{}2^y\sigma _\mathrm{}1^x+\sigma _{\mathrm{}+1}^x\sigma _{\mathrm{}+2}^y\sigma _{\mathrm{}+1}^y\sigma _{\mathrm{}+2}^x)].`$ (4)
Importantly, $`\widehat{j}_E`$ is an integral of motion, $`[\widehat{j}_E,\widehat{H}]=0`$, thus indicating that i) the transport of energy is singular in this system and ii) the states carrying fixed energy current can be obtained as stationary states of $`\widehat{H}_{\mathrm{xxz}}`$.
The $`XXZ`$ model can be described in terms of interacting fermions and has been solved using the Bethe Ansatz method. The same approach works in the presence of the driving term, $`\lambda \widehat{j}_E`$, as well, and the solution has been given in . An interesting feature of the solution is that the system displays rigidity against the drive, namely the ground state supports a nonzero energy current, $`\widehat{j}_Ej_E0`$, only if the coupling $`\lambda `$ exceeds some critical value $`\lambda _c(\gamma )`$. As we are interested in fixed energy currents, we simply choose sufficiently large values of $`\lambda `$. Furthermore, since in realistic situations $`j_E`$ turns out to be small, we concentrate on the region $`\lambda \lambda _c(\gamma )`$, in which case $`j_E(\lambda \lambda _c)`$.
Once the energy current flows, an important restructuring takes place in the ground state. The single Fermi sea characterizing the ground state without current splits into two Fermi seas as shown in Fig.1 for the simple case of the $`XX`$ limit ($`\mathrm{\Delta }=0`$) where a free-fermion description applies. There are now four Fermi wave vectors, $`\pm \pi /2`$ and $`\pi /2\pm \delta k`$, and the structure of the ground state immediately implies that there will be gapless excitations at wave vectors $`0,\delta k,2\delta k,\pi \delta k,\pi `$ and $`\pi +\delta k`$, which is readily confirmed by the exact solution at arbitrary $`\mathrm{\Delta }`$.
Thus an incommensurability characterized by $`\delta k`$ appears in the system. This can be seen readily in the ground-state correlations. Indeed, it has been shown that, in the $`XX`$ limit, the longitudinal correlations for small $`j_E0`$ can be expressed in a scaling form
$$\frac{\sigma _{\mathrm{}}^x\sigma _{\mathrm{}+n}^x_{j_E0}}{\sigma _{\mathrm{}}^x\sigma _{\mathrm{}+n}^x_{j_E=0}}=\mathrm{\Phi }(\delta kn)$$
(5)
where $`\mathrm{\Phi }(x0)=1`$ and the large argument asymptotics of the scaling function is given by
$$\underset{x\mathrm{}}{lim}\mathrm{\Phi }(x)\frac{1}{\sqrt{x}}(1+\mathrm{cos}x).$$
(6)
Since $`\sigma _{\mathrm{}}^x\sigma _{\mathrm{}+n}^x_{j_E=0}(1)^n/\sqrt{n}`$ equations (5,6) imply that, as the current is switched on, the static structure factor develops additional peaks at $`k=\pi \pm \delta k`$ (as it will turn out, $`\delta k`$ is small thus it is better to speak about the $`k=\pi `$ peak developing shoulders for $`j_E0`$).
In order to connect $`\delta k`$ to the current one determines both $`j_E`$ and $`\delta k`$ through $`\lambda `$ and then eliminates the Lagrange multiplier. The expressions are simple for the $`XX`$ limit
$$j_E=\frac{𝒥^2}{2\pi \mathrm{}}(1\frac{1}{\lambda ^2}),\mathrm{cos}\delta k=\lambda ^1$$
(7)
and, for small currents ($`\lambda \lambda _c=1`$), they yield
$$\delta k=\sqrt{\frac{j_E}{j_E^{(1)}}}.$$
(8)
where a ‘natural unit’ of the current, $`j_E^{(1)}=𝒥^2/h`$, has been introduced.
The above calculation can be carried out for any $`0\mathrm{\Delta }<1`$ and the result for $`\delta k`$ differs only in a prefactor of order unity
$$\delta k=\frac{2\gamma }{\pi \mathrm{sin}\gamma }\sqrt{\frac{j_E}{j_E^{(1)}}}.$$
(9)
As one can see, the largest $`\delta k`$ is obtained in the $`XX`$ limit ($`\gamma \pi /2`$).
In principle, if $`\delta k`$ is large enough then the extra peaks at $`\pi \pm \delta k`$ should be observable as Bragg peaks in an elastic neutron scattering experiment. In practice, however, the incommensurate modulations of distinct spin chains are not correlated and, as a consequence, the delta function of the Bragg peak would spread out into a plain and the effect would be unobservable.
## III Structure factor and an experimental setup
It is more promising to look for an experimental signature in an inelastic neutron scattering experiment where the excitations of the system are measured and no coherence among the chains is needed. Taking into account the facts that, for $`j_E=0`$, most of the spectral weight is concentrated on the region around the antiferromagnetic wave vector $`\pi `$, and furthermore that, for $`j_E0`$, there are gapless excitations at wave vectors $`\pi \pm \delta k`$, we expect that the presence of the current manisfests itself via the emergence of additional inelastic peaks at wave vector $`\pi \pm \delta k`$. This expectation can be put on a more solid base by calculating the dynamic structure factor and examining the relative weights at wave vectors $`\pi `$ and $`\pi \pm \delta k`$.
The simplest case is again the $`XX`$ limit where the calculation of the time-dependent transverse correlation function, $`\sigma _n^z(t)\sigma _0^z(0)`$, is straightforward. There is, however, a principal problem at the outset of the calculation. Namely, it is not clear whether the time-evolution of $`\sigma _n^z(t)`$ is governed by $`\widehat{H}_{\mathrm{xx}}`$ or by $`\widehat{H}_{\mathrm{xx}}\lambda \widehat{j}^E`$. We shall take the view that in reality the current-carrying state is formed as a result of boundary conditions. Thus the local perturbation caused by an incoming neutron evolves by the local hamiltonian i.e. by $`\widehat{H}_{\mathrm{xx}}`$ . Once $`\sigma _n^z(t)\sigma _0^z(0)`$ is known its Fourier transform in time and space gives the structure factor $`S_{zz}(k,\omega )`$ as shown in Fig.2.
As we can see, a large part of the weight of the $`j_E=0`$ peak of the structure factor at $`\pi `$ shifts to $`\pi \pm \delta k`$ for $`j_E0`$. Thus one can expect that even if $`\delta k`$ is small, the presence of a small energy current will result in a broadening by $`2\delta k`$ of the inelastic peak centered at wavevector $`\pi `$. It is this broadening that we propose as an experimental signature for the current-carrying state. The remaining question now is how to estimate $`\delta k`$.
As we can see from (9), an estimate of $`\delta k`$ requires the value of the energy current, $`j_E`$. Thus we should, in principle, calculate $`j_E`$ in a spin chain where the two ends are kept at different temperatures. We are unable to do this for any reasonable size system, and so we shall treat the energy flux as a parameter taken from experiments ($`j_Ej_E^{exp}`$). Then a thermodynamic measurement of $`j_E^{exp}`$ can be used to estimate the value of $`\delta k`$ in an independent neutron-scattering experiment. Below we shall show how to estimate $`j_E^{exp}`$ using parameters from a realistic experimental setup.
Let the sample be a cube of side $`l=10^2m`$ and let the spin chains be along $`x`$ direction with the distance between the neighboring chains being $`d=10^9m`$. Furthermore, let the sides of the cube perpendicular to the chains be at temperatures $`T`$ and $`T+\delta T`$ (see Fig.3). The temperature should be chosen to be low in order to minimize the phonon contribution to $`j_E^{exp}`$. However, $`T`$ cannot be too small since then the small coupling between the chains makes the system three-dimensional. The value of $`\delta T`$ should be, in principle, chosen large but the limitations of cryogenics of a realistic setup restrict the steady-state temperature-differences to $`\delta T0.1T`$. From the above considerations we arrive at the following ranges for the possible temperatures and temperature-differences
$$T=(110)^oK;\delta T=0.1T=(0.11)^oK.$$
(10)
The total current of heat across the sample can now be estimated as
$`j_E^{total}=\kappa l^2{\displaystyle \frac{\delta T}{l}}`$ (11)
provided we know the heat conductivity, $`\kappa `$. We note here that the finiteness of the experimental $`\kappa `$ is not in contradiction with the singular nature ($`\kappa ^{int}=\mathrm{}`$) of the heat conductivity of integrable spin chains. A macroscopic sample consist of spin chains of characteristic length $`\mathrm{}10^4cm`$ and the energy must also be transported between chains. This leads to the loss of ideal conductivity and results in a finite $`\kappa `$. Consequently, the estimate of energy flux using the experimental $`\kappa `$ does give an estimate of the energy flux through the chains provided the spin chains are the main channels of energy transport.
Unfortunately, $`\kappa `$, is not available for the materials we have in mind, , and another problematic issue is how much of the conductivity comes from the spin-chains. Since measurements of the magnetothermal conductivity of magnetic materials indicate that spin waves provide a significant fraction of the low-temperature thermal conductivity, we shall assume that an order of magnitude estimate of the energy current through the spin chains is given by $`j_E^{total}`$. Furthermore, we shall assume that, as a value of $`\kappa `$, we can take a characteristic value of this parameter in crystalline magnetic materials in the temperature range $`T=(110)^oK`$ :
$$\kappa (110)\frac{W}{m{}_{}{}^{o}K}.$$
(12)
We can then estimate $`j_E^{total}(10^310^1)W`$ and, since the number of spin-chains in the sample is $`𝒩=l^2/d^2=(10^2/10^9)^2=10^{14}`$, we obtain the energy flux per chain, $`j_E^{exp}`$, as
$$j_E^{exp}\frac{j_E^{total}}{𝒩}=(10^{17}10^{15})W.$$
(13)
As we have seen (8) the natural unit of energy current in a spin chain is $`j_E^{(1)}=J^2/h`$. Using a characteristic value of $`J=(110)^oK`$ for the spin coupling we find $`j_E^{(1)}(10^{12}10^{10})W`$ and obtain the following estimate for the shift of the wavenumber
$$\delta k\sqrt{\frac{j_E^{exp}}{j_E^{(1)}}}10^410^2.$$
(14)
This is our central result. Since $`\delta k10^2`$ is accessible in an inelastic neutron scattering experiment, the effect of shift in the wavenumber should be observable.
In summary, we have studied an integrable system which doesn’t obey Fourier’s law. We proposed that, under some simplifying assumptions, one can explore states of this system which carry current of energy and, furthermore, one can derive theoretical results verifiable in experiments.
## Acknowledgement
I thank T. Antal, R. Cowley, F. Essler, A. Rákos, L. Sasvári, G. M. Schütz, and A. Tsvelik for helpful discussions. This work has been supported by the by the Hungarian Academy of Sciences (Grant No. OTKA T 029792).
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# Rapidly convergent series for the Weierstrass zeta-function and the Kronecker function
## 1. Formulas
In this section we derive our formulas by classical means. In the next section we’ll explain how one can guess these formulas from the computation of certain triple products on elliptic curves.
### 1.1. Weierstrass zeta-function
Let $`L`$ be a lattice in $``$, $`\omega _1`$, $`\omega _2`$ be generators of $`L`$ such that $`\mathrm{Im}(\overline{\omega _1}\omega _2)>0`$. The Weierstrass zeta-function is defined by the series
$$\zeta (x,L)=\frac{1}{x}+\underset{\omega L0}{}(\frac{1}{x+\omega }\frac{1}{\omega }+\frac{x}{\omega ^2})$$
for $`xL`$. One has
$$\zeta (x+\omega ,L)\zeta (x,L)=\eta (\omega )$$
for any $`\omega L`$, where $`\eta (\omega )`$ is a constant. We denote $`\eta _i=\eta (\omega _i)`$ for $`i=1,2`$. It is well-known that
(1.1)
$$\eta _i=2\zeta (\frac{\omega _i}{2})$$
for $`i=1,2`$. Following Hecke (see ) for any $`x=x_1\omega _1+x_2\omega _2L`$ we set
$$Z(x,L)=\zeta (x,L)x_1\eta _1x_2\eta _2$$
(here $`x_1`$ and $`x_2`$ are real).
###### Theorem 1.
For any $`xL`$ one has the following identity
(1.2)
$$Z(x,L)=\underset{\omega L}{}\frac{\mathrm{exp}(\frac{\pi }{a(L)}|\omega +x|^2)}{\omega +x}\underset{\omega L0}{}\frac{\mathrm{exp}(\frac{\pi }{a(L)}|\omega |^2+2\pi iE_L(\omega ,x))}{\omega }$$
where $`a(L)=\mathrm{Im}(\overline{\omega _1}\omega _2)`$ is the area of $`/L`$,
$$E_L(x,y)=\frac{\mathrm{Im}(\overline{x}y)}{a(L)}=\frac{\overline{x}yx\overline{y}}{2ia(L)}$$
is the symplectic form on $``$ (considered as a real space) associated with the oriented lattice $`L`$.
Let $`\mathrm{SS}()`$ be the Schwarz space of $``$. For any $`\phi \mathrm{SS}()`$ we define its symplectic Fourier transform by the formula
$$\widehat{\phi }(y)=_x\phi (x)\mathrm{exp}(2\pi iE_L(y,x))d_Lx$$
where $`d_Lx`$ is the Haar measure on $``$ normalized by the condition $`_{/L}d_Lx=1`$.
We will need the following simple lemma.
###### Lemma 1.1.
For any $`\phi \mathrm{SS}()`$, any $`x,y`$ one has
$$\underset{\omega L}{}\phi (\omega +x)\mathrm{exp}(2\pi iE_L(\omega +x,y))=\underset{\omega L}{}\widehat{\phi }(\omega +y)\mathrm{exp}(2\pi iE_L(\omega ,x))$$
Proof. By Poincare summation formula the distribution $`\delta _L=_{\omega L}\delta _\omega `$ is Fourier self-dual. Since the translation by $`x`$ goes under Fourier transform to the multiplication by $`\mathrm{exp}(2\pi iE_L(\mathrm{?},x))`$ the result follows. ∎
Proof of theorem 1. Let us denote
$$f(x)=\underset{\omega L}{}\frac{\mathrm{exp}(\frac{\pi }{a(L)}|\omega +x|^2)}{\omega +x}\underset{\omega L0}{}\frac{\mathrm{exp}(\frac{\pi }{a(L)}|\omega |^2+2\pi iE_L(\omega ,x))}{\omega }$$
where $`xL`$. We claim that
$$\overline{}f(x)=\frac{\pi }{a(L)}.$$
Indeed, we have
$$\overline{}f(x)=\frac{\pi }{a(L)}\left(\underset{\omega L}{}\mathrm{exp}(\frac{\pi }{a(L)}|\omega +x|^2)\underset{\omega L0}{}\mathrm{exp}(\frac{\pi }{a(L)}|\omega |^2+2\pi iE_L(\omega ,x))\right),$$
so our claim follows from Lemma 1.1 since the function $`\mathrm{exp}(\frac{\pi }{a(L)}|x|^2)`$ goes to itself under the symplectic Fourier transform.
It follows that the function $`g(x)=f(x)+\frac{\pi }{a(L)}\overline{x}`$ is holomorphic on $`L`$. Furthermore, looking at the series for $`f`$ we immediately see that $`g`$ has simple poles at all the lattice points $`\omega L`$ with residues equal to $`1`$.
On the other hand, from the fact that the symplectic form $`E_L`$ takes integer values on $`L`$ one immediately derives that $`f(x+\omega )=f(x)`$ for any $`\omega L`$. Thus, we have
$$g(x+\omega )=g(x)+\frac{\pi }{a(L)}\overline{\omega }$$
for $`\omega L`$. The Legendre period relation
$$\eta _1\omega _2\eta _2\omega _1=2\pi i$$
implies that there exists a constant $`c`$ such that
(1.3)
$$\eta _i=c\omega _i+\frac{\pi }{a(L)}\overline{\omega _i}$$
for $`i=1,2`$. It follows that $`h(x)=g(x)\zeta (x,L)+cx`$ is a holomorphic function on $`CL`$, periodic with respect to $`L`$. Comparing the polar parts of $`g`$ and $`\zeta `$ at the lattice points we conclude that $`h`$ is holomorphic on $``$. Therefore, $`h`$ is constant and we have
$$f(x)\zeta (x,L)=\frac{\pi }{a(L)}\overline{x}cx+h.$$
From the definition of the constant $`c`$ we derive that
$$x_1\eta _1+x_2\eta _2=cx+\frac{\pi }{a(L)}\overline{x}$$
where $`x=x_1\omega _1+x_2\omega _2`$. Thus, we have
$$f(x)\zeta (x,L)=x_1\eta _1x_2\eta _2+h.$$
It is easy to see that if $`x\frac{1}{2}LL`$ then $`f(x)`$. Thus, substituting $`x=\frac{\omega _1}{2}`$ and using the identity (1.1) we derive that $`h=0`$. ∎
Remarks. 1. It is not obvious that the right hand side of (1.2) is holomorphic in $`\omega _1`$, $`\omega _2`$. This fact is equivalent to the identity
$$\underset{\omega L}{}(\omega +x)\mathrm{exp}(\frac{\pi }{a(L)}|\omega +x|^2)=\underset{\omega L}{}\omega \mathrm{exp}(\frac{\pi }{a(L)}|\omega |^2+2\pi iE_L(\omega ,x))$$
which can be easily deduced from Lemma 1.1.
2. It is well-known that the following Epstein’s zeta function
$$\phi _1(s,L,x)=\underset{\omega L}{}\frac{1}{(\omega +x)|\omega +x|^{2s1}}$$
defined for $`\mathrm{Re}(s)>1`$ extends to an entire function of $`s`$ (for fixed $`xL`$). It was shown by N. Katz (Cor. 3.2.24 of ) that
$$\phi _1(\frac{1}{2},L,x)=Z(x,L)$$
for $`xLL`$. It is not clear whether there exists an expression for $`\phi _1(s,L,x)`$ for arbitrary $`s`$ similar to the one in (1.2).
Differentiating the identity (1.2) we obtain the following series for the Weierstrass $`\mathrm{}`$-function $`\mathrm{}(x)=\zeta ^{}(x)`$:
(1.4)
$$\mathrm{}(x)=c+\underset{\omega L}{}\frac{(1+\frac{\pi }{a(L)}|\omega +x|^2)\mathrm{exp}(\frac{\pi }{a(L)}|\omega +x|^2)}{(\omega +x)^2}+\frac{\pi }{a(L)}\underset{\omega L0}{}\frac{|\omega |^2\mathrm{exp}(\frac{\pi }{a(L)}|\omega |^2+2\pi iE_L(\omega ,x))}{\omega ^2}$$
where the constant $`c`$ is determined from (1.3). Differentiating one more time we get
(1.5)
$$\mathrm{}^{}(x)=\underset{\omega L}{}\frac{(1+(1+\frac{\pi }{a(L)}|\omega +x|^2)^2)\mathrm{exp}(\frac{\pi }{a(L)}|\omega +x|^2)}{(\omega +x)^3}+\frac{\pi ^2}{a(L)^2}\underset{\omega L0}{}\frac{|\omega |^4\mathrm{exp}(\frac{\pi }{a(L)}|\omega |^2+2\pi iE_L(\omega ,x))}{\omega ^3}.$$
### 1.2. Kronecker function
Let us consider the following holomorphic function in $`3`$ variables $`\tau `$, $`x`$, $`y`$, where $`\mathrm{Im}(\tau )>0`$, $`0<\mathrm{Im}(x),\mathrm{Im}(y)<\mathrm{Im}(\tau )`$:
$$F(x,y;\tau )=\underset{(m+\frac{1}{2})(n+\frac{1}{2})>0}{}\mathrm{sign}(m+\frac{1}{2})\mathrm{exp}(2\pi imn\tau +2\pi imx+2\pi iny)$$
where $`m,n`$ are integers (our choice of sign is compatible with the notation in Zagier’s paper , but our variables $`x`$ and $`y`$ differ from those used in by the factor $`2\pi i`$). We call it the Kronecker function since Kronecker discovered (see ) the following remarkable identity:
(1.6)
$$F(x,y;\tau )=\frac{\theta _{11}^{}(0,\tau )}{2\pi i}\frac{\theta _{11}(x+y,\tau )}{\theta _{11}(x,\tau )\theta _{11}(y,\tau )}$$
where
$$\theta _{11}(x,\tau )=\underset{n}{}(1)^n\mathrm{exp}(\pi i(n+\frac{1}{2})^2\tau +2\pi i(n+\frac{1}{2})x),$$
$`\theta _{11}^{}`$ is the derivative of $`\theta _{11}(x,\tau )`$ with respect to $`x`$. In particular, this identity gives a meromorphic continuation of $`F`$ to $`\times ^2`$ with poles along the divisors $`xL_\tau `$, $`yL_\tau `$, where $`L_\tau =+\tau `$.
###### Theorem 2.
One has the following identity
(1.7)
$$\begin{array}{c}2\pi iF(x,y,\tau )=\mathrm{exp}(\frac{\pi }{\mathrm{Im}\tau }x(y\overline{y}))_{\omega L_\tau }\frac{\mathrm{exp}(\frac{\pi }{\mathrm{Im}\tau }|\omega +x|^22\pi iE(\omega ,y))}{\omega +x}+\hfill \\ \mathrm{exp}(\frac{\pi }{\mathrm{Im}\tau }y(x\overline{x}))_{\omega L_\tau }\frac{\mathrm{exp}(\frac{\pi }{\mathrm{Im}\tau }|\omega +y|^22\pi iE(\omega ,x))}{\omega +y}\hfill \end{array}$$
where $`E=E_{L_\tau }`$.
Proof. For a fixed $`\tau `$ let $`f(x,y)`$ be the function in the right hand side of (1.7). First we have to check that $`f(x,y)`$ is meromorphic in $`x`$ and $`y`$. Let us denote $`a=\mathrm{Im}(\tau )`$. We have
$`{\displaystyle \frac{}{\overline{x}}}f={\displaystyle \frac{\pi }{a}}\mathrm{exp}({\displaystyle \frac{\pi }{a}}x(y\overline{y})){\displaystyle \underset{\omega L_\tau }{}}\mathrm{exp}({\displaystyle \frac{\pi }{a}}|\omega +x|^22\pi iE_L(\omega ,y))+`$
$`{\displaystyle \frac{\pi }{a}}\mathrm{exp}({\displaystyle \frac{\pi }{a}}y(x\overline{x})){\displaystyle \underset{\omega L_\tau }{}}\mathrm{exp}({\displaystyle \frac{\pi }{a}}|\omega +y|^22\pi iE_L(\omega ,x)).`$
Thus, we have to prove that
$$\underset{\omega L_\tau }{}\mathrm{exp}(\frac{\pi }{a}|\omega +x|^22\pi iE_L(\omega +x,y))=\underset{\omega L_\tau }{}\mathrm{exp}(\frac{\pi }{a}|\omega +y|^22\pi iE_L(\omega ,x)).$$
But this follows easily from Lemma 1.1 since the function $`\mathrm{exp}(\frac{\pi }{a}|x|^2)`$ is Fourier self-dual.
Next we observe that for any $`\omega L`$ one has
$$f(x+\omega ,y)=\mathrm{exp}(\frac{\pi }{a}(\omega \overline{\omega })y)f(x,y).$$
Since $`f(x,y)=f(y,x)`$ we conclude that $`f`$ has the same quasi-periodicity equations as $`F`$. Hence, $`f/F`$ is periodic with respect to $`L_\tau `$ in both variables. Since both $`f`$ and $`F`$ have poles of the first order at $`xL_\tau `$ and $`yL_\tau `$ the only possible poles of $`f/F`$ can come from zeroes of $`F`$. But the only zeroes of $`F`$ are the zeroes of the first order along the divisor $`x+yL_\tau `$. On the other hand, one immediately checks that $`f(x,x)=0`$. Therefore, $`f/F`$ is holomorphic, so it should be constant. Now the identity follows by the comparison of the residues of $`F`$ and $`f`$ at $`x=0`$. ∎
From identities (1.2) and (1.7) one immediately deduces the following result.
###### Corollary 1.2.
One has
$$\left(2\pi iF(x,y,\tau )\frac{1}{y}\right)|_{y=0}=\zeta (x,L_\tau )x\eta _1$$
where as generators of $`L_\tau `$ we take $`\omega _1=1`$, $`\omega _2=\tau `$.
Another way to express the relation between the function $`F`$ and Weierstrass zeta-function is the following:
$$\pi i(F(x,y,\tau )+F(x,y,\tau ))|_{y=0}=\zeta (x,L_\tau )x\eta _1.$$
This can be deduced either from the above corollary or using (1.6) and the formula
$$\frac{\theta _{11}^{}(x)}{\theta _{11}(x)}=\zeta (x,L_\tau )x\eta _1$$
which can be seen from the decomposition of $`\theta _{11}`$ into an infinite product.
## 2. Explanation
Both the functions $`Z(x,L)`$ and $`F(x,y,\tau )`$ have nice modular properties. The modular forms $`Z(x_1\omega _1+x_2\omega _2,L)`$ for fixed $`x_1,x_2`$ were considered by Hecke in (he called them “Teilwerte” of the Weierstrass zeta-function). The modular equation for $`F(x,y,\tau )`$ can be found in . The formulas (1.2) and (1.7) provide an alternative explanation of modularity but each of the series in the right hand side is non-holomorphic (only the difference of two such series is). In this section we show that this is related to the computation of certain triple products on elliptic curve using non-holomorphic data (namely, hermitian metrics). However, since the result doesn’t depend on a choice of non-holomorphic data the resulting expressions are holomorphic in the modular parameter.
We refer to for the general discussion of higher products on elliptic curve. The triple products related to the two series considered in the previous section are of the following type. Let $`L`$, $`M_1`$, $`M_2`$ be hermitian line bundles of degree $`1`$ on a complex elliptic curve. Then one can consider a triple product
$$m_3:H^0(E,M_1)H^1(E,L^1)H^0(E,M_2)H^0(E,M_1M_2L^1).$$
Recall that it is given by the formula
(2.1)
$$m_3(s_1,e,s_2)=\mathrm{pr}(Q(s_1e)s_2s_1Q(es_2))$$
where $`s_iH^0(E,M_i)`$, $`i=1,2`$ the class $`eH^1(E,L^1)`$ is represented by a harmonic $`(0,1)`$-form, $`Q=\overline{}^{}G_\overline{}`$ where $`G_\overline{}`$ is the Green operator corresponding to the laplacian $`\mathrm{\Delta }_\overline{}=\overline{}\overline{}^{}+\overline{}^{}\overline{}`$ (where $`\overline{}^{}`$ is conjugate to $`\overline{}`$ with respect to the hermitian metric), $`\mathrm{pr}=\mathrm{id}G_\overline{}\mathrm{\Delta }_\overline{}`$ is the harmonic projector. The formula (2.1) shows that in fact our triple product depends only on the operator $`Q`$ acting on forms with values in $`M_1L^1`$ and $`M_2L^1`$. Both these line bundles are of degree zero so (up to switching $`M_1`$ and $`M_2`$) the following three possibilities can occur:
(a) $`M_iL^1\simeq ̸𝒪_E`$ for $`i=1,2`$. Then we have $`Q=(\overline{})^1`$ in the formula (2.1) so this triple product doesn’t depend on metrics. We will show that in this case $`m_3`$ is expressed in terms of the series appearing in the formula (1.7).
(b) $`M_1L^1𝒪_E`$, $`M_2L^1\simeq ̸𝒪_E`$. In this case the operator $`Q`$ depends on a hermitian metric on $`𝒪_E`$. However, there is a natural choice of a constant metric on $`𝒪_E`$ and it is easy to see that $`Q`$ doesn’t change if we rescale a metric by a constant. In this case we’ll express $`m_3`$ in terms of the series from the formula (1.2).
(c) $`M_1L^1M_2L^1𝒪_E`$. One can easily see that in this case $`m_3=0`$.
To compute triple products in the cases (a) and (b) we represent our elliptic curve in the form $`E=/+\tau `$. We choose $`L`$ to be the line bundle on $`E`$ such that the theta-function
$$\theta (z)=\theta (z,\tau )=\underset{n}{}\mathrm{exp}(\pi i\tau n^2+2\pi inz)$$
descends to a section of $`L`$. Thus, the pull-back of $`L`$ to $``$ is canonically trivialized. For $`u`$ let us denote by $`L(u)`$ the line bundle $`t_u^{}L`$ where $`t_u:EE`$ is the translation by $`u`$ (note that a choice of $`u`$ induces a trivialization of the pull-back of $`L(u)`$ to $``$). We define the hermitian metric on $`L(u)`$ by the formula
$$f,g_{L(u)}=_{C/+\tau }f(x)\overline{g(x)}\mathrm{exp}(2\pi a(x_2^2+2x_2u_2)dx_1dx_2$$
where we use real coordinates $`x_1,x_2`$ defined by $`x=x_1+x_2\tau `$, so that $`u=u_1+u_2\tau `$ and we denote $`a=\mathrm{Im}\tau `$. With respect to this metric one has
$$t_u^{}\theta ^2=\frac{1}{\sqrt{2a}}\mathrm{exp}(2\pi au_2^2).$$
As line bundles $`M_1`$ and $`M_2`$ we take $`L(u)`$, $`L(v)`$ for some $`u,v`$, so that we have natural choice of sections $`s_1=t_u^{}\theta `$, $`s_2=t_v^{}\theta `$. As a harmonic $`(0,1)`$-form representing a non-trivial class $`eH^1(E,L^1)`$ we take
$$\alpha =\frac{\pi \sqrt{2}}{\sqrt{a}}\overline{\theta (x)}\mathrm{exp}(2\pi ax_2^2)d\overline{x}.$$
Our computation will be based on the following formula which was proven in (it is equivalent to eq. (2.2.1) of , on the other hand, it can be deduced from Prop. 4.1 of ):
(2.2)
$$\begin{array}{c}\theta (x+y)\overline{\theta (x+z)}\mathrm{exp}(2\pi a(x_2^2+2x_2z_2))=\hfill \\ \frac{1}{\sqrt{2a}}_{m,n}(1)^{mn}\mathrm{exp}(\frac{\pi }{2a}(|m\tau n|^2+2(m\overline{\tau }n)y2(m\tau n)\overline{z}+(y\overline{z})^2))\phi _{yz,m,n}(x)\hfill \end{array}$$
where we denote
$$\phi _{w,m,n}(x)=\mathrm{exp}(2\pi i(mx_1+(nw)x_2)),$$
the summation is over $`(m,n)^2`$. Note that $`(\phi _{w,m,n})_{(m,n)^2}`$ descend to the orthonormal basis of sections on $`L(w)L^1`$. We have
$$\overline{}\phi _{w,m,n}=\frac{\pi }{a}(m\tau n+w)\phi _{w,m,n}d\overline{x}.$$
In the case $`w+\tau `$ this allows us to compute $`Q=(\overline{})^1`$ in terms of coefficients with the basis $`\phi _{w,m,n}d\overline{x}`$. In the case $`w=0`$ the operator $`Q`$ still coincides with $`(\overline{})^1`$ on $`\phi _{0,m,n}d\overline{x}`$ for $`(m,n)(0,0)`$. On the other hand, $`\phi _{0,0,0}=1`$ and we have
$$\overline{}^{}(d\overline{x})=0,$$
hence, $`Q(d\overline{x})=0`$.
Let us first consider the case (a). Then we have
$$m_3(t_u^{}\theta ,\alpha ,t_v^{}\theta )=\mathrm{pr}(h_ut_v^{}\theta h_vt_u^{}\theta )$$
where we denote $`h_w=Q(t_w^{}\theta \alpha )`$. Using formula (2.2) we get that for any $`w+\tau `$ one has
$$h_w=\underset{m,n}{}a_{m,n}(w)\phi _{w,m,n}$$
where
(2.3)
$$a_{m,n}(w)=\frac{(1)^{mn}\mathrm{exp}(\frac{\pi }{2a}(|m\tau n|^2+2(m\overline{\tau }n)w+w^2))}{m\tau n+w}.$$
By definition the above triple product is proportional to $`t_{u+v}^{}\theta `$ so the computation reduces to calculating the coefficient
(2.4)
$$\frac{m_3(t_u^{}\theta ,\alpha ,t_v^{}\theta ),t_{u+v}^{}\theta }{t_{u+v}^{}\theta ^2}=\sqrt{2a}\mathrm{exp}(2\pi a(u_2+v_2)^2)\left(h_ut_v^{}\theta ,t_{u+v}^{}\theta h_vt_u^{}\theta ,t_{u+v}^{}\theta \right).$$
Now by definition we have
$`h_ut_v^{}\theta ,t_{u+v}^{}\theta ={\displaystyle _{/+\tau }}h_u(x)\theta (x+v)\overline{\theta (x+u+v)}\mathrm{exp}(2\pi a(x_2^2+2x_2(u_2+v_2)))𝑑x_1𝑑x_2=`$
$`h_u,t_{u+v}^{}\theta \overline{t_v^{}\theta }\mathrm{exp}(2\pi a(x_2^2+2x_2v_2))`$
where the last scalar product is taken with respect to the metric on $`L(u)L^1`$. Applying (2.2) we get
$$t_{u+v}^{}\theta \overline{t_v^{}\theta }\mathrm{exp}(2\pi a(x_2^2+2x_2v_2))=\underset{m,n}{}b_{m,n}(u,v)\phi _{u,m,n}$$
where
(2.5)
$$b_{m,n}(u,v)=\frac{1}{\sqrt{2a}}(1)^{mn}\mathrm{exp}(\frac{\pi }{2a}(|m\tau n|^2+2(m\overline{\tau }n)(u+v)2(m\tau n)\overline{v}+(u+v\overline{v})^2)).$$
Since $`\phi _{u,m,n}`$ is an orthonormal system we derive
$$h_ut_v^{}\theta ,t_{u+v}^{}\theta =\underset{m,n}{}a_{m,n}(u)\overline{b_{m,n}(u,v)}.$$
Substituting the expressions (2.3) and (2.5) and simplifying we obtain
$$\sqrt{2a}\mathrm{exp}(2\pi a(u_2+v_2)^2)h_ut_v^{}\theta ,t_{u+v}^{}\theta =\mathrm{exp}(\frac{\pi }{a}u(v\overline{v}))\underset{\omega +\tau }{}\frac{\mathrm{exp}(\frac{\pi }{a}|\omega +u|^2+2\pi iE(\omega ,v))}{\omega +u}$$
where $`E`$ is the symplectic form associated with the oriented lattice $`+\tau `$. Substituting this into the equation (2.4) we obtain that the coefficient of $`m_3(t_u^{}\theta ,\alpha ,t_v^{}\theta )`$ with $`t_{u+v}^{}\theta `$ is equal to
$$\mathrm{exp}(\frac{\pi }{a}u(v\overline{v}))\underset{\omega +\tau }{}\frac{\mathrm{exp}(\frac{\pi }{a}|\omega +u|^2+2\pi iE(\omega ,v))}{\omega +u}\mathrm{exp}(\frac{\pi }{a}v(u\overline{u}))\underset{\omega +\tau }{}\frac{\mathrm{exp}(\frac{\pi }{a}|\omega +v|^2+2\pi iE(\omega ,u))}{\omega +v}.$$
According to the identity (1.7) this expression is equal to $`2\pi iF(u,v,\tau )`$. In it was shown that the series defining $`F(u,v,\tau )`$ appears as the corresponding triple product in the Fukaya category of the torus, so we can view the identity (1.7) as a manifestation of the homological mirror symmetry<sup>1</sup><sup>1</sup>1In we worked with the RHS of the identity (1.7) without presenting an explicit series for it..
In the case (b) we set $`u=0`$ and slightly modify the above computation. Namely, we have
$$h_0=Q(\theta \alpha )=\underset{(m,n)(0,0)}{}a_{m,n}(0)\phi _{0,m,n}$$
where $`a_{m,n}(0)`$ are still given by the formula (2.3). The formula (2.4) still holds for $`u=0`$, so we obtain that the coefficient of the triple product in this case is equal to
$$\underset{\omega +\tau ,\omega 0}{}\frac{\mathrm{exp}(\frac{\pi }{a}|\omega |^2+2\pi iE(\omega ,v))}{\omega }\underset{\omega +\tau }{}\frac{\mathrm{exp}(\frac{\pi }{a}|\omega +v|^2)}{\omega +v}.$$
According to the identity (1.2) this is equal to $`Z(v,+\tau )`$. Thus, considering triple products above one can discover the identity (1.2) as follows. One starts with the identity (1.7) which follows from the homological mirror symmetry for elliptic curve. Then one tries to pass to the limit as $`u0`$. From the expression of the Weierstrass zeta-function as the logarithmic derivative of sigma-function it is easy to guess that the limit of $`F(u,v)`$ should be related to $`\zeta (v,+\tau )`$. On the other hand, it should be modular, so one naturally arrives to considering $`Z(v,+\tau )`$.
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# 1 Introduction and summary
## 1 Introduction and summary
The Quicksort algorithm of Hoare is “one of the fastest, the best-known, the most generalized, the most completely analyzed, and the most widely used algorithms for sorting an array of numbers” . Quicksort is the standard sorting procedure in Unix systems, and Philippe Flajolet, a leader in the field of analysis of algorithms, has noted that it is among “some of the most basic algorithms—the ones that do deserve deep investigation” . Our goal in this introductory section is to review briefly some of what is known about the analysis of Quicksort and to summarize how this paper advances that analysis.
The Quicksort algorithm for sorting an array of $`n`$ numbers is extremely simple to describe. If $`n=0`$ or $`n=1`$, there is nothing to do. If $`n2`$, pick a number uniformly at random from the given array. Compare the other numbers to it to partition the remaining numbers into two subarrays. Then recursively invoke Quicksort on each of the two subarrays.
Let $`X_n`$ denote the (random) number of comparisons required (so that $`X_0=0`$). Then $`X_n`$ satisfies the distributional recurrence relation
$$X_n\stackrel{}{=}X_{U_n1}+X_{nU_n}^{}+n1,n1,$$
where $`\stackrel{}{=}`$ denotes equality in law (i.e., in distribution), and where, on the right, $`U_n`$ is distributed uniformly on the set $`\{1,\mathrm{},n\}`$, $`X_j^{}\stackrel{}{=}X_j`$, and
$$U_n;X_0,\mathrm{},X_{n1};X_0^{},\mathrm{},X_{n1}^{}$$
are all independent.
As is well known and quite easily established, for $`n0`$ we have
$$\mu _n:=𝐄X_n=2(n+1)H_n4n2n\mathrm{ln}n,$$
where $`H_n:=_{k=1}^nk^1`$ is the $`n`$th harmonic number and $``$ denotes asymptotic equivalence. It is also routine to compute explicitly the standard deviation of $`X_n`$ (see Exercise 6.2.2-8 in ), which turns out to be $`n\sqrt{7\frac{2}{3}\pi ^2}`$.
Consider the standardized variate
$$Y_n:=(X_n\mu _n)/n,n1.$$
Régnier showed using martingale arguments that $`Y_nY`$ in distribution, with $`Y`$ satisfying the distributional identity
(1.1)
$$Y\stackrel{}{=}UY+(1U)Z+g(U)=:h_{Y,Z}(U),$$
where
(1.2)
$$g(u):=2u\mathrm{ln}u+2(1u)\mathrm{ln}(1u)+1,$$
and where, on the right of $`\stackrel{}{=}`$ in (1.1), $`U`$, $`Y`$, and $`Z`$ are independent, with $`Z\stackrel{}{=}Y`$ and $`U\text{unif}(0,1)`$. Rösler showed that (1.1) characterizes the limiting law $`(Y)`$, in the precise sense that $`F:=(Y)`$ is the *unique* fixed point of the operator
$$G=(V)SG:=(UV+(1U)V^{}+g(U))$$
(in what should now be obvious notation) subject to
$$𝐄V=0,\mathrm{𝐕𝐚𝐫}V<\mathrm{}.$$
Thus it is clear that fundamental (asymptotic) probabilistic understanding of Quicksort’s behavior relies on fundamental understanding of the limiting distribution $`F`$. In this regard, Rösler showed that
(1.3) the moment generating function (mgf) of $`Y`$ is everywhere finite,
and Hennequin and Rösler showed how all the moments of $`Y`$ can be pumped out one at a time, though there is no known expression for the mgf nor for the general $`p`$th moment in terms of $`p`$. Tan and Hadjicostas proved that $`F`$ has a density $`f`$ which is almost everywhere positive, but their proof does not even show whether $`f`$ is continuous.
The main goal of this paper is to prove that $`F`$ has a density $`f`$ which is infinitely differentiable, and that each derivative $`f^{(k)}(y)`$ decays as $`y\pm \mathrm{}`$ more rapidly than any power of $`|y|^1`$: this is our main Theorem 3.1. In particular, it follows that each $`f^{(k)}`$ is bounded (cf. Theorem 3.3).
Our main tool will be Fourier analysis. We begin in Section 2 by showing (see Theorem 2.9) that the characteristic function $`\varphi `$ for $`F`$ has rapidly decaying derivatives of every order. Standard arguments reviewed briefly at the outset of Section 3 then immediately carry this result over from $`\varphi `$ to $`f`$. Finally, in Section 4 we will use the boundedness and continuity of $`f`$ to establish an integral equation for $`f`$ (Theorem 4.1). As a corollary, $`f`$ is everywhere positive (Corollary 4.2).
###### Remark 1.1.
(a) Our method is sufficiently computational that we will prove, for example, that $`f`$ is bounded by $`16`$. This is not sharp numerically, as Figure 4 of strongly suggests that the maximum value of $`f`$ is about $`2/3`$. However, in future work we will rigorously justify (and discuss how to obtain bounds on the error in) the numerical computations used to obtain that figure, and the rather crude bounds on $`f`$ and its derivatives obtained in the present paper are needed as a starting point for that more refined work.
(b) Very little is known rigorously about $`f`$. For example, the figure discussed in (a) indicates that $`f`$ is unimodal. Can this be proved? Is $`f`$ in fact *strongly* unimodal (i.e., log-concave)? What can one say about changes of signs for the derivatives of $`f`$?
(c) Knessl and Szpankowski purport to prove very sharp estimates of the rates of decay of $`f(y)`$ as $`y\mathrm{}`$ and as $`y\mathrm{}`$. Roughly put, they assert that the left tail of $`f`$ decays doubly exponentially (like the tail of an extreme-value density) and that the right tail decays exponentially. But their “results” rely on several unproven assumptions. Among these, for example, is their assumption (59) that
$$𝐄e^{\lambda Y}\mathrm{exp}(\alpha \lambda \mathrm{ln}\lambda +\beta \lambda +\gamma \mathrm{ln}\lambda +\delta )\text{as }\lambda \mathrm{}$$
for some constants $`\alpha (>0),\beta ,\gamma ,\delta `$. \[Having assumed this, they derive the values of $`\alpha `$, $`\gamma `$, and $`\delta `$ exactly, and the value of $`\beta `$ numerically.\]
## 2 Bounds on the limiting Quicksort characteristic function
We will in this section prove the following result on superpolynomial decay of the characteristic function of the limit variable $`Y`$.
###### Theorem 2.1.
For every real $`p0`$ there is a smallest constant $`0<c_p<\mathrm{}`$ such that the characteristic function $`\varphi (t):𝐄e^{itY}`$ satisfies
(2.1)
$$|\varphi (t)|c_p|t|^p\text{ for all }t𝐑\text{.}$$
These best possible constants $`c_p`$ satisfy $`c_0=1`$, $`c_{1/2}2`$, $`c_{3/4}\sqrt{8\pi }`$, $`c_14\pi `$, $`c_{3/2}<187`$, $`c_{5/2}<103215`$, $`c_{7/2}<197102280`$, and the relations
(2.2) $`c_{p_1}^{1/p_1}`$ $`c_{p_2}^{1/p_2},0<p_1p_2;`$
(2.3) $`c_{p+1}`$ $`2^{p+1}c_p^{1+(1/p)}p/(p1),p>1;`$
(2.4) $`c_p`$ $`2^{p^2+6p},p>0.`$
\[The numerical bounds are not sharp (except in the trivial case of $`c_0`$); they are the best that we can get without too much work, but we expect that substantial improvements are possible.\]
###### Proof.
The basic approach is to use the fundamental relation (1.1). We will first show, using a method of van der Corput , that the characteristic function of $`h_{y,z}(U)`$ is bounded by $`2|t|^{1/2}`$ for each $`y,z`$. Mixing, this yields Theorem 2.1 for $`p=1/2`$. Then we will use another consequence of (1.1), namely, the functional equation
(2.5)
$$\varphi (t)=_{u=0}^1\varphi (ut)\varphi ((1u)t)e^{itg(u)}𝑑u,t𝐑,$$
or rather its consequence
(2.6)
$$|\varphi (t)|_{u=0}^1|\varphi (ut)||\varphi ((1u)t)|𝑑u,$$
and obtain successive improvements in the exponent $`p`$.
We give the details as a series of lemmas, beginning with a standard calculus estimate . Note that it suffices to consider $`t>0`$ in the proofs because $`\varphi (t)=\overline{\varphi (t)}`$ and thus $`|\varphi (t)|=|\varphi (t)|`$. Note also that the best constants satisfy $`c_p=sup_{t>0}t^p|\varphi (t)|`$ (although we do not know in advance of proving Theorem 2.1 that these are finite), and thus $`c_p^{1/p}=sup_{t>0}t|\varphi (t)|^{1/p}`$, which clearly satisfies (2.2) because $`|\varphi (t)|1`$.
###### Lemma 2.2.
Suppose that a function $`h`$ is twice continuously differentiable on an open interval $`(a,b)`$ with
$$h^{}(x)c>0\text{ and }h^{\prime \prime }(x)0\text{ for }x(a,b)\text{.}$$
Then
$$\left|_{x=a}^be^{ith(x)}𝑑x\right|\frac{2}{ct}\text{ for all }t>0\text{.}$$
###### Proof.
By considering subintervals $`(a+\epsilon ,b\epsilon )`$ and letting $`\epsilon 0`$, we may without loss of generality assume that $`h`$ is defined and twice differentiable at the endpoints, too. Then, using integration by parts, we calculate
$`{\displaystyle _{x=a}^b}e^{ith(x)}𝑑x`$ $`={\displaystyle \frac{1}{it}}{\displaystyle _{x=a}^b}\left[{\displaystyle \frac{d}{dx}}e^{ith(x)}\right]{\displaystyle \frac{dx}{h^{}(x)}}`$
$`={\displaystyle \frac{1}{it}}\left\{{\displaystyle \frac{e^{ith(x)}}{h^{}(x)}}|_{x=a}^b{\displaystyle _{x=a}^b}e^{ith(x)}d\left({\displaystyle \frac{1}{h^{}(x)}}\right)\right\}.`$
So
$$\begin{array}{cc}\hfill \left|_{x=a}^be^{ith(x)}𝑑x\right|& \frac{1}{t}\left\{\left(\frac{1}{h^{}(b)}+\frac{1}{h^{}(a)}\right)+_{x=a}^b\right|d\left(\frac{1}{h^{}(x)}\right)|dx\}\hfill \\ & =\frac{1}{t}\left\{\left(\frac{1}{h^{}(b)}+\frac{1}{h^{}(a)}\right)+_{x=a}^b\left[d\left(\frac{1}{h^{}(x)}\right)\right]𝑑x\right\}\hfill \\ & =\frac{1}{t}\left\{\left(\frac{1}{h^{}(b)}+\frac{1}{h^{}(a)}\right)+\left(\frac{1}{h^{}(a)}\frac{1}{h^{}(b)}\right)\right\}\hfill \\ & =\frac{2}{th^{}(a)}\frac{2}{ct}.\mathit{}\hfill \end{array}$$
###### Lemma 2.3.
For any real numbers $`y`$ and $`z`$, the random variable $`h_{y,z}(U)`$ defined by (1.1) satisfies
$$|𝐄e^{ith_{y,z}(U)}|2|t|^{1/2}.$$
###### Proof.
We will apply Lemma 2.2, taking $`h`$ to be $`h_{y,z}`$. Observe that
$$h_{y,z}^{\prime \prime }(u)=2\left(\frac{1}{u}+\frac{1}{1u}\right)=\frac{2}{u(1u)}8\text{ for }u(0,1)$$
and that
$$h_{y,z}^{}(u)=0\text{ if and only if }u=\alpha _{y,z}:=\frac{1}{1+\mathrm{exp}\left(\frac{1}{2}(yz)\right)}(0,1).$$
Let $`t>0`$ and $`\gamma >0`$. If in Lemma 2.2 we take $`a:=\alpha _{y,z}+\gamma t^{1/2}`$ and $`b:=1`$, and assume that $`a<b`$, then note
$$h^{}(u)=h_{y,z}^{}(u)=_{x=\alpha _{y,z}}^uh_{y,z}^{\prime \prime }(x)𝑑x8(u\alpha _{y,z})8\gamma t^{1/2}\text{ for all }u(a,b)\text{.}$$
So, by Lemma 2.2,
$$\left|_{u=\alpha _{y,z}+\gamma t^{1/2}}^1e^{ith_{y,z}(u)}𝑑u\right|\frac{2}{t}[8\gamma t^{1/2}]^1=\frac{1}{4\gamma }t^{1/2}.$$
Trivially,
$$\left|_{u=\alpha _{y,z}}^{\alpha _{y,z}+\gamma t^{1/2}}e^{ith_{y,z}(u)}𝑑u\right|\gamma t^{1/2},$$
so we can conclude
$$\left|_{u=\alpha _{y,z}}^1e^{ith_{y,z}(u)}𝑑u\right|[(4\gamma )^1+\gamma ]t^{1/2}.$$
This result is trivially also true when $`a=\alpha _{y,z}+\gamma t^{1/2}b=1,`$ so it holds for all $`t,\gamma >0`$. The optimal choice of $`\gamma `$ here is $`\gamma _{\text{opt}}=1/2`$, which yields
$$\left|_{u=\alpha _{y,z}}^1e^{ith_{y,z}(u)}𝑑u\right|t^{1/2}\text{for all }t>0.$$
Similarly, for example by considering $`uh(1u)`$,
$$\left|_0^{\alpha _{y,z}}e^{ith_{y,z}(u)}𝑑u\right|t^{1/2}\text{for all }t>0,$$
and we conclude that the lemma holds for all $`t>0`$, and thus for all real $`t`$. ∎
###### Lemma 2.4.
For any real $`t`$, $`|\varphi (t)|2|t|^{1/2}`$.
###### Proof.
Lemma 2.3 shows that
$$\left|𝐄\left(e^{ith_{Y,Z}(U)}\right|Y,Z)\right|2|t|^{1/2}$$
and thus
$$|\varphi (t)|=\left|𝐄e^{ith_{Y,Z}(U)}\right|𝐄\left|𝐄\left(e^{ith_{Y,Z}(U)}\right|Y,Z)\right|2|t|^{1/2}.$$
The preceding lemma is the case $`p=1/2`$ of Theorem 2.1. We now improve the exponent.
###### Lemma 2.5.
Let $`0<p<1`$. Then
$$c_{2p}\frac{\left[\mathrm{\Gamma }(1p)\right]^2}{\mathrm{\Gamma }(22p)}c_p^2.$$
###### Proof.
By (2.6) and the definition of $`c_p`$,
$$|\varphi (t)|_{u=0}^1c_p^2|ut|^p|(1u)t|^p𝑑u=c_p^2|t|^{2p}_{u=0}^1u^p(1u)^p𝑑u,$$
and the result follows by evaluating the beta integral. ∎
In particular, recalling $`\mathrm{\Gamma }(1/2)=\sqrt{\pi }`$, Lemmas 2.4 and 2.5 yield
(2.7)
$$|\varphi (t)|\frac{4\pi }{|t|}.$$
This proves (2.1) for $`p=1`$, with $`c_14\pi `$, and thus by (2.2) for every $`p1`$ with $`c_p(4\pi )^p`$; applying Lemma 2.5 again, we obtain the finiteness of $`c_p`$ in (2.1) for all $`p<2`$. Somewhat better numerical bounds are obtained for $`1/2<p<1`$ by taking a geometric average between the cases $`p=1/2`$ and $`p=1`$: the inequality
$$|\varphi (t)|(2t^{1/2})^{22p}(4\pi t^1)^{2p1}=2^{2p}\pi ^{2p1}t^p,t>0,$$
shows that $`c_p2^{2p}\pi ^{2p1}`$, $`1/2p1`$. In particular, we have $`c_{3/4}\sqrt{8\pi }`$, and thus, by Lemma 2.5, $`c_{3/2}8\pi ^{1/2}\left[\mathrm{\Gamma }(1/4)\right]^2<186.4<187`$.
###### Lemma 2.6.
Let $`p>1`$. Then
$$c_{p+1}2^{p+1}c_p^{1+(1/p)}p/(p1).$$
###### Proof.
Assume that $`t2c_p^{1/p}`$. Then, again using (2.6),
$`|\varphi (t)|`$ $`{\displaystyle _{u=0}^1}\mathrm{min}({\displaystyle \frac{c_p}{(ut)^p}},1)\mathrm{min}({\displaystyle \frac{c_p}{[(1u)t]^p}},1)𝑑u`$
$`=2{\displaystyle _{u=0}^{c_p^{1/p}t^1}}{\displaystyle \frac{c_p}{[(1u)t]^p}}𝑑u+{\displaystyle _{u=c_p^{1/p}t^1}^{1c_p^{1/p}t^1}}{\displaystyle \frac{c_p^2}{[u(1u)t^2]^p}}𝑑u`$
$`{\displaystyle \frac{2}{\left[1c_p^{1/p}t^1\right]^p}}{\displaystyle \frac{c_p^{1+(1/p)}}{t^{p+1}}}+2{\displaystyle \frac{c_p^2}{t^{2p}}}{\displaystyle _{u=c_p^{1/p}t^1}^{1/2}}{\displaystyle \frac{du}{[u(1u)]^p}}`$
$`{\displaystyle \frac{2}{(1/2)^p}}c_p^{1+(1/p)}t^{(p+1)}+{\displaystyle \frac{2}{(1/2)^p}}{\displaystyle \frac{c_p^2}{t^{2p}}}{\displaystyle _{u=c_p^{1/p}t^1}^{1/2}}u^p𝑑u`$
$`2^{p+1}\left\{c_p^{1+(1/p)}t^{(p+1)}+{\displaystyle \frac{1}{p1}}c_p^2t^{2p}\left[c_p^{1/p}t^1\right]^{(p1)}\right\}`$
$`=2^{p+1}c_p^{1+(1/p)}{\displaystyle \frac{p}{p1}}t^{(p+1)}.`$
We have derived the desired bound for all $`t2c_p^{1/p}`$. But also, for all $`0<t<2c_p^{1/p}`$, we have
$$2^{p+1}c_p^{1+(1/p)}\frac{p}{p1}t^{(p+1)}\frac{p}{p1}1|\varphi (t)|,$$
so the estimate holds for all $`t>0`$. ∎
Lemma 2.6 completes the proof of finiteness of every $`c_p`$ in (2.1) (by induction), and of the estimate (2.3). The bound for $`c_{3/2}`$ obtained above now shows (using Maple) that $`c_{5/2}<103215`$, which then gives $`c_{7/2}<197102280`$.
We can rewrite (2.3) as
$`c_{p+1}^{1/(p+1)}`$ $`2c_p^{1/p}\left(1+{\displaystyle \frac{1}{p1}}\right)^{1/(p+1)}2c_p^{1/p}\mathrm{exp}\left({\displaystyle \frac{1}{(p1)(p+1)}}\right)`$
$`=2c_p^{1/p}\mathrm{exp}\left({\displaystyle \frac{1}{2(p1)}}{\displaystyle \frac{1}{2(p+1)}}\right).`$
Hence, by induction, if $`p=n+\frac{5}{2}`$ for a nonnegative integer $`n`$, then
$$c_p^{1/p}2^nc_{5/2}^{2/5}e^{(1/3)+(1/5)}=C2^p,$$
where $`C:=2^{5/2}e^{8/15}c_{5/2}^{2/5}<30.6<2^5`$, using the above estimate of $`c_{5/2}`$. Consequently, $`c_p^{1/p}<2^{p+5}`$ when $`p=n+\frac{5}{2}`$. For general $`p>3/2`$ we now use (2.2) with $`p_1=p`$ and $`p_2=p\frac{5}{2}+\frac{5}{2}`$, obtaining $`c_p^{1/p}<2^{p_2+5}<2^{p+6}`$; the case $`p3/2`$ follows from (2.2) and the estimate $`c_{3/2}^{2/3}<33<2^6`$. This completes the proof of (2.4) and hence of Theorem 2.1. ∎
###### Remark 2.7.
We used (1.1) in two different ways. In the first step we conditioned on the values of $`Y`$ and $`Z`$, while in the inductive steps we conditioned on $`U`$.
###### Remark 2.8.
A variety of other bounds are possible. For example, if we begin with the inequality (2.7), use (2.6), and proceed just as in the proof of Lemma 2.6, we can easily derive the following result in the case $`t8\pi `$:
(2.8)
$$|\varphi (t)|\frac{32\pi ^2}{t^2}\left(\mathrm{ln}\left(\frac{t}{4\pi }\right)+2\right)\frac{32\pi ^2\mathrm{ln}t}{t^2}\text{ for all }t1.72\text{.}$$
The result is trivial for $`1.72t<8\pi `$, since then the bounds exceed unity.
Since $`Y`$ has finite moments of all orders \[recall (1.3)\], the characteristic function $`\varphi `$ is infinitely differentiable. Theorem 2.1 implies a rapid decrease of all derivatives, too.
###### Theorem 2.9.
For each real $`p0`$ and integer $`k0`$, there is a constant $`c_{p,k}`$ such that
$$|\varphi ^{(k)}(t)|c_{p,k}|t|^p\text{ for all }t𝐑\text{.}$$
###### Proof.
The case $`k=0`$ is Theorem 2.1, and the case $`p=0`$ follows by $`|\varphi ^{(k)}(t)|𝐄|Y|^k`$. The remaining cases follows from these cases by induction on $`k`$ and the following calculus lemma. ∎
###### Lemma 2.10.
Suppose that $`g`$ is a complex-valued function on $`(0,\mathrm{})`$ and that $`A,B,p>0`$ are such that $`|g(t)|At^p`$ and $`|g^{\prime \prime }(t)|B`$ for all $`t>0`$. Then $`|g^{}(t)|2\sqrt{AB}t^{p/2}`$.
###### Proof.
Fix $`t>0`$ and let $`\theta =\mathrm{arg}(g^{}(t))`$. For $`s>t`$,
$$\mathrm{Re}(e^{i\theta }g^{}(s))\mathrm{Re}(e^{i\theta }g^{}(t))|g^{}(s)g^{}(t)||g^{}(t)|B(st)$$
and thus, integrating from $`t`$ to $`t_1:=t+(|g^{}(t)|/B)`$,
$$\begin{array}{cc}\hfill \mathrm{Re}\left(e^{i\theta }(g(t_1)g(t))\right)& _t^{t_1}\left(|g^{}(t)|B(st)\right)𝑑s\hfill \\ & =(t_1t)|g^{}(t)|\frac{1}{2}B(t_1t)^2=|g^{}(t)|^2/(2B).\hfill \end{array}$$
Consequently,
$$|g^{}(t)|^2/(2B)|g(t)|+|g(t_1)|2At^p,$$
and the result follows. ∎
In other words, the characteristic function $`\varphi `$ belongs to the class $`𝒮`$ of infinitely differentiable functions that, together with all derivatives, decrease more rapidly than any power. (This is the important class of test functions for tempered distributions, introduced by Schwartz ; it is often called the class of *rapidly decreasing* $`C^{\mathrm{}}`$ functions.)
## 3 The limiting Quicksort density $`f`$ and its derivatives
We can now improve the result by Tan and Hadjicostas on existence of a density $`f`$ for $`Y`$. It is an immediate consequence of Theorem 2.1, with $`p=0`$ and $`p=2`$, say, that the characteristic function $`\varphi `$ is integrable over the real line. It is well-known—see, e.g., \[3, Theorem XV.3.3\]—that this implies that $`Y`$ has a bounded continuous density $`f`$ given by the Fourier inversion formula
(3.1)
$$f(x)=\frac{1}{2\pi }_{t=\mathrm{}}^{\mathrm{}}e^{itx}\varphi (t)𝑑t,x𝐑.$$
Moreover, using Theorem 2.1 with $`p=k+2`$, we see that $`t^k\varphi (t)`$ is also integrable for each integer $`k0`$, which by a standard argument (cf. \[3, Section XV.4\]) shows that $`f`$ is infinitely smooth, with a $`k`$th derivative ($`k0`$) given by
(3.2)
$$f^{(k)}(x)=\frac{1}{2\pi }_{t=\mathrm{}}^{\mathrm{}}(it)^ke^{itx}\varphi (t)𝑑t,x𝐑.$$
It follows further that the derivatives are bounded, with
(3.3)
$$\underset{x}{sup}|f^{(k)}(x)|\frac{1}{2\pi }_{t=\mathrm{}}^{\mathrm{}}|t|^k|\varphi (t)|𝑑t\text{(}k0\text{)},$$
and these bounds in turn can be estimated using Theorem 2.1. Moreover, as is well known , \[13, Theorem 7.4\], an extension of this argument shows that the class $`𝒮`$ discussed at the end of Section 2 is preserved by the Fourier transform, and thus Theorem 2.9 implies that $`f𝒮`$:
###### Theorem 3.1.
The Quicksort limiting distribution has an infinitely differentiable density function $`f`$. For each real $`p0`$ and integer $`k0`$, there is a constant $`C_{p,k}`$ such that
$$|f^{(k)}(x)|C_{p,k}|x|^p\text{ for all }x𝐑\text{.}$$
For numerical bounds on $`f`$, we can use (3.3) with $`k=0`$ and Theorem 2.1 for several different $`p`$ (in different intervals); for example, using $`p=0`$, $`1/2`$, $`1`$, $`3/2`$, $`5/2`$, $`7/2`$, and taking $`t_1=4`$, $`t_2=4\pi ^2`$, $`t_3=(187/(4\pi ))^2`$, $`t_4=103215/187`$, $`t_5=197102280/103215`$,
(3.4)
$$\begin{array}{cc}\hfill f(x)& \frac{1}{2\pi }_{t=\mathrm{}}^{\mathrm{}}|\varphi (t)|𝑑t=\frac{1}{\pi }_{t=0}^{\mathrm{}}|\varphi (t)|𝑑t\hfill \\ & \frac{1}{\pi }_{t=0}^{\mathrm{}}\mathrm{min}(1,2t^{1/2},4\pi t^1,187t^{3/2},103215t^{5/2},197102280t^{7/2})𝑑t\hfill \\ & =\frac{1}{\pi }(_{t=0}^{t_1}dt+_{t=t_1}^{t_2}2t^{1/2}dt+_{t=t_2}^{t_3}4\pi t^1dt+_{t=t_3}^{t_4}187t^{3/2}dt\hfill \\ & +_{t=t_4}^{t_5}103215t^{5/2}dt+_{t=t_5}^{\mathrm{}}197102280t^{7/2}dt)\hfill \\ & 18.2.\hfill \end{array}$$
###### Remark 3.2.
We can do somewhat better by using the first bound in (2.8) over the interval $`(103.18,1984)`$ instead of (as above) Theorem 2.1 with $`p=1`$, $`3/2`$, $`5/2`$, $`7/2`$ over $`(103.18,t_3)`$, $`(t_3,t_4)`$, $`(t_4,t_5)`$, $`(t_5,1984)`$, respectively. This gives
$$f(x)<15.3.$$
Similarly, (3.3) with $`k=1`$ and the same estimates of $`|\varphi (t)|`$ as in (3.4) yield
$$|f^{}(x)|\frac{1}{2\pi }_{t=\mathrm{}}^{\mathrm{}}|t||\varphi (t)|𝑑t=\frac{1}{\pi }_{t=0}^{\mathrm{}}t|\varphi (t)|𝑑t<3652.1,$$
which can be reduced to $`2492.1`$ by proceeding as in Remark 3.2. The bound can be further improved to $`2465.9`$ by using also $`p=9/2`$.
Somewhat better bounds are obtained by using more values of $`p`$ in the estimates of the integrals, but the improvements obtained in this way seem to be slight. We summarize the bounds we have obtained.
###### Theorem 3.3.
The limiting Quicksort density function $`f`$ satisfies $`\mathrm{max}_xf(x)<16`$ and $`\mathrm{max}_x|f^{}(x)|<2466`$. ∎
The numerical bounds obtained here are far from sharp; examination of Figure 4 of suggests that $`\mathrm{max}f<1`$ and $`\mathrm{max}|f^{}|<2`$. Our present technique cannot hope to produce a better bound on $`f`$ than $`4/\pi >1.27`$, since neither Lemma 2.3 nor (2.6) can improve on the bound $`|\varphi (t)|1`$ for $`|t|4`$. Further, *no* technique based on (3.3) can hope to do better than the actual value of $`(2\pi )^1_{t=\mathrm{}}^{\mathrm{}}|\varphi (t)|𝑑t`$, which from cursory examination of Figure 6 of appears to be about $`2`$.
## 4 An integral equation for the density $`f`$
Our estimates are readily used to justify rigorously the following functional equation.
###### Theorem 4.1.
The continuous limiting Quicksort density $`f`$ satisfies (pointwise) the integral equation
$$f(x)=_{u=0}^1_{y𝐑}f(y)f\left(\frac{xg(u)(1u)y}{u}\right)\frac{1}{u}𝑑y𝑑u,x𝐑,$$
where $`g()`$ is as in (1.2).
###### Proof.
For each $`u`$ with $`0<u<1`$, the random variable
(4.1)
$$uY+(1u)Z+g(u)$$
\[with notation as in (1.1)\] has the density function
(4.2)
$$f_u(x):=_{z𝐑}f(z)f\left(\frac{xg(u)(1u)z}{u}\right)\frac{1}{u}𝑑z,$$
where the integral converges for each $`x`$ since, using Theorem 3.3, the integrand is bounded by $`f(z)(\mathrm{max}f)/u16f(z)/u`$; dominated convergence using the continuity of $`f`$ and the same bound shows further that $`f_u`$ is continuous.
This argument yields the bound $`f_u(x)16/u`$, and since $`f_u=f_{1u}`$ by symmetry in (4.1), we have $`f_u(x)16/\mathrm{max}(u,1u)32`$. This uniform bound, (1.1), and dominated convergence again imply that $`_0^1f_u(x)𝑑u`$ is a continuous density for $`Y`$, and thus equals $`f(x)`$ for every $`x`$. ∎
It was shown in that $`f`$ is positive almost everywhere; we now can improve this by removing the qualifier “almost.”
###### Corollary 4.2.
The continuous limiting Quicksort density function is everywhere positive.
###### Proof.
We again use the notation (4.2) from the proof of Theorem 4.1. Fix $`x𝐑`$ and $`u(0,1)`$. Since $`f`$ is almost everywhere positive , the integrand in (4.2) is positive almost everywhere. Therefore $`f_u(x)>0`$. Now we integrate over $`u(0,1)`$ to conclude that $`f(x)>0`$. ∎
Alternatively, Corollary 4.2 can be derived directly from Theorem 4.1, without recourse to . Indeed, if $`f(y_0)>0`$ and $`u_0(0,1)`$, set $`x=y_0+g(y_0)`$; then the integrand in the double integral for $`f(x)`$ in Theorem 4.1 is postive for $`(u,y)`$ equal to $`(u_0,y_0)`$, and therefore, by continuity, also in some small neighborhood thereof. It follows that $`f(y_0+g(u_0))>0`$. Since $`u_0`$ is arbitrary and the image of $`(0,1)`$ under $`g`$ is $`((2\mathrm{ln}21),1)`$, an open interval containing the origin, Corollary 4.2 follows readily.
###### Remark 4.3.
In future work, we will use arguments similar to those of this paper, together with other arguments, to show that when one applies the method of successive substitutions to the integral equation in Theorem 4.1, the iterates enjoy exponential-rate uniform convergence to $`f`$. This will settle an issue raised in the third paragraph of Section 3 in .
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# Dynamical growth of the hadron bubbles during the quark-hadron phase transition
## I Introduction
The phenomena of phase transition has attracted many researchers from diverse areas due to many interesting and common features that occur near the transition point. Recently, a considerable amount of attention is being paid to the study of relativistic heavy ion collisions where a phase transition is expected from the normal nuclear matter to a deconfined state of quarks and gluons . The quark gluon plasma (QGP), if formed, would expand hydrodynamically and would cool down until it reaches a critical temperature $`T_c`$ where a phase transition to hadron phase begins. Although the order of such a phase transition remains an unsettled issue, a considerable amount of work has been carried out to understand the dynamics assuming it to be of first order and also assuming that the homogeneous nucleation is applicable . In the ideal Maxwell construction, the temperature of the plasma remains fixed at $`T_c`$ during the phase transition until the hadronization gets completed. However, if the hadronization proceeds through nucleation, it will not begin at $`T=T_c`$ due to the large nucleation barrier. The nucleation of the hadron bubbles can begin only from a supercooled metastable state. If the amount of supercooling is small, the nucleation rate is computed from $`I=A\mathrm{exp}(\mathrm{\Delta }F/T)`$ which gives the probability per unit time per unit volume to nucleate a region of the stable phase (the hadron phase) within the metastable phase (the QGP phase). Here $`\mathrm{\Delta }F`$ is the minimum energy needed to create a critical bubble and the prefactor $`A`$ is the product of statistical and dynamical factors. The statistical factor $`\mathrm{\Omega }_0`$ is a measure of both the available phase space as the system goes over the saddle and of the statistical fluctuations at the saddle relative to the equilibrium states. The dynamical prefactor $`\kappa `$ gives the exponential growth rate of the bubble or droplet sitting on the saddle.
In an earlier work, Langer and Turski derived the dynamical growth rate ($`\kappa `$) of the liquid droplet based on a non-relativistic formalism. Subsequently, $`\kappa `$ was derived both by Turski-Langer and Kawasaki for a liquid-gas phase transition near the critical point, to be
$`\kappa `$ $`=`$ $`{\displaystyle \frac{2\lambda \sigma T}{l^2n_l^2R^3}},`$ (1)
which involves the thermal conductivity $`\lambda `$, the surface free energy $`\sigma `$, the latent heat per molecule $`l`$ and the density of the molecules in the liquid phase $`n_l`$. The interesting physics in this expression is the thermal conductivity which appears as an essential ingredient for the transportation of the latent heat away from the surface region so that the droplet can grow. For a relativistic system like a baryon free quark gluon plasma which has no net conserved charge, the thermal conductivity vanishes. Hence, the above formula obviously can not be applied to such systems. Therefore, Csernai and Kapusta re-derived $`\kappa `$ for a baryon free plasma using earlier formalism of Langer-Turski , but extending their work to the relativistic domain . In the work of Langer and Turski , $`\kappa `$ was derived by solving a set of linearized hydrodynamic equations in the liquid, vapor and the interfacial regions. However, in the relativistic formalism, Csernai-Kapusta mostly concentrated in the interfacial region. Their primary motivation was to know the velocity profile in the surface region which was then used to estimate the energy flow across the surface. Then they used the condition, that the energy flux which is to be transported outwards should be balanced by the viscous heat dissipation as follows
$`\mathrm{\Delta }\omega {\displaystyle \frac{dR}{dt}}=(4\eta /3+\zeta )vdv/dr,`$ (2)
where $`R`$ is the radius of the hadron bubble and $`v(r)`$ is the flow velocity just outside the surface of the bubble. Accordingly, they obtained an expression for $`\kappa `$ given by
$`\kappa `$ $`=`$ $`{\displaystyle \frac{4\sigma (\frac{4}{3}\eta +\zeta )}{(\mathrm{\Delta }\omega )^2R^3}},`$ (3)
where $`\eta `$ and $`\zeta `$ are the shear and bulk viscosity coefficients respectively and $`\mathrm{\Delta }\omega `$ is the difference in the enthalpy densities of the plasma and the hadronic phases, $`\omega =e+p`$. The above expression implies that energy flow $`\omega 𝐯`$ is provided by the viscous effects. There will be no bubble growth in the case of an ideal plasma with zero viscosity. Thus, the viscosity plays the same role as thermal conductivity in case of a relativistic fluid like the quark gluon plasma with zero baryon density. This approach has been extended by Venugopalan and Vischer for the case of baryon rich QGP where both viscous damping and thermal dissipation are significant. On the contrary, Ruggeri and Friedman (RF) argued that the energy flow does not vanish in absence of any heat conduction or viscous damping. Since the change of the energy density $`e`$ in time is given in the low velocity limit by the conservation equation , $`e/t=.(\omega 𝐯)`$ which implies that the energy flow $`\omega 𝐯`$ is always present. Therefore, following a different approach, Ruggeri and Friedman derived an expression for $`\kappa `$ which does not vanish in the absence of viscosity, given by
$`\kappa =\sqrt{{\displaystyle \frac{2\sigma }{R^3}}{\displaystyle \frac{\omega _q}{(\mathrm{\Delta }\omega )^2}}}.`$ (4)
The viscous effects cause only small perturbations to the above equation. This result is in contradiction with the expression given in Eq. (3) according to which the hadron bubble will not grow in the limit of vanishing viscosity. The difference between Csernai-Kapusta (CK) and Ruggeri-Friedman (RF) results are due to the technical differences in the treatment of the pressure gradients and it needs further investigation. Motivated by this, we re-derive $`\kappa `$ using Csernai-Kapusta formalism which is a relativistic generalization of Langer-Turski (LT) procedure . However, unlike Csernai-Kapusta, we solve the linearized hydrodynamic equations in all regions namely the exterior quark region, the interior hadron region as well as the interfacial or the surface region. We found that in the limit of zero viscosity, our prefactor $`\kappa `$ depends only on two scale parameters, the correlation length $`\xi `$ and the critical radius of the hadron bubble $`R`$. We have also obtained the prefactor for a viscous medium where it can be written in a simple way as the sum of a viscous and a non-viscous terms. Interestingly, using certain assumption for the velocity of sound in the medium around the saddle configuration, the viscous and non-viscous components are found to be similar to the results as obtained by Csernai-Kapusta \[Eq. (3)\] and Ruggeri-Friedman \[for zero viscosity, Eq. (4)\] respectively.
####
The paper is organized as follows. We begin with a brief review of the Csernai-Kapusta and Turski-Langer formalism describing the energy-momentum conserving equations of motion in section II. In section III, we solve these equations to derive the dynamical prefactor. Finally, the numerical calculations and the conclusions are presented in sections IV and V respectively.
## II The relativistic hydrodynamics for baryon free plasma
In the case of relativistic hydrodynamics, we consider the energy density $`e(𝐫,t)`$ and the flow velocity $`𝐯(𝐫,t)`$ of the fluid as two independent variables that describe the dynamics of the system. The equations of motion can be obtained from the local conservation laws:
$$_\mu T^{\mu \nu }=_\mu n^\mu =0.$$
(5)
Here $`T^{\mu \nu }`$ is the energy momentum tensor and $`n^\mu `$ represents the baryon four vector. In the presence of viscosity, the energy-momentum tensor $`T^{\mu \nu }`$ and baryon four vector current $`n^\mu `$ can be decomposed into an ideal and a viscous part
$$T^{\mu \nu }=[(e+p)u^\mu u^\nu pg^{\mu \nu }]+\tau ^{\mu \nu },$$
(6)
$$n^\mu =nu^\mu +\nu ^\mu .$$
(7)
Here $`e`$, $`p`$ and $`n`$ are the energy density, pressure and particle number density. The fluid four velocity is given by $`u^\mu =\gamma (1,𝐯)`$ and $`\tau ^{\mu \nu }`$ and $`\nu ^\mu `$ are the dissipative corrections. The form of the dissipative terms $`\tau ^{\mu \nu }`$ and $`\nu ^\mu `$ depend on the definition of what constitutes the local rest frame of the fluid. The four velocity $`u^\mu `$ should be defined in such a way that in a proper frame of any given fluid element, the energy and the number densities are expressible in terms of other thermodynamic quantities by the same formulae, when dissipative processes are not present. It is also necessary to specify whether $`u^\mu `$ is the velocity of energy transport or particle transport. Accordingly, there exist two definitions for the rest frame; one due to Landau and other due to Eckart. In the Landau approach, $`u^\mu `$ is taken as the velocity of energy transport so that energy three flux $`T^{0i}`$ vanishes in a comoving frame . In the Eckart definition, $`u^\mu `$ is taken as the velocity of the particle transport and the particle three current, rather than the energy three flux vanishes in the fluid rest frame . So in the Eckart definition of rest frame, the particle four vector can be written as $`n^\mu =(n,\mathrm{𝟎})`$, whereas in the Landau definition of rest frame $`n^\mu =(n,\nu )`$. Therefore, the two frames are related by a Lorentz transformation with a boost velocity $`\nu /n`$. It is found that due to ill defined boost velocity , the energy three flux in the Eckart frame (which involves heat conductivity $`\lambda `$) is not well defined as $`\lambda `$ diverges in the limit of chemical potential $`\mu `$0. On the other hand, in the Landau definition heat conduction enters as a correction to baryon flux. It was shown that inspite of the divergence of $`\lambda `$, the correction to the baryon flux $`\nu ^\mu `$ is finite . Therefore, we will use the Landau definition for the subsequent study and also we will assume a baryon free plasma for simplicity. We can now write the equations of motion from the conservation law $`_\mu T^{\mu \nu }=0`$ using Landau definition
$$_te=.(\omega 𝐯)+\mathrm{O}(𝐯^\mathrm{𝟐}),$$
(8)
$`_t(\omega 𝐯)`$ $`=`$ $`.(\omega 𝐯𝐯)p`$ (10)
$`+(({\displaystyle \frac{4}{3}}\eta +\zeta ).𝐯)+\mathrm{O}(𝐯^\mathrm{𝟐}).`$
Here $`\omega =(e+p)`$ and $`\eta `$ and $`\zeta `$ are the shear and bulk viscosity coefficients respectively. We have also assumed the low speed limit where $`\gamma 1`$. Although, the fluid velocity is small, the velocity of individual particles is large. Thus, the expressions for the energy density, pressure etc. are taken to be same as that in the relativistic case.
Following Ref. , Eq. (10) can also be written in terms of the Helmholtz free energy $`f(e)`$ and the usual gradient energy $`\frac{1}{2}K(e)^2`$ as
$`_t(\omega 𝐯)`$ $`=`$ $`.(\omega 𝐯𝐯)p^{}`$ (12)
$`+(({\displaystyle \frac{4}{3}}\eta +\zeta ).𝐯)+\mathrm{O}(𝐯^\mathrm{𝟐}),`$
where
$$p^{}=K(^2e)e+\frac{f}{e}e.$$
(13)
The constant $`K`$ is related to the surface tension $`\sigma `$ as
$`\sigma =K{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑r\left({\displaystyle \frac{de}{dr}}\right)^2.`$ (14)
It can be noted by comparing Eq. (12) with Eq. (10) that $`p^{}`$ is not simply a pressure but a combination of $`f`$ and a force term $`K(^2e)e`$ which is related to the surface tension given by Eq. (14). The pressure inside the interface differs from that outside so there is necessarily a pressure gradient at the interface. The term given by Eq. (13) is needed to balance the differential pressure otherwise the Euler or the Navier Stokes equation would require a changing fluid velocity even in a stationary configuration, which is unphysical.
## III Solution of the relativistic hydrodynamic equations
The above hydrodynamic equations \[Eq. (10) or Eq. (12)\] can be solved after linearizing around the saddle configuration. The saddle point corresponds to the stationary solution when $`e(𝐫,t)=\overline{e}(r)`$ and $`𝐯(𝐫,t)=0`$ and also $`\overline{e}`$ satisfies,
$$K(^2\overline{e})+\frac{f}{\overline{e}}=0.$$
(15)
We can now write the equations of motion for small deviations about the stationary configuration by defining $`e=\overline{e}(r)+\nu (𝐫,t)`$ where $`\nu `$ is a small fluctuation in energy density and $`𝐯=0+𝐯(𝐫,t)`$ and linearizing Eqs. (8) and (12) around this configuration
$`_t\nu (𝐫,t)`$ $`=`$ $`.(\omega 𝐯(𝐫,t)),`$ (16)
$`_t(\overline{\omega }𝐯(𝐫,t))`$ $`=`$ $`\overline{e}\left(K^2+f^{\prime \prime }\right)\nu (𝐫,t)`$ (18)
$`+(({\displaystyle \frac{4}{3}}\eta +\zeta ).𝐯(𝐫,t)).`$
Here $`f^{\prime \prime }=^2f/e^2`$, evaluated around the stationary configuration. The dynamical prefactor $`\kappa `$ is determined with the radial perturbations of the form
$`\nu (𝐫,t)=\nu (r)e^{\kappa t},`$ (19)
$`𝐯(𝐫,t)=𝐯(r)e^{\kappa t}.`$ (20)
It can be seen from Eqs. (16) and (18) that the radial deviations are governed by the equations of motion of the following form
$`\kappa \nu (r)`$ $`=`$ $`.\left(\overline{\omega }𝐯(r)\right),`$ (21)
$`\kappa \overline{\omega }𝐯(r)`$ $`=`$ $`\overline{e}\left(K^2+f^{\prime \prime }\right)\nu (r)`$ (23)
$`+(({\displaystyle \frac{4}{3}}\eta +\zeta ).𝐯(r)).`$
Following Langer and Turski, we find the solution for $`\nu (r)`$ in each of the three regions (i) the interior region of hadron phase, $`rR\xi `$ (ii) the exterior region of QGP phase, $`rR+\xi `$ and (iii) the interface region, $`R\xi rR+\xi `$, where $`R`$ is the radius of the hadron bubble with origin at $`r=0`$. The interfacial region has a thickness of the order of the correlation length $`\xi `$. Further, it is assumed that, everywhere outside the droplet, the energy density $`\overline{e}(r)`$ has the value $`e_q`$, the quark density. Within the droplet, $`\overline{e}(r)`$ is equal to the hadron density $`e_h`$. Thus, $`\overline{e}(r)`$ describes a smooth interfacial profile at $`r=R`$ going from $`e_h`$ to $`e_q`$ within a region of roughly of the order of the correlation length $`\xi `$. We then evaluate the relative amplitudes in the above three regions by matching the values at the boundaries. Finally, we evaluate $`\kappa `$ applying the condition
$`{\displaystyle _0^{\mathrm{}}}r^2\nu (r)𝑑r=0.`$ (24)
which is the conservation law implied by Eq. (16)
Before we proceed further, it may be noted here that the above set of linear equations \[Eqs. (21)\] is obtained from Eq. (12) which contains $`f(e)`$ and $`K`$ explicitly. As will be shown subsequently, this form is suitable for the interfacial region where $`\overline{e}`$ is nonzero. Such an equation has been used in Ref. to evaluate the velocity profile at the surface region in order to evaluate the energy loss due to dissipation. In another approach, Ruggeri-Friedman linearize Eq. (10) only in the exterior region with the assumption $`p=c_s^2e`$ where $`c_s`$ is the velocity of sound in the medium. We will discuss about the validity of the above assumption particularly near the stationary configuration as we proceed further. However, the advantage of using Eq.(10) along with the above relation is that one of the variable (say $`p`$) can be eliminated so that the linear equation becomes a simple wave equation. Therefore, we adopt both the approaches as in the following. We solve the relativistic hydrodynamic equations by linearising Eqs. (8) and (10) in the interior and exterior regions whereas we use the linear Eq. (21) in the interfacial region.
### A The interior and exterior region
In these regions, $`\overline{e}(r)`$ is varying so slowly that the gradient energy can be ignored so that Eq. (13) is consistent with
$$e\frac{f}{e}=fp$$
(25)
Since $`f/\overline{e}=0`$ at the saddle point \[see Eq. (15)\], the above relation would imply $`\overline{p}`$ is independent of the energy density at the stationary configuration. Like energy density $`\nu `$, we also consider a small fluctuation in pressure so that we have $`p(r)=\overline{p}+\beta (r,t)`$ (recall, $`e(r)=\overline{e}(r)+\nu (r,t)`$). We assume that the corresponding fluctuations satisfy the relation $`\beta =c_s^2\nu `$ where $`c_s^2`$ is a constant ($`c_s`$ could be the velocity of sound in the medium around the saddle configuration). Therefore, in the interior and exterior regions, we solve the equation
$$\kappa ^2\nu (r)=c_s^2^2\nu (r)+\frac{\kappa }{\overline{\omega }}\left(\frac{4}{3}\eta +\zeta \right)^2\nu (r),$$
(26)
obtained from Eqs. (8) and (10) after linearizing around the stationary configuration and also using the relation $`^2\beta =c_s^2^2\nu `$. Such a relation has also been used in ref, but we differ in our interpretation of the pressure gradient.
Assuming spherically symmetric solutions of the form
$`\nu (r)={\displaystyle \frac{\mathrm{Constant}}{r}}e^{\pm qr},`$ (27)
we get the relation
$$\overline{\kappa }^2=c_s^2q^2,$$
(28)
where
$$\overline{\kappa }^2=\frac{\kappa ^2}{\left(1+\frac{\kappa }{\overline{\omega }c_s^2}\left(\frac{4}{3}\eta +\zeta \right)\right)}.$$
(29)
The above relation holds both for QGP and the hadron regions, except for the fact that the viscosity coefficients are different in two phases. The interior and exterior solutions, therefore, are
$`\nu (r)`$ $`=`$ $`{\displaystyle \frac{A}{r}}\mathrm{sinh}(qr)\mathrm{for}0rR\xi `$ (30)
and
$`\nu (r)`$ $`=`$ $`{\displaystyle \frac{B}{r}}e^{q(rR)}\mathrm{for}rR+\xi .`$ (31)
If $`\kappa `$ is small, the solution will be the one in which $`\nu (r)`$ varies slowly over a distance of the order of correlation length $`\xi `$ so that $`q\xi <<1`$. Since $`\kappa `$ is related to $`q`$, next we proceed to estimate it by solving the linear hydrodynamic equation in the interfacial region and matching it at the boundary.
### B Interfacial region
As will be shown subsequently, the velocity varies as $`vr^2`$ in this region so that $`r^2v`$ remains constant. As a consequence, $`.𝐯=r^2d(r^2v)/dr`$ vanishes at the surface region. Therefore, ignoring the viscous term and eliminating $`𝐯`$ from Eqs. (21), an equation for $`\nu (r)`$ is obtained as
$`\kappa ^2\nu (r)`$ $`=`$ $`.\left[\overline{e}\left(K^2+f^{\prime \prime }\right)\nu (r)\right].`$ (32)
Further, $`\kappa `$ is assumed to be small so that in the first approximation, we can completely neglect the terms containing $`\kappa ^2`$. Thus, to a good approximation in the interfacial region, $`\nu (r)`$ satisfies
$`.\left[{\displaystyle \frac{d\overline{e}}{dr}}\left(K^2+f^{\prime \prime }\right)\nu (r)\right]=0.`$ (33)
We follow the procedure described in Ref. to get the solution of the above equation in the interfacial region (see appendix A for detail),
$`\nu (r)`$ $``$ $`{\displaystyle \frac{a(R)R^2\mathrm{\Delta }\overline{e}}{2\sigma r}}{\displaystyle \frac{d\overline{e}}{dr}},`$ (34)
where $`\mathrm{\Delta }\overline{e}=e_qe_h`$. This solution is quite similar to that found in Ref. , with $`\overline{n}`$ replaced by $`\overline{e}`$ and $`a`$ which is now a function of $`r`$ evaluated at $`R`$. Note that in the quark and the hadron regions (far away from the surface region), $`\overline{e}(r)`$ is nearly constant. Thus, Eq. (33) becomes undefined in theses regions. Therefore, a different set of equations has been used in the exterior-interior regions as discussed in the previous sections.
We can also get an expression for velocity $`v(r)`$ from the relation,
$$\kappa \nu (r)=\frac{1}{r^2}\frac{d}{dr}[r^2\overline{\omega }v(r)],$$
(35)
Substituting $`\nu (r)`$ from Eq. (34), we get
$$v(r)=\frac{D}{r^2\overline{\omega }}_0^rr𝑑r\frac{d\overline{e}}{dr},$$
(36)
where $`D`$ is a constant. The above equation can be integrated to give
$$v(r)\frac{D}{\overline{\omega }_q}\frac{R}{r^2}.$$
(37)
Recall that this result is consistent with our assumption that $`r^2v`$ is constant in the surface region.
### C Dynamical Prefactor
For the interfacial region, the solution is given in the appendix A. It remains now only to apply Eq. (24) to compute $`q`$ (or $`\kappa `$). As in , we can neglect the contribution coming from the interior region $`(r<R)`$ and the terms of order $`qR\sqrt{\xi /R}`$ in the exterior region. The contribution coming from the interfacial region is $`aR^3(\mathrm{\Delta }\overline{e})^2/2\sigma `$ where $`a`$ is a function of $`r`$ related to the constant $`B`$ and the second derivative of $`f`$ w.r.t. $`\overline{e}_q`$. The exterior region contribution is $`B/q^2`$. Combining both the terms, we get
$`q=\sqrt{{\displaystyle \frac{2\sigma B}{a(R)R^3(\mathrm{\Delta }\overline{e})^2}}}.`$ (38)
Assuming $`a(R+\xi )=a(R\xi )a(R)`$ and using the relation $`(^2f/\overline{e}_q^2)^1`$ for $`B/a`$ \[see Eq. (A7)\], we obtain
$`q=\sqrt{{\displaystyle \frac{2\sigma }{R^3(\mathrm{\Delta }\overline{e})^2}}{\displaystyle \frac{1}{(^2f/\overline{e}_{q}^{}{}_{}{}^{2})}}}.`$ (39)
We can also eliminate $`^2f/\overline{e}_q^2`$ by using the relation
$`{\displaystyle \frac{1}{K}}{\displaystyle \frac{^2f}{\overline{e}_q^2}}={\displaystyle \frac{1}{\xi _q^2}},`$ (40)
where $`\xi _q`$ is the correlation length and $`K`$ is related to the surface tension $`\sigma `$ given by Eq. (14). The choice of $`K`$ depends on the energy density profile $`\overline{e}(r)`$. Following , $`\sigma `$ can be related to $`K`$ in the planar interface approximation at $`T_c`$ as
$`\sigma =K(\mathrm{\Delta }\overline{e})^2/6\xi _q,`$ (41)
which will result in
$`q=\sqrt{{\displaystyle \frac{\xi _q}{3R^3}}}.`$ (42)
Therefore, in the case of a non-viscous plasma, we get a very simple relation for $`\kappa `$ given by
$`\kappa =c_s\sqrt{{\displaystyle \frac{\xi _q}{3R^3}}}=\xi _q^1\sqrt{{\displaystyle \frac{c_s^2x^3}{3}}}=\xi _q^1f(x),`$ (43)
where $`x=\xi _q/R`$. This can be viewed as the critical behavior of $`\kappa `$ that scales as $`\xi _q^1f(x)`$. However, this scaling law is different from the dynamical scaling law that one finds in the case of a non-relativistic liquid-vapor transition where $`\kappa `$ scales as $`\xi ^0R^3`$. While this needs further investigation, one of the reason for this discrepancy could be unlike the static scaling, the dynamical scaling depends on the dynamical behavior of the system which is definitely different depending on whether the medium is relativistic or non-relativistic. The above result is also valid in the case of the viscous plasma, only instead of $`\kappa `$, $`\overline{\kappa }`$ will scale as $`\xi _q^1f(x)`$. Therefore, for viscous quark gluon plasma, this scaling results in a quadratic equation in $`\kappa `$ with solution given by
$`\kappa ={\displaystyle \frac{\alpha q^2}{2}}+c_sq\sqrt{1+{\displaystyle \frac{\alpha ^2q^2}{4c_s^2}}},`$ (44)
where $`\alpha =(4\eta /3+\zeta )/\overline{\omega }`$. Since in the first approximation $`q\kappa `$, we can neglect the second term under the square root which are higher order in $`q^2`$ and $`\alpha ^2`$ (viscosity). Finally, we get
$`\kappa ={\displaystyle \frac{q^2}{2\overline{\omega }}}\left({\displaystyle \frac{4}{3}}\eta +\xi \right)+c_sq.`$ (45)
Using Eq. (42) for $`q`$, we can obtain a general expression for $`\kappa `$ for a viscous QGP as
$`\kappa =c_s\sqrt{{\displaystyle \frac{\xi _q}{3R^3}}}+{\displaystyle \frac{\xi _q}{6R^3}}{\displaystyle \frac{1}{\overline{\omega }_q}}\left({\displaystyle \frac{4}{3}}\eta _q+\zeta _q\right).`$ (46)
Therefore, the prefactor $`\kappa `$ can be written as the sum of two terms having a non-viscous ($`\kappa _0`$) and a viscous ($`\kappa _v`$) component. However, both $`\kappa _0`$ and $`\kappa _v`$ have simple dependence on the correlation length $`\xi _q`$ and the bubble radius $`R`$. As can be seen from Eq. (46), the first term is more dominating as compared to the second one particularly when $`T`$ is close to $`T_c`$. However, as temperature decreases, the viscous contribution competes with that of the non-viscous one.
We can also express the above equation (46) in a different way by assuming $`c_s^2`$ as
$$c_s^2=\overline{\omega }_q\frac{^2f}{\overline{e}_q^2},$$
(47)
The above relation is analogous to the non-relativistic expression for velocity of sound in the medium which has a similar relation with $`\omega `$ and $`e`$ replaced by $`n`$ (density) . Then from Eq. (39) we get
$`q=c_s^1\sqrt{{\displaystyle \frac{2\sigma }{R^3}}{\displaystyle \frac{\overline{\omega }_q}{(\mathrm{\Delta }\overline{\omega })^2}}},`$ (48)
where we have used the approximation $`\mathrm{\Delta }\overline{\omega }\mathrm{\Delta }\overline{e}`$ since the pressure difference is negligible as compared to the difference in energy density. Now using the above $`q`$ in Eq. (45), the prefactor $`\kappa `$ can be written as
$`\kappa =\sqrt{{\displaystyle \frac{2\sigma }{R^3}}{\displaystyle \frac{\overline{\omega }_q}{(\mathrm{\Delta }\overline{\omega })^2}}}+{\displaystyle \frac{1}{c_s^2}}{\displaystyle \frac{\sigma }{R^3(\mathrm{\Delta }\overline{\omega })^2}}\left({\displaystyle \frac{4}{3}}\eta _q+\zeta _q\right).`$ (49)
As can be seen, the first term in the above equation is same as Eq. (4) as obtained by Ruggeri and Friedman corresponding to the case of a non-viscous plasma. The second term is similar to the result obtained by Csernai and Kapusta except with a minor difference, i.e., instead of 4, we have a factor of $`c_s^2`$ in the numerator \[see Eq. (3)\]. However this is a small difference which can be removed by redefining $`K`$ \[see Eq. (41)\].
It may be mentioned here that the relation $`\beta =c_s^2\nu `$ has been used to obtain $`\kappa `$ as given by Eq. (46). This is the main result of our work. It is also satisfying to note that we can recover the result of Csernai-Kapusta and Ruggeri-Friedman under the assumption for $`c_s^2`$ given by Eq. (47). In analogy with the non-relativistic case, we may interpret $`c_s`$ as the velocity of sound around the saddle configuration ( Recall that saddle point is the configuration where Eq. (15) is satisfied). Although, the above interpretation needs further justification, it is sufficient to say that the results of Eq. (49) can be recovered using $`c_s^2`$ as given by Eq. (47).
### D Result and discussion
We should point out here that there are several reasonable assumptions that have been made in our derivation of the dynamical prefactor. As discussed in the text, most of these approximations are same as that of used in the original work of Langer-Turski since we use the same procedure except that the equations follow relativistic hydrodynamics. An important aspect where we differ from both Ref and Ref is the use of the relation $`\beta =c_s^2\nu `$ assumed to be valid in the quark and hadron regions, ($`\beta `$ and $`\nu `$ are the radial deviations of the pressure and energy density from the stationary solution ). Although, we differ in our interpretation, such a relation has also been used by Ruggeri and Friedman to eliminate one of the hydrodynamic variables. However, we do not make any such assumption in the interfacial region. As a result, the linearized equation used in the interfacial region is different from the one used in the exterior-interior regions. Within above formalism, we have derived an expression for the dynamical prefactor $`\kappa `$. Two important aspects of our result are (a) the prefactor $`\kappa `$ can be written as a linear sum of a non-viscous $`(\kappa _0)`$ and a viscous ($`\kappa _v`$) component and (b) the non-viscous component $`(\kappa _0)`$ which depends on two parameters $`R`$ and $`\xi `$ is finite in the limit of zero viscosity. The present result on $`\kappa _0`$ is also in agreement with the view point of Ruggeri-Friedman that the viscosity is not essential for the dynamical growth of the hadron bubbles. This fact is also evident from Eq. (8) which implies a non-vanishing energy flow $`\omega 𝐯`$ even in the absence of viscosity. Only terms second order in $`𝐯`$ appear in the energy equation in the presence of viscosity and the momentum equation contains a term linear in $`𝐯`$. This means that viscosity terms are relatively unimportant in the energy transport for small value of $`𝐯`$. The momentum equation, however, indicates that viscosity influences the time evolution of $`𝐯`$. Thus, viscosity can serve to disrupt the energy flow and generate entropy but cannot be the only mechanism for energy removal.
The above aspect apparently is in contradiction to the general expectation that transport coefficients like viscosity and thermal conductivity are essential for the removal of the latent heat and hence for the growth of the hadron bubbles . This point needs further clarification. The formation of the hadron bubbles can be interpreted as the thermal fluctuation of the new phase within a correlated volume of radius $`R`$ and surface thickness $`\xi `$. According to the fluctuation-dissipation theory, this fluctuation would mean certain amount of heat dissipation. We can understand the origin of this heat dissipation as follows. In the absence of thermal conductivity, the dissipative losses that occur in a fluid are due to the coefficients of shear and bulk viscosity that depend on the gradient and divergence of the velocity field respectively. In case of an incompressible fluid, the losses are due to the shear stress alone since the bulk viscosity that provides resistance to the expansion (or contraction) does not exist. However, due to the nucleation of the critical size hadron bubble, the pressure or the tension in the fluid is no longer uniform, the pressure inside the hadron bubble being more than the outside. Due to this pressure difference, the hadron bubbles will keep expanding with a non-zero wall velocity. Thus, the fluid medium outside will exert a frictional force on the bubble wall (causing heat dissipation) whose magnitude depends on the pressure difference between the two phases . In the field theoretical language, this dissipation corresponds to the coupling of the order parameter $`\varphi `$ to the fluid which acts as heat bath. Estimates for it in the context of electroweak theory have been given in Refs. . Therefore, our non-viscous part of the prefactor corresponds to a dissipation of dynamical nature which does not depend on any transport coefficients like viscosity or thermal conductivity. This dissipation basically arises due to non-uniform pressure across the interface. Following a different approach, Ignatius had also derived $`\kappa `$ in the limit of zero fluid velocity to be $`2/(\eta R_c^2)`$ where $`\eta `$ is a phenomenological friction parameter (not to be confused with the shear viscosity of the plasma) responsible for the energy transportation between the order parameter and the fluid. Recently, Alamoudi et al have also studied the dynamical viscosity and the growth rate of the nucleating bubble where the viscosity effects arise due to the interaction of the unstable co-ordinate with the stable fluctuations. They estimate a growth rate which depends on $`R`$, $`\xi `$ and the self coupling $`\lambda `$ $`(\kappa \frac{\sqrt{2}}{R}[1.003\lambda T\xi (\frac{R}{\xi })^2])`$. In the limit of weak coupling, the above growth rate scales as $`R^1`$ which is also consistent with our result .
Finally, we conclude this section with the comment that in case of relativistic heavy ion collisions, appreciable amount of nucleation begins from a super cooled metastable QGP phase at which the radius of the critical hadron bubble is of the same order as the width of the bubble interface. At such point, the homogeneous nucleation theory may break down. However, the system moves out of this problematic region quickly due to the release of latent heat which heats up the medium again towards $`T_c`$. In the other application, in case of cosmology, this problem is not serious, although there the actual value of the dynamical growth rate may be of less important. In either case, this study has significance for homogeneous nucleation under the thin-walled bubble approximation.
## IV Numerical results
In the following, we compare $`\kappa `$ obtained from different methods. In the case of a second order phase transition, the correlation length $`\xi `$ scales in the proximity of the critical point as $`\xi (T)=\xi (0)(1T_0/T)^\nu `$ where $`\nu =0.63`$ . However, in the case of a first order phase transition, the transition temperature $`T_0`$ is smaller than $`T_c`$ and approaches $`T_c`$ only in the limit when strength of the transition becomes weak. Therefore, unlike the second order case, $`\xi _q`$ at $`T_c`$ will be finite and which, in the present context, represents the thickness of the interfacial region such that $`R>>\xi _q`$. Further, we ignore the temperature dependence of $`\sigma `$ and $`\xi _q`$ and treat them as constant parameters. This assumption can be justified when the amount of supercooling is small and the medium returns to $`T_c`$ due to the release of latent heat .
Figure 1 shows the temperature dependence of $`\kappa `$ given by Eq. (46) along with viscous ($`\kappa _v`$) and non-viscous ($`\kappa _0`$) components at two different values of $`\sigma `$. Following , we take $`\eta _q`$ as $`2.5T^3`$ and set $`\zeta _q`$ to zero. With decreasing temperature as well as with decreasing $`\sigma `$, the value of the critical radius \[which is obtained from Laplace formulae, see Eq. (52)\] decreases. Therefore, the $`\kappa `$, $`\kappa _0`$ and $`\kappa _v`$ increase with decreasing temperature and also they have higher values for smaller $`\sigma `$, as expected. The behavior of $`\kappa _v`$ is quite different from that of $`\kappa _0`$. Initially, near $`TT_c`$, the $`\kappa _v`$ has small value, but it exceeds $`\kappa _0`$ as temperature comes down particularly at smaller $`\sigma `$ values. Figure 2 shows the similar plot as that of Figure 1 where we have used Eq. (49) to estimate $`\kappa `$. As seen from the figures, both the estimates have similar behavior although Eq.(49) yields slightly higher values for $`\kappa `$ as compared to Eq.(46). The above studies also suggest that the effect of viscosity is negligible at higher $`\sigma `$ values and also for small amount of supercooling. However, its effect can not be ignored at much lower temperature particularly when $`\sigma `$ is small. In figure 3. we have also compared only the non-viscous part ($`\kappa _0`$) of the prefactor as obtained from Eqs.(46) and (49) at two different values of $`\xi _q`$. Within the present set of parameters, the non-viscous parts of the prefactor obtained by both the methods behave similar way.
We also study the dynamics of the nucleation and super cooling by computing the nucleation rate as
$`I={\displaystyle \frac{\kappa }{2\pi }}{\displaystyle \frac{\mathrm{\Omega }_0}{V}}e^{F_\mathrm{C}/T},`$ (50)
where $`F_\mathrm{C}`$ is free energy needed to form a critical bubble in the metastable (supercooled) background. The dynamical prefactor $`\kappa `$ is estimated using Eq. (46) whereas the the statistical prefactor $`\mathrm{\Omega }_0`$ is taken from the previous works as
$`{\displaystyle \frac{\mathrm{\Omega }_0}{V}}={\displaystyle \frac{2}{3}}\left({\displaystyle \frac{\sigma }{3T}}\right)^{3/2}\left({\displaystyle \frac{R}{\xi _q}}\right)^4,`$ (51)
where $`R`$ is the radius of the critical bubble. Under thin-wall approximation $`F_C`$ and $`R`$ for a spherical bubble are given by
$`F_C={\displaystyle \frac{4\pi }{3}}\sigma R^2,R={\displaystyle \frac{2\sigma }{p_hp_q}}.`$ (52)
From the nucleation rate $`I(T)`$, the fraction of space which has been converted to hadron phase can be calculated. If the system cools to $`T_c`$ at a proper time $`\tau _c`$, then at some later time $`\tau `$ the fraction $`h`$ of space which has been converted to hadronic gas is
$`h(\tau )={\displaystyle _{\tau _c}^\tau }𝑑\tau ^{}I(T(\tau ^{}))[1h(\tau ^{})]V(\tau ^{},\tau ).`$ (53)
Here $`V(\tau ^{},\tau )`$ is the volume of a bubble at time $`\tau `$ which had been nucleated at an earlier time $`\tau ^{}`$; this takes into account the bubble growth. The factor $`\left[1h(\tau ^{})\right]`$ is the available space for new bubbles to nucleate. The model for bubble growth is simply taken as
$`V(\tau ^{},\tau )={\displaystyle \frac{4\pi }{3}}\left(R(T(\tau ^{}))+{\displaystyle _\tau ^{}^\tau }𝑑\tau ^{\prime \prime }v(T(\tau ^{\prime \prime }))\right)^3,`$ (54)
where $`v(T)=3[1T/T_c]^{3/2}`$ is the velocity of the bubble growth at temperature $`T`$ . (Recall that this velocity is different from the velocity of the nucleated bubble surface as used in the previous sections). The evolution of the energy momentum in 1+1 dimension is given by
$`{\displaystyle \frac{de}{d\tau }}+{\displaystyle \frac{\omega }{\tau }}={\displaystyle \frac{\frac{4}{3}\eta +\zeta }{\tau ^2}}.`$ (55)
In this work, we use the bag equation of state for QGP. The energy and enthalpy densities in pure QGP and hadron phases are taken as
$`e_q(T)=3a_qT^4+B,\omega _q(T)=4a_qT^4,`$ (56)
$`e_h(T)=3a_hT^4,\omega _h(T)=4a_hT^4.`$ (57)
Here, $`a_q`$ and $`a_h`$ are related to the degrees of freedom operating in two phases and $`B`$ is the bag pressure. The quark phase is assumed to consist of massless gas of $`u`$, $`d`$ quarks and gluon while the hadron phase contains massless pions. Thus the coefficients $`a_q=37\pi ^2/90`$ and $`a_h=3\pi ^2/90`$. In the transition region, the energy density at a time $`\tau `$ can be written in terms of hadronic fraction $`h(\tau )`$ as
$`e(\tau )`$ $`=`$ $`e_q(T)+(e_h(T)e_q(T))h(\tau ).`$ (58)
Following , the viscosity coefficients for the QGP and the hadron phases are chosen as $`\eta _q=2.5T^3`$, $`\zeta _q=0`$, $`\eta _h=1.5T^3`$ and $`\zeta _h=T^3`$. The other parameters are $`T_c`$ = 160 MeV and $`\sigma =30`$ MeV/fm<sup>3</sup>. With the above set of parameters, Eqs. (55) and (53) are solved to get $`h`$ and $`T`$ as a function of time $`\tau `$ .
Figures 4(a) and 4(b) show the plot of $`s\tau `$\- the rate of entropy production and $`T/T_c`$ \- the rate of supercooling as a function of $`\tau `$ both for ideal hydrodynamic (IHD) and viscous hydrodynamic (VHD) expansions of the system. The system cools below $`T_c`$ until nucleation rate becomes significant. Afterwards, bubble nucleation and growth reheats the system due to the release of latent heat. This behavior is similar to what has been studied earlier in Refs. using a prefactor which explicitly depends on the viscosity coefficient of the plasma. In the present work, since $`\kappa `$ has both viscous and non-viscous components, we study the supercooling and extra entropy production both with $`\kappa _0`$ and $`\kappa `$ particularly when the medium is non-viscous. First we consider only the ideal hydrodynamic expansion. The short-dashed curve and the long-dashed curves are obtained using Eq. (46) for $`\kappa _0`$ (no viscosity) and $`\kappa `$ ($`\kappa _v`$ included) respectively. Since there is supercooling with $`\kappa _0`$, extra entropy is generated even without viscosity. As shown earlier (see Figures 1 and 2) , the effect of viscosity on $`\kappa `$ is not significant with a reasonable choice of $`\eta _q=2.5T^3`$ particularly for small amount of supercooling. Therefore, inclusion of $`\kappa _v`$ does not effect the supercooling much (see the long-dashed curve). The supercooling (hence the entropy production) comes down only by about $``$ 1% due to viscosity, with the present set of parameters. As mentioned before, eventhough we use $`\kappa `$, we do not include viscosity in the hydrodynamical evolution just to bring out the additional effect due to the use of $`\kappa `$ instead of $`\kappa _0`$ in the prefactor. However, when the plasma is viscous, the VHD should be used for consistency (i.e. when $`\kappa _v`$ is included). The use of VHD reduces the supercooling by about 10 $`\%`$ as shown by the solid curve. Although, the amount of supercooling reduces, the entropy production goes up. Since the effect of viscosity on $`\kappa `$ is insignificant, the reduction in supercooling is purely due to the viscous heating of the medium. As a result extra entropy is generated in addition to the entropy that is produced due to supercooling.
## V Conclusion
To summarize, we have derived an expression for the dynamical prefactor which governs the initial growth of critical size bubbles nucleated in first order phase transition. In the case of a non-viscous plasma, the dynamical growth rate is found to depend only on the correlation length and the size of the hadron bubble which are two meaningful scale parameters to describe the critical phenomena at the transition point. The correction to the dynamical prefactor due to viscosity is found to be additive and does not affect the growth process significantly though additional entropy is generated due to viscous heating of the medium. Since the prefactor does not vanish in the limit of zero viscosity, extra entropy is produced during the process of nucleation even when the fluid is non-viscous. Nearly similar conclusions are also drawn by Ruggeri and Friedman who had derived dynamical prefactor by solving relativistic hydrodynamics following a different approach. However, unlike their result, the present prefactor can be written as the sum of viscous and non-viscous terms. Interestingly, using an assumption for velocity of sound in the medium (around the saddle configuartion) which has a form analogous to what is used for non-relativistic plasma, the viscous and the non-viscous parts are found to be similar to the results as obtained by Csernai-Kapusta and Ruggeri-Friedman respectively.
In the present work we solve relativistic hydrodynamic equations both in the interior-exterior, i.e., quark-hadron regions and surface regions. The linear hydrodynamic equation used in the quark-hadron region is obtained after eliminating one of the variable using the relation $`\beta =c_s^2\nu `$ which is not valid in the surface region. Therefore, a different equation is used for the surface region which involves the extra gradient energy. This is where we differ from the Csernai and Kapusta method. Further, Csernai and Kapusta derived $`\kappa `$ by equating the flow of the outward energy flux with the dissipative loss due to viscosity of the medium and the contribution due to a dynamical dissipation was not included. On the other hand, Ruggeri and Friedman solve the hydrodynamic equation only in the quark region and use a set of boundary conditions with certain assumptions. In this context, the present formalism is more general as we solve the linearized hydrodynamic equations in all space and obtain an expression for the prefactor by matching the solutions at the boundary of the interface. Moreover, our result is different in the sense that it has a very simple dependence on the correlation length and radius of the hadron bubble although the CK and RF results can be obtained from it under certain assumption.
## A
We find out the solution for $`\nu (r)`$ that satisfies the radial equation (see Eq. 33),
$`{\displaystyle \frac{d}{dr}}\left[r{\displaystyle \frac{d\overline{e}}{dr}}\left(K{\displaystyle \frac{d^2}{dr^2}}+f^{\prime \prime }\right)\chi (r)\right]=0,`$ (59)
where $`\nu (r)=\chi (r)/r`$. Finally we solve for $`\chi (r)`$ from the equation
$`\left(K{\displaystyle \frac{d^2}{dr^2}}+{\displaystyle \frac{^2f}{\overline{e}^2}}\right)\chi (r)=a(r).`$ (60)
The above equation is quite identical to the one used by Langer-Turski in the surface region. The only difference is that the constant $`a`$ now depends on $`r`$ as $`a(r)(r\overline{e})^1`$. (Note that $`\overline{e}`$ peaks at $`rR`$). Since $`R>>\xi `$ and $`\overline{e}(r)`$ varies sharply in the range $`R\xi rR+\xi `$, $`a(r)`$ mostly depends on the $`\overline{e}(r)`$ variation. Therefore, to a good approximation we can write $`a(r)(R\overline{e})^1`$. The general form of the solution of Eq. (60) is given by
$`\chi (r)={\displaystyle 𝑑r^{}G(r,r^{})a(r^{})},`$ (61)
where $`G`$ is the Green’s function satisfying
$`\left(K{\displaystyle \frac{d^2}{dr^2}}+{\displaystyle \frac{^2f}{\overline{e}^2}}\right)G(r,r^{})=\delta (rr^{}).`$ (62)
On either side of the interface $`^2f/\overline{e}^2`$ is nearly constant. Using the relation $`^2\overline{e}=0`$, it is easy to verify that
$`\chi (r)`$ $``$ $`a(r)\left({\displaystyle \frac{^2f}{\overline{e}^2}}\right)^1`$ (63)
is an approximate solution of Eq. (60) at the interface boundary. Matching the solution in the interfacial region given by Eq. (63) with the solution in the interior region given by Eq. (30) at $`R\xi `$ and with the solution in the exterior region Eq. (31) at $`R+\xi `$, give the following conditions
$`A\mathrm{sinh}(q_hR)=a(R\xi )\left({\displaystyle \frac{^2f}{\overline{e_h}^2}}\right)^1`$ (64)
and
$`B=a(R+\xi )\left({\displaystyle \frac{^2f}{\overline{e}_{q}^{}{}_{}{}^{2}}}\right)^1.`$ (65)
Here the condition $`q\xi <<1`$ has been used. To get a solution inside the interface, we follow the same procedure as that of Ref., i.e., we use the spectral decomposition of $`G`$ as
$`G(r,r^{})={\displaystyle \underset{n}{}}{\displaystyle \frac{\chi _n(r)\chi _n(r^{})}{\overline{\lambda }_n}},`$ (66)
where $`\overline{\lambda }_n`$ are the s-wave eigenvalues and $`\chi _n`$ are the corresponding eigenfunctions. For value of $`r`$ near $`R`$, the sum will be dominated by the first term. This is because $`\overline{\lambda }_1\frac{2K}{R^2}`$, vanishes as $`R`$ becomes large. Since $`\chi _1(r)(\frac{K}{\sigma })^{1/2}(d\overline{e}/dr)`$ is sharply peaked at interface, using Eqs. (61) and (66) we get
$`\chi (r)`$ $``$ $`{\displaystyle \frac{a(R)R^2\mathrm{\Delta }\overline{e}}{2\sigma }}{\displaystyle \frac{d\overline{e}}{dr}},`$ (67)
where $`\mathrm{\Delta }\overline{e}=e_qe_h`$. we can also estimate the variation of $`a`$ in the range $`R\xi rR+\xi `$,
$$\frac{a(R+\xi )}{a(R)}\frac{\overline{e}(R)}{\overline{e}(R+\xi )}.$$
(68)
Assuming $`\xi `$ to be the half width of the full maxima, the above ratio could be $`\sqrt{2}`$.
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# Non–monotonous crossover between capillary condensation and interface localisation/delocalisation transition in binary polymer blends
## 1 Introduction
We study the phase behaviour of a symmetric binary mixture confined into a slit–like pore. The shift of the critical point for symmetric film surfaces (capillary condensation) has been studied extensively. The unmixing transition approaches the critical temperature of the infinite system for thick films. Below the critical temperature the coexisting phases differ in their composition at the center of the film. More recently, novel types of phase transitions have been studied in the case of antisymmetric surface fields,\[2-6\] i.e., one surface attracts the $`A`$ component of a symmetric mixture in exactly the same way the opposite surface attracts the $`B`$ component. Close to the critical temperature in the bulk, enrichment layers of the components gradually form at the surfaces and stabilise an interface at the center of the film (soft–mode phase). It is only close to the temperature of the second order wetting transition in the semi–infinite system that the symmetry is spontaneously broken and the interface is located at either surface. Upon increasing the film thickness the temperature of this interface localisation/delocalisation transition converges to the wetting temperature rather than the critical temperature of the bulk, and the interpretation of the transition in terms of a thin film critical critical point or a wetting transition in the limit of large film thickness has been discussed. The effect of an external field (i.e., gravity) has been explored. Analytical approaches and simulations have considered systems with second order wetting transitions in the semi–infinite system or the behaviour at bulk coexistence only. This has excluded a possible interplay between phase behaviour and prewetting, which might alter the topology of the phase diagram in thin films, from consideration.
In this Letter we revisit the interface localisation/delocalisation transition for a first order wetting transition in the corresponding semi–infinite system and explore the crossover from the interface localisation/delocalisation transition to capillary condensation upon varying the interaction $`\mathrm{\Lambda }_2`$ at one surface. The interaction $`\mathrm{\Lambda }_1`$ at the other surface is kept constant. This is the first systematic theoretical study of boundary conditions which are neither strictly symmetric nor antisymmetric. It is clearly pertinent to the interpretation of experiments as the idealized limiting cases are never strictly realized. Our calculations yield information about the range of asymmetry where the behaviour characteristic of the symmetric and antisymmetric boundary conditions is observable, and we explore the dependence of the topology of the phase diagrams on the surface interactions. Both the phase diagram as a function of the intensive variables incompatibility and chemical potential as well as the binodals are discussed.
## 2 Self–consistent field calculations
We calculate the phase behaviour of a confined $`AB`$ mixture within the self-consistent field theory of Gaussian polymers.\[12-14\]Polymer mixtures in confined geometries are interesting for many applications (e.g., coating, lubrication), and consequentially have attracted abiding recent attention. In these systems the soft mode phase mentioned above has been experimentally studied. Thus we focus on these systems, although we believe that our results carry over to other confined mixtures at least qualitatively. We consider a film with volume $`V_0=\mathrm{\Delta }_0\times L\times L`$. $`\mathrm{\Delta }_0`$ denotes the film thickness, while $`L`$ is the lateral extension of the film. The density at the film surfaces decreases to zero in a boundary region of width $`\mathrm{\Delta }_w`$ according to
$$\mathrm{\Phi }_0(x)=\{\begin{array}{cc}\frac{1\mathrm{cos}\left(\frac{\pi x}{\mathrm{\Delta }_w}\right)}{2}\hfill & \text{}0x\mathrm{\Delta }_w\hfill \\ 1\hfill & \text{}\mathrm{\Delta }_wx\mathrm{\Delta }_0\mathrm{\Delta }_w\hfill \\ \frac{1\mathrm{cos}\left(\frac{\pi (\mathrm{\Delta }_0x)}{\mathrm{\Delta }_w}\right)}{2}\hfill & \text{}\mathrm{\Delta }_0\mathrm{\Delta }_wx\mathrm{\Delta }_0\hfill \end{array}$$
(1)
where $`\mathrm{\Phi }_0`$ denotes the ratio of the monomer density and the value $`\rho `$ in the middle of the film. The thickness $`\mathrm{\Delta }`$ of an equivalent film with constant monomer density $`\rho `$ is $`\mathrm{\Delta }=\mathrm{\Delta }_0\mathrm{\Delta }_w`$. We choose $`\mathrm{\Delta }_w=0.15R_e`$, where $`R_e`$ is the end–to–end distance. Both surfaces interact with the monomer species via a short range potential $`H`$:
$$H(x)=\{\begin{array}{cc}\frac{4\mathrm{\Lambda }_1R_e\left\{1+\mathrm{cos}\left(\frac{\pi x}{\mathrm{\Delta }_w}\right)\right\}}{\mathrm{\Delta }_w}\hfill & \text{}0x\mathrm{\Delta }_w\text{ }\hfill \\ 0\hfill & \text{}\mathrm{\Delta }_wx\mathrm{\Delta }_0\mathrm{\Delta }_w\hfill \\ \frac{4\mathrm{\Lambda }_2R_e\left\{1+\mathrm{cos}\left(\frac{\pi (\mathrm{\Delta }_0x)}{\mathrm{\Delta }_w}\right)\right\}}{\mathrm{\Delta }_w}\hfill & \text{}\mathrm{\Delta }_0\mathrm{\Delta }_wx\mathrm{\Delta }_0\hfill \end{array}$$
(2)
$`H>0`$ is attractive for the $`A`$ monomers and repulsive for the $`B`$ species. The normalisation of the surface fields $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$, which act on the monomers close to the left and the right surface, is chosen such that the integrated interaction energy between the surface and the monomers is independent of the width of the boundary region $`\mathrm{\Delta }_w`$.
$`A`$ and $`B`$ polymers contain $`N`$ monomers and are structurally symmetric. The polymer conformations $`\{𝐫_\alpha (\tau )\}`$ determine the microscopic $`A`$ monomer density $`\widehat{\mathrm{\Phi }}_A(𝐫)=\frac{N}{\rho }_{\alpha =0}^{n_A}_0^1\mathrm{d}\tau \times `$
$`\delta \left(𝐫𝐫_\alpha (\tau )\right)`$, where the sum runs over all $`n_A`$ $`A`$ polymers in the system and $`0\tau 1`$ parameterises the contour of the Gaussian polymer. A similar expression holds for $`\widehat{\mathrm{\Phi }}_B(𝐫)`$. With this definition the semi–grandcanonical partition function takes the form:
$`𝒵`$ $`{\displaystyle \underset{n_A=1}{\overset{n}{}}}{\displaystyle \frac{e^{+\mathrm{\Delta }\mu n_A/2k_BT}}{n_A!}}{\displaystyle \frac{e^{\mathrm{\Delta }\mu n_B/2k_BT}}{n_B!}}`$ $`{\displaystyle 𝒟_A[𝐫]𝒫_A[𝐫]𝒟_B[𝐫]𝒫_B[𝐫]\delta \left(\mathrm{\Phi }_0\widehat{\mathrm{\Phi }}_A\widehat{\mathrm{\Phi }}_B\right)}`$ (3)
$`\times \mathrm{exp}(\rho {\displaystyle }\mathrm{d}^3𝐫\{\chi \widehat{\mathrm{\Phi }}_A\widehat{\mathrm{\Phi }}_BH(\widehat{\mathrm{\Phi }}_A\widehat{\mathrm{\Phi }}_B\})`$
where $`n=n_A+n_B`$ and $`\mathrm{\Delta }\mu `$ represents the exchange potential between $`A`$ and $`B`$ polymers. The functional integral $`𝒟`$ sums over all conformations of the polymers, and $`𝒫[𝐫]\mathrm{exp}\left(\frac{3}{2R_e^2}_0^1d\tau \left(\frac{\mathrm{d}𝐫}{\mathrm{d}\tau }\right)^2\right)`$ denotes the statistical weight of a non–interacting Gaussian polymer. The second factor enforces the monomer density profile to comply with Eq.(1) (incompressibility). The Boltzmann factor in the partition function incorporates the thermal repulsion between unlike monomers and the interactions between monomers and surfaces. The strength of the repulsion is described by the Flory–Huggins parameter $`\chi `$.
In mean field approximation the free energy is obtained as the extremum of the functional:
$`{\displaystyle \frac{𝒢[W_A,W_B,\mathrm{\Phi }_A,\mathrm{\Phi }_B,\mathrm{\Xi }]}{nk_BT}}+\mathrm{ln}{\displaystyle \frac{n}{V_0}}\mathrm{ln}𝒬+{\displaystyle \frac{1}{V}}{\displaystyle \mathrm{d}^3𝐫\chi N\mathrm{\Phi }_A\mathrm{\Phi }_B}HN\left\{\mathrm{\Phi }_A\mathrm{\Phi }_B\right\}`$
$`{\displaystyle \frac{1}{V}}{\displaystyle \mathrm{d}^3𝐫\left\{W_A\mathrm{\Phi }_A+W_B\mathrm{\Phi }_B\right\}}+\mathrm{\Xi }\left\{\mathrm{\Phi }_0\mathrm{\Phi }_A\mathrm{\Phi }_B\right\}`$ (4)
with respect to its five arguments. $`𝒬_A`$ denotes the single chain partition function:
$$𝒬_A[W_A]=\frac{1}{V_0}𝒟_1[𝐫]𝒫_1[𝐫]e^{_0^1d\tau W_A(𝐫(\tau ))}$$
(5)
a similar expression holds for $`𝒬_B`$, and $`𝒬=\mathrm{exp}(\mathrm{\Delta }\mu /2k_BT)𝒬_A+\mathrm{exp}(\mathrm{\Delta }\mu /2k_BT)𝒬_B`$.
The values of $`W_A,W_B,\mathrm{\Phi }_A,\mathrm{\Phi }_B,\mathrm{\Xi }`$ which extremize the free energy functional are denoted by lower–case letters and satisfy the self-consistent set of equations
$$w_A(𝐫)=\chi N\varphi _B(𝐫)H(𝐫)N+\xi (𝐫)\text{and}\varphi _A(𝐫)=\frac{V}{𝒬}\frac{𝒟𝒬_A}{𝒟w_A(𝐫)}$$
(6)
$`\mathrm{\Phi }_0(𝐫)=\varphi _A(𝐫)+\varphi _B(𝐫)`$, and similar expressions for $`\varphi _B`$ and $`w_B`$.
To calculate the monomer density we employ the end segment distribution $`q_A(𝐫,t)`$
$$q_A(𝐫,\tau )=_0^\tau 𝒟_1[𝐫]𝒫_1[𝐫]\delta (𝐫𝐫(\tau ))e^{_0^\tau dtw_a(𝐫(t))}$$
(7)
It satisfies the diffusion equation:
$$\frac{q_A(𝐫,\tau )}{\tau }=\frac{R_e^2}{6}\mathrm{}q_A(𝐫,\tau )w_A(𝐫)q_A(𝐫,\tau )$$
(8)
Then, the $`A`$ monomer density and the single chain partition can be calculated via
$$\varphi _A(𝐫)=\frac{Ve^{\mathrm{\Delta }\mu /2k_BT}}{V_0𝒬}_0^1d\tau q_A(𝐫,\tau )q_A(𝐫,1\tau )\text{and}𝒬_A=\frac{1}{V_0}\mathrm{d}^3𝐫q_A(𝐫,1)$$
(9)
We expand the spatial dependence of the densities and fields in a set of orthonormal functions $`\{f_k(x)=\sqrt{2}\mathrm{sin}(\pi kx/\mathrm{\Delta }_0)\}`$ with $`k=1,2,\mathrm{}`$. This procedure results in a set of non–linear equations which are solved by a Newton–Raphson–like method. We use up to 80 basis functions and achieve a relative accuracy $`10^4`$ in the free energy. For symmetric boundary fields only components with an odd index $`k`$ are non–zero. Substituting the extremal values of the densities and fields into the free energy functional (4) we calculate the free energy $`G`$ of the different phases. At coexistence the two phases have equal semi–grandcanonical potential.
## 3 Results
The phase boundaries as a function of the surface fields are presented in Fig.1. First, we consider the antisymmetric boundary condition $`\mathrm{\Lambda }_1N=\mathrm{\Lambda }_2N=0.5`$. One half of the phase diagram is enlarged in the inset. If we reduce the temperature along the symmetry axis of the phase diagram $`\varphi =1/2`$, an interface is stabilized at the center of the film for $`\chi N<\chi _{\mathrm{triple}}N`$, and at $`\chi _{\mathrm{triple}}N`$ we encounter a first order interface localisation/delocalisation transition at which the interface jumps discontinuously to one of the two surfaces. At this triple point phases with compositions $`\varphi _{\mathrm{triple}}`$, $`1/2`$ and $`1\varphi _{\mathrm{triple}}`$ coexist. This behaviour for $`\mathrm{\Delta }\mu =0`$ has been confirmed by computer simulations of Ising models.
For all other compositions in the range $`[\varphi _{\mathrm{triple}}:1\varphi _{\mathrm{triple}}]`$, however, one encounters two transitions upon cooling. Consider cooling at $`\varphi _{\mathrm{triple}}<\varphi <1/2`$ as indicated by the arrow in the inset. At high temperatures enrichment layers gradually form at the surfaces. Slightly below the critical temperature of the bulk an interface is stabilized which separates a thin $`A`$–rich layer at one surface from a thicker $`B`$–rich layer at the opposite surface (cf. Fig.1b upper panel). Laterally, the system is homogenous (soft–mode phase) and an $`AB`$ interface runs parallel to the surfaces. Upon cooling, we first encounter a phase separation into lateral regions with a thin and a thick $`A`$–rich layer(cf. Fig.1b middle panel). This phase coexistence is the analogy of the prewetting line in a semi–infinite system. Since the coexistence involves only finite layer thicknesses it also persists in a thin film, provided the film thickness is sufficiently large. The different layer thicknesses correspond to different compositions of the film and give rise to two symmetric coexistence regions. This prediction has not yet been observed in computer simulations or experiments. Upon further cooling, we encounter a second transition at $`\chi _{\mathrm{triple}}N`$ where the two phase coexistences merge and the system laterally segregates into $`A`$–rich and $`B`$–rich regions(cf. Fig.1b lower panel). For temperatures far below the triple temperature the interface between the coexisting phases runs straight across the film. The angle between the interface and the surface corresponds to the contact angle $`\mathrm{\Theta }`$ of droplets in the semi–infinite system. Upon increasing the film thickness the triple temperature approaches the wetting transition temperature, the two coexistence regions become narrower and converge to the prewetting lines, and the critical temperature of the film tends to the prewetting critical temperature of the semi–infinite system.
The phase coexistence in terms of incompatibility $`\chi N`$ and exchange potential $`\mathrm{\Delta }\mu `$ is shown in Fig.2. Although the Hamiltonian of the system is symmetric with respect to the exchange $`AB`$, phase coexistence occurs at $`\mathrm{\Delta }\mu =0`$ only below the triple temperature. At $`\chi _{\mathrm{triple}}N`$ the coexistence curve bifurcates and the branches move away from the symmetry axis. The two coexistence lines resemble the prewetting lines which correspond to the first order wetting transitions at the two surfaces.
At low temperatures $`\chi \chi _{\mathrm{triple}}`$ the two coexisting phases are almost pure and the coexistence value of the chemical potential is given by $`\mathrm{\Delta }\mu _{\mathrm{coex}}/k_BT=4(\mathrm{\Lambda }_1+\mathrm{\Lambda }_2)NR_e/\mathrm{\Delta }`$. Upon increasing $`\mathrm{\Lambda }_2`$ this value shifts to more negative values. Only for large film thickness the two phase region on the $`B`$–rich side of the phase diagram remains. For thin films, however, $`\mathrm{\Delta }\mu _{\mathrm{coex}}`$ is shifted to negative values which are eventually smaller than the chemical potential of the prewetting critical point at the surface which attracts $`A`$ (cf. Fig.2). For the parameters chosen this occurs roughly around $`\mathrm{\Lambda }_2N0.2`$. In this case the two phase region on the $`B`$–rich side as well as the associated critical point disappears upon increasing $`\mathrm{\Lambda }_2`$ further.
The critical point on the $`A`$–rich side of the phase diagram shifts to higher temperatures upon increasing $`\mathrm{\Lambda }_2`$ similar to the prewetting at the corresponding surface. Note that $`\chi _{\mathrm{wet}}N24(\mathrm{\Lambda }_2N)^2`$ for strong segregation. Although the film surfaces favor (on average) the $`A`$ component of the mixture, the critical point occurs at $`\mathrm{\Delta }\mu >0`$ for $`\mathrm{\Lambda }_2N\stackrel{<}{}0.25`$. Around $`\varphi =1/2`$ or $`\mathrm{\Delta }\mu _{\mathrm{coex}}=0`$ the $`B`$–rich binodal shows a concave curvature, which is the remanent of the wetting transition at the surface which favours $`B`$. Upon increasing $`\mathrm{\Lambda }_2`$ to positive values this feature disappears, of course.
For an almost neutral surface $`\mathrm{\Lambda }_2N0`$ the critical temperature passes through a maximum and the critical temperature of the capillary condensation is approached from above. The surface field dependence of the critical temperature and critical density is investigated in Fig.3. This non-monotonic behaviour of the critical temperature and density can be rationalised as follows: For $`\mathrm{\Lambda }_2=\mathrm{\Lambda }_1`$ the composition profiles of the two coexisting phases close to the critical point are symmetric (c.f. Fig.3 inset). It is the difference in the composition at the center which distinguishes between the phases and vanishes at the critical point. Decreasing $`\mathrm{\Lambda }_2`$ we reduce the influence of the surfaces, and the critical temperature increases and the critical composition decreases, respectively, as both tend towards their bulk values. For large negative values of $`\mathrm{\Lambda }_2`$ the transition is qualitatively different. It is associated with the prewetting transition, i.e., the profiles of the coexisting phases resemble a profile across an interface and it is the location of the interface which distinguishes the two phases (c.f. Fig.3 inset). Decreasing $`\mathrm{\Lambda }_2`$ shifts the wetting transition to lower temperatures and higher $`A`$ content. The gradual crossover between the two distinct behaviours occurs around $`\mathrm{\Lambda }_20`$. The corresponding profiles are included into Fig.3 (inset). In both phases the composition at the surface favouring $`A`$ is almost saturated and the composition decreases to a value larger or smaller than 1/2 at the almost neutral surface. Upon approaching the critical point the composition difference at the neutral surface vanishes. This is similar to the capillary condensation mechanism for $`\mathrm{\Lambda }_2>0`$. Both profiles, however, can also be perceived as very broad interfaces (with a width comparable to the film thickness) with different centers. This resembles the prewetting–like behaviour for $`\mathrm{\Lambda }_2<0`$.
## 4 Discussion
In summary, we have explored the unusual dependence of the phase behaviour of a confined symmetric mixture on the surface fields. For nearly antisymmetric surface fields which are strong enough to produce a first order wetting transition our mean field calculations predict two coexistence regions which correspond to the prewetting transition at either surface. We vary the surface interactions from antisymmetric to symmetric and the critical temperature(composition) passes through a maximum(minimum) when one surface is almost neutral. We expect the qualitative dependence of the phase diagram on the surface fields to be generic. Fluctuations, which are not included in the mean field calculations, impart 2D Ising rather than mean field behaviour to the critical points. Their importance is, however, restricted to a narrow region around the critical point in the limit of strong interpenetration $`\rho R_e^3/N1`$. Capillary waves increase the range of the effective interaction between interface and surface by a factor $`1+\omega /2`$ where $`\omega =k_BT/4\pi \xi ^2\sigma `$. $`\xi `$ and $`\sigma `$ denote the correlation length and interfacial tension between the coexisting bulk phases, respectively. Simulations indicate that $`1/\omega \sqrt{\chi N}f(\chi N)\rho R_e^3/N`$, where the scaling function $`f(\chi N)`$ approaches a constant value for $`\chi N\mathrm{}`$. Since the wetting transition in binary polymer blends is typically first order and occurs at strong segregation, the capillary parameter $`\omega `$ is small. Experiments on polymer blends have explored the phase behaviour in confined geometry, and these systems as well as computer simulations might prove convenient for testing our predictions.
We are currently exploring the thickness dependence of the phase diagram for antisymmetric surfaces within the self–consistent field theory and a Landau–Ginzburg approach. Contrary to the behaviour close to a second order wetting transition, we find that the triple point approaches the temperature of the first order wetting transition from above upon increasing the film thickness. The transition between first and second order interface localisation/delocalisation transition and the behaviour at the tricritical point will be elucidated.
\***
It is a pleasure to thank E. Reister and F. Schmid for helpful discussions. Financial support was provided by the DFG under grant Bi314/17 within the Priority Program “Wetting and Structure Formation at Interfaces”, the VW Stiftung, and PROALAR2000.
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# NEW PHENOMENA II: RECENT RESULTS FROM THE FERMILAB TEVATRON
## 1 Introduction
In these, and other$`^\mathrm{?}`$, proceedings we summarize some recent results of searches for new physics at the Fermilab Tevatron. Specifically, we review new results on R-parity violating Supersymmetry (SUSY) and $`Z^{}`$/Technicolor models in the $`ee`$ and $`t\overline{t}`$ channels. We conclude by introducing Sherlock, a new quasi-model-independent search strategy.
## 2 R-parity Violating SUSY
Motivated in part by the interpretation$`^\mathrm{?}`$ of the reported HERA excess of high $`Q^2`$ events$`^\mathrm{?}`$, and in part for theoretical reasons, many recent searches for SUSY have focused on the possibility that R-parity is not conserved. In such a scenario at the Tevatron, pairs of SUSY particles will typically be produced and then decay to standard model particles and two neutralinos. However, instead of leaving the detector without depositing any energy (as in R-parity conserving models), R-parity violating terms, $`\lambda _{ijk}\mathrm{}_i\stackrel{~}{\mathrm{}}_j\overline{e}_k+\lambda _{ijk}^{}\mathrm{}_i\stackrel{~}{q}_j\overline{d}_k+\lambda _{ijk}^{\prime \prime }\overline{u}_i\overline{\stackrel{~}{d}_j}\overline{d}_k`$, allow each neutralino to decay via three body decay to standard model particles. For example in a scenario with a non-zero all-leptonic coupling term, a neutralino could decay via a $`\nu _e`$ and a virtual $`\stackrel{~}{\nu }_e`$ with the $`\stackrel{~}{\nu }_e`$ decaying via a $`\lambda _{121}`$ coupling and producing $`\mu ^{}e^+`$. A similar decay of the other neutralino in the event allows the event to produce a total of 4 charged leptons as well as two neutrinos each of which can contribute to missing transverse energy. The CDF and DØ collaborations have both searched for this type of production and decay.
At DØ the search for R-parity violating SUSY is performed by looking in the $`eee,ee\mu ,e\mu \mu `$ and $`\mu \mu \mu `$ channels for excesses of events with missing transverse energy. Since the search is identical to that for $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_2^0`$ production and decay in mSUGRA, identical data sets and selection criteria are employed$`^\mathrm{?}`$. The results, along with the luminosity for each sample, are shown in Table 1 with no candidates in the data. A similar search for R-parity violating SUSY was performed at CDF in 87.5 pb<sup>-1</sup> of data. It is complementary in that it requires four leptons, but does not require missing transverse energy. Using all combinations of electrons and muons in the final state, there is one candidate event in the data, consistent with the background prediction of 1.3$`\pm `$0.4 events.
Limits on R-parity violating models, which can be set as a function of the mSUGRA SUSY model parameters ($`m_0,m_{1/2},A_0,tan\beta `$ and $`\mu `$), are shown Figures 2 and 2 for different values of the $`\lambda `$ coupling constant and $`tan\beta `$. In both cases the regions of 95% C.L. exclusion correspond to the space below the dark solid lines. The lighter curves on Figure 2 indicate the value of $`\lambda `$ such that the average decay length of the LSP is less than 1 cm. Since the search is not sensitive to displaced decays of the neutralino, the region below the curves labeled with $`\lambda _{121}`$ and above the $`10^3`$ line is excluded if $`\lambda _{121}>10^3`$.
## 3 Technicolor and Neutral Heavy Vector Bosons
In Technicolor models$`^\mathrm{?}`$ and models which have additional neutral heavy vector bosons$`^\mathrm{?}`$, new heavy particles (labeled $`\rho _T,\omega _T`$ or $`Z^{}`$ respectively) are produced and can decay via $`e^+e^{}`$. In both cases, a straight-forward search for a resonance, or excess at high mass, in the $`ee`$ invariant mass spectra could easily illuminate a signal. At D$`Ø`$ this analysis uses the same data set as the recently published paper setting limits on quark and lepton compositeness$`^\mathrm{?}`$. The luminosity for this data set is 120.9 pb<sup>-1</sup>. The backgrounds are dominated by standard model production of $`Z/\gamma ^{}ee`$, and instrumental backgrounds (fakes). The invariant mass spectrum for the data and backgrounds is shown in Figure 4. There is no evidence for resonant production in the data or for an excess at high mass. A similar search using both $`ee`$ and $`\mu \mu `$ final states by CDF was recently published$`^\mathrm{?}`$. In both searches there is no evidence for new physics.
Figure 4 shows the DØ preliminary 95% C.L. cross section upper limits on production of $`\rho _T/\omega _Te^+e^{}`$. Also shown on the plot are theoretical production curves for two different scenarios in which the decay $`\rho _TW\pi _T`$ is allowed or not allowed, affecting the branching ratio to $`ee`$. Assuming the $`\rho _T`$ and $`\omega _T`$ have the same mass, the mass limits are M$`{}_{\rho _T}{}^{}>225`$ GeV if the decay $`\rho _TW\pi _T`$ is kinematically disallowed or if M$`{}_{\mathrm{T}}{}^{}>200`$ GeV which suppresses the decay $`\omega _T\gamma \pi _T`$. While CDF does not set specific limits on Technicolor particles in this channel, their sensitivity is comparable$`^\mathrm{?}`$.
Similarly, CDF and DØ also set limits on neutral heavy vector bosons which decay via $`Z^{}ee`$. Assuming the couplings to known fermions are the same as in the standard model, DØ sets a mass limit of M$`{}_{Z^{}}{}^{}>670`$ GeV in the electron only channel at 95% C.L. CDF’s combined 95% C.L. limit from $`ee`$ and $`\mu \mu `$ is M$`{}_{Z^{}}{}^{}>`$690 GeV. CDF goes further and sets limits on other $`Z^{}`$ models with extended gauge groups$`^\mathrm{?}`$.
Another search for neutral heavy vector bosons was carried out by the CDF collaboration in the $`t\overline{t}`$ channel$`^\mathrm{?}`$<sup>,</sup>$`^\mathrm{?}`$. The analysis, representing 106 pb<sup>-1</sup> of data, uses the ‘lepton+jets’ data set ($`e`$ and $`\mu `$) which was used for the top quark mass measurements$`^\mathrm{?}`$. Using the same fitting techniques which are used to measure the mass, the final state objects are fit to the $`t\overline{t}`$ hypothesis constraining the mass of top quark mass to be 175 GeV. The best fit is then used to calculate the invariant mass of the $`\overline{t}t`$ system which is searched for resonant structure or excesses. There is no evidence for new physics seen. While this channel is not competitive with $`Z^{}`$ limits from $`ee`$ and $`\mu \mu `$ searches (assuming standard model couplings), Technicolor models with leptophobic topcolor production$`^\mathrm{?}`$ could make this channel highly produced at the Tevatron. Experimental 95% C.L. cross section upper limits are at the few picobarn level and are compared with theoretical production cross sections in Figure 6.
## 4 Sherlock: A New Quasi-Model-Independent Method for Searches for New Physics
Finally, we introduce Sherlock, a new quasi-model-independent search method developed at DØ. This method provides a prescription for searching for new physics by systematically looking for excesses in multi-dimensional data distributions. Since we assume that the physics responsible for electroweak symmetry breaking occurs at mass scales large compared to standard model backgrounds, we currently add in the assumption that the new physics is characterized by high P<sub>T</sub> final state particles. It is this feature which makes it quasi-model-independent.
The method consists of a three part prescription and algorithm. The first part is to pick a data set (such as the inclusive $`e\mu `$+X sample) and categorize each event according to its observed final state particles (number of electrons, jets, photons etc.). For each category of events, the kinematic variables for the sample are uniquely specified by an a priori prescription. A region, R, is then defined in the multi-variable space surrounding one or more of the data points. By giving a precise definition of the region for an event, or set of events, an amount of parameter space is determined and the probability for the background in that region to fluctuate up to or above the number of observed events in the region gives a quantitative measure of the degree of interest of the region. The algorithm then searches for the most interesting region (largest excess relative to background) in all of variable space including the high P<sub>T</sub> region. Once this most interesting region is found, it is compared to the most interesting region found in a large sample of hypothetical similar experiments (HSEs) drawn from the parent distributions of the backgrounds (according to statistical and systematic uncertainties). In this way, the true degree of interest of the region of largest excess is quantified in terms of the fraction of HSEs which give regions which are more interesting than the one observed in our data (again due simply to statistical fluctuations or systematic misunderstanding of the data).
Sherlock is run on 108 pb<sup>-1</sup> of inclusive high P<sub>T</sub> $`e\mu `$+X data taken at DØ. As a test and an illustration of the sensitivity of the method, we have used this algorithm on a set of mock experiments drawn from the background estimations. The $`e\mu X`$ sample has the advantage of having two known ‘signals’ which give high P<sub>T</sub> physics in the final state: WW and $`t\overline{t}`$ production. Figure 8 shows the results of running a series of mock experiments in which the ‘known’ backgrounds include only $`Z/\gamma \tau \tau `$ and fakes and the ‘data’ is drawn only from the background distributions. The results are shown in the figure with the dashed lines which shows the significance in standard deviations. As expected, the results are basically Gaussian and centered on 0 with unit width. When ‘signal’ events of WW and $`t\overline{t}`$ are added to the ‘data’ (according to expected standard model production expectations, smeared by statistical and systematic uncertainties) but not to the background expectations, as shown by the solid line in Figure 8, Sherlock reports that in the over 50% of the mock experiments it finds a statistically significant excess at or greater than 2$`\sigma `$<sup>a</sup><sup>a</sup>aWe again note that Sherlock doesn’t know anything about $`WW`$ or $`t\overline{t}`$ and that it is in no way optimized for finding them. It is simply looking for an excess of events in the high-P<sub>T</sub> region.. Running Sherlock on the data itself, using on $`Z/\gamma \tau \tau `$ and fakes as the background, Sherlock picks out significant excesses ($`>2\sigma `$) in both the $`e\mu \overline{)}\mathrm{E}_\mathrm{T}`$ ($`WW`$) and $`e\mu \overline{)}\mathrm{E}_\mathrm{T}jj`$ ($`t\overline{t}`$) data correctly indicating the presence of both WW and $`t\overline{t}`$ in the data.
Figure 8 shows the results when WW are added to the background expectations so that we may see the sensitivity to $`t\overline{t}`$ alone. Again, in an ensemble of mock experiments Sherlock picks out an excess in the data at greater than or equal to the 2$`\sigma `$ level in over 25% of the cases. Running on the data we observe an excess at the 1.9$`\sigma `$ level in the $`e\mu \overline{)}\mathrm{E}_\mathrm{T}jj`$ data, correctly identifying $`t\overline{t}`$ in the data.
Including all known standard model sources in the backgrounds and running Sherlock yields no evidence of an excess in any $`e\mu X`$ channel indicating agreement with the standard model. Specifically, in 71% of hypothetical similar experiments we expect to see an excess more interesting than the most interesting region of excess than is observed in our data.
The Sherlock method is a novel approach to searching for new physics in the data. While a dedicated search is always best for a specific signal hypothesis (e.g. SM Higgs with M<sub>H</sub>=130 GeV), the number of possible models to search for is very large. Sherlock is a very powerful method for being sensitive to a large number of new physics models by being much more model independent. While we have only applied the method to a single data set, it is generally applicable and should prove to be an immensely valuable tool in Run II to complement the dedicated searches.
## 5 Conclusions
We have presented the results of a number of recent searches for new physics from the Fermilab Tevatron based on the data taken during 1992-1996. While the next run with upgraded detectors and high luminosity is just over a year away, we continue to make progress in searching for new physics as well as setting limits on important theoretical models. The future at the Tevatron appears bright, and the lab should continue to be an interesting and exciting place to search for new physics in the coming years.
## Acknowledgments
The author would like to thank Bruce Knuteson, Juan Valls, Maxwell Chertok, Meenakshi Narain, Elemer Nagy, Gustaaf Brooimans, Andre Turcot, Greg Graham, Ray Culbertson and Kaori Maeshima for their assistance with this talk and proceedings.
## References
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# Cooling of a single atom in an optical trap inside a resonator
## I Introduction
A recent experiment succeeded in trapping a single atom with single photons inside an optical cavity and in monitoring the atomic motion with the resolution approaching the standard quantum limit for position measurements. Yet a second experiment has likewise reported single-atom trapping at the few-photon level, although in this case the trapping potential and diffusion are in fact well approximated by a free-space semiclassical theory .
One future objective for such experiments is to use atoms trapped in cavities for quantum communication purposes, with atoms serving as quantum memories and photons as the transporters of quantum information . While the single-photon trapping experiments provide a new paradigm for quantum measurement and control, they are nevertheless not entirely suitable for the purpose of distributed quantum networks where qubits will be communicated among quantum nodes. The reason is the short trapping life time of the atoms as well as limited operation flexibility. A better strategy might be to use the cavity QED field for quantum state entanglement and distribution while an additional (external) trapping mechanism provides the necessary confinement of the atomic center-of-mass motion. For instance, in another recent experiment from the Caltech group , mean trapping times of $`28`$ms (as compared to mean trapping times of $`<1`$msec in the experiments ) were achieved by employing a far-off resonant trapping (FORT) beam along the cavity axis. In that experiment the trapping lifetime was limited due to intensity fluctuations of the intracavity FORT beam . Here we consider the situation of current improved experiments in which a single atom is held inside an optical cavity in a stable FORT beam of minimum intensity fluctuations.
Several mechanisms for cooling inside optical resonators have been discussed before . Here we discuss in detail how the combination of an external trapping potential and the cavity QED field adds flexibility in predetermining where and to what degree atoms will be trapped and cooled. Moreover, our calculations go beyond the weak driving limit discussed in . That is, we allow the “probe” field driving the cavity to be so strong as to appreciably modify the dynamical behavior of, rather than merely probe, the atom-cavity system.
This paper is organized as follows. In Section II we describe the physical situation of an atom trapped in an optical potential and strongly interacting with a cavity QED field. We give the evolution equations for both internal and external atomic degrees of freedom and for the quantized cavity mode. Section III contains an exposition on how we calculated friction and diffusion coefficients from the forces acting on the atom. Section IV contains the main results of this paper: we discuss simple pictures for cooling mechanisms, based on the dressed state structure of the atom-cavity system, and give numerical results for the typical cooling and diffusion rates, and hence “temperatures” for single atoms under various trapping conditions. We also study the saturation behavior under strong driving conditions and perform simulations of the full 3-D motion of atoms trapped in particular wells that show how the probe field transmission is correlated with the atomic motion and how trapping times can be prolonged by strong cooling. Section IV F concludes with a brief discussion of a slightly different trapping scheme. The summary highlights the main results.
## II Description of problem
We consider a single two-level atom coupled to a single quantized cavity mode and coupled to a (classical) far-off resonant trapping (FORT) beam. In most of the paper we assume that the FORT shifts the atomic excited state $`|e`$ up and the ground state $`|g`$ down by an amount $`S_F(\stackrel{}{r})`$ (i.e., the energy of the ground state is $`E_gS_F`$, that of the excited state $`E_e+S_F`$), as this is the situation pursued in previous and current experiments . In Section IV F, however, we will also study the different situation where both ground and excited states are shifted down by $`S_F`$ (see e.g. ). The FORT beam coincides with one of the longitudinal modes of the cavity and its wavelength $`\lambda _F`$ is longer than that of the main cavity mode of interest for cavity QED, $`\lambda _0`$. In fact, in the experiments the cavity length $`L`$ is $`104\lambda _0/2=102\lambda _F/2`$.
The position-dependent AC-Stark shift due to the FORT field is of the form
$$S_F(\stackrel{}{r})=S_0\mathrm{sin}^2(k_Fz)\mathrm{exp}(2\rho ^2/w_0^2),$$
(1)
with $`S_0>0`$ the maximum shift, $`k_F=2\pi /\lambda _F`$ the wave vector of the FORT field, $`w_0`$ the size of the Gaussian mode of the cavity, while $`z`$ and $`\rho `$ give the coordinate along, and the distance perpendicular to, the cavity axis, respectively. The quantized cavity mode is assumed to have the same transverse dimensions $`w_0`$<sup>*</sup><sup>*</sup>*It is in fact the Rayleigh ranges of the beams that are identical, so that $`w_0^{\mathrm{FORT}}/w_0^{\mathrm{cav}}=\sqrt{\lambda _0/\lambda _F}0.99`$ so that the atom-cavity coupling is determined by
$$g(\stackrel{}{r})=g_0\mathrm{sin}(kz)\mathrm{exp}(\rho ^2/w_0^2),$$
(2)
with $`g_0`$ the maximum coupling rate and $`k=2\pi /\lambda _0`$ the wave vector of the cavity mode. Under conditions where the cavity is not driven too strongly, the atom will be trapped around the anti-nodes of the red-detuned FORT field. Thanks to the fact that $`\lambda _0\lambda _F`$, the atom will experience a different coupling strength to the cavity mode in each different well. Figure 1 shows the axial pattern arising from the FORT and cavity fields. For illustrative purposes we choose here (and in the rest of this paper) a cavity of length $`L=16\lambda _0=15\lambda _F`$. This does not influence the basic physics involved: in particular we note that the precise value of $`\lambda _F`$ is largely irrelevant on the time scales considered here, as the FORT field is detuned far from atomic resonance. The choice of $`L=16\lambda _0=15\lambda _F`$ just means that only 8 wells out of 30 are qualitatively and quantitatively different.
This is illustrated in Fig. 2 where we plot the value of the cavity QED coupling $`g`$ at the anti-nodes of the FORT (i.e., the bottom of the trapping potential). In particular, there are 2 anti-nodes in which $`g=0`$, and 4 in which $`|g|`$ attains it maximum.
The cavity is driven by an external classical field $`(t)=_0\mathrm{exp}(i\omega _pt)`$, at a frequency $`\omega _p`$, which is used to probe the atom-cavity system and which may cool the atom at the same time. In the following, the strength of the driving field is indicated by the number of cavity photons $`N_e`$ that would be present if there were no atom in the cavity, rather than by $`_0`$. This closely follows the experimental procedure for determining the driving strength. The relation between the two is
$$N_e=\frac{_0^2}{\kappa ^2+\mathrm{\Delta }_c^2},$$
(3)
with $`\mathrm{\Delta }_c=\omega _c\omega _p`$ the detuning of the probe from the cavity frequency $`\omega _c=kc`$. The Hamiltonian for the internal atomic degrees of freedom and the quantized cavity mode is, in a frame rotating at the probe frequency $`\omega _p`$, given by
$`H`$ $`=`$ $`\mathrm{}\mathrm{\Delta }_ca^+a+\mathrm{}\mathrm{\Delta }_a\sigma ^+\sigma ^{}+2\mathrm{}S_F(\stackrel{}{r})(\sigma ^+\sigma ^{}1/2)`$ (5)
$`+\mathrm{}_0(a^++a)+\mathrm{}g(\stackrel{}{r})(a^+\sigma ^{}+\sigma ^+a).`$
Here $`\mathrm{\Delta }_a=\omega _a\omega _p`$ is the detuning of the atomic resonance from the probe frequency. In all numerical examples given below the cavity frequency is chosen to coincide with the atomic frequency, so that $`\mathrm{\Delta }_c=\mathrm{\Delta }_a`$. The quantity $`\mathrm{\Delta }_p\mathrm{\Delta }_a`$ is then referred to as the probe detuning. Note here that without a FORT the optimum cavity and atom detunings are not equal . In our case, however, the FORT effectively changes the atomic frequency in a position-dependent way and thus the precise value of the atomic detuning relative to the cavity detuning is largely irrelevant. Indeed optimum cooling conditions will exist in certain wells but not in others, which is one feature that allows one to distinguish various wells.
Coupling the atom and the cavity to the remaining modes of the electro-magnetic field leads by a standard procedure to the master equation for the density operator of the coupled atom-cavity system,
$`{\displaystyle \frac{\mathrm{d}\rho }{\mathrm{d}t}}=i[H,\rho ]/\mathrm{}\kappa \{a^+a,\rho \}+2\kappa a\rho a^++`$ (6)
$`{\displaystyle \frac{\mathrm{\Gamma }}{2}}\{\sigma ^+\sigma ^{},\rho \}+{\displaystyle \frac{3\mathrm{\Gamma }}{8\pi }}{\displaystyle }\mathrm{d}^2\widehat{k}{\displaystyle \underset{\widehat{ϵ}}{}}(\widehat{d}\widehat{ϵ})^2\mathrm{exp}(i\stackrel{}{k}\stackrel{}{r})\times `$ (7)
$`\times \sigma ^{}\rho \sigma ^+\mathrm{exp}(i\stackrel{}{k}\stackrel{}{r}),`$ (8)
with $`\mathrm{\Gamma }`$ the spontaneous decay rate and $`\kappa `$ the cavity decay rate. We are mainly interested in the strong-coupling regime, where $`g_0\mathrm{\Gamma },\kappa `$.
We treat the external (center-of-mass) degrees of freedom of the atom classically, an approximation justified at the end of Section IV C. For a discussion of various interesting effects arising from the quantized external motion of an atom in a cavity QED field, we refer the reader to .
In the quasiclassical approximation (i.e., where we retain the full quantum character of the internal degrees of freedom and of the cavity mode; see for a full discussion of this approximation), the integral in (6) can be evaluated to give the simpler result
$`{\displaystyle \frac{\mathrm{d}\rho }{\mathrm{d}t}}`$ $`=`$ $`i[H,\rho ]/\mathrm{}\kappa \{a^+a,\rho \}+2\kappa a\rho a^+`$ (10)
$`{\displaystyle \frac{\mathrm{\Gamma }}{2}}\{\sigma ^+\sigma ^{},\rho \}+\mathrm{\Gamma }\sigma ^{}\rho \sigma ^+.`$
The force acting on the atom consists of two parts, one due to spontaneous emission, whose mean vanishes on average, and the other part is represented by the operator
$`\stackrel{}{F}`$ $``$ $`\stackrel{}{}H`$ (11)
$`=`$ $`2\mathrm{}\stackrel{}{}S_F(\sigma ^+\sigma ^{}1/2)\mathrm{}\stackrel{}{}g(a^+\sigma ^{}+\sigma ^+a),`$ (12)
which has contributions arising from the FORT potential and from the interaction with the cavity mode. It was only the latter part that was considered in and that leads to 1-D cooling to temperatures of the order $`k_BT\mathrm{min}(\mathrm{}\kappa ,\mathrm{}\mathrm{\Gamma }/2)`$. See also Refs. for similar calculations on single atoms moving in cavity QED field, and Refs. for calculations of diffusion of atoms in optical traps in free space.
It can be shown starting from a fully quantized description, that the semiclassical motion of the atom is described by a Fokker-Planck equation for the Wigner distribution function containing (position-dependent) friction and diffusion coefficients. Equivalently, we may use stochastic equations for the classical atomic position and velocity variables $`\stackrel{}{r}`$ and $`\stackrel{}{v}`$ of the form
$`\mathrm{d}\stackrel{}{r}`$ $`=`$ $`\stackrel{}{v}\mathrm{d}t,`$ (13)
$`\mathrm{d}\stackrel{}{v}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{F}}{m}}\mathrm{d}t\beta \stackrel{}{v}\mathrm{d}t+B\mathrm{d}\stackrel{}{W},`$ (14)
where $`.`$ denotes an expectation value, $`\beta `$ is the friction tensor (with dimensions of a rate), $`m`$ the mass of the atom, $`B`$ is a tensor such that $`D=BB^T/2`$ is the velocity diffusion tensor (with dimension m<sup>2</sup>/s<sup>3</sup>), and $`\mathrm{d}\stackrel{}{W}`$ is a 3-dimensional Wiener process which satisfies $`\mathrm{d}W_i\mathrm{d}W_j=\delta _{ij}\mathrm{d}t`$ . Starting with the expression (11) for the force operator, we can calculate $`\beta `$ and $`D`$ by the procedure outlined in the next Section.
## III Friction and diffusion
Refs. employ Heisenberg equations of motion for various field and atomic operators to find friction and diffusion coefficients. These equations are not closed and, consequently, an approximation has to be made in order to find solutions. The natural assumption is to consider the weak driving limit (i.e., $`_0\kappa `$) and truncate the available Hilbert space to that part containing no more than a single cavity photon. This allows one to close the Heisenberg equations . Here we employ a different method (using the density matrix equations) to calculate friction and diffusion coefficients that does not require us to stay within the weak driving limit, but in addition we used Ref. ’s procedure here to obtain results in the weak driving limit for verification purposes. In any case, it is still true that the most interesting regime is where only one or few photons are involved. Note that given the strong coupling between atom and cavity field, even a single photon is sufficient to lead to regimes far beyond the weak driving limit. In our examples we truncated the Hilbert space to photon numbers of around 4 or smaller. We refer to for an exposition on how to represent operators in truncated Hilbert spaces of precisely this form in a numerically convenient manner.
The master equation (10) is written as
$$\frac{\mathrm{d}\rho }{\mathrm{d}t}=\rho .$$
(15)
Numerically, the Liouvillian superoperator $``$ is converted into a pre-multiplication operator by methods explained in . In order to find friction and diffusion coefficients we apply a simple procedure, which yields these coefficients at zero velocity: this is sufficient for our purposes as the atom we are interested in, Cs, is relatively heavy. More precisely, the relevant dimensionless parameters determining the velocity dependence of friction and diffusion coefficients are $`kv/\mathrm{\Gamma }`$ and $`kv/\kappa `$ (see for instance ), and both are very small in all our simulations. In particular, $`\mathrm{\Gamma }/k4.3`$ m/s and $`\kappa /k3.4`$ m/s, while velocities in the trapping regime we are interested in (where atoms are localized in wells at low temperatures for times $`\kappa ^1,\mathrm{\Gamma }^1`$) are around the Doppler limit velocity
$$v_D=\sqrt{\frac{\mathrm{}\mathrm{\Gamma }/2}{m}}8.8\mathrm{cm}/\mathrm{s}.$$
(16)
Also note that the standard procedure of continued fractions to calculate the full velocity dependence is not directly applicable to the present case, as the potential through which the atom is moving is not periodic ($`\lambda _F\lambda _0`$).
For an atom moving at velocity $`\stackrel{}{v}`$ we write
$$\frac{\mathrm{d}}{\mathrm{d}t}=\frac{}{t}+\stackrel{}{v},$$
(17)
and expand (15) in powers in $`\stackrel{}{v}`$ and solve for the steady state. The zeroth-order solution is then the steady state $`\rho _0`$ at zero velocity:
$$\rho _0=0,$$
(18)
while the first-order term $`\rho _1`$ is determined by
$$\rho _1=\stackrel{}{v}\rho _0.$$
(19)
The zeroth-order force is the steady-state force for an atom at rest, and is given by
$$\stackrel{}{F}_0=\mathrm{Tr}(\rho _0\stackrel{}{}H).$$
(20)
Similarly, the friction coefficients follow from the first-order term in the force
$$\stackrel{}{F}_1=\mathrm{Tr}(\rho _1\stackrel{}{}H),$$
(21)
by identifying
$$\stackrel{}{F}_1\beta m\stackrel{}{v},$$
(22)
where $`\beta `$ is a 3-by-3 tensor. In our case (), the gradients along the cavity axis are larger in magnitude than those in the transverse directions by roughly a factor $`kw_0150`$ (and around the cavity axis where the atoms spend most of their time the radial gradients are even smaller, of course). Since the friction coefficient scales with the product of two gradients (cf. Eqs (19) and (21)), the largest element of the tensor $`\beta `$ is the $`zz`$ component. Next largest in magnitude are the off-diagonal components such as $`\beta _{xz}`$ and $`\beta _{zx}`$. Their effects, however, can be safely neglected in our case: firstly, the force in the $`z`$ direction proportional to $`\beta _{zx}v_x`$ is smaller than the friction force $`\beta _{zz}v_z`$ by roughly a factor $`kw_0`$. Secondly, the force in the $`x`$ direction $`\beta _{xz}v_z`$ is not a friction force (as it is not proportional to $`v_x`$), and its contribution is averaged out because the oscillations in $`v_z`$ are faster than those in the $`x`$ direction by another factor $`kw_0`$. Finally, the purely radial friction rates such as $`\beta _{xx}`$ are too small ($`1`$s<sup>-1</sup> on average) to have any influence on the time scales considered here. Thus we take only $`\beta _{zz}`$ into account.
The diffusion coefficient, again at zero velocity, is calculated as follows. The standard method is to use the quantum regression theorem, and a particularly useful (for numerical purposes) interpretation of that theorem is given in . The momentum diffusion tensor $`D_p`$ is given by
$$D_p=\underset{t\mathrm{}}{lim}\mathrm{Re}_0^{\mathrm{}}d\tau \stackrel{}{F}(t)\stackrel{}{F}(t\tau )\stackrel{}{F}(t)\stackrel{}{F}(t\tau ),$$
(23)
and its relation to the velocity diffusion tensor is $`D=D_p/m^2`$. Before eliminating any degrees of freedom, the total system in fully quantized form is described by a time-independent Hamiltonian, which we denote by $`H_{\mathrm{tot}}`$. In that case the time evolution of all operators is determined by $`\mathrm{exp}(iH_{\mathrm{tot}}t)`$, and two-time averages of the form $`A(t)B(t\tau )`$ as appearing in (23) can be written as
$$A(t)B(t\tau )=\mathrm{Tr}\left[A\mathrm{exp}(iH_{\mathrm{tot}}\tau )B\rho _{\mathrm{tot}}(t)\mathrm{exp}(iH_{\mathrm{tot}}\tau )\right],$$
(24)
with $`\rho _{\mathrm{tot}}`$ the density matrix of the total system. This expression formally contains the evolution of a density matrix over a time interval $`\tau `$ starting from an initial density matrix $`\rho _{\mathrm{init}}B\rho _{\mathrm{tot}}(t)`$. The quantum regression theorem now states that (24) is still valid for the reduced density matrix that evolves under the Liouvillian $``$. That is, instead of (24) we may use
$$A(t)B(t\tau )=\mathrm{Tr}\left[A\mathrm{exp}(\tau )B\rho (t)\right].$$
(25)
In our case, $``$ is a time-independent operator and hence the right-hand side of (25) can be evaluated by expanding $`\mathrm{exp}(\tau )`$ in an exponential time series, as in the methods developed in . This then is the method we use here to evaluate the friction and diffusion tensors, and the results have been checked in the low-intensity limit by applying the different methods from to the same problem.
Diffusion due to spontaneous emission is not obtained this way (as the bath of vacuum modes has been eliminated already), but can be obtained by standard methods and gives an independent additional three components $`(D_p)_{ii}^{SE}=N_i\mathrm{}^2k^2\mathrm{\Gamma }/2\sigma ^+\sigma ^{}_0`$ for $`i=x,y,z`$, with $`._0`$ denoting a steady-state value and with the dimensionless factor $`N_i`$ depending on polarization. When the two-level system is formed by two Zeeman levels that are connected by circularly polarized light propagating in the $`z`$ direction, we have $`N_z=2/5`$, and $`N_x=N_y=3/10`$.
Since the diffusion coefficients, just as the friction coefficients, scale as the square of a gradient, the largest component is $`D_{zz}`$. Off-diagonal elements such as $`D_{xz}`$ and $`D_{zx}`$ are, again, smaller by roughly a factor $`kw_0150`$, while the diagonal radial components such as $`D_{xx}`$ are in fact largely determined by spontaneous emission, and are of similar or larger magnitude than the off-diagonal elements. The proper way to take into account the off-diagonal elements of the diffusion tensor $`D`$ is to diagonalize $`D`$, and consider 3 independent diffusion processes along the axes of the basis that diagonalizes $`D`$ with the eigenvalues of $`D`$ as diffusion coefficients. Using the fact that $`D_{zz}`$ is large we can calculate both eigenvalues and eigenbasis perturbatively. The eigenvalues to first order are given by
$`D_{x^{}x^{}}`$ $`=`$ $`D_{xx}{\displaystyle \frac{D_{xz}D_{zx}}{D_{zz}}}+\mathrm{}`$ (26)
$`D_{z^{}z^{}}`$ $`=`$ $`D_{zz}+{\displaystyle \frac{D_{xz}D_{zx}}{D_{zz}}}+\mathrm{},`$ (27)
where the $`\mathrm{}`$ stands for terms of higher order in $`1/(kw_0)`$, while the axes change as
$`\widehat{z}^{}`$ $`=`$ $`\widehat{z}+\widehat{x}{\displaystyle \frac{D_{xz}}{D_{zz}}}+\mathrm{}`$ (28)
$`\widehat{x}^{}`$ $`=`$ $`\widehat{x}+\widehat{z}{\displaystyle \frac{D_{zx}}{D_{zz}}}+\mathrm{}.`$ (29)
The fact that $`\widehat{z}^{}`$ is slightly tilted towards the $`x`$ direction implies that a small part of the large diffusion coefficient $`D_{z^{}z^{}}`$ will contribute to diffusion in the $`x`$ direction. This increase, however, is almost exactly compensated for by the decrease in $`D_{x^{}x^{}}`$. In particular, the velocity in the $`x`$ direction undergoes the following Wiener process:
$`\mathrm{d}v_x=\sqrt{2D_{x^{}x^{}}+2D_{z^{}z^{}}{\displaystyle \frac{D_{zx}^2}{D_{zz}^2}}+\mathrm{}}\mathrm{d}W.`$ (30)
In our case it turns out that $`D_{xx}D_{zz}D_{zx}^2`$ (see Fig. 3), so that effects due to the off-diagonal elements of the diffusion tensor can in fact be neglected. The figure also shows that the previous considerations about the relative sizes of the various components of $`D`$ do not just hold on average, but also locally.
Thus, friction is appreciable only along the cavity axis, while diffusion has two main contributions: from spontaneous emission in all three directions, and a large diffusion along the cavity axis from fluctuations in the FORT and cavity QED forces.
## IV Numerical results
The following results pertain to a Cs atom, with the ground state given by $`|6S_{1/2};F=4;m_F=4`$ and the excited state by $`|6P_{3/2};F=5;m_F=5`$, so that $`\lambda _0=852.4`$nm and $`\mathrm{\Gamma }/(2\pi )=5.2`$MHz. The cavity parameters are $`\kappa /(2\pi )=4`$MHz and $`g_0/(2\pi )=30`$MHz, and $`w_0=20\mu `$, which are typical for the experiments discussed in . Furthermore, the values for $`S_0`$ examined here are $`S_0/(2\pi )=10,50`$MHz. Both of these values are close to those explored in the actual experiment , and they contrast the behavior of atoms in shallow ($`S_0<g_0`$) and deep ($`S_0>g_0`$) wells. Typical values for $`N_e`$ range from $`10^3`$ to 0.1.
### A Dressed state structure
We first focus on the atomic motion along the cavity axis. The simplest way to get a feeling for the results for $`\beta _{zz}`$ and $`D_{zz}`$ as a function of the probe detuning $`\mathrm{\Delta }_p`$ is to first consider the eigenenergies of the dressed atom-cavity states. When we neglect dissipation for the moment, and take the limit of no driving ($`N_e=0`$), we can easily find the energies of the lower dressed states $`|\psi _\pm `$ containing at most one excitation: the state containing no excitation is the ground state with an energy of $`E_0=\mathrm{}S_F(\stackrel{}{r})`$, while the energies of the two dressed states in the manifold of states containing a single excitation are
$$E_\pm =\mathrm{}\omega _a\pm \mathrm{}\sqrt{g(\stackrel{}{r})^2+S_F(\stackrel{}{r})^2},$$
(31)
if the atom and cavity are on resonance. The excited dressed states are given by
$$|\psi _{}=\mathrm{sin}\theta |g,1+\mathrm{cos}\theta |e,0,$$
(32)
with
$`\mathrm{sin}\theta `$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{g^2+(\sqrt{g^2+S_F^2}S_F)^2}}},`$ (33)
$`\mathrm{cos}\theta `$ $`=`$ $`{\displaystyle \frac{S_F\sqrt{g^2+S_F^2}}{\sqrt{g^2+(\sqrt{g^2+S_F^2}S_F)^2}}}.`$ (34)
In Figures 4 (10MHz FORT) and 5 (50MHz FORT) we plot the transition frequencies (relative to $`\omega _a`$) from the ground state to these two excited states as functions of position, i.e.,
$$\mathrm{\Delta }_\pm =S_F(\stackrel{}{r})\pm \sqrt{g(\stackrel{}{r})^2+S_F(\stackrel{}{r})^2}.$$
(35)
This expression along with the figures explicitly shows that the main features of the atom-cavity system are determined by the ratio $`S_0/g_0`$. It furthermore shows an important difference with the situation of trapping with a FORT in free space. The fact that the excited state shifts up while the ground state shifts down implies that ground and excited states are trapped in different positions in free space. In the presence of the quantized cavity field, however, both the lower excited dressed state and the ground state are now shifted down. This may improve trapping and cooling conditions, as detailed below.
### B Cooling mechanisms
We now take a closer look at cooling mechanisms. In the regime of weak driving, we will find that the friction coefficient $`\beta _{zz}`$ is positive (corresponding to cooling) when the probe field is tuned slightly (by an amount $`\kappa ,\mathrm{\Gamma }/2`$) below the transition to the relevant dressed state while for blue detuning the friction coefficient is negative, leading to exponential heating of the atom’s velocity. This can be understood by analogy with Doppler cooling: by tuning below resonance, the process of stimulated absorption followed by spontaneous emission leads to a loss of energy, while the maximum cooling rate is achieved by maximizing the product of excitation rate and detuning. Now looking back to Figs. 4 and 5 one sees that the variation of $`\mathrm{\Delta }_+`$ with position is larger than that of $`\mathrm{\Delta }_{}`$, because both the ground state and the lower excited dressed state shift down, while the upper excited dressed state shifts upward. Generally speaking, for cooling purposes it is better to tune to the lower excited state so as to have smaller spatial variations in cooling rates. More importantly, the upper excited state energies decrease with increasing radial distance, whereas the lower excited state energy increases. Thus, for the upper state the probe detuning changes from red to blue, so that an atom cooled on axis will in fact be heated if it moves away radially. For the lower dressed state the probe detuning becomes more red, so that an atom that is optimally cooled on axis will still be cooled away from the axis, but at a lower rate.
The most popular explanation for intra-cavity cooling exploits analogies with Sisyphus cooling , although another explanation for cavity-based cooling based on asymmetries in coherent scattering was recently put forward in . Here we illustrate the Sisyphus picture for cooling inside optical wells within an optical resonator, using a very simple dressed-state picture, that makes use of only the lower dressed state and the ground state, relevant in the low-excitation limit. We choose one particular well, from $`z=2.0\lambda _F`$ to $`z=2.5\lambda _F`$, and one particular set of parameters given in the caption of Fig. 6. In that figure we plot the decay rate $`\gamma _{}`$ of the lower dressed state and the excitation rate from ground to the dressed state, $`\mathrm{\Omega }_{}`$, as functions of position. In the weak driving limit the decay rate is given by
$`\gamma _{}`$ $`=`$ $`\psi _{}|\kappa a^+a+\mathrm{\Gamma }\sigma ^+\sigma ^{}/2|\psi _{}`$ (36)
$`=`$ $`\mathrm{sin}^2\theta \kappa +\mathrm{cos}^2\theta \mathrm{\Gamma }/2,`$ (37)
and the excitation rate by
$`\mathrm{\Omega }_{}`$ $`=`$ $`|g,0|_0(a^++a)|\psi _{}|`$ (38)
$`=`$ $`_0|\mathrm{sin}\theta |.`$ (39)
These two quantities, together with the detuning of the probe field from (dressed-state) resonance determine the steady-state population in the lower dressed state, according to
$$n_{}=\frac{\mathrm{\Omega }_{}^2}{(\mathrm{\Delta }_{}\mathrm{\Delta }_p)^2+\gamma _{}^2}.$$
(40)
The population $`n_{}`$ is plotted in Fig. 7, along with the transition frequency $`\mathrm{\Delta }_{}`$. These two quantities are sufficient to understand the Sisyphus cooling mechanism.
Since an atom in the ground state is moving in a conservative potential around the equilibrium position $`z=2.25\lambda _F`$, the following Sisyphus picture should be taken as to apply to the motion of the atom in addition to that conservative motion (see (41)). Suppose, for example, that the atom is at position $`z=2.2\lambda _F`$ and moving towards the right (cf. Fig. 7). The probability to be in the excited state now decreases (according to the lower part of Fig. 7), while the energy of the excited state relative to the ground state is increasing: in other words, an atom in the excited state is climbing uphill (again, in relation to the ground state), but will likely make the down transition to the ground state, thus leading to cooling at that particular position. Similarly, at $`z=2.4\lambda _F`$ an atom moving to the left is going uphill while having an increased chance of decaying to the ground state, again leading to cooling. This picture in fact shows that the cooling rate is expected to be proportional to the gradient of $`n_{}`$ and the gradient of $`\mathrm{\Delta }_{}`$. More precisely, the force on the atom at position $`z`$ is approximately given by
$`F_z`$ $``$ $`\mathrm{}{\displaystyle \frac{dS_F}{dz}}\mathrm{}n_{}(zv/\gamma _{}){\displaystyle \frac{d\mathrm{\Delta }_{}}{dz}}`$ (41)
$``$ $`\mathrm{}{\displaystyle \frac{dS_F}{dz}}\mathrm{}n_{}(z){\displaystyle \frac{d\mathrm{\Delta }_{}}{dz}}+{\displaystyle \frac{\mathrm{}v}{\gamma _{}}}{\displaystyle \frac{dn_{}}{dz}}{\displaystyle \frac{d\mathrm{\Delta }_{}}{dz}},`$ (42)
where the argument of $`n_{}`$ indicates the lag between the atom reaching a position $`z`$ and reaching its steady state, with the lag time scale determined by the inverse decay rate from the dressed state. From the second line we see that the friction coefficient $`\beta _{zz}`$ is approximated by
$$R\frac{\mathrm{}}{m\gamma _{}}\frac{dn_{}}{dz}\frac{d\mathrm{\Delta }_{}}{dz}.$$
(43)
Indeed, Fig. 8 shows the similar behavior of $`\beta _{zz}`$ and $`R`$ as functions of position.
### C Friction, diffusion and equilibrium rms velocities
In Figures 910 we give examples of friction and diffusion coefficients for both the 10 and 50 MHz FORTs, as functions of the atomic position. They illustrate the point that in the low-excitation limit red (blue) detuning leads to cooling (heating) (cf. Figs. 4 and 5). They moreover clearly show how all wells are quantitatively different, with cooling rates and diffusion strengths differing by orders of magnitude over the various wells, and with $`\beta _{zz}`$ being negative in some wells, and always positive in others. This of course also implies that the temperatures reached by atoms in thermal equilibrium vary with position.
For the case of the shallow FORT we consider weak driving ($`N_e=0.001`$), whereas for the deeper FORT the driving field is taken to be stronger by an order of magnitude. The stronger driving field increases cooling rates while the fact that deeper wells trap the atoms better means that correspondingly larger diffusion rates still can be tolerated.
The stable equilibrium points $`z_n^e`$ are located around the maxima of $`S_F`$, i.e. around $`z_n=(n1/2)\lambda _F/2`$ for integer $`n`$, because it is the FORT that gives the main contribution to the total force (even for the smallest value of $`S_0=2\pi \times 10`$MHz considered here). The cavity QED field gives only a small correction to the force and hence to the equilibrium position. In each equilibrium point, we can define a measure for the expected rms velocity of the atom along the $`z`$ axis in thermal equilibrium by considering averages over local wells
$$v_{\mathrm{rms}}^z=\sqrt{\frac{\overline{D}_{zz}}{\overline{\beta }_{zz}}}\mathrm{if}\overline{\beta }_{zz}>0,$$
(44)
in terms of the friction and diffusion coefficients. This averaging procedure gives a sensible measure for the rms velocity only if the atom indeed samples the whole well. This condition is fulfilled for the relatively shallow wells originating from $`S_0=2\pi \times 10`$MHz, and Fig. 11 uses this averaging procedure. For the 50MHz FORT, however, we averaged over only part of the well, namely a region of size $`\lambda _F/10`$ symmetrically around the equilibrium point. This choice is rather arbitrary, and thus Fig. 12 just gives an indication of what rms velocities to expect for atoms trapped in the corresponding wells, although the simulations in fact do confirm these values.
We see here that depending on the probe detuning, the atom will be cooled to low temperatures either in all wells, or only in wells where $`g`$ is large in the equilibrium point, or only in wells where $`g`$ is small. This shows the flexibility that a FORT beam adds: one can predetermine to a certain degree in which well the atom will be trapped (and cooled) for longer times and in which it will not be.
Under the current conditions $`\kappa >\mathrm{\Gamma }/2`$ the lowest temperatures achievable are determined by the Doppler velocity $`v_D`$. More precisely, the lower limit on rms velocities along the cavity axis is expected to be
$$v_D^z=\sqrt{0.7\frac{\mathrm{}\mathrm{\Gamma }}{2m}},$$
(45)
where the factor $`0.7=(1+2/5)/2`$ comes from the fact that in our case the diffusion due to spontaneous emission in the $`z`$ direction is 2/5th of the full 3-D value. We tested that for smaller $`\kappa `$ the rms velocities indeed do become even smaller, now determined by $`\sqrt{\mathrm{}\kappa /m}`$, thus confirming predictions of .
Finally, we note that the quasiclassical approximation used throughout this paper is justified as neither the recoil limit is reached nor the resolved-sideband limit, i.e.
$`\mathrm{}\mathrm{\Gamma }/2`$ $``$ $`(\mathrm{}k)^2/m,`$ (46)
$`\mathrm{}\mathrm{\Gamma }/2`$ $``$ $`h\nu _{\mathrm{osc}},`$ (47)
with $`\nu _{\mathrm{osc}}`$ the oscillation frequency of the atom in a well (see below), although in some cases the latter condition is only marginally fulfilled, namely when $`\nu _{\mathrm{osc}}=600`$kHz, which is only a factor 4 smaller than $`\mathrm{\Gamma }/(4\pi )`$.
### D Saturation behavior
We now briefly turn to the question of the nonlinear behavior of the atom-cavity system with increasing excitation. In the absence of saturation effects, both friction and diffusion coefficients would increase linearly with $`N_e`$. For the same parameters as Fig. 9, Figure 13 shows nonlinearities setting in around $`N_e=0.01`$. The friction coefficient even starts to decrease around $`N_e=0.1`$ as a result of the local values of $`\beta _{zz}`$ becoming negative where they were positive in the weak driving limit. The concomitant effect on the $`v_{\mathrm{rms}}^z`$ is shown as well.
### E Simulations
We also performed Monte-Carlo simulations of the 3-D motion in given wells by solving the Langevin equations (13) for position and velocity (see also ). The experimental procedure switches the FORT field on only when an atom has been detected and when it consequently has partly fallen through the cavity already . We accordingly fix initial conditions as follows: We start the atom on the cavity axis, and we fix the downward velocity to be $`v_x=10`$ cm/s. Furthermore, we chose $`v_z=0`$ cm/s, and the initial position along the $`z`$ axis to be $`\lambda _F/8`$ away from the equilibrium point. The initial position and velocity were fixed so that all variations in trapping times and rms velocities are solely due to the random fluctuations of the forces acting on the atom, rather than from random initial conditions. Experimentally these two are mixed of course.
Since atoms with these initial conditions do not possess angular momentum around the $`z`$ axis, this in some sense represents a favorable case (although the atoms are not put in the bottom of the well). However, in the course of their evolution the atoms do acquire angular momentum so that this is in fact not a severe restriction. For more detail see below (Figure 22).
In Figure 14 we plot the results of simulations of 1000 trajectories for an atom in the shallow well of 10 MHz. We plot the average rms velocity along the cavity axis as a function of trapping time for each trajectory. Here we defined the “trapping time” as the time spent by the atom in one particular given well of size $`\lambda _F/2`$. The actual trapping time inside the cavity may be longer, obviously, as the atom may subsequently get trapped in different wells. For very short trapping times, $`v_{rms}`$ is determined by the initial condition, but for longer times lower temperatures corresponding to those calculated in Fig.11 are reached. Note however that the simulations were done in 3-D, and as such do not necessarily give the same temperatures as predicted for on-axis (1-D) motion in Figures 11 and 12. Nevertheless, the effect of the atoms’ radial motion is apparently not strong, and in fact atoms leave the well while still being trapped radially. This is partly due to the fact that all (especially heating) rates in the radial direction are smaller by a factor $`kw_0150`$ than those in the axial direction.
About half of the atoms is basically not trapped at all. The remaining atoms have a probability $`P(T)`$ to be trapped longer than a time $`T`$, with $`P(T)`$ decaying exponentially with $`T`$. The average trapping time for these parameters is found to be $`\tau `$ 25ms, as shown in Fig. 14.
In Figures 15 and 16 we plot for the same 10MHz FORT an example of a single trajectory, after the atom has spent 4ms in the trap. The oscillation frequencies along the $`z`$ and the radial directions differ by two orders of magnitude (since $`kw_0150`$): in the $`z`$ direction the oscillation rate is $``$ 200kHz, in the radial direction $``$ 2.2kHz. The photon transmission follows both these oscillations so that in principle the atomic motion in both axial and radial direction is detectable. Experimentally, though, the oscillations along the cavity axis may be too fast to be accessible. In particular, the average rate at which photons leaking out through one end of the cavity are detected is at most (the efficiency is less than 100%) equal to the cavity decay rate multiplied by the average number of photons inside the cavity. For the parameters of Fig. 15 this amounts to a rate $`0.01\times \kappa 2.5\times 10^5`$/sec, which corresponds to just about one photon per oscillation period.
The figures show that when the atom is in a position where it is not coupled to the cavity ($`g=0`$), the number of photons in the cavity drops to $`N_e=0.001`$. Similarly, when the atom moves away radially, the transmission drops.
To make a direct comparison with the trapping times achieved in the experiment , we now turn to the case of a 50 MHz FORT. We plot rms velocities vs trapping times for 300 trajectories for an atom trapped in the well ranging from $`z=2\lambda _F`$ to $`z=2.5\lambda _F`$.
For the parameters of Fig. 17 the atom is either trapped for long times ($`>10`$ms) or only for a short time ($`<1`$ms), both with about $`50\%`$ probability. In the latter case the rms velocity is determined just by the (arbitrarily chosen) initial condition and is around 30cm/s, but for longer trapping times the effects of cooling are visible. Thermal equilibrium is reached with $`v_{rms}8`$cm/s, thus confirming the results of Fig. 12. The distribution of trapping times again follows an exponential law, and the average trapping time, as determined from the tail of the distribution, is $`\tau 250`$ms, which is ten times longer than for the (fluctuating) 50MHz FORT used in . This shows the great potential of holding single atoms in the cavity for extended periods of time if the intensity fluctuations of the FORT beam can be minimized. Experimental efforts along this path are currently underway.
Also for this case we plot snapshots for a single trajectory, taken after the atom has spent 25ms in the trap. Compared to the 10MHz FORT, the oscillations of the atom along the cavity axis and in the radial direction become faster by about a factor of 3. The axial oscillation frequency is about 600kHz, while along the radial direction the oscillations occur at a rate 6.2kHz, i.e., again slower by two orders of magnitude. In this case, the photon transmission still follows directly the axial oscillations but no longer follows the radial excursions of the atom, as now the fluctuations in the magnitude of $`g`$ at the atom’s position along the cavity axis are in fact larger than those due to the radial excursions of the atom. This is partly due to the fact that in the simulations here the driving field is stronger than for the 10MHz example above so that fluctuations in the atomic motion occur at a shorter time scale, and partly simply because the radial excursions are small. Fig. 19(c) shows that it is primarily the axial fluctuations that determine the variations in the numbers of photons inside the cavity.
Generally speaking, the axial excursions determine (local) minimum and maximum transmission levels (as in Fig. 15). When these minima and/or maxima depend on the radial position, then the radial motion could in principle be visible in the cavity transmission level. This depends in turn on whether the axial fluctuations on the time scale of the transverse motion are sufficiently small so as not to hide the radial dependence. There seems to be no simple general rule how this interplay between radial and axial motions depends on detunings, driving strength, and the particular well.
In contrast, in a different well, the one ranging from $`z=0.0`$ to $`z=0.5\lambda _F`$, the photon number in the cavity does follow the radial motion, as the radial excursions become larger. Perhaps more importantly, the average transmission level is higher by more than a factor 2 compared to the previous case, as a result of $`g`$ being larger in this well (cf. Fig. 2). This shows how, in principle, different wells may be experimentally distinguished via the transmission of the probe field through the cavity.
We also simulated the motion of an atom trapped under more adverse conditions, namely for an atom in the well $`[z=\lambda _F1.5\lambda _F]`$ at a probe detuning $`\mathrm{\Delta }_p/(2\pi )=5`$MHz. According to Fig. 12, the atom is not cooled on axis under these conditions (i.e. the average friction coefficient around the equilibrium point on the $`z`$ axis is negative). This is confirmed by Fig. 21: the mean trapping time for an atom starting at $`z=1.125\lambda _F`$ is now very short, about 1.6ms, while the average rms velocity is $`v_{\mathrm{rms}}^z28`$ cm/s, as determined essentially by the initial condition.
Finally, we consider the influence of different initial conditions on trapping and cooling. All the results so far were obtained by considering atoms that initially are moving on axis. Thus, they have no angular momentum along the $`z`$ axis, nor any radial potential energy. Figure 22 shows a plot of rms velocities vs. trapping times for atoms trapped under the same conditions as for Fig. 17) (i.e., in the well from $`z=2\lambda `$ to $`z=2.5\lambda `$, for $`\mathrm{\Delta }_p=10\times 2\pi `$MHz, $`N_e=0.01`$, and $`S_0=50\times 2\pi `$MHz), but with different (nonzero) values for the initial angular momentum.
Obviously, the more initial potential energy the atom has, the less likely it is to be trapped. In fact, the angular momentum does not play any role here, as confirmed by similar calculations with initial conditions chosen such that the atoms have no initial angular momentum but have the same potential energy. The results are the same in that case. For atoms starting at $`y=0.2w_0`$ the trapping times and rms velocities are basically not affected, and the trapping time is still around 250ms. But for atoms starting at $`y=0.5w_0`$ the effect of their increased potential energy leads to clearly shorter trapping times (by roughly a factor of 2), and for atoms starting at $`y=w_0`$ this effect is even more pronounced with a decrease in trapping time of about a factor of 10.
### F A different trapping structure
We now consider a different case where the atomic excited state is assumed to be shifted down by the FORT field, just as the ground state is (see, for instance ). This can be achieved by using a FORT that is (red) detuned in such a way that the excited atomic state is relatively closer to resonance with a higher-lying excited state than with the ground state. This situation at first sight looks even more appealing for trapping purposes, as now both excited and ground state will be trapped in the same positions. Moreover, fluctuations in the force due to the FORT are dimished.
We consider only the 50MHz FORT here, and compare this case to the previous 50MHz FORT case, and in particular we refer the reader back to Figs. 5, 10, and 17. For ease of comparison we keep $`\lambda _F`$ the same, and assume for simplicity that the excited state is shifted down by an amount $`S_F`$, so that the shifts of the ground and excited state are in fact identical.
The fact that ground and excited state see the same potential, implies that the transition frequencies to the dressed states are simply periodic in space with period $`\lambda _0`$, as shown in Fig. 23, rather than aperiodic as in Fig. 5.
Similarly, the fluctuations in the force due to the FORT now vanish, as both ground and excited state undergo the same shift, so that the diffusion coefficient is periodic with period $`\lambda _0`$. Also the friction force arises only from the cavity QED part and is periodic. Yet, the different wells are not equivalent. The forces are, just as before, driven by both cavity QED field and the FORT, and the value of $`g`$ at the antinode of the FORT still varies over the different wells. This is illustrated in Fig. 24 where the rms velocities in the 8 different wells are shown, along with friction and diffusion coefficients. Since in this example the probe field is detuned below the lower dressed state, one has cooling everywhere in space.
The simulations show that the mean trapping time is smaller, although the rms velocities are just as small as before. The reason is the less favorable cooling condition away from the cavity axis. In particular, for the parameters used here the expected rms velocity $`v_{\mathrm{rms}}^z`$ steadily increases to 90cm/s at a radial distance $`\rho =2w_0`$, while for the simulations of Figs. 17, $`v_{\mathrm{rms}}^z`$ is increasing only slowly to 12cm/s. This large difference can be understood by noting the difference in dressed state structures between the two cases. For the case of Fig. 5, the transition frequency to the lower dressed state around the equilibrium position $`z2.25\lambda _F`$ does not change much with increasing radial distance, so that the probe field in that trapping region is always detuned below resonance by an amount that stays more or less constant. For the dressed state structure of Fig. 23, however, the probe detuning increases from $`5`$MHz to $`35`$MHz below resonance, thus leading to much worse cooling conditions. In other words, the presence of opposite level shifts due to the FORT makes the spatial variation of the transition frequency to the lower dressed state smaller: compare $`\mathrm{\Delta }_{}=S_F\sqrt{S_F^2+g^2}`$ to $`\mathrm{\Delta }_{}=g`$, especially when $`gS_F`$.
The alternative trapping potential is, therefore, not necessarily more favorable for trapping purposes. On the other hand, all atoms are captured now and are trapped for at least 10ms. This can be understood from the simple fact that here the friction coefficient is positive in the entire well.
## V Summary
We analyzed cooling limits and trapping mechanisms for atoms trapped in optical traps inside optical cavities. The main distinguishing feature from previous discussions on cooling of atoms inside cavities is the presence of the external trapping potential with a different spatial periodicity as compared to the cavity QED field. This not only provides better cooling and trapping conditions but the different spatial period makes the various potential wells qualitatively different. Atoms can be trapped in regions of space where the coupling to the cavity QED field is maximum, minimum or somewhere in between. Depending on the laser detuning, cooling may take place only in wells where the atom is minimally coupled to the cavity QED field, or where it is maximally coupled. This allows one in principle to distinguish to a certain degree the different atomic positions along the cavity axis, namely, by comparing
* the average transmission level
* the fluctuations of the cavity transmission
* the total trapping time
which reflect, respectively, the average atom-cavity coupling, the temperature of the atom and under certain conditions the radial motion, and the overall cooling and trapping conditions. This is an important additional tool useful for eventual control of coherent evolution of the atomic center-of-mass degrees of freedom, as relevant to performing quantum logic operations.
## Acknowledgements
We thank Andrew Doherty, Klaus Mølmer and David Vernooy for helpful discussions and comments. This work was funded by DARPA through the National Science Foundation and the QUIC (Quantum Information and Computing) program administered by the US Army Research Office and the Office of Naval Research.
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# Embeddings of relatively free groups into finitely presented groups
## 1 Introduction
By the Higman embedding theorem every finitely generated recursively presented group can be isomorphically embedded into some finitely presented group. Moreover the embedding can preserve the solvability of the word problem or even the solvability of the conjugacy problem . The embedding can be quasi-isometric and can at the same time dramatically improve the isoperimetric function of the group . Moreover if the word problem for a finitely generated group $`G`$ is solvable in non-deterministic time $`T(n)`$ by a Turing machine then $`G`$ is a quasi-isometric subgroup of a finitely presented group $`H`$ whose Dehn function is polynomially equivalent to $`T(n)`$ . Thus the word problem of a finitely generated group $`G`$ is in NP if and only if $`G`$ is a quasi-isometric subgroup of a finitely presented group with polynomial Dehn function .
There exists (see, for example, Valiev ) a relatively simple finite presentation of a group containing all recursively presented groups, but the embedding of a concrete finitely generated group in this universal group is not really effective. The embeddings mentioned in the previous paragraph are not really effective either. In both case one needs to know a precise description of a Turing machine (or a Diophantine equation or an $`S`$-machine ) “solving” the word problem in the original group in order to find the presentation of the bigger group or to find a copy of the given group inside the universal group.
The main goal of our paper is to show that very often an embedding of $`G`$ into $`H`$ can be given in a more explicit and straightforward way. In particular, the description of the set of defining words for $`H`$ is explicit. (More precisely the $`S`$-machines used in these constructions are so small that it is easy to write down all the commands, or simply forget about machines altogether and write down all the relations of the presentation of $`H`$.) We can also say a lot about the structure of $`H`$, and about the way $`G`$ is embedded into $`H`$.
In this paper, we consider two classes of groups. The first class consists of relatively free groups of finite ranks in varieties of groups. For example, we show that finitely generated free solvable, or free Burnside groups of sufficiently large exponents can be easily embedded into finitely presented groups with polynomial isoperimetric functions. Notice that the isoperimetric functions of these relatively free groups are not even defined because these groups are (as a rule) infinitely presented. In order to deal with relatively free groups, we introduce the so called verbal isoperimetric function of a relatively free group. This function seems to be interesting in itself.
The second class of groups consists of one relator metabelian Baumslag-Solitar groups $`BS_{k,1}=a,b|a^b=a^k`$. Here $`x^y`$ stands for $`y^1xy`$. It is well known that the Dehn function of the group $`BS_{k,1}`$ is exponential . It is also well known that these groups are very “stubborn”: they “resist” being embedded into groups with small isoperimetric functions. For example, there is a conjecture that the Dehn function of every 1-related group containing $`BS_{k,1}`$, $`k2`$ is exponential. In this paper, we present simple finite presentations of groups with polynomial isoperimetric functions which contain $`BS_{k,1}`$.
Let us give the necessary definitions. Let $`H`$ be a group given by a finite presentation $`P=x_1,\mathrm{},x_k|r_1,\mathrm{},r_{\mathrm{}}`$. A non-decreasing function $`f:𝐍𝐍`$ is called an isoperimetric function of this presentation if any word $`w=w(x_1,\mathrm{},x_k)`$ of length $`n`$, which is equal to $`1`$ in $`H`$, can be written in the free group $`F(x_1,\mathrm{},x_k)`$ as a product of at most $`f(n)`$ conjugates of the relators of $`H`$. In other words, $`f(n)`$ is an isoperimetric function of the presentation $`P`$ if every loop of length $`n`$ in the Cayley complex corresponding to the presentation $`P`$ has area at most $`f(n)`$. Isoperimetric functions of different finite presentations of the same group are equivalent in some natural sense , so one can talk about isoperimetric functions of the group $`H`$ forgetting about its presentations. We do not distinguish equivalent isoperimetric functions in this paper. Following Gersten, the smallest isoperimetric function of a group $`H`$ is called the Dehn function of $`H`$.
Now let us define the verbal isoperimetric functions. It is easy to see that any variety of groups $`𝒱`$ that can be defined by a finite set of identities, can be also defined by a single law $`v=1`$ for some word $`v=v(x_1,\mathrm{}x_k)`$ from the (absolutely) free group $`F`$ of infinite rank. The verbal subgroup $`VF`$ consists of all words vanishing in all groups of the variety $`𝒱`$. Any word $`wV`$ is freely equal to a product
$$\underset{i=1}{\overset{N}{}}u_iv(X_{i1},\mathrm{},X_{im})^{\pm 1}u_i^1$$
$`(1.1)`$
for some words $`u_i`$ and $`X_{ij}`$. We call a non-decreasing function $`f_v:𝐍𝐍`$ a verbal isoperimetric function of the word $`v`$, if for any word $`wV`$ there is a representation (1.1), where $`_{ij}|X_{ij}|f_v(|w|)`$. The smallest verbal isoperimetric function will be called the verbal Dehn function. Note that the verbal Dehn function exists because in the definition of verbal isoperimetric functions one may restrict oneself to words in variables $`x_1,\mathrm{},x_n`$ and there are only finitely many such words of any given length.
If an identity $`v^{}(x_1,\mathrm{},x_m^{})=1`$ is equivalent to the identity $`v=1`$ then every value of $`v^{}`$ is a product of a fixed number of values of the word $`v`$, and vice versa. Thus the verbal Dehn function of the variety $`V`$ does not depend on the choice of a defining law if one identifies functions which are $`\mathrm{\Theta }`$-equivalent. Recall that two functions $`f,g:𝐍𝐍`$ are $`\mathrm{\Theta }`$-equivalent if there exist positive constants $`c,d`$ such that for every $`n𝐍`$ we have:
$$f(n)cg(n)\text{ and }g(n)df(n).$$
Assume that a word $`w`$ is a product of two words $`w^{},w^{\prime \prime }`$ which depend on disjoint sets of variables, $`S^{}`$ and $`S^{\prime \prime }`$. Then a representation
$$w=u_iv(X_{i1},\mathrm{},X_{im})u_i^1$$
induces similar representations of $`w^{}`$ and $`w^{\prime \prime }`$ if we substitute 1 for the variables from $`S^{\prime \prime }`$ or from $`S^{}`$. For the corresponding words $`X_{ij}^{}`$ and $`X_{ij}^{\prime \prime }`$ one has the obvious inequality $`|X_{ij}^{}|+|X_{ij}^{\prime \prime }||X_{ij}|`$. Hence the following statement holds.
###### Proposition 1.1
Any verbal Dehn function $`f`$ is superadditive, i.e. $`f(n_1+n_2)f(n_1)+f(n_2)`$.
Remark. Recall that the superadditivity of any (usual) Dehn function is still an open problem (see Guba and Sapir ).
In this paper we produce explicite finite presentations of groups $`H(v,m)`$ to prove the following
###### Theorem 1.2
Let $`f(n)`$ be the verbal isoperimetric function of a group variety $`𝒱`$ defined by an identity $`v=1`$. Then the free group $`F_m(𝒱)`$ of rank $`m`$ in the variety $`𝒱`$ can be isomorphically embedded into a finitely presented group $`H=H(v,m)`$ with an isoperimetric function $`n^2f(n^2)^2`$. The presentation of $`H`$ is explicitely constructed given the word $`v`$, and the embedding is quasi-isometric.
Recall that a finitely generated subgroup $`G`$ of a finitely generated group $`H`$ is quasi-isometric (in other terminology, undistorted) if there is a positive constant $`C`$ such that $`|g|_GC|g|_H`$ for any $`gG`$. Here $`|g|_G`$, $`|g|_H`$ are the lengths of $`g`$ in the fixed finite sets of generators of $`G`$ and $`H`$ respectively. This notion is well defined, but the constant $`C`$ does depend on the choice of the finite generating sets.
Remark. Without the word “explicit” and with slightly worse estimate for the isoperimetric function of the bigger group this theorem is a corollary of the main result of . Indeed, it is easy to see that if $`f`$ is a verbal isoperimetric function of a group variety $`𝒱`$ then the word problem in any relatively free group in this variety can be solved by a non-deterministic Turing machine in time $`f(n)`$. Thus by the main result in , this group can be embedded into a finitely presented group with isoperimetric function $`n^2f(n^2)^4`$ (compare with the estimate $`n^2f(n^2)^2`$ in our theorem).
For example, let $`𝒱=_n`$ be the Burnside variety consisting of all groups satisfying the identity $`x^n=1`$ for a fixed large odd number $`n`$. Then the function $`f(s)=s^4`$ is a verbal isoperimetric function for $`𝒱`$. One can refer, for instance, to Storozhev’s argument in , $`\mathrm{\S }28.2`$: a word $`w`$ of length $`s`$ that is equal to 1 in the free Burnside group $`B(m,n)`$, is a product of $`n`$-th powers of some words $`A_i`$ where the lengths of $`A_i`$ are bounded by a linear function of $`s`$, and the number of factors is bounded by a cubic function. Thus the function $`s^{18}`$ is an isoperimetric function of $`H(v,m)`$ in this case. It can be proved by refining Storozhev’s argument that the function $`s^{1+ϵ(n)}`$ is a verbal isoperimetric function for $`_n`$ where $`ϵ(n)0`$ for $`n\mathrm{}`$ (R. Mikhajlov, unpublished). Therefore for very large odd $`n`$ the group $`H(x^n,m)`$ satisfies a verbal isoperimetric inequality with $`f(s)=s^{8+ϵ}`$ for a small $`ϵ>0`$.
The Schreier – Reidemeister rewriting shows by induction, that the variety of all solvable groups of derived length at most $`d`$ has a polynomial verbal isoperimetric function. It is interesting to calculate (up to the equivalence) the verbal Dehn function of these varieties, or the varieties of all nilpotent groups of given nilpotency classes, etc. Results of Yu. G. Kleiman show that there exists an identity defining a solvable group variety with non-recursive verbal Dehn function. Indeed, Kleiman constructed a finitely based solvable variety with undecidable identity problem: given a word $`w`$, it is impossible to decide whether the law $`w=1`$ follows from the defining laws of his variety. Clearly this implies that Kleiman’s variety cannot have a recursive verbal isoperimetric function.
Similar questions can be raised for the function which counts just the minimal number of factors $`N`$ in (1.1). Notice that from the decidability of the Diophantine theory of the free group (Makanin ), it follows that this function is recursive if and only if the corresponding verbal Dehn function is recursive.
###### Theorem 1.3
The group $`BS_{k,1}=a,b|b^1ab=a^k`$ is quasi-isometrically embeddable into a finitely presented group $`H_{k,1}`$ with isoperimetric function $`n^{10}`$.
This theorem gives an upper bound for the Dehn function of an appropriate finitely presented group containing $`BS_{k,1}`$ as a subgroup, and in contrast to , the presentation of $`H_{k,1}`$ is explicit and $`S`$-machines are used only in an implicite way.
## 2 Defining relations
Let $`𝒱`$ be the variety of groups defined by an identity $`v(x_1,\mathrm{},x_k)=1`$ where $`v`$ is a cyclically reduced word, that can be written as $`y_1\mathrm{}y_{M1}`$ with $`y_sx_{i_s}^{\pm 1}`$.
First we define an auxiliary group $`G(v,m)`$ which plays the same role as the group $`G_N(𝒮)`$ in (and bares some similarity with the Boone group from ). The set of generators of $`G(v,m)`$ consists of the letters $`a_1,\mathrm{},a_m`$; $`q_1,\mathrm{},q_M`$; $`k_1,\mathrm{},k_N`$ (where $`N29`$); $`r_{1,1},\mathrm{},r_{m,k}`$.
We present the set of relations of the group $`G(v,m)`$ in three different ways. First we just list all the relations. Then we will present these relations in terms of an $`S`$-machine, and finally we shall present a picture which will contain all the relations, and also show the main corollaries of the relations.
2.A. The list of relations. The finite set of defining relators for $`G(v,m)`$ is given below by equalities (2.1) - (2.6).
$$r_i^1q_{s+1}r_i=a_jq_{s+1}$$
$`(2.1)`$
if $`i=(j,l)`$ and $`y_sx_{\mathrm{}}`$;
$$r_i^1q_sr_i=q_sa_j^1$$
$`(2.2)`$
if $`i=(j,l)`$ and $`y_sx_{\mathrm{}}^1`$;
$$r_i^1q_sr_i=q_s$$
$`(2.3)`$
for all combinations of $`i`$ and $`s`$ which are not considered in (2.1) or (2.2);
$$r_i^1a_jr_i=a_j,iI,1jm;$$
$`(2.4)`$
$$r_i^1k_jr_j=k_j,iI,1jN;$$
$`(2.5)`$
$$k_1(q_1q_2\mathrm{}q_M)\mathrm{}k_N(q_1q_2\mathrm{}q_M)=1.$$
$`(2.6)`$
2.B. $`S`$-machines. Let us give a precise definition of $`S`$-machines . Let $`k`$ be a natural number. Consider now a language of admissible words. It consists of words of the form
$$q_1u_1q_2\mathrm{}u_kq_{k+1}$$
where $`q_i`$ are letters from disjoint sets $`Q_i`$, $`i=1,\mathrm{},k+1`$, $`u_i`$ are reduced group words in an alphabet $`Y_i`$ ($`Y_i`$ are not necessarily disjoint), the sets $`\overline{Y}=Y_i`$ and $`\overline{Q}=Q_i`$ are disjoint.
Notice that in every admissible word, there is exactly one representative of each $`Q_i`$ and these representatives appear in this word in the order of the indices of $`Q_i`$.
If $`0ijk`$ and $`W=q_1u_1q_2\mathrm{}u_kq_{k+1}`$ is an admissible word then the subword $`q_iu_i\mathrm{}q_j`$ of $`W`$ is called the $`(Q_i,Q_j)`$-subword of $`W`$ ($`i<j`$).
An $`S`$-machine is a rewriting system . The objects of this rewriting system are all admissible words.
The rewriting rules, or $`S`$-rules, have the following form:
$$[U_1V_1,\mathrm{},U_mV_m]$$
where the following conditions hold:
Each $`U_i`$ is a subword of an admissible word starting with a $`Q_{\mathrm{}}`$-letter and ending with a $`Q_r`$-letter (where $`\mathrm{}=\mathrm{}(i)`$ must not exceed $`r=r(i)`$, of course).
If $`i<j`$ then $`r(i)<\mathrm{}(j)`$.
Each $`V_i`$ is also a subword of an admissible word whose $`Q`$-letters belong to $`Q_{\mathrm{}(i)}\mathrm{}Q_{r(i)}`$ and which contains a $`Q_{\mathrm{}}`$-letter and a $`Q_r`$-letter.
If $`\mathrm{}(1)=1`$ then $`V_1`$ must start with a $`Q_1`$-letter and if $`r(m)=k+1`$ then $`V_n`$ must end with a $`Q_{k+1}`$-letter (so tape letters are not inserted to the left of $`Q_1`$-letters and to the right of $`Q_{k+1}`$-letters).
To apply an $`S`$-rule to a word $`W`$ means to replace simultaneously subwords $`U_i`$ by subwords $`V_i`$, $`i=1,\mathrm{},m`$. In particular, this means that our rule is not applicable if one of the $`U_i`$’s is not a subword of $`W`$. The following convention is important:
After every application of a rewriting rule, the word is automatically reduced. We do not consider reducing of an admissible word a separate step of an $`S`$-machine.
We also always assume that an $`S`$-machine is symmetric, that is for every rule of the $`S`$-machine the inverse rule (defined in the natural way) is also a rule of this $`S`$-machine.
Notice that virtually any $`S`$-machine is highly nondeterministic.
Among all admissible words of an $`S`$-machine we fix one word $`W_0`$. If an $`S`$-machine $`𝒮`$ can take an admissible word $`W`$ to $`W_0`$ then we say that $`𝒮`$ accepts $`W`$. We can define the time function of an $`S`$-machine as usual. If $`ZZ_1\mathrm{}Z_n=W_0`$ is an accepting computation of the $`𝒮`$-machine $`𝒮`$ then $`|Z|+|Z_1|+\mathrm{}+|Z_n|`$ is called the area of this computation. This allows us to define the area function of an $`S`$-machine.
In , it is showed how to associate a group with any $`S`$-machine. Our group $`G(v,m)`$ is the group associated with the following simple $`S`$-machine $`𝒮(v)`$.
Its language of admissible words coincides with the set of words of the form $`q_1u_1q_2u_2\mathrm{}q_M`$ where $`u_i`$ are any group words in the alphabet $`\{a_1,\mathrm{},a_m\}`$. Each command of this $`S`$-machine corresponds to a variable of $`v`$ and a letter from $`\{a_1,\mathrm{},a_m\}`$. Let $`x`$ be one of the variables of $`v`$, which occurs with exponent $`+1`$ at positions $`i_1,\mathrm{},i_s`$ of $`v`$ and occurs with exponent $`1`$ at positions $`j_1,\mathrm{},j_t`$ and let $`a`$ be any letter from $`\{a_1,\mathrm{},a_m\}`$. Then the corresponding command multiplies $`q_{i_1+1},\mathrm{},q_{i_s+1}`$ by $`a`$ on the left, multiplies $`q_{j_1},\mathrm{},q_{j_t}`$ by $`a^1`$ on the right, and does not change other $`q`$.
For example, if $`v=x^{M1}`$ then $`v`$ contains only one variable, and commands of the $`S`$-machine $`𝒮(v)`$ are indexed by letters from $`\{a_1,\mathrm{},a_m\}`$ and each command has the form $`[q_2a_iq_2,\mathrm{},q_Ma_iq_M]`$.
Let $`W_0=q_1\mathrm{}q_M`$. Here is the main property of the $`S`$-machine $`𝒮(v)`$. This statement can be easily proved by induction.
The main property of $`𝒮(v)`$. An admisible word $`q_1u_1q_2\mathrm{}u_{M1}q_M`$ is accepted by the $`S`$-machine $`𝒮(v)`$ if and only if $`u_iu_jy_i=y_j`$. In other words acceptable words are obtained from the words of the form $`v(u_1,\mathrm{},u_k)=u_{i_1}\mathrm{}u_{i_{M1}}`$ by inserting $`q_1,\mathrm{},q_M`$ between the factors $`u_{i_j}`$.
For example, if $`v=x^{M1}`$ then accepted words have the form $`q_1uq_2u\mathrm{}uq_M`$ where $`u`$ is an arbitrary word in the alphabet $`\{a_1,\mathrm{},a_n\}`$.
The construction (essentialy from ) which simulates an $`S`$-machine in a group is the following.
Let $`\mathrm{\Theta }`$ be the set of rules of $`𝒮`$. Let us call one of each pair of mutually inverse rules from $`\mathrm{\Theta }`$ positive and the other one negative. The set of all positive rules will be denoted by $`\mathrm{\Theta }_+`$.
Let $`N`$ be any positive integer ($`29`$). Let $`A`$ be the set of all letters occurring in the admissible words of $`𝒮`$ union with the set $`\{k_j|j=1,\mathrm{},N\}\mathrm{\Theta }_+.`$
Our group is generated by the set $`A`$ subject to the set $`𝒫_N(𝒮)`$ of relations described below.
1. Transition relations. These relations correspond to elements of $`\mathrm{\Theta }_+`$.
Let $`r\mathrm{\Theta }_+`$, $`r=[U_1V_1,\mathrm{},U_pV_p]`$. Then we include relations $`U_1^r=V_1,\mathrm{},U_p^r=V_p`$ into $`𝒫_N(𝒮)`$. If some set $`Q_j`$ does not have a representative in any of the words $`U_i`$ then we include all the commutativity relations $`q^r=q`$, $`qQ_j`$.
2. Auxiliary relations.
These are all possible relations of the form $`rx=xr`$ where $`x`$ is one of the letters in $`\{a_1,\mathrm{},a_m,k_1,\mathrm{},k_N\}`$, $`r\mathrm{\Theta }^+`$.
3. The hub relation.
$$k_1W_0k_2W_0\mathrm{}k_NW_0=1$$
It is easy to see that the group presentation associated with our $`S`$-machine $`S(v)`$ coincides with the presentation constructed above in section A.
2.C. A picture. For simplicity, let us take $`v=x^n`$: our construction does not depend much on the word $`v`$ anyway.
Further simplifying the situation (and the future picture) let us take $`n=3`$. The construction really does not depend much on $`n`$ either, so we shall sometimes write $`n`$ instead of $`3`$. Figure 1 shows the van Kampen diagram (below it will be called a disc) with boundary label $`\mathrm{\Sigma }(q_1uq_2uq_3uq_4)=k_1q_1uq_2uq_3uq_4k_2q_1uq_2uq_3uq_4k_3\mathrm{}k_Nq_1uq_2uq_3uq_4`$.
Fig. 1.
On the boundary of this diagram we can read the word $`\mathrm{\Sigma }(q_1uq_2uq_3uq_4)`$. The words on each of the concentric circles is labeled by $`\mathrm{\Sigma }(q_1u_iq_2u_iq_3u_iq_4)`$ where $`u_i`$ is a prefix of $`u`$ of length $`i1`$. The word written on the innermost circle is the hub, $`\mathrm{\Sigma }(q_1q_2q_3q_4)`$. The edges connecting the circles are labeled by letters $`r_1,\mathrm{},r_m`$ corresponding to the letters of $`u`$. The cells tessellating the space between the circles have labels
* $`q_i^{r_j}=a_jq_i`$, $`i=2,3,4`$, $`j=1,\mathrm{},m`$; $`q_1^{r_j}=q_1`$.
* $`ar=ra`$, $`a\{a_1,\mathrm{},a_m\}`$, $`r\{r_1,\mathrm{},r_m\}`$
* $`kr=rk`$, $`k\{k_1,\mathrm{},k_N\}`$, $`r\{r_1,\mathrm{},r_m\}`$.
These are exactly the relations of our group $`G(v,m)`$.
The following two lemmas summarize two main features of the presentation of the group $`G(v,m)`$. For any reduced words $`X_1,\mathrm{}X_k`$ in the alphabet $`a_1^{\pm 1},\mathrm{},a_m^{\pm 1}`$ we define the word $`\mathrm{\Lambda }(X_1,\mathrm{},X_k)`$ to be the word $`u(X_1,\mathrm{},X_k)`$ with letters $`q_1,\mathrm{},q_M`$ inserted between factors $`X_{i_j}`$ (recall that by the main property of $`𝒮(v)`$ these are all acceptable words of the $`S`$-machine $`𝒮(v)`$.) More precisely
$$\mathrm{\Lambda }(X_1,\mathrm{},X_k)q_1X_{i_1}q_2X_{i_2}\mathrm{}X_{i_{M1}}q_M,$$
provided $`v(x_1,\mathrm{},x_k)x_{i_1}x_{i_2}\mathrm{}x_{i_{M1}}`$. The following claim is an immediate corollary of the relations (2.1)-(2.5) (and is evident from Figure 1).
###### Lemma 2.1
Assume $`i=(j,\mathrm{})`$. Then in view of relations (2.1) - (2.4), the word $`r_i^1\mathrm{\Lambda }(X_1,\mathrm{},X_k)r_i`$ (the word $`r_i\mathrm{\Lambda }(X_1,\mathrm{},X_k)r_i^1`$) is equal to the word $`\mathrm{\Lambda }(X_1^{},\mathrm{},X_k^{})`$, where in the free group $`X_u^{}=X_u`$ for $`u\mathrm{}`$, and $`X_{\mathrm{}}^{}=X_{\mathrm{}}a_j`$ ($`X_{\mathrm{}}^{}=X_{\mathrm{}}a_j^1`$).
$`\mathrm{}`$
Now set
$$\mathrm{\Sigma }(X_1,\mathrm{},X_k)k_1\mathrm{\Lambda }(X_1,\mathrm{},X_k)\mathrm{}k_N\mathrm{\Lambda }(X_1,\mathrm{},X_k).$$
In particular $`\mathrm{\Sigma }(1,\mathrm{},1)`$ is the left-hand side of the relation (2.6). Since any $`k`$-tuple $`(X_1,\mathrm{},X_k)`$ is a result of iterated multiplications of the components of the $`k`$-tuple $`(1,\mathrm{},1)`$ by letters $`a_j^{\pm 1}`$, Lemma 2.1 and relations (2.6) imply
###### Lemma 2.2
For any reduced words $`X_1,\mathrm{},X_k`$, the word $`\mathrm{\Sigma }(X_1,\mathrm{},X_k)`$ is conjugate to the word $`\mathrm{\Sigma }(1,\mathrm{},1)`$ in virtue of relations (2.1) - (2.6). In particular, $`\mathrm{\Sigma }(X_1,\mathrm{},X_k)=1`$ in the group $`G(v,m)`$.
$`\mathrm{}`$
This lemma is also evident from Figure 1, because Figure 1 is the van Kampen diagram over the presentation of $`G(v,m)`$ with boundary label
$$\mathrm{\Sigma }(X_1,\mathrm{},X_k).$$
Now let us define the presentation of the group $`H=H(v,m)`$ (the construction is similar to the Aanderaa construction from ). Let us add new letters $`\rho ,d,b_1,\mathrm{},b_m`$ to the above presentation of $`G(v,m)`$, and add the following relations:
$$\rho ^1k_1\rho =k_1d^1,\rho ^1k_2\rho =dk_2;$$
$`(2.7)`$
$$\rho ^1k_j\rho =k_j,\mathrm{\hspace{0.33em}3}jN;$$
$`(2.8)`$
$$\rho ^1q_j\rho =q_j,\mathrm{\hspace{0.33em}1}jM;$$
$`(2.9)`$
$$\rho ^1a_j\rho =a_j,\mathrm{\hspace{0.33em}1}jm;$$
$`(2.10)`$
$$d^1a_jd=a_jb_j,\mathrm{\hspace{0.33em}1}jm;$$
$`(2.11)`$
$$d^1q_jd=q_j,\mathrm{\hspace{0.33em}1}jM;$$
$`(2.12)`$
$$b_ja_{\mathrm{}}=a_{\mathrm{}}b_j,\mathrm{\hspace{0.33em}1}j,lm;$$
$`(2.13)`$
$$b_jq_{\mathrm{}}=q_{\mathrm{}}b_j,\mathrm{\hspace{0.33em}1}jm,1lM;$$
$`(2.14)`$
$$w(b_1,\mathrm{},b_m)=1$$
$`(2.15)`$
for any cyclically reduced word $`w=w(b_1,\mathrm{},b_m)`$ which is equal to 1 in the relatively free group $`F_m(𝒱)`$ with the basis $`(b_1,\mathrm{},b_m)`$.
It is easy to see that these relations do not depend on the structure of the word $`v`$ (only on the length of $`v`$). All these relations can be put in one picture, the following Figure 2. These relations together with the relations of $`G(v,m)`$ form the presentation of $`H(v,m)`$.
Fig. 2.
This is an annular diagram over the presentation of $`H(v,m)`$ (as above we assume for simplicity that $`v=x^3`$). It is obtained in the following way. Take the disc $`\mathrm{\Delta }`$ on Figure 1. Since $`\rho `$ commutes with all generators of $`G(v,m)`$ except $`k_1`$ and $`k_2`$, and $`k_1^\rho =k_1d^1,k_2^\rho =dk_2`$, we can form an annulus of $`\rho `$-cells with the inner boundary labeled by the same word as the boundary of $`\mathrm{\Delta }`$, and the outer boundary labeled by the same word with $`d^1`$ inserted next to the right of $`k_1`$ and $`d`$ inserted next to the left of $`k_2`$. Glue in the disc $`\mathrm{\Delta }`$ inside this annulus. We obtain the part of the diagram on Figure 2 formed by the disc and the $`\rho `$-annulus enveloping the disc. Let us call this part $`\mathrm{\Delta }_1`$. Now $`d`$ commutes with all $`q`$’s, and we have that $`a_i^d=a_ib_i`$. Also take into account that $`b_i`$ commutes with all the $`a`$’s and $`q`$’s. This implies that if $`U=q_1uq_2uq_3uq_4`$ is the word written between $`k_1`$ and $`k_2`$ on the boundary of $`\mathrm{\Delta }`$ (read clockwise) then
$$U^d=q_1uu_bq_2uu_bq_3uu_bq_4=q_1uq_2uq_3uq_4u_b^3=Uu_b^3$$
(all equalities hold modulo the presentation of $`H(v,m)`$). Here $`u_b`$ is the word $`u`$ rewritten in the alphabet $`\{b_1,\mathrm{},b_m\}`$. The corresponding diagram over $`H(v,m)`$ can be attached to $`\mathrm{\Delta }_1`$ along the arc labeled by $`d^1Wd`$. Let the resulting diagram be denoted by $`\mathrm{\Delta }_2`$. Now we can get the diagram on Figure 2 by identifying the ends of the arc labeled by $`u_b^3`$ on the boundary of $`\mathrm{\Delta }_2`$.
Notice that the outer boundary of the diagram on Figure 2 is labeled by the same word as the boundary of the disc $`\mathrm{\Delta }`$. Thus all the relations $`u_b^3=1`$ (relations (2.15) above) follow from the other relations from the presentation of $`H(x^3,m)`$. The next statement generalizes these observations to an arbitrary $`v`$.
###### Lemma 2.3
Let $`X_1,\mathrm{},X_k`$ be words in $`a_1^{\pm 1},\mathrm{},a_m^{\pm 1}`$, and assume that $`Y_1,\mathrm{},Y_k`$ are their copies in $`b_1^{\pm 1},\mathrm{},b_m^{\pm 1}`$ (obtained by replacing every $`a_j`$ by $`b_j`$ in $`X_1,\mathrm{},X_k`$). Then the relation
$$d^1\mathrm{\Lambda }(X_1,\mathrm{},X_k)d=\mathrm{\Lambda }(X_1,\mathrm{},X_k)v(Y_1,\mathrm{},Y_k)$$
follows from relations (2.11) - (2.14); relations (2.15) follows from relations (2.1) - (2.14). In particular, $`H`$ is a finitely presented group.
$`\mathrm{}`$ The first claim is an immediate corollary of the relations (2.11) – (2.14) and the definition of the word $`\mathrm{\Lambda }\mathrm{\Lambda }(X_1,\mathrm{},X_k)`$. As for the second statement, it suffices to prove it for all words of the type $`w(b_1,\mathrm{},b_m)v(Y_1,\mathrm{},Y_k)`$ only. In order to do that, we will apply Lemma 2.2, relations (2.7) – (2.10), and the first claim of Lemma 2.3:
$$1=\rho ^1\mathrm{\Sigma }(X_1,\mathrm{},X_k)\rho =k_1d^1\mathrm{\Lambda }dk_2\mathrm{\Lambda }\mathrm{}k_N\mathrm{\Lambda }=$$
$$k_1\mathrm{\Lambda }v(Y_1,\mathrm{},Y_k)k_2\mathrm{\Lambda }\mathrm{}k_N\mathrm{\Lambda }.$$
By Lemma 2.2 the last product remains being equal to 1 after erasing the factor $`v(Y_1,\mathrm{},Y_k)`$. Hence this factor vanishes itself. $`\mathrm{}`$
## 3 Bands and annuli.
Consider a simply connected van Kampen diagram $`\mathrm{\Delta }`$ over the presentation of the group $`H`$ (see or ; we assume that any edge of any van Kampen diagram is labeled by one letter, as in ). If a face $`\mathrm{\Pi }`$ of $`\mathrm{\Delta }`$ corresponds to a relation containing letters $`x`$ and $`y`$ then $`\mathrm{\Pi }`$ is said to be a $`(x,y)`$-cell. Thus we can talk about $`(\rho ,a)`$-cells, $`(a_j,\rho )`$-cell, $`(r,q)`$-cells, etc. Similarly if the relation contains letter $`x`$ then we shall call the corresponding cell an $`x`$-cell.
The boundary of a $`\rho `$-cell $`\mathrm{\Pi }`$ has exactly two $`\rho `$-edges labeled by $`\rho ^{\pm 1}`$. These labels are inverses of each other when one reads the boundary label of $`\mathrm{\Pi }`$. This gives us an opportunity to construct “bands” of several $`\rho `$-cells.
A $`\rho `$-band of length 0 has no faces and consists of one $`\rho `$-edge. A $`\rho `$-band of length 1 is just a single $`\rho `$-cell. Assume by induction, that we have a $`\rho `$-band $`T^{}=[\mathrm{\Pi }_1,\mathrm{},\mathrm{\Pi }_{s1}]`$ of length $`s1`$ constructed of $`s1`$ distinct $`\rho `$-cells $`\mathrm{\Pi }_1,\mathrm{},\mathrm{\Pi }_{s1}`$, and the boundary of $`T^{}`$ has a $`\rho `$-edge $`e`$, which is a common edge of the boundary $`(\mathrm{\Pi }_{s1})`$ and the boundary of a $`\rho `$-cell $`\mathrm{\Pi }_s`$ which is distinct from $`\mathrm{\Pi }_1,\mathrm{},\mathrm{\Pi }_{s1}`$. Then we are able to construct a $`\rho `$-band $`T=[\mathrm{\Pi }_1,\mathrm{},\mathrm{\Pi }_s]`$ of length $`s`$ whose boundary $`T`$ is the union of the the boundaries $`T^{}`$ and $`\mathrm{\Pi }_s`$ minus the edge $`e`$. A $`\rho `$-band $`T`$ is maximal if is not contained in a $`\rho `$-band of a greater length.
Thus, the boundary of a $`\rho `$-band $`T`$ has the form $`e_1pe_2q`$, where $`e_1`$ and $`e_2`$ are $`\rho `$-edges (we call them ends of $`T`$), and the paths $`p,q`$ (sides of the band $`T`$) consist of $`a`$\- and $`q`$-edges, but contain no $`\rho `$-edges.
If the ends of $`T`$ coincide, one may identify them and the annular subdiagram $`T`$ is called a $`\rho `$-annulus . For example the diagram on Figure 2 contains a $`\rho `$-annulus enveloping the disc.
Similarly we define $`d`$-bands (and annuli) that by definition can be constructed of $`(d,a)`$-, $`(d,t)`$-, and $`(d,q)`$-cells.
$`b`$-bands can be constructed of $`(b,a)`$\- and $`(b,q)`$-cells.
$`r`$-bands are constructed of $`(r,q)`$-, $`(r,a)`$-, and $`(r,k)`$-cells.
$`q`$-bands are constructed of $`(q,r)`$-, $`(q,\rho )`$-, $`(q,d)`$-, and $`(q,b)`$-cells.
$`a`$-bands are created of $`(a,r)`$-, $`(a,\rho )`$-, $`(a,d)`$-, and $`(a,b)`$-cells.
Notice that $`(\rho ,d)`$-cells of type (2.7) cannot be included in a $`d`$-band but they can be terminal for $`d`$-bands, i.e. a maximal $`d`$-band can end only on the contour of a $`(\rho ,k)`$-cell or on the contour of the diagram $`\mathrm{\Delta }`$. Similarly, $`(r,q)`$-cells are terminal for $`a`$-bands, hubs, corresponding to relation (2.6), are terminal for $`k`$\- and $`q`$-bands, $`(d,a)`$-cells are terminal for $`b`$-bands. Also a $`b`$-band can terminate on the contour of a $`G_b`$-cell (by definition, a $`G_b`$-cell corresponds to a relation (2.15)).
Now consider an $`r`$-band $`T=[\pi _0,\pi _1,\mathrm{},\pi _{\mathrm{}},\pi _{\mathrm{}+1}]`$ and a $`q`$-band
$`T^{}=[\pi _0,\gamma _1,\mathrm{},\gamma _s,\pi _{s+1}]`$ which have no common faces except for $`\pi _0`$ and $`\pi _{\mathrm{}+1}`$, and with all ends of $`T`$ and $`T^{}`$ lying on the outer boundary of the annulus $`S`$ formed by $`T`$ and $`T^{}`$. Then this annulus is called an $`(r,q)`$annulus. It consists of the $`r`$-part $`T`$ and the $`q`$-part $`T^{}`$. The faces $`\pi _0`$ and $`\pi _{\mathrm{}+1}`$ are its corner cells.
The definitions of $`(\rho ,a)`$-, $`(q,b)`$-annuli, etc. are quite similar.
A diagram $`\mathrm{\Delta }`$ is called minimal in this section if there exists no other diagram $`\mathrm{\Delta }^{}`$ such that (1) $`\mathrm{\Delta }^{}`$ has the same boundary label as $`\mathrm{\Delta }`$, and (2) the number of faces of each of the types (2.1) - (2.15) in $`\mathrm{\Delta }^{}`$ does not exceed the similar number for $`\mathrm{\Delta }`$, (3) the total number of faces in $`\mathrm{\Delta }^{}`$ is smaller than the number of faces in $`\mathrm{\Delta }`$. In the next two sections we shall a stronger definition of minimality.
The main lemma of this section claims that there are no annuli of various kinds in minimal diagrams without hubs.
###### Lemma 3.1
Let $`\mathrm{\Delta }`$ be a minimal diagram over $`H`$ containing no hubs. Then $`\mathrm{\Delta }`$ has no
(1) $`\rho `$-annuli,
(2) $`r`$-annuli,
(3) $`(r,q)`$-annuli,
(4) $`q`$-annuli,
(5) $`(r,k)`$-annuli,
(6) $`k`$-annuli,
(7) $`(\rho ,k)`$-annuli,
(8) $`(a,b)`$-annuli,
(9) $`(d,a)`$-annuli,
(10) $`(\rho ,a)`$-annuli,
(11) $`(r,a)`$-annuli,
(12) $`d`$-annuli,
(13) $`b`$-annuli,
(14) $`a`$-annuli,
(15) $`(\rho ,q)`$-annuli,
(16) $`(d,q)`$-annuli,
(17) $`(q,b)`$-annuli
$`\mathrm{}`$ To prove statements (1) - (17) we use a simultaneous induction on the number of faces in the minimal subdiagram $`\mathrm{\Delta }_S`$ containing a conjectural counterexample, i.e. an annulus $`S`$. This means that we may assume that $`\mathrm{\Delta }`$ has no annulus $`S^{}`$ of types (1) - (17) such that the subdiagram $`\mathrm{\Delta }_S^{}`$ has fewer faces than $`\mathrm{\Delta }_S`$.
(1) Let $`S`$ be a $`\rho `$-annulus. Assume that $`S`$ has a $`(\rho ,k_j)`$-cell. Then this cell belongs to a $`k_j`$-band $`T`$, which must intersect $`S`$ at least twice, because by the lemma condition $`\mathrm{\Delta }`$ contains no hubs (terminal cells for $`k`$-bands). In such a case $`T`$ and a subband of $`S`$ form a smaller $`(\rho ,k)`$-annulus $`S^{}`$ than $`S`$, contrary to claim (7) of the lemma. The only case when $`S^{}`$ is not smaller than $`S`$ is when $`S`$ consists just of two $`(\rho ,k_j)`$-cells with a common $`\rho `$-edge and with “mirror” labels. Such a pair of mirror faces is impossible in a minimal diagram. For more details on the cell cancellation, see .
Therefore $`S`$ contains only $`(\rho ,a)`$\- and $`(\rho ,q)`$-cells. Consequently the outer and the inner boundaries of $`S`$ have identical labels. This makes it possible to delete the interior of $`S`$ and then identify the sides of $`S`$. Such a surgery does not change the boundary label of $`\mathrm{\Delta }`$, contrary to the minimality of $`\mathrm{\Delta }`$.
(3) Assume that $`S`$ is a $`(r,q)`$-annulus. Let $`T`$ be the $`q`$-part and $`T^{}`$ be the $`r`$-part. $`S`$ has no $`k`$-cells, because otherwise a smaller $`(r,k)`$-annulus appears, contradicting claim (5). Analogously, by (3) (for smaller annuli) there are no non-corner $`(r,q)`$-cells in $`S`$. The same argument shows that there are no other $`(r,q)`$-cells in the subdiagram $`\mathrm{\Delta }`$.
The corner $`(r,q)`$-cells are included in the same $`r`$\- and $`q`$-bands. This implies that they correspond to the same relation and simultaneously have or have no $`a_j`$-edges (for the same $`j`$) on the inner border of $`S`$ with opposite directions. Only these corner cells can be terminal for $`a`$-bands crossing $`S`$.
Therefore if there exist non-corner cells in $`S`$, then by (11) there exist exactly 2 such cells, they must be neighbors in the $`r`$-part of $`S`$, and must have “mirror” labels, contrary to the minimality of $`\mathrm{\Delta }`$.
Hence the $`r`$-part of $`S`$ has no non-corner cells. Then the standard cancelation argument can be applied to the corner cell. This contradicts the minimality assumption again.
(2), (4) - (17) The proofs of all these statements are similar to the two proofs of (1) and (3) given above. The reader could examine them as an exercise or read similar explanation for Lemma 6.1 (claims (1) - (20)) or Section 7 of . $`\mathrm{}`$
## 4 Hubs and spokes
A spoke is a maximal $`k`$-band having an end on a hub (or on a disc in the next section). Obviously, another end lies on a hub too, or on the boundary $`\mathrm{\Delta }`$.
It will be convenient to restrict the notion of a minimal diagram used in previous section as follows.
A type $`\tau =\tau (\mathrm{\Delta })`$ of a diagram $`\mathrm{\Delta }`$ over $`H`$ is the 4-tuple $`\tau =(\tau _1,\tau _2,\tau _3,\tau _4)`$ where $`\tau _1`$ is the number of hubs in $`\mathrm{\Delta }`$ (or, in the next section, the number of discs), $`\tau _2`$ is the number of $`(\rho ,k)`$\- and $`(r,k)`$-cells, $`\tau _3`$ is the number of all other faces except $`(b,q)`$-, and $`(b,a)`$\- cells, $`\tau _4`$ is the number of $`(b,t)`$-, $`(b,q)`$-, and $`(b,a)`$-cells in $`\mathrm{\Delta }`$. Set $`\tau <\tau ^{}`$ if $`\tau _1<\tau _1^{}`$, or $`\tau _1=\tau _1^{}`$, but $`\tau _2<\tau _2^{}`$, and so on. Further a diagram is said to be minimal if it has the minimal type among all diagrams with the same boundary label. Clearly, this notion of minimality is stronger than that in Section 3, that is a diagram which is minimal under this definition is also minimal under the definition in the previous section.
Since our construction of the group $`G(v,m)`$ is essentially the same as the construction of the group $`G_N(𝒮)`$ of , the next lemma follows from Lemma 11.1 of . Nevertheless we present a direct proof here.
###### Lemma 4.1
Let $`\mathrm{\Delta }`$ be a minimal diagram over $`G(v,m)`$ containing two hubs $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$. Assume that the hubs have two common consecutive spokes $`T_1`$ and $`T_2`$ such that the subdiagram $`\mathrm{\Delta }_0`$ bounded by these spokes and the hubs contains no hubs. Then $`\mathrm{\Delta }`$ is not minimal diagram: by passing to another diagram with the same boundary label one can decrease the number of hubs by 2.
$`\mathrm{}`$ Let $`\mathrm{\Delta }_1`$ be the subdiagram constructed of $`\mathrm{\Delta }_0`$ and the spokes $`T_1`$, $`T_2`$. Without loss of generality we may assume that $`\mathrm{\Delta }_1`$ is a minimal diagram. By Lemma 3.1 (5) for $`\mathrm{\Delta }_1`$, every $`r`$-band crossing $`T_1`$, must cross $`T_2`$ as well. Therefore each of the 4 sides of the bands $`T_1`$, $`T_2`$ must have the same label $`V=V(r_1,r_2,\mathrm{})`$. Thus the boundary label of $`\mathrm{\Delta }_0`$ is the commutator $`[V,\mathrm{\Lambda }_1]`$ for the word $`\mathrm{\Lambda }_1\mathrm{\Lambda }(1,\mathrm{},1)q_1tq_2t\mathrm{}q_M`$, i.e. $`V`$ commutes with $`\mathrm{\Lambda }_1`$ modulo all the defining relations of $`G(v,m)`$ excluding the hub. Being a word in the alphabet $`r_1^{\pm 1},r_2^{\pm 1},\mathrm{}`$, $`V`$ commutes with letters $`k_j`$ (see relations (2.5)). Consequently, $`V`$ commutes (modulo (2.1) - (2.5)) with the word $`W`$, which is a cyclic permutation of the left-hand side of the hub relation (2.6) written on $`\mathrm{\Pi }_1`$, $`\mathrm{\Pi }_2`$ in opposite directions starting with the ends of the band $`T_1`$.
Thus, if we cut out the hubs $`\mathrm{\Pi }_1`$, $`\mathrm{\Pi }_2`$ from $`\mathrm{\Delta }`$ and then make a cut along the band $`T_1`$ border, we get a hole labeled by the word $`[V,W]`$ equal to 1 by (2.1) – (2.5). Now by the van Kampen lemma we are able to insert a diagram of a type $`(0,,,)`$ in this hole, reducing the number of hubs in $`\mathrm{\Delta }`$ by 2. $`\mathrm{}`$
###### Lemma 4.2
Let $`\mathrm{\Delta }`$ be a diagram over $`H(v,m)`$ containing two hubs $`\mathrm{\Pi }_1`$, $`\mathrm{\Pi }_2`$ with common consecutive $`k_j`$\- and $`k_s`$-spokes $`T_1`$ and $`T_2`$ where $`\{j,s\}\{1,2\}`$. Assume that the subdiagram $`\mathrm{\Delta }_0`$ bounded by these spokes and the hubs contains no hubs. Then the diagram $`\mathrm{\Delta }`$ is not minimal: By passing to another diagram with the same boundary label, one can decrease the number of hubs by 2.
$`\mathrm{}`$ The proof is similar to the proof of Lemma 4.1, but now the word $`V`$ may contain $`\rho ^{\pm 1},r_1,r_2,\mathrm{}`$. So, the difference is that $`V`$ may not commute with the letters $`k_1`$, $`k_2`$ occurring in $`W`$. However $`k_1,k_2`$ occur just in the subword $`k_1\mathrm{\Lambda }_1k_2`$ of the left-hand side in (2.7), which is also a subword of $`W`$ in view of the condition $`\{k_j,k_s\}\{1,2\}`$. Therefore it suffices to check that $`V`$ commutes with $`k_{\mathrm{}}`$ for $`\mathrm{}1,2`$, with $`\mathrm{\Lambda }_1`$, and with $`k_1\mathrm{\Lambda }_1k_2`$ modulo relations (2.1) - (2.5), (2.7) - (2.15). The first property follows from (2.5) and (2.9). The second one can be explained exactly as in Lemma 4.1.
To prove the third commutativity, notice first of all that by relations (2.15) and Lemma 2.3 $`d`$ commutes with $`\mathrm{\Lambda }(X_1,\mathrm{},X_k)`$ for any words $`X_1,\mathrm{},X_k`$ in the alphabet $`\{a_1^{\pm 1},\mathrm{},a_m^{\pm 1}\}`$. Then $`\rho `$ commutes with both $`\mathrm{\Lambda }(X_1,\mathrm{},X_k)`$ (see (2.9), (2.10)) and $`k_1\mathrm{\Lambda }(X_1,\mathrm{},X_k)k_2`$ since by (2.7) $`\rho ^1k_1=k_1d^1\rho ^1`$ and $`k_2\rho =\rho dk_2`$.
Also recall that by (2.5) and Lemma 2.2
$$r_i^{\pm 1}k_1\mathrm{\Lambda }(X_1,\mathrm{},X_k)k_2r_i^1=k_1r_i^{\pm 1}\mathrm{\Lambda }(X_1,\mathrm{},X_k)r_i^1k_2=k_1\mathrm{\Lambda }(X_1^{},\mathrm{},X_k^{})k_2$$
for some $`X_1^{},\mathrm{},X_k^{}`$.
Therefore permuting the words $`V^1`$ and $`V`$ (these are word in in $`\rho ^{\pm 1},r_1^{\pm 1},r_2^{\pm 1},\mathrm{}`$) letter-by-letter with $`k_1`$ and $`k_2`$, we get
$$V^1k_1\mathrm{\Lambda }_1k_2V=k_1V^1\mathrm{\Lambda }_1Vk_2.$$
The right-hand side is equal to $`k_1\mathrm{\Lambda }_1k_2`$, as was mentioned earlier. $`\mathrm{}`$
With any minimal diagram $`\mathrm{\Delta }`$ over $`H(v,m)`$ we associate the following graph $`\mathrm{\Gamma }=\mathrm{\Gamma }_\mathrm{\Delta }`$. One vertex of $`\mathrm{\Gamma }`$ (exterior vertex) is taken outside $`\mathrm{\Delta }`$ on the plane. Every interior vertex is chosen inside a hub. The edges between the vertices are drawn along the ”medians” of the spokes. The exterior edges are incident to the exterior vertex. The other edges are interior. Finally to complete the definition of $`\mathrm{\Gamma }`$, we erase the interior $`k_1`$-spoke for every pair of hubs connected by both $`k_1`$\- and $`k_2`$-spokes.
By Lemma 4.2 there are no bigons formed by interior edges of $`\mathrm{\Gamma }`$. Also there are no loops in $`\mathrm{\Gamma }`$ because every letter $`k_j`$ occurs once in the left-hand side of (2.7). Therefore in standard way the Euler formula implies that $`\mathrm{\Gamma }`$ has many exterior edges. (The restriction $`N16`$ would be enough here; the stronger condition $`N128`$ will be useful for Sections 5-7.) The following statement (see Lemma 2.13 in , or Lemmas 3.2, 3.3 in , or Lemma 11.5 in ) will be sufficient for our purpose.
###### Lemma 4.3
Let $`\mathrm{\Delta }`$ be a minimal diagram containing at least one hub. Then there exists a hub $`\mathrm{\Pi }`$ in $`\mathrm{\Delta }`$ such that at least $`N4`$ consecutive spokes starting on $`\mathrm{\Pi }`$, have their ends on the boundary $`\mathrm{\Delta }`$, and moreover there are no other hubs between the spokes of this set. The number of $`k`$-edges in $`\mathrm{\Delta }`$ is at least 3 times greater than the number of hubs in $`\mathrm{\Delta }`$.
$`\mathrm{}`$
###### Lemma 4.4
The natural homomorphism of the group $`F_m(𝒱)`$ onto $`H(v,m)`$ (well defined in view of relations (2.15)) is injective.
$`\mathrm{}`$ Assume that a word $`ww(b_1,\mathrm{},b_2)`$ vanishes under the homomorphism. Then there exists a minimal van Kampen diagram $`\mathrm{\Delta }`$ with the boundary label $`w`$. By Lemma 4.3 there are no hubs in $`\mathrm{\Delta }`$ because no letter $`k_j`$ occurs in $`w`$. Then by Lemma 3.1 (1) $`\mathrm{\Delta }`$ has no $`\rho `$-annuli, and consequently, it has no $`\rho `$-cells at all. Quite similarly, Lemma 3.1 allows us to exclude $`r`$-, $`q`$-, $`k`$-, $`d`$-, and $`a`$-cells from $`\mathrm{\Delta }`$ consequently. For example, there are no $`d`$-cells because maximal $`d`$-bands could terminate on $`\rho `$-cells only. Thus, $`\mathrm{\Delta }`$ has $`G_b`$-cells only, that correspond to relations (2.15). Hence $`w=1`$ in $`G`$ as desired. $`\mathrm{}`$
## 5 The band structure of disc-based diagrams
It will be convenient to extend the list of relations of the group $`H(v,m)`$ by adding the relations $`\mathrm{\Sigma }(X_1,\mathrm{},X_k)=1`$ for all $`k`$-tuples of reduced words $`X_1,\mathrm{},X_k`$ in $`a_1^{\pm 1},\mathrm{},a_m^{\pm 1}`$. Such an enlargement does not change $`H(v,m)`$ by Lemma 2.2. Any face of a diagram over this presentation of $`H(v,m)`$, that corresponds to some relation $`\mathrm{\Sigma }(X_1,\mathrm{},X_k)=1`$, will be called a disc.
The notion of a minimal diagram will be further restricted by replacing discs for hubs in the definition of a minimal diagram from the previous section. With any minimal diagram $`\mathrm{\Delta }`$ we associate a graph $`\mathrm{\Gamma }(\mathrm{\Delta })`$. Its definition repeats the definition of the graph $`\mathrm{\Gamma }_\mathrm{\Delta }`$ given in Section 4 where discs replace hubs. Notice that by Lemma 2.2 every disc can be replaced by a hub (which is a disc too) and a number of faces of smaller ranks (but the resulted diagram may be not minimal). The possibility of such a replacement and Lemma 4.2 show that the graph $`\mathrm{\Gamma }(\mathrm{\Delta })`$ has no bigons as well. Therefore the statement of Lemma 4.3 is also true for $`\mathrm{\Gamma }(\mathrm{\Delta })`$.
Let $`S`$ be a $`\rho `$\- or $`r`$-band that consequently intersects at $`(\rho ,k)`$\- or $`(r,k)`$-cells a series of consequent spokes $`T_1,\mathrm{},T_{\mathrm{}}`$ starting on a disc $`D`$. We say that the band $`S`$ envelopes disc $`D`$ if $`\mathrm{}>N/2`$, and there are no other discs in the sectors formed by $`S`$, $`T_j`$ and $`T_{j+1}`$ for $`j=1,\mathrm{},\mathrm{}1`$.
###### Lemma 5.1
A minimal diagram $`\mathrm{\Delta }`$ over $`H(v,m)`$ has no bands which envelope discs.
$`\mathrm{}`$ The proof is completely similar to the proofs of Lemma 8.4 or Lemma 2.17 in . Therefore we give just a brief explanation below referring for details to or .
Arguing by contradiction, one can choose the closest to $`D`$ band $`S`$ that envelopes it. Then the intersection cells $`\pi _1,\mathrm{},\pi _{\mathrm{}}`$ of $`S`$ and $`T_1,\mathrm{},T_{\mathrm{}}`$ have common $`k`$-edges with $`D`$. It suffices to prove that the diagram $`\mathrm{\Delta }_0`$ consisting of $`D`$ and $`\pi _1,\mathrm{},\pi _{\mathrm{}}`$ and considered separately from $`\mathrm{\Delta }`$, is not minimal.
For this purpose we attach auxiliary $`(r_s,k)`$-cells $`\pi _{\mathrm{}+1},\mathrm{},\pi _N`$ to $`\mathrm{\Delta }_0`$ along the $`N\mathrm{}`$ free $`k`$-edges of the disc $`D`$ (if $`S`$ is a $`r_s`$ band) so that all the cells $`\pi _1,\mathrm{},\pi _N`$ would be attached to $`D`$ uniformly, i.e. their $`r`$-edges would be directed all ’to’ or all ’from’ $`D`$. Then adding several faces of smaller ranks we can get a diagram $`\mathrm{\Delta }_1`$ with a label $`\mathrm{\Sigma }(X_1^{},\mathrm{},X_k^{})`$ by Lemma 2.2.
Therefore, conversely, one can construct a diagram $`\mathrm{\Delta }_2`$, with the same boundary label as $`\mathrm{\Delta }_0`$, consisting of a disc $`D^{}`$ labeled by $`\mathrm{\Sigma }(X_1^{},\mathrm{},X_k^{})`$, mirror copies of $`(k,r)`$-cells $`\pi _{\mathrm{}+1},\mathrm{},\pi _N`$, and faces of smaller ranks. But this contradicts to the minimality of $`\mathrm{\Delta }_0`$ since $`N\mathrm{}<\mathrm{}`$.
If $`S`$ is a $`\rho `$-band, the proof is similar, but the boundary label of the disc $`D^{}`$ coincides with the label of $`D`$, since $`\rho `$ commutes with $`\mathrm{\Sigma }(X_1,\mathrm{},X_k)`$ as was explained in Lemma 4.2. $`\mathrm{}`$
###### Lemma 5.2
A minimal diagram $`\mathrm{\Delta }`$ has no annuli of types (1) - (17) from Lemma 3.1 (even if hubs or discs occur in $`\mathrm{\Delta }`$).
$`\mathrm{}`$ This is similar to Lemmas in Section 4 of . Let us show, for example, that statement (3) of this lemma can be deduced from statement (3) of Lemma 3.1.
Let $`\mathrm{\Delta }_S`$ be a minimal subdiagram of $`\mathrm{\Delta }`$, containing a $`(r,q)`$-annulus $`S`$. By Lemma 3.1(3) it contains a disc. By Lemma 4.3 there exists a disc $`D`$ in $`\mathrm{\Delta }_S`$ such that its consecutive spokes $`T_1,\mathrm{},T_{N4}`$ intersect the $`r`$-part $`R`$ of $`S`$ (since the $`q`$-part has no $`k`$-cells at all). If there are other discs in $`\mathrm{\Delta }_S`$ and the $`q`$-part $`Q`$ of $`S`$ occurs between some $`T_j`$ and $`T_{j+1}`$, then again as in Lemma 4.3 (but for the graph obtained by erasing the vertex in $`D`$ and the edges incident to it), we get another disc $`D^{}`$ and spokes $`T_1^{},\mathrm{},T_{N5}^{}`$ starting on it, such that $`R`$ intersects them consecutively , and there are neither hubs nor $`Q`$ between the spokes. But this contradicts Lemma 5.1 because $`N5N/2`$. $`\mathrm{}`$
###### Lemma 5.3
Any two distinct maximal bands $`T`$ and $`T^{}`$ have at most one common face in a minimal diagram over $`H(v,m)`$.
$`\mathrm{}`$ Basically the statement follows from Lemma 5.2. However we have to remember that a-priori, a multiple intersection of two bands $`T`$ and $`T^{}`$ does not imply that they form even one annulus, because one or both ends of the band can be inside the “annulus”. Figure 3 shows a spiral multiple intersection.
Fig. 3.
In fact such spirals cannot occur, and the reader can find details in Lemmas 5.1 and 5.8 from . Here we just explain the idea for the particular case when $`T`$ is a $`d`$-band and $`T^{}`$ is an $`a`$-band. In this case a terminal cell $`\mathrm{\Pi }`$ for $`T`$ must be inside the subdiagram $`\mathrm{\Delta }_S`$ bounded by $`S`$. $`\mathrm{\Pi }`$ is a $`(\rho ,k)`$-cell (see relations (2.7)). Therefore by Lemma 4.3 and Lemma 3.1 (6) the boundary $`\mathrm{\Delta }_S`$ must be crossed by at least one $`k`$-band. But both $`d`$-band $`T`$ and $`a`$-band $`T^{}`$ have no $`k`$-cells at all, a contradiction. $`\mathrm{}`$
The following statement is similar to Lemma 4.36 in .
###### Lemma 5.4
There is no $`b`$-band $`T`$ in a minimal diagram $`\mathrm{\Delta }`$ such that both ends of $`T`$ belong to $`G_b`$-cells.
$`\mathrm{}`$ First assume that both ends of $`T`$ belong to the boundary of one $`G_b`$-cell $`\mathrm{\Pi }`$. Then the boundary of $`T`$ and a subpath of $`\mathrm{\Pi }`$ form the boundary of a subdiagram $`\mathrm{\Delta }_0`$ with the boundary label $`UV`$ where $`U`$ is a word in $`b_1^{\pm 1},\mathrm{},b_m^{\pm 1}`$, and $`V`$ is a word in $`a_1^{\pm 1},\mathrm{},a_m^{\pm 1}`$, $`q_1^{\pm 1},\mathrm{},q_M^{\pm 1}`$. By lemmas 4.3 and 3.1 $`\mathrm{\Delta }_0`$ has no discs and $`k`$-cells. By lemma 3.1 the band $`T`$ has no $`a`$\- and $`q`$-cells because otherwise we would get $`(a,b)`$\- or $`(q,b)`$-bands. Thus the band $`T`$ has length 0. Hence there is a loop in $`\mathrm{\Pi }`$ whose label must be equal to 1 in the group $`G`$ by Lemma 4.4. Then one can replace the interior of this loop and $`\mathrm{\Pi }`$ by a single $`G_b`$-cell, contrary to the minimality of $`\mathrm{\Delta }`$.
Assume now that the ends of $`T`$ belong to distinct $`G_b`$-cells $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$. Recall that the boundary label of $`T`$ commutes with any word in $`b_1^{\pm 1},\mathrm{},b_m^{\pm 1}`$ by relations (2.13) and (2.14). Therefore the diagram $`\mathrm{\Delta }`$ is not minimal: the subdiagram consisting of the two $`G_b`$-cells connected by the border of $`T`$, can be replaced by a diagram consisting of one $`G_b`$-cell and several cells of smaller ranks that correspond to relations (2.13) and (2.14). $`\mathrm{}`$
## 6 Upper bounds for the number of cells in minimal diagrams
Recall that by Lemma 2.3 the group $`H(v,m)`$ has a finite presentation.
###### Lemma 6.1
For any word $`w`$ of length $`n`$ in the generators of the group $`H(v,m)`$, that is equal to 1 in $`H(v,m)`$, there is a diagram over the finite presentation (2.1) - (2.14) with $`n^2O(f(O(n^2))^2`$ cells, where $`f`$ is the function given in Theorem 1.
$`\mathrm{}`$ First consider a minimal diagram $`\mathrm{\Delta }`$ with the boundary label $`w`$ over the disc-based presentation.
It follows from the statement of Lemma 4.3 for discs, that the number of discs in $`\mathrm{\Delta }`$ is $`O(n)`$. Therefore by Lemma 5.2 the number of maximal $`\rho `$-, $`r`$-, $`q`$-, and $`k`$-bands in $`\mathrm{\Delta }`$ is $`O(n)`$ because the boundary of any disc has no $`\rho `$\- and $`r`$-edges, and has $`O(1)`$ $`q`$-, and $`k`$-edges.
Therefore by Lemma 5.3 the number of $`(r,q)`$-cells (the intersections of $`r`$\- and $`q`$-bands) is $`O(n^2)`$. Similarly, the number of $`(r,k)`$-, $`(\rho ,k)`$-, and $`(\rho ,q)`$-cells is $`O(n^2)`$.
By Lemma 5.2 the number of maximal $`d`$-bands in $`\mathrm{\Delta }`$ is $`O(n^2)`$ since only $`(\rho ,k)`$-cells can be terminal for $`d`$-bands.
A similar argument shows that the number $`N_1`$ of all maximal $`a`$-bands terminating on $`(r,q)`$-cells or on $`\mathrm{\Delta }`$ is $`O(n^2)`$. To obtain similar upper bound for the number of all $`a`$-bands, it suffices to explain why the number $`N_2`$ of maximal $`a`$-bands having both ends on discs, is not greater than $`N_1`$. Indeed the number of common spokes of two discs is at most 2. Consequently the number of common $`a`$-bands for 2 discs is at most $`3/N<1/9`$ of the number of maximal $`a`$-bands starting on each of the disc. This implies that $`N_2N_1`$.
The upper bounds for the numbers of maximal $`k`$-, $`q`$-, and $`a`$-bands show that the sum of perimeters of all discs in $`\mathrm{\Delta }`$ is $`O(n^2)`$.
Now by Lemma 5.3 we are able to conclude that the number of $`(\rho ,a)`$-, $`(r,a)`$-, $`(d,q)`$-, $`(d,t)`$-cells is $`O(n^3)`$, and the number of $`(d,a)`$-cells is $`O(n^4)`$.
Only $`(d,a)`$-cells can be terminal for $`b`$-bands. So the above estimate of the number of $`(d,a)`$-cells together with Lemma 5.2 imply that the number of maximal $`b`$-bands is $`O(n^4)`$. This implies, first of all, that the number of $`(b,a)`$-cells is $`O(n^6)`$, and the number of $`(b,q)`$\- and $`(b,t)`$-cells is $`O(n^4)`$. Second, this implies that the sum of perimeters of all $`G_b`$-cells is $`O(n^4)`$. By the definition of the function $`f`$ every $`G_b`$-cell with boundary length $`s`$ can be tiled by cells with boundary labels of the form $`v(Y_1,\mathrm{},Y_k)`$ and with $`|Y_{ij}|`$ at most $`f(s)`$. By the superadditivity of the function $`f`$ (see Proposition 1.1) the set of all $`G_b`$-cells in $`\mathrm{\Delta }`$ can be replaced by the set of $`G_b`$-cells with labels of the form $`v(Y_1,\mathrm{},Y_k)`$ and with $`|Y_{ij}|`$ at most $`f(O(n^4))`$.
By Lemma 2.3 every $`G_b`$-cell $`\mathrm{\Pi }`$ labeled by a word $`v(Y_1,\mathrm{},Y_k)`$, can be replaced by a some number $`N_\mathrm{\Pi }`$ of cells corresponding to relations (2.1) - (2.14). The straightforward computation shows that $`N_\mathrm{\Pi }=O(|Y_1|+\mathrm{}+|Y_k|)^2`$. Therefore the set of all cells corresponding to relations (2.1) - (2.14), that replace all $`G_b`$-cells of $`\mathrm{\Delta }`$ is $`O(f(O(n^4))^2)`$.
Finally, every disc which is not a hub, can be replaced by a hub and several cells corresponding to relations (2.1) - (2.6). We need at most $`O(l^2)`$ such cells for every disc of perimeter $`\mathrm{}`$ which can be easily derived from the proof of Lemma 2.2. Therefore all discs in $`\mathrm{\Delta }`$ can be replaced by $`O(n^4)`$ cells of types (2.1) - (2.5).
Summing all the upper bounds for the numbers of cells of different types, we obtain an isoperimetric function of the group $`H(v,m)`$ that is equal to $`O(f(O(n^4))^2)`$. To achieve the better upper bound of $`n^2O(f(n^2)^2)`$ claimed in the Theorem, we prove that the perimeter of every $`G_b`$-cell in $`\mathrm{\Delta }`$ is in fact bounded by $`O(n^2)`$. The explanation is completely similar to the proof of Lemma 5.15 in . It is based on the fact that the number of the maximal $`b`$-bands starting on the same $`G_b`$-cell and terminating on $`a`$-cells lying in the same $`a`$-band, cannot exceed 3 (see Lemma 4.34 in ). We leave details to the reader. $`\mathrm{}`$
Proof of Theorem 1. $`\mathrm{}`$ The first statement of the theorem follows from Lemmas 2.3, 4.4 and 6.1. To prove that the constructed embedding of $`G`$ into $`H(v,m)`$ is undistorted, we need a longer and more complicated argument. The reasoning essentially coincides with that in Sections 10-12 of or in Section 7 in . Therefore we will not repeat it here referring the reader to and . $`\mathrm{}`$
## 7 The embeddings of the groups $`BS_{k,1}`$
It is more convenient to change the names of generators of $`BS_{k,1}`$ to $`b_1,b_2`$. So $`BS_{k,1}=b_1,b_2|b_2^{b_1}=b_2^k`$. Obviously the words $`W_n=W_n(b_1,b_2)(b_1^nb_2b_1^n)b_2(b_1^nb_2^1b_1^n)b_2^1`$ represent the identity in the group $`B=BS_{k,1}`$, and it is known that their “areas” exponentially increase depending on $`n`$ if $`k2`$. To prove Theorem 2 we are going to embed $`BS_{k,1}`$ into a finitely presented group $`H=H_{k,1}`$ such that areas of the words $`W_n`$ with respect of defining relations of $`H`$ have quadratic growth (a relatively easier task), and then we will prove a polynomial isoperimetric inequality for all words vanishing in $`H`$ (a harder job).
The embedding is similar to the one used for relatively free groups (see the proof of Theorem 1 above). But in order to show flexibility of our approach, we modify the second step of the embedding slightly. The main difference of the construction that we are about to present and the construction presented above is the absence of the letter $`d`$. Instead, we have different letters in different sectors of the discs.
An auxiliary group $`G=G_{k,1}`$ is given by generators
$$a_1,a_2,c,q_1,q_2,q_3,q_4,k_1,\mathrm{},k_N$$
where $`N29`$ and by defining relations
$$rq_1r^1=q_1a_1,rq_2r^1=q_2a_1^1,rq_3r^1=q_3a_1,rq_4r^1=q_4a_1^1,rq_5r^1=q_5c;$$
$`(7.1)`$
$$rk_jr^1=k_j,\mathrm{\hspace{0.33em}1}jN;$$
$`(7.2)`$
$$ra_1r^1=a_1,ra_2r^1=a_2,rcr^1=c;$$
$`(7.3)`$
$$q_1a_2q_2a_2q_3a_2^1q_4a_2^1k_1q_5\mathrm{}k_Nq_5=1.$$
$`(7.4)`$
The relation (7.4) is called the hub relation.
Note that it is easy to draw the disc corresponding to these relations and write the rules of the corresponding $`S`$-machine. It is a good exercise for a reader who wants to learn how to draw van Kampen diagrams and write programs for $`S`$-machines.
The words $`\mathrm{\Sigma }_s=\mathrm{\Sigma }(a_1^s,a_2)`$ are defined as follows:
$$\mathrm{\Sigma }_sq_1a_1^sa_2q_2a_1^sa_2q_3a_1^sa_2^1q_4a_1^sa_2^1k_1q_5c^s\mathrm{}k_Nq_5c^s.$$
An obvious analog of Lemma 2.2 says that $`r\mathrm{\Sigma }_sr^1=\mathrm{\Sigma }_{s+1}`$ in view of relations (7.1) – (7.3), and the deduction takes $`O(n^2)`$ of application of relations (7.1) – (7.3).
The group $`H=H_{k,1}`$ is defined by adding to the presentation of the group $`G`$ new generators $`b_1,b_2,\rho `$, and by adding new relations
$$\rho a_j\rho ^1=a_jb_j,j=1,2;$$
$`(7.5)`$
$$\rho c\rho ^1=c;$$
$`(7.6)`$
$$\rho q_j\rho ^1=q_j,j=1,2,3,4,5;$$
$`(7.7)`$
$$\rho k_j\rho ^1=k_j,\mathrm{\hspace{0.33em}1}jN;$$
$`(7.8)`$
$$b_ja_i=a_ib_j,i,j=1,2;$$
$`(7.9)`$
$$b_jq_i=q_ib_j,i=1,2,3,4,j=1,2;$$
$`(7.10)`$
$$b_1b_2b_1^1=b^k.$$
$`(7.11)`$
By definition, a $`B`$-cell in a diagram over $`H`$ corresponds to any cyclically reduced word $`w(b_1,b_2)`$ such that $`w=1`$ follows from (7.11). The following analog of Lemma 2.3 can be verified immediately (and similar to the proof of Lemma 2.3).
###### Lemma 7.1
For any $`n0`$ the relation
$$\mathrm{\Sigma }_nq_1a_1^na_2q_2a_1^na_2q_3a_1^na_2^1q_4a_1^na_2^1k_1q_5c^n\mathrm{}k_Nq_5c^n=1$$
can be obtained by application of $`O(n^2)`$ relations (7.1) – (7.3) to the hub, and the relation $`W_n(b_1,b_2)=1`$ can be deduced from (7.1) – (7.10) in $`O(n^2)`$ steps.
$`\mathrm{}`$
The proof of the next claim is absolutely similar to the proof of Lemma 3.1.
###### Lemma 7.2
A minimal hub-free diagram over the presentation of the group $`H_{k,1}`$ has no $`\rho `$\- $`r`$-, $`(r,q)`$-, $`q`$-, $`(r,k)`$-, $`k`$-, $`(\rho ,k)`$-, $`(a,b)`$-, $`(\rho ,a)`$-, $`(\rho ,c)`$-, $`(r,a)`$-, $`(r,c)`$-, $`b`$-, $`a`$-, $`(\rho ,q)`$\- or $`(q,b)`$-annuli.
$`\mathrm{}`$
The statement of Lemma 4.1 is also true for the group $`G=G_{k,1}`$, because the boundary label of a spoke has the form $`r^{\mathrm{}}`$, but in the HNN-extension of a free group with the stable letter $`r`$ defined by relations (7.1)–(7.3), the word $`r^{\mathrm{}}`$ can commute with a subword of the left-hand side of (7.4) written between neighbor $`k`$-letters only if $`\mathrm{}=0`$. In view of minimality of $`\mathrm{\Delta }`$, this means that the spokes have length 0, and the two hubs form a mirror pair of cells, that cancel.
###### Lemma 7.3
The statement of Lemma 4.2 is true for diagrams over $`H=H_{k,1}`$. Hence Lemmas 4.3, 4.4 are also valid for $`H`$.
$`\mathrm{}`$ As in the proof of Lemma 4.2, we have to consider the boundary label $`V=V(\rho ^{\pm 1},r^{\pm 1})`$ of the spokes $`T_1,T_2`$ that must commute with $`q_5`$ in view of relations (7.1)–(7.3), (7.4)–(7.11). This is possible if and only if the exponent sum for $`r`$ is zero in $`V`$, since the group defined by (7.1)–(7.3), (7.4)–(7.11) has the retraction preserving $`r,q_5`$ and $`c`$ and mapping the other generators to 1.
As in Lemma 4.2 we have to prove that $`V`$ commutes with the cyclic permutation $`W`$ of the left-hand side of (7.4) beginning with $`k_j`$. The word $`V`$ commutes with any letter $`k_j`$ by relations (7.2), (7.8). Therefore it suffices to prove that $`V`$ commutes with the word $`\mathrm{\Lambda }_0q_1a_2q_2a_2q_3a_2^1q_4a_2^1`$.
It is clear that the equality $`r\mathrm{\Lambda }_nr^1=\mathrm{\Lambda }_{n+1}`$ follows from (7.1)–(7.3) for
$`\mathrm{\Lambda }_sq_1a_1^sa_2q_2a_1^sa_2q_3a_1^sa_2^1q_4a_1^sa_2^1`$ and any integer $`n`$. The equalities $`\rho \mathrm{\Lambda }_s\rho ^1=\mathrm{\Lambda }_s`$ follows from (7.5), (7.7), (7.9)–(7.10) and Lemma 7.1. Since the exponent sum for the occurrences of $`r`$ in $`V`$ is equal to 0, $`V\mathrm{\Lambda }_0V^1=V_0`$, as desired. $`\mathrm{}`$
Let us mention another similarity of the groups $`H(v,m)`$ and $`H_{k,1}`$: the claims of Lemmas 5.15.4 are true for $`H=H_{k,1}`$ as well, and the proofs are quite analogous to those in Section 5.
To complete the proof of Theorem 2 we need
###### Lemma 7.4
Let $`w=w(b_1^{\pm 1},b_2^{\pm 1})`$ be a word of length $`n`$ such that $`w=1`$ in $`H`$. Then this equality can be derived from the trivial one by application of $`O(n^4)`$ of relations (7.1)–(7.11).
$`\mathrm{}`$ Since $`w=1`$ in $`BS_{k,1}`$, the exponential sum over the occurrences of $`b_1`$ in $`w`$ is equal to $`0`$. Therefore the word $`w`$ is freely equal to a product $`_iv_i`$ where $`v_ib_1^{s_i}b_2^{\pm 1}b_1^{s_i}`$, $`|s_i|<n/2`$ and the number of factor do not exceed $`n`$. Passing to a freely conjugate word, we may assume that $`0s_i<n`$ for every $`i`$.
Let $`s=\mathrm{min}s_i`$, and assume there are factors $`v_ib_1^sb_2b_1^s`$ and $`v_jb_1^sb_2^1b_1^s`$. Then by Lemma 7.1 we can transpose $`v_i`$ with a neighbor factor $`v_{i\pm 1}`$ by applying the defining relations at most $`O(n^2)`$ times, and so after $`O(n^3)`$ applications of the relations the factors $`v_i`$ and $`v_j`$ cancel.
Now assume that all the $`v_i`$’s with $`s_i=s`$ are equal, and $`\mathrm{}`$ is the number of such $`v_i`$’s. Since $`b_1^tb_2b_1^t=b_2^{k^t}`$ for any $`ts`$ in $`G`$, and the product of the words $`v_i`$ is equal to 1 in $`G`$, the number $`\mathrm{}`$ is a multiple of $`k`$.
Again, by Lemma 7.1 we can collect the $`\mathrm{}=km`$ factors (with minimal $`s_i`$’s) at the end of the word $`w`$ applying the relations at most $`\mathrm{}O(n^3)`$ times. This suffix is (in the free group) the product of $`m`$ factors $`u_ib_1^sb_2^{\pm k}b_1^s`$, and applying relator (7.11) $`m`$ times we can rewrite it as the product of $`m`$ factors $`b_1^{s+1}b_2^{\pm 1}b_1^{s1}`$.
Thus, we need $`O(n^3)`$ relations to reduce the word $`w`$ to 1, or $`\mathrm{}O(n^3)`$ relations to decrease the number of factors $`v_i`$ by $`\mathrm{}(k1)`$. The lemma is proved. $`\mathrm{}`$
Now an isoperimetric inequality for the group $`H=H_{k,1}`$ can be obtained by almost the same reasoning as in Lemma 6.1 for $`H=H(v,m)`$ with the only essential difference that the function $`f(n^2)^2`$ should be replaced by $`(n^2)^4`$, since by Lemma 7.4 the “area” of a $`B`$-cell of perimeter $`s`$ (in the relators (7.1)–(7.11)) is at most $`O(s^4)`$. The property that the constructed in the proof of Theorem 2 embedding is undistorted, can be justified in the same manner as for Theorem 1.
Alexander Yu. Olshanskii
Department of Mathematics
Vanderbilt University
olsh@math.vanderbilt.edu
and
Department of Higher Algebra
MEHMAT
Moscow State University
olsh@nw.math.msu.su
Mark V. Sapir
Department of Mathematics
Vanderbilt University
http://www.math.vanderbilt.edu/$``$msapir
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# Schwinger Pair Production via Instantons in Strong Electric Fields
## I Introduction
Strong electromagnetic fields lead to two physically important phenomena: the pair-production and vacuum polarization. A strong electric field makes the quantum electrodynamics (QED) vacuum unstable which decays by emitting significantly boson or fermion pairs sau ; hei ; sch . The vacuum fluctuations of an external electromagnetic field also result in an effective action of the nonlinear Maxwell equations hei ; sch ; wei . As its long history, there have been developed many different methods such as the proper time method sch ; dew , canonical method dew , etc., to derive the QED effective action in external electromagnetic fields. Also there have been applications to various physical problems gre . The proper time method by Schwinger sch and DeWitt dew has widely been employed to compute the effective action. The real part of the effective action leads to the vacuum polarization and the imaginary part to the pair-production. Though that method is conceptually well-defined and technically rigorous, it is sometimes difficult to apply the method to some concrete physical problems such as inhomogeneous electromagnetic fields and others. On the other hand, the canonical method dew proves quite efficient in calculating the pair-production rate of bosons and fermions by static or time-dependent uniform electric fields in many physical contexts.
In canonical approach the most frequently used gauge for the electromagnetic potential is the time-dependent gauge. In that gauge the Klein-Gordon equation for bosons or the Dirac equation for fermions in a uniform electric field, when appropriately mode-decomposed, takes the form of time-dependent Schrödinger equations. Now the pair-production by the external electric field is analogous to the particle production by a time-dependent metric of a curved spacetime par ; par2 ; par3 . In both problems one imposes the same boundary condition that an incident, positive frequency component in the past infinity is scattered by a potential barrier into a superposition of positive and negative frequency components in the future infinity. It is the complex conjugate of the boundary condition for scattering problems in quantum mechanics. The coefficients determine the Bogoliubov transformation and, in particular, the coefficient of the negative frequency component gives the number created bosons or fermions per mode. The pair-production rates were calculated for time-varying electric fields bre ; nar ; pop ; pop2 ; cor ; bal ; klu ; dun . Using both canonical and path integral methods, the pair-production in a uniform electric field was studied in the time-dependent gauge pad ; pad2 and in Rindler coordinates gab . The pair-production was also studied for a uniform electric field confined to a finite region, an inhomogeneous field wan ; mar ; mar2
A shortcoming of the time-dependent gauge is that except for uniform fields, the gauge potential and thereby the Klein-Gordon equation involve both the space and time coordinates at the same time. So it is technically difficult to apply the Bogoliubov transformation for inhomogeneous fields. On the other hand, in the space-dependent (Coulomb) gauge for a static electric field, each mode of the Klein-Gordon equation for bosons or the Dirac equation for fermions takes the form of a time-independent Schrödinger equation for quantum tunnelling through a potential barrier. In that space-dependent gauge there is no direct interpretation of wave components in terms of positive and negative frequencies. However, in the case of the static uniform electric field, Brezin and Itzykson explained the dominant contribution to the pair-production rate by quantum tunnelling through the potential barrier bre ; itz , and Casher et al rederived the Schwinger’s pair-production rate by semiclassical tunnelling calculation neu ; neu2 ; neu3 . Nikishov found the pair-production rate in scattering matrix formalism for a uniform field and an inhomogeneous field of Sauter type gauge potential nik . Hansen and Ravndal showed that the transmission probability through the barrier of a uniform electric field gives the probability for pair-production for bosons and fermions han , solving the Klein paradox kle ; hun ; sau . Also Padmanabhan pad2 suggested that the reflection probability of the scattering problem gives the correct relative probability for the vacuum-to-vacuum transition for bosons. The role of tunnelling solutions for pair-production was also noticed in Refs.pop ; ste ; bro ; bro2 ; pare .
The purpose of this paper is to interpret and derive the boson or fermion pair-production rate by strong static uniform or inhomogeneous electric fields in terms of instantons through potential barriers in the space-dependent gauge in any spacetime dimensions. This formula in terms of the instanton action may provide a simple way to estimate the pair-production rate by inhomogeneous electric fields, for instance, from charged black holes, neutron stars, or astrophysical objects ruf ; ruf2 . For these static inhomogeneous fields it is easier to apply the space-dependent gauge than the time-dependent gauge. We propose that the single- and multi-instantons for quantum tunnelling determine somehow the single- and multi-pair production. In particular, we show that all the contributions from multi-instantons and anti-instantons yield exactly the total tunnelling probability for the static uniform electric field, and thereby determine the relative vacuum-to-vacuum transition and the boson pair-production rate. We further show that the instanton interpretation together with the Pauli blocking gives correctly the fermion production rate by the static uniform electric field. Using the formula in terms of the instanton action, we find the pair-production rates for bosons and fermions which are asymptotically valid for extremely strong electric fields. Also the pair-production rates for bosons and fermions by a static inhomogeneous electric field are calculated using WKB (adiabatic) approximation for the instantons. Finally we show that according to the instanton interpretation a static localized magnetic field does not lead to any pair-production, confirming the result from the proper time method.
The organization of this paper is as follows. In Sec. II we show that the tunnelling probability by instantons gives correctly the pair-production rates for bosons and fermions by a static uniform electric field. We calculate the pair-production rates in any spacetime dimensions and find their asymptotic form for extremely strong field and compare them with those from other methods. In Sec. III we extend the instanton interpretation of pair-production to an inhomogeneous electric field and find the pair-production rates in terms of the instanton action. In Sec. IV we apply the idea to a static magnetic field to show that any pair of boson or fermion are not produced. This resolves some of the puzzling issue in the canonical method on the pair-production by a static localized magnetic field.
## II Uniform Electric Field
We consider a charged boson in a static uniform electric field in a $`(d+1)`$-dimensional Minkowski spacetime. It satisfies the Klein-Gordon equation (in units of $`\mathrm{}=c=1`$)
$$\left[\eta ^{\mu \nu }\left(\frac{}{x^\mu }+iqA_\mu \right)\left(\frac{}{x^\nu }+iqA_\nu \right)+m^2\right]\mathrm{\Phi }(t,𝐱)=0,$$
(1)
where $`q`$ is the charge and $`m`$ the mass of the boson. In the space-dependent (Coulomb) gauge, the vector potential for the uniform electric field in the $`x_{}`$-direction is given by
$$A_\mu (t,𝐱)=(E_0x_{},0,\mathrm{},0).$$
(2)
Each Fourier-mode of the boson field
$$\mathrm{\Phi }(t,𝐱)=e^{i(𝐤_{}𝐱_{}\omega t)}\varphi _{\omega ,𝐤_{}}(x_{}),$$
(3)
satisfies the one-dimensional equation
$$\left[\frac{1}{2}\frac{d^2}{dx_{}^2}\frac{1}{2}\left(\omega +qE_0x_{}\right)^2\right]\varphi _{\omega ,𝐤_{}}(x_{})=\frac{1}{2}(m^2+𝐤_{}^2)\varphi _{\omega ,𝐤_{}}(x_{}).$$
(4)
Now one may interpret Eq. (4) as a Schrödinger-like equation for a unit mass moving in the inverted harmonic potential with the center at $`x_{,c}=\omega /(qE_0)`$ and the energy $`ϵ=(m^2+𝐤_{}^2)/2`$. As the energy is negative $`(ϵ<0)`$, Eq. (4) indeed describes a tunnelling problem for all transverse momenta $`𝐤_{}`$. The wave function describing the tunnelling process is given by the complex parabolic cylindrical function abr
$$\varphi _{\omega ,𝐤_{}}(\xi )=cE(a_𝐤_{},\xi ),$$
(5)
where $`c`$ is a complex number, and
$$\xi =\sqrt{\frac{2}{qE_0}}(\omega +qE_0x_{}),a_𝐤_{}=\frac{m^2+𝐤_{}^2}{2qE_0}.$$
(6)
It has the asymptotic forms in two regimes
$`\varphi _{\omega ,𝐤_{}}(\xi )`$ $`=`$ $`A\phi _{\omega ,𝐤_{}}(\xi )B\phi _{\omega ,𝐤_{}}^{}(\xi ),(\xi 2\sqrt{a_𝐤_{}}),`$
$`\varphi _{\omega ,𝐤_{}}(\xi )`$ $`=`$ $`C\phi _{\omega ,𝐤_{}}^{}(\xi ),(\xi 2\sqrt{a_𝐤_{}}),`$ (7)
where
$$\phi _{\omega ,𝐤_{}}(\xi )=\sqrt{\frac{2}{|\xi |}}e^{\frac{i}{4}\xi ^2}.$$
(8)
Here the coefficients are given by
$$A=ic\sqrt{1+e^{2\pi a_𝐤_{}}},B=ice^{\pi a_𝐤_{}},C=c.$$
(9)
In the region $`\xi 2\sqrt{a_𝐤_{}}`$, the components $`\phi _{\omega ,𝐤_{}}e^{i\omega t}`$ describes an incoming particle, from and $`\phi _{\omega ,𝐤_{}}^{}e^{i\omega t}`$ an outgoing particle, to $`\xi =\mathrm{}`$, whereas in the region $`\xi 2\sqrt{a_𝐤_{}}`$ the component $`\phi _{\omega ,𝐤_{}}^{}e^{i\omega t}`$ describes an incoming anti-particle from $`\xi =+\mathrm{}`$. Hansen and Ravndal showed that the transmission probability $`|C/A|^2`$ gives the probability for one-pair production han . Also Padmanabhan suggested that the reflection probability $`|B/A|^2`$ gives the relative probability for the vacuum-to-vacuum transition pad ; pad2 . His interpretation implies that $`\phi _{\omega ,𝐤_{}}(\mathrm{})e^{i\omega t}`$ and $`\phi _{\omega ,𝐤_{}}^{}(\mathrm{})e^{i\omega t}`$ correspond to the incoming and outgoing vacuum state, respectively. Extending their arguments to any static field, we further propose that the single- and multi-pair production of bosons are related to the single- and multi-instantons of potential barrier in such a way that the tunnelling probability $`P^\mathrm{t}`$ gives the probability for the pair-production and therefore the relative probability for the vacuum-to-vacuum transition is given by the probability for the no-pair production $`P^{\mathrm{n}\mathrm{p}}=1P^\mathrm{t}`$. Here and from now on we restrict the tunnelling probability to the transmission probability through potential barrier but exclude any nonzero transmission probability above a potential barrier or a potential well. Further we shall assume that the tunnelling probability is accurately given by the instanton action or with its higher corrections.
To see how the instanton interpretation works for the uniform electric field, we calculate the tunnelling probability from the asymptotic form (7) and compare it with the result from the instanton calculation. As the negative energy for all momenta $`𝐤_{}`$ is below the potential barrier in Eq. (4), the tunnelling probability is given by the transmission probability
$$P_𝐤_{}^{\mathrm{b}.\mathrm{t}}=\left|\frac{C}{A}\right|^2=\frac{1}{e^{2\pi a_𝐤_{}}+1}.$$
(10)
Likewise, the probability for the no-pair production, i.e., the vacuum-to-vacuum transition, given by the reflection probability
$$P_𝐤_{}^{\mathrm{b}.\mathrm{n}\mathrm{p}}=1P_𝐤_{}^{\mathrm{b}.\mathrm{t}}=\frac{1}{1+e^{2\pi a_𝐤_{}}}=\left|\frac{B}{A}\right|^2,$$
(11)
as a consequence of the flux conservation. Hence, what is needed in finding the probability for the no-pair production (vacuum-to-vacuum transition) even in a general electric field is the corresponding total tunnelling probability via the single- and multi-instantons.
Now let us interpret the tunnelling probability (10) in terms of multi-instantons and anti-instantons of tunnelling process. In instanton physics col , the leading contribution to the tunnelling probability
$$P_𝐤_{}^\mathrm{t}=e^{2S_𝐤_{}},$$
(12)
is determined by the single-instanton action
$$S_𝐤_{}=_x_{}^{x^+}𝑑x_{}\sqrt{m^2+𝐤_{}^2\left(\omega +qE_0x_{}\right)^2}=\pi a_𝐤_{},$$
(13)
where $`x_\pm =\pm \sqrt{m^2+𝐤_{}^2}\omega `$ are the classical turning points. We propose that the single-instanton and multi-instantons may be related in a certain way with one-pair and multi-pair production, whereas multi-anti-instantons with the annihilation of created boson pairs. As there is no limitation from the Pauli blocking for the multi-pair production of bosons, the correct total tunnelling probability should take into account both multi-instantons and anti-instantons
$$P_𝐤_{}^{\mathrm{b}.\mathrm{t}}=\underset{n=1}{\overset{\mathrm{}}{}}(1)^{n+1}e^{2nS_𝐤_{}}=\frac{1}{e^{2S_𝐤_{}}+1},$$
(14)
where instantons contribute positively and anti-instantons negatively. Similarly, the relative probability for the no-pair production (vacuum-to-vacuum transition) is given by
$$P_𝐤_{}^{\mathrm{b}.\mathrm{n}\mathrm{p}}=\underset{n=0}{\overset{\mathrm{}}{}}(1)^ne^{2nS_𝐤_{}}=\frac{1}{1+e^{2S_𝐤_{}}},$$
(15)
These results agree with Eqs. (10) and (11). The physical interpretation of the alternating signs is that only the instantons of even repeated periodic motions in the inverted potential contribute positively (creating pairs) to the tunnelling probability, whereas the anti-instantons of odd repeated periodic motions contribute negatively (annihilating created-pairs) to the tunnelling probability.
The vacuum means the absence of any particle for possible physical states. So the vacuum-to-vacuum transition, i.e., the vacuum persistence, is the total relative probability for the no-pair production:
$$|0,\mathrm{out}|0,\mathrm{in}|^2=\underset{\mathrm{all}\mathrm{states}}{}P_𝐤_{}^{\mathrm{b}.\mathrm{n}\mathrm{p}}=\mathrm{exp}\left[\underset{\mathrm{all}\mathrm{states}}{}\mathrm{ln}(1+e^{2S_𝐤_{}})\right].$$
(16)
On the other hand, the vacuum-to-vacuum transition is given by the imaginary part of the effective action for boson
$$|0,\mathrm{out}|0,\mathrm{in}|^2=\mathrm{exp}\left[2VT\mathrm{Im}_{\mathrm{eff}.}^\mathrm{b}\right],$$
(17)
where $`V`$ and $`T`$ are the relevant volume and the duration of time. Therefore, the pair-production rate per unit time per unit volume is twice of the imaginary part of the effective action:
$$w^\mathrm{b}=2\mathrm{I}\mathrm{m}_{\mathrm{eff}.}^\mathrm{b}=\frac{1}{VT}\underset{\mathrm{all}\mathrm{states}}{}\mathrm{ln}(1+e^{2S_𝐤_{}}).$$
(18)
Then the pair-production rate for bosons is explicitly given by
$`w^\mathrm{b}`$ $`=`$ $`{\displaystyle \frac{(2s+1)V_{}}{V}}{\displaystyle \frac{d\omega d𝐤_{}^{d1}}{(2\pi )^d}\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^{n+1}}{n}e^{\frac{\pi n}{qE_0}𝐤_{}^2}e^{\frac{\pi m^2}{qE_0}n}}`$ (19)
$`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^d}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^{n+1}\left({\displaystyle \frac{qE_0}{n}}\right)^{(d+1)/2}e^{\frac{\pi m^2}{qE_0}n},`$
where $`s`$ is the spin of the boson. Here we used $`𝑑\omega =(qE_0)V_{}`$, where $`V_{}`$ is the longitudinal extension of the field, and $`V=V_{}V_{}`$, $`V_{}`$ being the transverse volume nik . It should be noted that Eq. (19) recovers the standard result for the boson pair-production in any dimension in Ref. gus .
The fermion pair-production can be understood similarly. The created fermion pair blocks the multi-pair production. So the total tunnelling probability for the fermion pair-production per each mode is just
$$P_𝐤_{}^{\mathrm{f}.\mathrm{t}}=e^{2S_𝐤_{}}.$$
(20)
Therefore, the relative probability for the no-pair production of fermions is now given by
$$P_𝐤_{}^{\mathrm{f}.\mathrm{n}\mathrm{p}}=1e^{2S_𝐤_{}}.$$
(21)
Finally, the fermion pair-production rate per unit time per unit volume is found to be
$$w^\mathrm{f}=2\mathrm{I}\mathrm{m}_{\mathrm{eff}.}^\mathrm{f}=\frac{1}{VT}\underset{\mathrm{all}\mathrm{states}}{}\mathrm{ln}(1e^{2S_𝐤_{}}),$$
(22)
and takes the form
$`w^\mathrm{f}`$ $`=`$ $`{\displaystyle \frac{(2s+1)V_{}}{V}}{\displaystyle \frac{d\omega d𝐤_{}^{d1}}{(2\pi )^d}\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}e^{\frac{\pi n}{qE_0}𝐤_{}^2}e^{\frac{\pi m^2}{qE_0}n}}`$ (23)
$`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^d}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{qE_0}{n}}\right)^{(d+1)/2}e^{\frac{\pi m^2}{qE_0}n}.`$
Also Eq. (23) recovers the standard result for the fermion pair-production in Ref. gus .
Though the production rate (19) for bosons and (23) for fermions are well defined for all electric fields, the series converge strongly for weak electric fields because all higher terms are exponentially suppressed. But for extremely strong electric fields the exponential terms approach to unity and the series are approximated by the Riemann eta function $`\eta (2)`$ for bosons and the Riemann zeta function $`\zeta (2)`$ for fermions. So, for strong electric fields, instead of using a special resummation of the series, we adopt directly the pair-production formula (18) and (22) and evaluate properly the integrals suitable for strong fields. In four dimensions $`(d=3)`$, the boson pair-production rate (18) becomes
$`w^\mathrm{b}`$ $`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^3}}(qE_0){\displaystyle _0^{\mathrm{}}}(2\pi )𝑑k_{}k_{}\mathrm{ln}\left(1+e^{\frac{\pi (m^2+k_{}^2)}{qE_0}}\right)`$ (24)
$`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^3}}(qE_0)^2\left\{{\displaystyle _0^{\mathrm{}}}𝑑y\mathrm{ln}(1+e^y){\displaystyle _0^{\pi m^2/qE_0}}𝑑y\mathrm{ln}(1+e^y)\right\},`$
where
$$y=\frac{\pi }{qE_0}k_{}^2+\frac{\pi m^2}{qE_0}.$$
(25)
Using the integral pbm
$$_0^{\mathrm{}}𝑑y\mathrm{ln}(1+e^y)=\frac{\pi ^2}{12},$$
(26)
and expanding the exponential and then the logarithmic function to any desired order
$$\mathrm{ln}(1+e^y)=\mathrm{ln}2\frac{1}{2}y+\frac{1}{8}y^2+\frac{1}{96}y^4+𝒪(y^5),$$
(27)
we obtain the pair-production rate
$`w^\mathrm{b}`$ $`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^3}}\{{\displaystyle \frac{\pi ^2}{12}}(qE_0)^2(\mathrm{ln}2)\pi m^2qE_0+{\displaystyle \frac{1}{4}}(\pi m^2)^2`$
$``$ $`{\displaystyle \frac{1}{24}}{\displaystyle \frac{(\pi m^2)^3}{qE_0}}{\displaystyle \frac{1}{480}}{\displaystyle \frac{(\pi m^2)^5}{(qE_0)^3}}+𝒪\left({\displaystyle \frac{(\pi m^2)^6}{(qE_0)^4}}\right)\}.`$ (28)
Similarly, the fermion pair-production rate (22) for strong fields takes the form
$`w^\mathrm{f}`$ $`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^3}}(qE_0){\displaystyle _0^{\mathrm{}}}(2\pi )𝑑k_{}k_{}\mathrm{ln}\left(1e^{\frac{\pi (m^2+k_{}^2)}{qE_0}}\right)`$ (29)
$`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^3}}(qE_0)^2\left\{{\displaystyle _0^{\mathrm{}}}𝑑y\mathrm{ln}(1e^y){\displaystyle _0^{\pi m^2/qE_0}}𝑑y\mathrm{ln}(1e^y)\right\}.`$
Using the integral pbm
$$_0^{\mathrm{}}𝑑y\mathrm{ln}(1e^y)=\frac{\pi ^2}{6},$$
(30)
and expanding the exponential function and then the logarithmic function
$$\mathrm{ln}(1e^y)=\mathrm{ln}y\frac{1}{2}y+\frac{1}{24}y^2+\frac{11}{720}y^4+𝒪(y^5),$$
(31)
we finally obtain the fermion pair-production rate
$`w^\mathrm{f}`$ $`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^3}}\{{\displaystyle \frac{\pi ^2}{6}}(qE_0)^2\pi m^2qE_0(\mathrm{ln}\left({\displaystyle \frac{qE_0}{\pi m^2}}\right)+1){\displaystyle \frac{1}{4}}(\pi m^2)^2`$
$`+`$ $`{\displaystyle \frac{1}{72}}{\displaystyle \frac{(\pi m^2)^3}{qE_0}}+{\displaystyle \frac{11}{3600}}{\displaystyle \frac{(\pi m^2)^5}{(qE_0)^3}}+𝒪\left({\displaystyle \frac{(\pi m^2)^6}{(qE_0)^4}}\right)\}.`$ (32)
The fermion pair-production rate (32) for strong electric fields confirms the result obtained from different methods in Refs. dit ; hey ; sol .
A comment is in order. The Schwinger pair-production by a static uniform electric field is an ideal calculation in which one neglects the pair-production due to the interactions of the created pairs with the electric field background and among the created pairs. For instance, a single-pair can produce another pair through the interaction with the electric field, whose rate is proportional to $`(qE_0/m^2)^2(q/m)^2`$ jau and can be larger than the multi-pair production rate, $`e^{(\pi m^2n)/(qE_0)}(qE_0)^2/n^2`$, from multi-instantons for all sufficiently large $`n`$ even for an extremely strong electric field $`E_0`$. However, we shall not consider this complicated real situation but rather focus on the ideal calculation without the back-reaction of produced pairs.
## III Inhomogeneous Electric Fields
We now consider the pair-production by a static inhomogeneous electric field. Without loss of generality, the electric field is assumed to be localized in the $`x_{}`$-direction and to have the gauge potential
$$A_\mu (t,𝐱)=(A_0(x_{}),0,\mathrm{},0),$$
(33)
where $`E(x_{})=dA_0(x_{})/dx_{}`$. We restrict only to the case where all produced particles $`(q>0)`$ and antiparticles reach the asymptotic regions $`x=+\mathrm{}`$ and $`x=\mathrm{}`$, respectively, without being bounded by the electric field. This requires that $`qA_0(\mathrm{})qA_0(+\mathrm{})2m`$. The mode-decomposed Klein-Gordon equation then takes the form
$$\left[\frac{1}{2}\frac{d^2}{dx_{}^2}\frac{1}{2}\left(\omega qA_0(x_{})\right)^2\right]\varphi _{\omega ,𝐤_{}}(x_{})=\frac{1}{2}(m^2+𝐤_{}^2)\varphi _{\omega ,𝐤_{}}(x_{}).$$
(34)
We can still interpret Eq. (34) as a one-dimensional quantum system of a unit mass with the potential $`(\omega qA_0(x_{}))^2/2`$ and the energy $`ϵ=(m^2+𝐤_{}^2)/2`$. In the WKB (adiabatic) approximation the asymptotic form for the tunnelling probability for each mode $`𝐤_{}`$ is given by fro ; fro2 ; ben
$$P_𝐤_{}^{\mathrm{b}.\mathrm{t}}=\frac{1}{e^{2S_𝐤_{}}+1},$$
(35)
where
$$S_𝐤_{}=\underset{n=0}{\overset{\mathrm{}}{}}S_𝐤_{}^{(2n)}.$$
(36)
Here the leading contribution to $`S_𝐤_{}`$ is given by the instanton action
$$S_𝐤_{}^{(0)}=_x_{}^{x_+}𝑑x_{}\left[Q_𝐤_{}(x)\right]^{1/2},$$
(37)
and the next-to-leading term by
$$S_𝐤_{}^{(2)}=_x_{}^{x_+}𝑑x_{}\left[\frac{1}{8}\frac{Q_𝐤_{}^{\prime \prime }(x)}{Q_𝐤_{}^{3/2}(x)}\frac{5}{32}\frac{Q_𝐤_{}^{}_{}{}^{}2(x)}{Q_𝐤_{}^{5/2}(x)}\right],$$
(38)
where
$$Q_𝐤_{}(x)=m^2+𝐤_{}^2\left(\omega qA_0(x_{})\right)^2.$$
(39)
Hence the relative probability for the no-pair production (vacuum-to-vacuum transition) of bosons is given by
$$P_𝐤_{}^{\mathrm{b}.\mathrm{n}\mathrm{p}}=\frac{1}{1+e^{2S_𝐤_{}}},$$
(40)
and for fermions by
$$P_𝐤_{}^{\mathrm{f}.\mathrm{n}\mathrm{p}}=1e^{2S_𝐤_{}}.$$
(41)
A few comments are in order. First, if the electric field extends over all the space as in the uniform field case and has the potential $`|A_0(\pm \mathrm{})|=\mathrm{}`$, then the potential barrier decreases indefinitely at both $`\pm \mathrm{}`$. Therefore, there are always instantons for all $`𝐤_{}`$. Second, if the electric field is localized or has finite values of the potential at $`\pm \mathrm{}`$, then pairs are produced only when $`\omega qA_0(+\mathrm{})m`$ and $`\omega qA_0(\mathrm{})m`$. So there is a change of the sign of $`\omega qA_0(x_{})`$ implying a potential barrier. Thus only those modes belonging to $`|𝐤_{}|k_{,\mathrm{max}}`$ have finite instantons and lead to pair-production, where the upper limit is given by the minimum of two asymptotic values
$$k_{,\mathrm{max}}^2=\mathrm{Min}\{\left(\omega qA_0(+\mathrm{})\right)^2m^2,\left(\omega qA_0(\mathrm{})\right)^2m^2\}.$$
(42)
In the inhomogeneous electric field, we obtain the boson pair-production rate per unit time per unit volume
$`w^\mathrm{b}`$ $`=`$ $`2\mathrm{I}\mathrm{m}_{\mathrm{eff}.}^\mathrm{b}`$ (43)
$`=`$ $`{\displaystyle \frac{(2s+1)}{VT}}{\displaystyle \underset{\mathrm{all}\mathrm{allowed}\mathrm{states}}{}}\mathrm{ln}\left(1+e^{2S_𝐤_{}}\right)`$
$`=`$ $`{\displaystyle \frac{(2s+1)}{(2\pi )^dV_{}}}{\displaystyle \frac{(d1)\pi ^{(d1)/2}}{\mathrm{\Gamma }(\frac{d+1}{2})}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{n+1}}{n}}{\displaystyle _{qA_0(+\mathrm{})+m}^{qA_0(\mathrm{})m}}𝑑\omega {\displaystyle _0^{k_{,\mathrm{max}}}}𝑑k_{}k_{}^{d2}e^{2nS_𝐤_{}},`$
and the fermion pair-production rate
$`w^\mathrm{f}`$ $`=`$ $`2\mathrm{I}\mathrm{m}_{\mathrm{eff}.}^\mathrm{f}`$ (44)
$`=`$ $`{\displaystyle \frac{(2s+1)}{2VT}}{\displaystyle \underset{\mathrm{all}\mathrm{allowed}\mathrm{states}}{}}\mathrm{ln}\left(1e^{2S_𝐤_{}}\right)`$
$`=`$ $`{\displaystyle \frac{(2s+1)}{2(2\pi )^dV_{}}}{\displaystyle \frac{(d1)\pi ^{(d1)/2}}{\mathrm{\Gamma }(\frac{d+1}{2})}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n}}{\displaystyle _{qA_0(+\mathrm{})+m}^{qA_0(\mathrm{})m}}𝑑\omega {\displaystyle _0^{k_{,\mathrm{max}}}}𝑑k_{}k_{}^{d2}e^{2nS_𝐤_{}}.`$
As an exactly solvable model we consider a localized electric field $`E(x_{})=E_0\mathrm{sech}^2(x_{}/L)`$ with the Sauter type gauge potential sau ; nik
$$A_0(x_{})=E_0L\mathrm{tanh}(\frac{x_{}}{L}).$$
(45)
In the limit of $`L\mathrm{}`$ the gauge potential (45) reduces to the uniform electric field in Sec. II. Since the gauge potential (45) is a more general case including the uniform field as a special case, it is worthy to apply the instanton interpretation to the pair-production and compare the result with the exact one. Bosons gain an additional contribution to momenta from the acceleration by the localized electric field and have asymptotic values at $`x_{}\pm \mathrm{}`$:
$$k_{}^2(\mathrm{})=(qE_0L+\omega )^2m^2𝐤_{}^2,k_{}^2(\mathrm{})=(qE_0L\omega )^2m^2𝐤_{}^2.$$
(46)
In the large $`L`$ limit the instanton action (37) is given by
$$S_𝐤_{}=\pi \frac{m^2+𝐤_{}^2}{2qE_0}\left[1+\frac{\omega ^2}{(qE_0L)^2}+\frac{m^2+𝐤_{}^2}{4(qE_0L)^2}+𝒪(\frac{1}{L^4})\right].$$
(47)
The exact wave function describing the tunnelling process is found
$$\varphi _{\omega ,𝐤_{}}(x_{})=Ce^{\mu \frac{x_{}}{L}}\mathrm{sech}^\nu (\frac{x_{}}{L})F(\alpha ,\beta ;\gamma ;\zeta ),$$
(48)
where $`F`$ is the hypergeometric function and
$`\mu `$ $`=`$ $`i{\displaystyle \frac{L}{2}}\left(k_{}(\mathrm{})+k_{}(\mathrm{})\right),\nu =i{\displaystyle \frac{L}{2}}\left(k_{}(\mathrm{})k_{}(\mathrm{})\right),`$
$`\alpha `$ $`=`$ $`\nu +{\displaystyle \frac{1}{2}}+i\sqrt{(qE_0L^2)^2{\displaystyle \frac{1}{4}}},\beta =\nu +{\displaystyle \frac{1}{2}}i\sqrt{(qE_0L^2)^2{\displaystyle \frac{1}{4}}},`$
$`\gamma `$ $`=`$ $`1iLk_{}(\mathrm{}),\zeta ={\displaystyle \frac{1}{2}}\left(1\mathrm{tanh}({\displaystyle \frac{x_{}}{L}})\right).`$ (49)
In the limit of $`x_{}L`$, Eq. (48) has the asymptotic form
$$\varphi _{\omega ,𝐤_{}}(x_{})=2^\nu Ce^{ik_{}(\mathrm{})x_{}}.$$
(50)
It describes a wave function after tunnelling (an incoming anti-particle). In the limit of $`x_{}L`$, we may use another form for Eq. (48)
$`\varphi _{\omega ,𝐤_{}}(x_{})=Ce^{\mu \frac{x_{}}{L}}\mathrm{sech}^\nu ({\displaystyle \frac{x_{}}{L}})[{\displaystyle \frac{\mathrm{\Gamma }(\gamma )\mathrm{\Gamma }(\gamma \alpha \beta )}{\mathrm{\Gamma }(\gamma \alpha )\mathrm{\Gamma }(\gamma \beta )}}F(\alpha ,\beta ;\gamma ;1\zeta )`$
$`+{\displaystyle \frac{\mathrm{\Gamma }(\gamma )\mathrm{\Gamma }(\alpha +\beta \gamma )}{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\beta )}}(1\zeta )^{\gamma \alpha \beta }F(\alpha ,\beta ;\gamma ;1\zeta )].`$ (51)
In this limit the first term of Eq. (51) describes an incident wave (an incoming particle) having the asymptotic form
$$\varphi _{\omega ,𝐤_{}}(x_{})=2^\nu C\frac{\mathrm{\Gamma }(\gamma )\mathrm{\Gamma }(\gamma \alpha \beta )}{\mathrm{\Gamma }(\gamma \alpha )\mathrm{\Gamma }(\gamma \beta )}e^{ik_{}(\mathrm{})x_{}}.$$
(52)
Therefore, from Eqs. (50) and (52) we can find the probability for tunnelling
$`P_\mathrm{k}_{}^{\mathrm{b}.\mathrm{t}}`$ $`=`$ $`{\displaystyle \frac{k_{}(\mathrm{})}{k_{}(\mathrm{})}}\left|{\displaystyle \frac{\mathrm{\Gamma }(\gamma \alpha )\mathrm{\Gamma }(\gamma \beta )}{\mathrm{\Gamma }(\gamma )\mathrm{\Gamma }(\gamma \alpha \beta )}}\right|^2`$ (53)
$`=`$ $`{\displaystyle \frac{\mathrm{sinh}\pi \left(Lk_{}(\mathrm{})\right)\mathrm{sinh}\pi \left(Lk_{}(\mathrm{})\right)}{\mathrm{cosh}\pi \left(\frac{L}{2}(k_{}(\mathrm{})+k_{}(\mathrm{}))+Q\right)\mathrm{cosh}\pi \left(\frac{L}{2}(k_{}(\mathrm{})+k_{}(\mathrm{}))Q\right)}},`$
where $`Q=\sqrt{(qE_0L^2)^21/4}`$. In the large $`L`$ limit we obtain approximately the probability for tunnelling
$$P_\mathrm{k}_{}^{\mathrm{b}.\mathrm{tun}.}=\frac{1}{1+e^{2S_𝐤_{}}}.$$
(54)
Here, we used the binomial expansion
$`k_{}(\pm \mathrm{})`$ $`=`$ $`(qE_0L\pm \omega )\left[1{\displaystyle \frac{m^2+𝐤_{}^2}{(qE_0L)^2\left(1\pm \frac{\omega }{(qE_0L)^2}\right)^2}}\right]^{1/2}`$ (55)
$`=`$ $`qE_0L\pm \omega {\displaystyle \frac{m^2+𝐤_{}^2}{2qE_0L}}\left[1{\displaystyle \frac{\omega }{qE_0L}}+{\displaystyle \frac{\omega ^2}{(qE_0L)^2}}\right]`$
$`{\displaystyle \frac{(m^2+𝐤_{}^2)^2}{8(qE_0L)^3}}\left[1{\displaystyle \frac{3\omega }{qE_0L}}+{\displaystyle \frac{6\omega ^2}{(qE_0L)^2}}\right]+\mathrm{}.`$
Therefore, using instanton action (47) we are able to obtain the pair-production rate for bosons and fermions according to Eqs. (43) and (44). Thus we have shown that in the space-dependent gauge the instanton interpretation for wave function gives correctly the pair-production rates for bosons and fermions for two exactly solvable models.
## IV Magnetic Fields
A static uniform magnetic field leads only to a real effective action and thus implies no-pair production sch . Recently it has also been shown that any static magnetic field, having no imaginary part, does not lead to the pair-production dun ; sri ; dun2 . On the other hand, in canonical method, each mode of the Klein-Gordon or Dirac equation has a nonzero reflection probability for static magnetic fields. This issue has been raised and discussed to interpret the reflection probability as the pair-production by static localized magnetic fields in Ref. sri . In this section we resolve this issue from the view point of the instanton interpretation.
Let us consider a static magnetic field in a 4-dimensional spacetime with the gauge potential
$$A_\mu (t,𝐱)=(0,A_1(x_2),0,0).$$
(56)
The magnetic field is given by $`𝐁=(dA_1(x_2)/dx_2)\widehat{𝐱}_3`$. The Klein-Gordon equation has the form
$$\left[\frac{^2}{t^2}\left(\frac{}{x_1}+iqA_1(x_2)\right)^2\frac{^2}{x_2^2}\frac{^2}{x_3^2}+m^2\right]\mathrm{\Phi }(t,𝐱)=0.$$
(57)
As in the case of the electric field, each mode of the field
$$\mathrm{\Phi }(t,𝐱)=e^{i(k_1x_1+k_3z_3\omega t)}\varphi _{\omega ,k_1,k_3}(x_2)$$
(58)
leads to a Schrödinger-like equation
$$\left[\frac{1}{2}\frac{d^2}{dx_2^2}+\frac{1}{2}\left(k_1qA_1(x_2)\right)^2\right]\varphi _{\omega ,k_1,k_3}(x_2)=\frac{1}{2}(\omega ^2m^2k_3^2)\varphi _{\omega ,k_1,k_3}(x_2).$$
(59)
As a one-dimensional quantum system, Eq. (59) has the potential $`(k_1qA_1(x_2))^2/2`$ and the energy $`(\omega ^2m^2k_3^2)/2`$.
In the case of a uniform magnetic field, the gauge potential $`A_1(x_2)=B_0x_2`$ is indefinitely unbounded at $`x_2=\pm \mathrm{}`$. Then the potential of Eq. (59) is exactly that of a harmonic oscillator and the energy is quantized
$$ϵ_n=qB_0(2n+1).$$
(60)
The quantized energy has been used to calculate the effective action in the uniform magnetic field hei . From the view point of instanton interpretation, there is no pair-production since there are no finite instantons at all. All would-be instantons from one spatial infinity to another are infinite and do not contribute to the tunnelling probability. This result agrees with that obtained from the proper time and other methods.
We now consider a static localized magnetic field $`𝐁(x_2)=B_0\mathrm{sech}^2(x_2/L)\widehat{𝐱}_3`$. The gauge potential is given by $`A_1(y)=B_0L\mathrm{tanh}(x_2/L)`$. The gauge potential in Eq. (59) has two asymptotic values $`(k_1\pm qB_0L)^2/2`$ at $`x_2=\pm \mathrm{}`$, respectively, and a minimum value in-between. There is a potential well instead of a potential barrier, so the reflection probability may be nonzero, though the tunnelling probability via instantons is zero. Therefore, according to the instanton interpretation, there is no pair-production. The instanton interpretation thus resolves the contradiction between the effective action and the canonical method raised in Ref. sri . Our result also agrees with that by Dunne and Hall, who showed that the imaginary part of the effective action vanishes for the static magnetic field considered above and therefore neither bosons nor fermions are produced in pairs dun2 ; dun3 .
## V Discussion and Conclusion
In this paper we have studied the pair-production of bosons and fermions by static uniform or inhomogeneous electric field. For these fields we used the space-dependent (Coulomb) gauge and solved the Klein-Gordon equation. For strong electric fields the mode-decomposed Klein-Gordon equations have potential barriers from the gauge potential. The set of wave functions describing pair-production in quantum field theory is the same of the standard scattering problem through potential barriers in quantum mechanics pad ; pad2 ; han , in contract with the time-dependent gauge. Together with the fact that the most dominant contribution to pair-production rate is the single instanton bre ; itz ; neu ; neu2 ; neu3 , we further propose that all multi-instantons contribute to the pair-production and anti-multi-instantons to the annihilation of created pairs and that the total tunnelling probability from all multi-instantons and anti-multi-instantons is related with pair-production and the probability for the vacuum-to-vacuum transition is the probability for no-pair production. Based on this we derived the pair-production formula for bosons (18) and fermions (22).
This instanton interpretation means that the single-instanton is related in a certain way with the single-pair production, multi-instantons with the multi-pair production and anti-multi-instantons with the annihilation of the created pairs. In fact, when the instanton action is large, the single-instanton is the dominant contribution to one-pair production and multi-instantons are the dominant contribution to the multi-pair production. Also it implies the no-pair production when there is not any tunnelling instanton. In the case of a uniform electric field, when all the contributions from multi-instantons and anti-instantons are taken into account, the pair-production rate for bosons calculated according to the instanton interpretation recovers the well-known result from the proper time method. By taking the Pauli-blocking into account the pair-production rate for fermions is found also to agree with the standard result. Further the pair-production formula from instanton action yields the correct forms (28) and (32) for extremely strong electric fields, confirming the consistency of the formula with other methods. Using the instantons obtained in the WKB (adiabatic) approximation we are also able to provide the formula for the pair-production rate of bosons and fermions by inhomogeneous electric fields.
As a by-product we are able to show that any static (localized) magnetic field cannot produce pairs of bosons or fermions. In the case of magnetic fields the space-dependent gauge reduces the Klein-Gordon equation to time-independent Schrödinger equations with potential wells instead of potential barriers of the electric field case. Since potential wells cannot have finite instantons and possible infinite instantons from either side of potential wells give the zero probability for pair-production, the pair-production of bosons or fermions cannot proceed. As the nonzero transmission probability through potential wells is the result without any finite instanton, the instanton interpretation excludes the possibility of pair-production by a static localized magnetic field in the canonical method raised in Ref. sri . Therefore we conclude that any static magnetic field does not lead to pair-production and the canonical method equipped with the instanton interpretation is compatible with the effective action method dun ; sri ; dun2 .
###### Acknowledgements.
We would like to thank F. C. Khanna and L. Sriramkumar for many useful discussions, G. Dunne for comments on no-pair production by a static magnetic field and H. Neuberger for useful information. S.P.K. would like to express his appreciation for the warm hospitality of the Theoretical Physics Institute, University of Alberta. The work of S.P.K. was supported by the Korea Science and Engineering Foundation under Grant No. 1999-2-112-003-5 and the work of D.N.P. by the National Sciences and Engineering Research Council of Canada.
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# Wide Field Imaging of the Hubble Deep Field-South Region I: Quasar Candidates 1footnote 11footnote 1Based on observations obtained at Cerro Tololo Inter-American Observatory, a division of the National Optical Astronomy Observatories, which is operated by the Association of Universities for Research in Astronomy, Inc. under cooperative agreement with the National Science Foundation.
## 1. Introduction
The recent Hubble Deep Field-South campaign (HDF–S; Williams et al. 2000) is distinguished in part from the hugely successful Hubble Deep Field-North (HDF–N; Williams et al. 1996) by ultraviolet (UV) spectroscopy (Ferguson et al. 2000) of the bright, radio-quiet quasar J2233–6033 (HDF–S QSO; $`V17.5`$, $`z=2.24`$; Boyle 1997), using the Space Telescope Imaging Spectrograph (STIS; Woodgate et al. 1998). Medium- and high-resolution STIS spectroscopy provides detailed coverage of the Ly$`\alpha `$ forest from $`1.2<z<1.93`$. Optical spectroscopy from the ground (Outram et al. 1999, Sealey et al. 1998, Savaglio et al. 1998, Cristiani et al. 2000) extends coverage of the Ly$`\alpha `$ forest to $`z=`$ 2.24 and also covers various metal lines along the QSO line of sight. The UV spectrum of the HDF–S QSO provides a unique window on the universe over this redshift interval where very little is known due to observational difficulties.
The HDF–S QSO spectra have revealed several features in the distribution of Ly$`\alpha `$ gas clouds (Savaglio et al. 1999), including a “void” spanning $`z=1.3831.460`$ (94 co-moving $`\mathrm{h}_{65}^1`$ Mpc)<sup>2</sup><sup>2</sup>2We assume H = 65 $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$, q = 0.5, $`\mathrm{\Lambda }=0`$, and co-moving distances throughout this paper., and two overdensities at $`z1.335`$ and $`z=1.921.99`$ (covering 64 $`\mathrm{h}_{65}^1`$ Mpc). Both Ly$`\alpha `$ overdensities have associated metal systems: one tentatively identified at $`z=1.336`$ (Ferguson et al. 2000), and two at $`z=1.929`$ and $`z=1.943`$ (Savaglio et al. 1998). There is a QSO, J2233–6032, at a redshift $`z=1.335`$ and a projected distance of 44″ (0.7 $`\mathrm{h}_{65}^1`$ Mpc) from the HDF–S QSO (Tresse et al. 1998). In addition, there are three galaxies with redshifts of $`z=`$ 1.268, 1.315 and 1.218 at a distance $`2\mathrm{h}_{65}^1`$ Mpc from the HDF–S QSO line-of-sight (Glazebrook et al. 2000). A concentrated ground based search for other objects associated with these features will be required to better understand their nature.
We have performed a multi-color ($`uBVRI`$+narrow-band) imaging survey of 1/2 square degree centered on the HDF–S (Palunas et al. 2000). The narrow-band filter (NB) was centered at 3958Å to detect Ly$`\alpha `$ at the redshift of the HDF–S QSO. Our goals for the survey include determining photometric redshifts of galaxies over a wide field around the HDF–S , searching for emission-line objects near the HDF–S QSO, and identifying other quasars in the field. Here we present candidate quasars in our catalog. We confirm the presence of Ly$`\alpha `$ nebulosity extending $`12\mathrm{}`$ around the HDF–S QSO, reported by Bergeron et al. (1999). We detect 10 point-like objects in emission through the NB filter. Of these, 7 satisfy our QSO color selection criteria. We report the discovery of a $`B20`$ quasar at $`z=1.56`$ just 6.7′ (6.8 $`\mathrm{h}_{65}^1`$ Mpc) from the HDF–S QSO line of sight and $``$12″ from the western edge of the WFPC2 field. Images and catalogs from the survey are available on the world wide web at http://hires.gsfc.nasa.gov/$``$gardner/hdfs.
## 2. Observations and Reductions
We observed the HDF–S region using the Big Throughput Camera (BTC) on the CTIO 4m during 1998 September. A complete description of the reductions can be found in Palunas et al. (2000). Images were taken in the Sloan Digital Sky Survey $`u`$, Johnson $`B`$ and $`V`$, and Cousins $`R`$ and $`I`$ filters. We also imaged the field using a narrow-band interference filter centered at 3958Å and 50 Å wide tuned to Ly$`\alpha `$ at the redshift of the HDF–S QSO. The BTC is a prime-focus mosaic camera consisting of 4 CCDs, each 2048$`\times `$2048 pixels (Wittman et al. 1998). The camera has a 0.43″ pixel size and covers an overall area of (34.8′)<sup>2</sup> with 5.4′ gaps between the CCDs. Individual exposures were dithered to fill in the gaps between the CCDs, which resulted in a contiguous field with non-uniform coverage. The field extends $`44.0\mathrm{}\times 44.0\mathrm{}`$; the total area covered, excluding regions around bright stars, is 1716 $`\mathrm{}^{\mathrm{}}`$. The NB imaging covers less area, 1352$`\mathrm{}^{\mathrm{}}`$.
The conditions were photometric. The photometry was calibrated using Landolt standards, and fixed to the Johnson $`UBV`$ and Cousins $`RI`$ photometric system as defined by Landolt (1992). Comparison of the Johnson $`U`$ and Sloan $`u`$ filters demonstrates that small color corrections are sufficient to place our observations onto the same system as Landolt.
The seeing was 1.2″–1.5″. Bias, dark subtraction and flatfielding of the images using superflats, were done using IRAF<sup>3</sup><sup>3</sup>3IRAF is distributed by NOAO, which is operated by AURA Inc., under contract to the NSF.. Astrometric corrections for the wide-field optical distortions were derived and the corrected images were combined using *btcproc* (Wittman et al. 1998). Catalogs of individual sources were compiled using SExtractor (Bertin & Arnouts 1996). The limiting magnitudes vary across the images. The 5$`\sigma `$ detection limits for point-like objects in each color range between $`24<u<25`$, $`25.6<B<26.6`$, $`25<V<26`$, $`25<R<25.8`$, $`23.5<I<24.4`$, with about 50% of the area of each image at the brightest limits. For the NB image the 5$`\sigma `$ limits range between 0.5—1.2 $`\times 10^{16}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$.
## 3. Quasar candidates
Most quasars are easily distinguished from other point sources on the basis of their unusual colors, particularly in the UV (e.g. Hall et al. 1996a, 1996b). We transform our colors from the Landolt to Cousins photometric system that Hall et al. use (Bessel 1991), and denote the transformed colors with a subscript C. The corrections are small and do not affect the results. Figure 6 shows the $`(uB)_\mathrm{C}`$ vs. $`(BV)_\mathrm{C}`$ colors of the point sources in our catalog. Our QSO selection criteria are given in Table 1. Hall et al.’s selection criteria are designed to exclude the dense stellar locus running diagonally across the diagram and compact emission-line galaxies which populate the upper right of the diagram. Their criteria are adjusted as a function of magnitude to account for broadening of the stellar locus due to photometric errors. Our criteria differ slightly at fainter magnitudes because our photometry is deeper and therefore the stellar locus exhibits less broadening.
Point sources in our data were selected on the basis of the FWHM and “class star” (cl) parameters from SExtractor. The selection criteria were determined by comparing to the sources classified in the HDF–S flanking fields (Lucas, R. et al. 2000). The following parameters maximized the number of point sources selected without excessive contamination by extended sources: for $`R20`$, FWHM $`<1.8`$″; for $`20<R23`$, FWHM $`<`$ 1.8″ and cl $`>0.95`$; and for $`R>23`$, FWHM $`<1.9\mathrm{}`$ and cl $`>.92`$. Full discussion will be provided elsewhere (Palunas et al. 2000, in preparation).
Using the above criteria we find 4 $`B<19`$, 13 $`B<20`$, 30 $`B<21`$, 71 $`B<22`$ and 154 $`B<23`$ QSO candidates in our field. Figure 6 shows the distribution of QSO candidates in the field. The QSO selection criteria identify quasars at $`z<2.3`$. Colors for QSOs at higher redshifts lie within the stellar locus. Hall et al. find a 60% QSO selection efficiency for objects with these colors at $`B<22`$. In Figure 6 we compare the integrated number counts of the QSO candidates we have selected, assuming 60% selection efficiency, to the QSO counts of Hartwick & Shade (1990).
## 4. The HDF–S QSO
Bergeron et al. (1999) report large scale extended Ly$`\alpha `$ nebulosity around the HDF–S QSO, detected in narrow-band images taken at the VLT. The nebulosity is reported to extend approximately 9.2″ $`\times `$ 12.1″; the outer parts are patchy but distributed symmetrically around the QSO. They find no evidence of a continuum component associated with the nebulosity. Emission-line nebulae of this size and luminosity are common around radio-loud quasars, but not around radio-quiet quasars such as the HDF-S QSO. However, the nebulae around radio-loud QSOs are generally asymmetric and roughly aligned with the radio jets.
We confirm the existence of extended Ly$`\alpha `$ emission around the HDF-S QSO. Figure 6 shows the narrow-band luminosity profile of the HDF–S QSO compared with the profiles of 8 stars within $``$3′ of the QSO. The stellar profiles are normalized to show the same flux as the QSO within a 4″ radius aperture. The images suffer from low-level ghost reflections off the NB filter; the manufacturer dropped a completed filter shortly before our observing run and the replacement did not have an AR coating. The ghosts are well modeled by a Gaussian which we subtract from the profiles. The peak amplitude of the ghost is $`1.75\times 10^4`$ times the stellar flux within a 4″ radius and its width is $`\sigma =3.6\mathrm{}`$. The flux from the ghost peaks at a radius of $`9.5`$arcsec. The QSO profile exhibits significant excess flux for $`r>4\mathrm{}`$ when compared to all 8 stellar profiles even where the flux from the model ghost is negligible.
The integrated Ly$`\alpha `$ flux in an annulus with 3″$`<r<`$10″ centered on the QSO is $`(2.9\pm 0.5)\times 10^{15}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$. Because the low flux levels and wide area make absolute flux measurement difficult we estimate the error by scaling the 1$`\sigma `$ deviation of the normalized relative flux of the 8 stars integrated over the same area as the QSO. Bergeron et al. (1999) report a flux of 3.2$`\times 10^{15}\mathrm{erg}\mathrm{s}^1\mathrm{cm}^2`$.
## 5. Emission-line objects
To identify point-like emission-line objects in the field we compare the NB magnitudes to an estimate of the continuum around the NB filter. We estimate the continuum as the average $`u`$\- and $`B`$-band flux, $`uB`$. In Fig. 6 we plot the $`uB`$ NB color-magnitude diagram for the point sources. The wide spread in $`uB`$ NB colors is due to the 4000Å break in spectra of the stars.
To set limits for emission-line detection we have estimated the upper 99.5% confidence limit on the $`uB`$ \- NB colors of non-detections using Monte Carlo simulations (Teplitz et al. 1998). A Gaussian distribution of errors is applied to the photometry for simulated line-free objects, and the confidence intervals are measured as a function of broad-band magnitude. We used the global average of the errors over the entire frame for the broad and narrow band data.
We detect 10 spatially unresolved, emission-line sources in our NB filter including the HDF–S QSO. Of these, 7 are QSO candidates by our criteria. Table 2 lists the emission-line objects. Figure 6 shows their position in the field. The observed equivalent widths are lower limits since part of the broad-line emission may fall outside of the narrow-band filter. Figure 6 plots the spectral energy distributions of the objects in AB magnitudes. We include JHK photometry (da Costa et al. 1998) for objects for which it is available.
### 5.1. QSO B
In addition to the emission we detect at 3958Å, Glazebrook et al. (2000) detect a broad emission line at 7174Å for QSO candidate B (see Table 2). The object is unresolved in an HDF–S flanking-field image. The emission lines are identified as $`[`$C IV$`]`$ $`\lambda `$1549 at $`z=1.56\pm 0.02`$ and Mg II $`\lambda `$2800 at $`z=1.56\pm 0.01`$. We adopt a redshift of $`1.56\pm 0.01`$. Relatively small systematic biases in the redshifts determined from these lines (e.g. Tytler & Fan 1992, Laor et al. 1995) cannot be resolved. The absolute $`B`$-band magnitude is $`\mathrm{M}(B)=25`$ using the Veron & Veron (1998) prescription. The rest-frame equivalent width of the C IV emission doublet is at least 60.5 Å. The rest-frame equivalent width of the Mg II lines is 47 Å. The FWHM of the Mg II line is 82 Å in the rest frame, compared to an instrumental resolution of 3 Å establishing it as a broad line typical of QSOs. There is an excess H-band flux which is probably due to H$`\alpha `$ emission. The equivalent width required to give the observed excess above the J and K bands is $``$240 Å in the rest frame. This is near the minimum observed equivalent width in quasars (Espey et al. 1989). The object is detected at 7$`\mu `$m (5$`\sigma `$) and 14$`\mu `$m bands (3$`\sigma `$) with the Infrared Space Observatory (ISO) (Oliver et al. 2000), but flux calibrated magnitudes are not yet available. The QSO is radio quiet; it is detected with a flux of 163 $`\mu `$Jy at 4.9 GHz and 111 $`\mu `$Jy at 8.6 GHz in observations at the Australia Telescope National Facility (ATNF) HDF–S (Norris et al. 2000).
QSO B is 6.7′ (6.8 $`\mathrm{h}_{65}^1`$ Mpc) from the line of sight to the HDF–S QSO and $``$12″ from the western edge of the WFPC2 field. There is no statistically significant excess of absorption lines in the HDF–S QSO associated with QSO B. Since it is so close to the HDF–S, this quasar provides an important new line-of-sight to study the $`z<1.56`$ features revealed in the HDF–S HST images and QSO spectrum.
### 5.2. Other Objects
QSO emission lines which could be detected in our NB images are Ly$`\alpha `$ at $`z2.26`$, C IV $`\lambda `$1549 at $`z1.55`$, C III $`\lambda `$1909 at $`z1.07`$ or Mg II $`\lambda `$2791 at $`z0.42`$. The NB images sample 6.24$`\times `$, 5.10$`\times `$, 3.69$`\times `$ and 1.24$`\times 10^4\mathrm{h}_{65}^3\mathrm{Mpc}^3`$ at each of these redshifts, respectively. The effective volume sampled could be as much as 50% higher depending on how broad the emission lines are. Compact emission-line galaxies could be detected in O II at $`z0.06`$. To estimate the space density of QSOs we integrate the (Hawkins & Veron 1995) luminosity function, with a constant logarithmic slope of 0.63, down to M$`{}_{B}{}^{}=23`$ which is the border between objects classified as QSOs or Seyfert 1 nuclei (Schmidt & Green 1983) and about our detection limit for QSOs at $`z=2.2`$. The expected space density of QSOs in units of $`\mathrm{h}_{65}^3\mathrm{Mpc}^3`$ is 2.4$`\times 10^5`$ at $`2.2<z<3.2`$, 1.5$`\times 10^5`$ at $`1.2<z<1.7`$ and 0.5$`\times 10^5`$ at $`0.7<z<1.2`$ in our cosmology. At $`z<1.5`$, the amplitude of the luminosity function falls as $`(1+\mathrm{z})^6`$. Under this assumption the expected space density at redshift z$`=0.4`$ is roughly 30 times less than that at z$`=1.5`$ or $`5\times 10^7\mathrm{h}_{65}^3\mathrm{Mpc}^3`$. The number of QSOs expected in our detection windows is about 2.5: 1.5 in Ly$`\alpha `$, 0.8 in C IV and 0.2 in C III and 0.005 in Mg II. We detect 7 emission-line objects which we identify as QSO candidates excluding the HDF–S QSO. The total number of detections agrees with the expected number to within the 3$`\sigma `$ Poisson errors (Gehrels 1986). Due to small number statistics there is no evidence of an overdensity in any of the redshift windows.
## 6. Summary
We have presented a sample of QSO candidates in the HDF–S region, including one confirmed QSO, at redshift $`z=1.56`$ just 6.7′ (6.8 $`\mathrm{h}_{65}^1`$ Mpc) from the HDF–S QSO. We confirm the existence of a large, 24″ diameter (490 $`\mathrm{h}_{65}^1`$ kpc) nebular emission surrounding the HDF–S QSO. The existence of such a cloud around a radio-quiet QSO is highly unusual. QSOs in the HDF–S region will provide an important tool to map structure over a much larger area than that covered by the HST data. Absorption line spectra of the brighter QSOs will allow detection of low column density gas clouds along lines of sight other than the HDF–S QSO.
We thank Dave Wittman and Steve Kraemer for useful discussions. We warmly acknowledge the excellent support of the CTIO observing staff. This work was funded by the National Research Council, by NASA grant NRA–98–03–UVG–011, and by NASA funding to the STIS Instrument Definition Team (459–10–60). Finally, we thank the anonymous referee for a rapid response and very helpful comments.
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# Primordial Fluctuations of the metric in the Warm Inflation Scenario
## I Introduction
The standard inflation scenario describes a quasi-de Sitter expansion in a supercooled scenario. This model separates expansion and reheating into two distinguished time periods. This theory assumes a second order phase transition of the inflaton field, followed by a localized mechanism that rapidly distributes the vacuum energy into thermal energy. Reheating after inflation occurs due to particle production by the oscillating inflaton field.
Warm Inflation is a theory of the early universe that explains the abrupt expansion and the exchange of energy beetwen the inflaton field and the thermal bath. This theory was proposed by A. Berera a few years ago, and generalized in another works. Quantum to classical transition and the power spectrum of the primordial fluctuations was studied in the framework of a stochastic approach for the warm inflation scenario. In this thermal scenario the rapid cooling followed by rapid heating is replaced by a smoothened dissipative mechanism. The warm inflation formalism predics energy fluctuations $`\frac{\delta \rho }{\rho }`$ that decreases with time in a power - law expansion, when the universe grows very rapidly. In this formalism a semiclassical expansion for the inflaton field $`\phi `$ was proposed $`\phi (\stackrel{}{x},t)=\varphi _c(t)+\varphi (\stackrel{}{x},t)`$, where $`\stackrel{}{x}`$ are the spatial cordinates, $`t`$ is the time, $`\varphi _c(t)`$ the spatially homogeneous field and $`\varphi (\stackrel{}{x},t)`$ are the quantum fluctuations. For consistency one requires $`<\phi (\stackrel{}{x},t)>=\varphi _c(t)`$, and $`<\varphi (\stackrel{}{x},t)>=<\dot{\varphi }(\stackrel{}{x},t)>=0`$. In the warm inflation scenario the rapid expansion of the universe is produced in presence of a thermal component. The kinetic energy density ($`\rho _{kin}=\rho _r+\frac{1}{2}\dot{\phi }^2`$ ) must be small with respect to the vacuum energy density: $`\rho (\phi )\rho _m(\phi )V(\phi )\rho _{kin}`$, where $`V(\phi )`$ is the potential associated with the field of matter $`\phi `$. Furthermore $`\rho _m`$ and $`\rho _r=\frac{\tau (\phi )}{8H(\phi )}\dot{\phi }^2`$ are the matter and radiation energy densities. This scenario provides thermal fluctuations compatible with de COBE data if the thermal equilibrium becomes near the minimum of the potential $`V(\phi )`$. Furthermore, particles are created during the expansion of the universe, and it is not necessary a further reheating era. The field $`\phi `$ interacts with other particles, which there are in the thermal bath at temperature $`T_r<T_{GUT}10^{15}`$ GeV.
The Lagrangian density that describes the warm inflation scenario is
$$(\phi ,\phi _{,\mu })=\sqrt{g}\left[\frac{R}{16\pi }+\frac{1}{2}g^{\mu \nu }\phi _{,\mu }\phi _{,\nu }+V(\phi )\right]+_{int},$$
(1)
where $`R`$ is the scalar curvature, $`g^{\mu \nu }`$ the metric tensor, $`g`$ is the metric and $`_{int}`$ takes into account the interaction of the field $`\phi `$ with other fields of the thermal bath. In this work I consider a perturbed Friedmann - Robertson - Walker (FRW)
$$ds^2=dt^2+a^2(t)[1+h(\stackrel{}{x},t)]d\stackrel{}{x}^2.$$
(2)
Here, $`a(t)`$ is the scale factor and $`h(\stackrel{}{x},t)`$ denote the fluctuations around the FRW metric with zero global curvature ($`k=0`$). The perturbations $`h(\stackrel{}{x},t)`$ are assumed to be small. Expanding $`H[\phi (\stackrel{}{x},t)]`$ around $`\varphi _c`$, one obtains $`H(\phi )H_c(\varphi _c)+H^{}(\varphi _c)\varphi (\stackrel{}{x},t)`$, at first order on $`\varphi `$. The perturbations of the metric, for small $`\varphi `$, are
$$h(\stackrel{}{x},t)2H_c^{}(\varphi _c)\varphi (\stackrel{}{x},t)𝑑t.$$
(3)
The equation of motion for the scalar field $`\phi `$ is
$$\ddot{\phi }\frac{1}{a^2}^2\phi +\left[3H(\phi )+\tau (\phi )\right]\dot{\phi }+V^{}(\phi )=0,$$
(4)
where $`\tau (\phi )\dot{\phi }`$ describes the dissipation due to the interaction of the field $`\phi `$ with the fields of the thermal bath. We write the semiclassical Friedmann equation for a globally flat FRW metric, which describes a globally isotropic and homogeneous universe: $`H^2(\phi )=\frac{8\pi }{3}G\left(\rho _m+\rho _r\right)`$. Here $`G=M_p^2`$ is the gravitational constant and $`M_p`$ the Planckian mass. The matter and radiation energy densities are $`\rho _m(\phi )=\frac{\dot{\phi }^2}{2}+\frac{1}{a^2}\left(\stackrel{}{}\phi \right)^2+V(\phi )`$ and $`\rho _r(\phi )=\frac{\tau (\phi )}{8H(\phi )}\dot{\phi }^2`$.
## II The classical field $`\varphi _c`$
As in previous works, I consider the following relation between the friction parameter $`\tau _c`$, and the Hubble one
$$\tau _c(\varphi _c)=\gamma H_c(\varphi _c),$$
(5)
where $`\gamma `$ is a dimensionless constant. The classical field is defined as a solution of the equation of motion
$$\ddot{\varphi }_c+[3H_c(\varphi _c)+\tau _c(\varphi _c)]\dot{\varphi }_c+V^{}(\varphi _c)=0,$$
(6)
on the unperturbed FRW metric: $`ds^2=dt^2+a^2(t)[1+h(\stackrel{}{x},t)]d\stackrel{}{x}^2`$ $`=dt^2+a^2(t)d\stackrel{}{x}^2`$.
The classical Hubble parameter is
$$H_c^2(\varphi _c)=\frac{4\pi }{3M_p^2}\left[\left(1+\frac{\tau _c}{4H_c}\right)\dot{\varphi }_c^2+2V(\varphi _c)\right],$$
(7)
and thus, the scalar potential is
$$V(\varphi _c)=\frac{3M_p^2}{8\pi }\left[H_c^2(\varphi _c)\frac{M_p^2}{12\pi }\left(H_c^{}\right)^2\left(1+\frac{\tau _c}{4H_c}\right)\left(1+\frac{\tau _c}{3H_c}\right)^2\right],$$
(8)
where one assumes that $`H(\phi )=H(\varphi _c)H_c`$ and $`\tau (\phi )=\tau (\varphi _c)\tau _c`$. The eq. (8) is obtained using the following equations that describe the dynamics of the classical field and the Hubble parameter: $`\dot{\varphi }_c=\frac{M_p^2}{4\pi }H_c^{}\left(1+\frac{\tau _c}{3H_c}\right)^1`$, and $`\dot{H}_c=\frac{M_p^2}{4\pi }(H_c^{})^2\left(1+\frac{\tau _c}{3H_c}\right)^1`$, where the prime denote the derivative with respect to the field. With the equation for $`\dot{\varphi }_c`$, one obtains the expression for the radiation energy density
$$\rho _r(\varphi _c)=\frac{\tau _c}{8H_c}\left(\frac{M_p^2}{4\pi }\right)^2(H_c^{})^2\left(1+\frac{\tau _c}{3H_c}\right)^2.$$
(9)
When the thermal equilibrium holds, the temperature of the bath is $`T_r\left(\frac{\tau _c(\varphi _c)}{8H_c(\varphi _c)}\dot{\varphi }_c^2\right)_{t1}^{1/4}`$. In this formalism, the expansion of the universe and the interaction of the inflaton with the bath are produced by the classical field $`\varphi _c`$.
## III The quantum fluctuations
For simplicity, I consider the quantum fluctuations on the expectation value of the metric
$$ds^2=dt^2+a(t)d\stackrel{}{x}^2.$$
(10)
However, a consistent treatment must consider the interaction of the quantum fluctuations with the metric. Here, the simplification $`H[\phi (\stackrel{}{x},t)]=H_c(\varphi _c)`$ will be assumed. The equation of motion for the quantum fluctuations \[with the simplification (10)\], is
$$\ddot{\varphi }\frac{1}{a^2}^2\varphi +[3H_c+\tau _c]\dot{\varphi }+V^{\prime \prime }(\varphi _c)\varphi =0.$$
(11)
The structure of this equation can be simplified with the mapp $`\chi =e^{3/2{\scriptscriptstyle (H_c+\tau _c/3)𝑑t}}\varphi `$, and one obtains
$$\ddot{\chi }a^2^2\chi \frac{k_o^2}{a^2}\chi =0,$$
(12)
where $`k_o^2(t)`$ $`=a^2\left[\frac{9}{4}\left(H_c+\frac{\tau _c}{3}\right)^2V^{\prime \prime }(\varphi _c)+\frac{3}{2}\left(\dot{H}_c+\frac{\dot{\tau }_c}{3}\right)\right]`$. The field that describes the quantum fluctuations can be written as a Fourier expansion of the modes $`\xi _k(t)e^{i\stackrel{}{k}.\stackrel{}{x}}`$
$$\chi (\stackrel{}{x},t)=\frac{1}{(2\pi )^{3/2}}d^3k[a_ke^{i\stackrel{}{k}.\stackrel{}{x}}\xi _k(t)+h.c.].$$
(13)
Since in the metric (10) there are not taken into account the quantum fluctuations, the field (13) can be consider as free. However, in a more consistent formalism must be considered the interaction of the quantum fluctuations with the fluctuations of the metric. I denote with $`a_k^{}`$ and $`a_k`$, the creation and annihilation operators. These operators satisfy the commutation relations $`[a_k,a_k^{}^{}]`$ $`=\delta ^{(3)}(\stackrel{}{k}\stackrel{}{k}^{})`$ and $`[a_k,a_k^{}]`$ $`=[a_k^{},a_k^{}^{}]=0`$. The commutation relation between the operators $`\dot{\chi }`$ and $`\chi `$ is
$$[\chi (\stackrel{}{x},t),\dot{\chi }(\stackrel{}{x},t)]=i\delta (\stackrel{}{x}\stackrel{}{x}^{}),$$
(14)
which is satisfied for $`\xi _k\dot{\xi }_k^{}\dot{\xi }_k\xi _k^{}=i`$.
I am interested in the study of the universe on a scale greater than the observable universe. Thus, I consider the field $`\chi `$, but taking into account only the modes with wavelength of size $`l\frac{1}{ϵk_o}`$, where $`ϵ1`$ is a dimensionless constant. This field is given by
$$\chi _{cg}(\stackrel{}{x},t)=\frac{1}{(2\pi )^{3/2}}d^3k\theta (ϵk_ok)[a_ke^{i\stackrel{}{k}.\stackrel{}{x}}\xi _k+h.c.].$$
(15)
The coarse - grained field (15) satisfy the following stochastic equation
$$\ddot{\chi }_{cg}\frac{k_o^2}{a^2}\chi _{cg}=ϵ\left(\frac{d}{dt}\left(\dot{k}_o\eta \right)+2\dot{k}_o\kappa \right),$$
(16)
with the noises
$`\eta (\stackrel{}{x},t)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle }d^3k\delta (ϵk_ok)[a_ke^{i\stackrel{}{k}.\stackrel{}{x}}\xi _k+h.c.],`$ (17)
$`\kappa (\stackrel{}{x},t)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{3/2}}}{\displaystyle }d^3k\delta (ϵk_ok)[a_ke^{i\stackrel{}{k}.\stackrel{}{x}}\dot{\xi }_k+h.c.].`$ (18)
The equation (16) is an second order operatorial stochastic equation, and also can be obtained in the standard inflation formalism (for $`V(\phi )=0`$).
Complex to real transition of the modes $`\xi _k`$ was studied in a previous article. When this transition occurs one obtains $`\xi _k(t)\xi _k^{}(t^{})\xi _k(t)\xi (t^{})`$, and the commutators $`[\chi _{cg},\eta ]`$, $`[\chi _{cg},\kappa ]`$, $`[\eta ,\kappa ]`$, and (14) are null. Hence, the self - correlation of the coarse - grained field $`\chi _{cg}`$ becomes
$$\chi _{cg}(t)\chi _{cg}(t^{})=\frac{1}{(2\pi ^{3/2})}_{ϵk_o(t)}^{ϵk_o(t^{})}d^3k\xi _k(t)\xi _k(t^{}).$$
(19)
The radiation energy fluctuations are $`\delta \rho _r(t)=\left|2\left(\frac{\gamma }{8}\right)\left(1+\frac{\gamma }{3}\right)^2H_c^{}H_c^{\prime \prime }\right|\varphi _{cg}^2^{1/2}`$, and the radiation energy density is $`\rho _r(t)`$ $`=\left(\frac{\gamma }{8}\right)\left(1+\frac{\gamma }{3}\right)^2\left(H_c^{}\right)^2`$. When the radiation energy fluctuations $`\frac{\delta \rho _r}{\rho _r}`$, decreases with time, the thermal equilibrium holds for sufficiently large times. In this case, the Fourier transform of (19) gives the spectral density for the quantum fluctuations in the infrared sector
$$S[\chi _{cg};\omega _k]=\frac{1}{\pi }_0^{\mathrm{}}𝑑t^{\prime \prime }cos[\omega t^{\prime \prime }]_{ϵk_o(t)}^{ϵk_o(t+t^{\prime \prime })}𝑑kk^2\xi _k(t)\xi _k(t+t^{\prime \prime }).$$
(20)
Here, $`\omega _k`$ is the frecuency of oscillation for each mode with wavenumber $`k`$. In a semiclassical representation to the warm inflation scenario, the fluctuations of the metric are dues to the quantum fluctuations $`\varphi _{cg}=e^{3/2{\scriptscriptstyle (H_c+\tau _c/3)𝑑t}}\chi _{cg}`$, on a scale much greater than the scale of the observable universe. The perturbed metric for the FRW universe is given by eq. (2), with $`h(\stackrel{}{x},t)`$ given by eq. (3). The temporal evolution of $`h(\stackrel{}{x},t)`$ is given by
$$h^2(t)^{1/2}a^2\left[2H^{}(\varphi _c)<\varphi _{cg}^2(\stackrel{}{x},t)>^{1/2}𝑑t\right].$$
(21)
Quantum to classical transition of the perturbations of the metric holds due to the quantum to classical transition of the coarse - grained quantum field $`\varphi _{cg}`$. It is produced when all of the modes of $`\chi _{cg}`$ become real, i.e.: $`\frac{1}{N(t)}_{k=0}^{k=ϵk_o}\alpha _k(t)1`$. Here $`N(t)`$ is the time dependent number of degrees of freedom of the infrared sector, which increases with time. The function $`\alpha _k(t)`$ was defined in a previous work, and is: $`\alpha _k=\left|\frac{v_k(t)}{u_k(t)}\right|`$, for $`\xi _k(t)=u_k(t)+iv_k(t)`$. This effect is due to the time evolution of the superhorizon with size $`l(t)\frac{1}{ϵk_o}`$. Thus, quantum to classical transition of the coarse - grained field occurs when $`\alpha _{k=ϵk_o}0`$.
## IV An example: fluctuations of the metric in power - law inflation
.
In this section I estime the fluctuations of the metric in the the example for power - law inflation. Here, the scale factor and the Hubble parameter are $`a(t)=H_o^1(t/t_o)^p`$, and $`H_c(t)=\frac{p}{t}`$, respectively. The temporal evolution of the scalar field is $`\varphi _c(t)=\varphi _om\mathrm{ln}\left[\frac{H_o}{p}t\right]`$. Hence, the scalar potential $`V(\varphi _c)`$ is given by
$$V(\varphi _c)=\frac{3M_p^2H_o^2}{8\pi }e^{2\varphi _c/m}\left[1\frac{M_p^2}{12\pi m^2}\left(1+\frac{\gamma }{4}\right)\left(1+\frac{\gamma }{3}\right)^2\right].$$
(22)
The temporal evolution of the quantum fluctuations is
$$<\varphi _{cg}^2(t)>^{1/2}|_{t1}t^{2\nu p(1\nu /pp/\nu )+3/2\gamma /2},$$
(23)
with $`\nu =\frac{\sqrt{1+9p^2(1+\gamma )^26p[(1+\gamma )+1]+9(1+\frac{\gamma }{4})(1+\frac{\gamma }{3})^2}}{2(p1)}`$. These fluctuations decrease with time when the universe grows very rapidly (for $`p1`$). So, the fluctuations of the metric are given (as $`H_c^{}(t)t^1`$) by
$$h^2(t)^{1/2}t^{2\nu p(1\nu /pp/\nu )+3/2\gamma /2},$$
(24)
which, since $`<\varphi _{cg}^2>^{1/2}`$, decreases with time for $`p1`$. The power spectral density for the matter field fluctuations is $`S\left[\chi _{cg},\omega _k=\frac{2k}{t^p}\right]\left|\omega _k\right|^n`$, with $`n=[4\nu (p1)2p+4]`$. The fluctuations of radiation energy density are
$$\frac{\delta \rho _r}{\rho _r}|_{t1}t^{2\nu p(1\nu /p(5/4)p/\nu )+1/2\gamma /2},$$
(25)
which decreases for $`p`$ sufficiently large. Thus, for $`p1`$ the thermal equilibrium holds due to $`\frac{d[\delta \rho _r/\rho _r]}{dt}<0`$.
## V Final Comments
In this work was studied the evolution for the fluctuations of the metric with a semiclassical representation of the inflaton field $`\phi `$, in the warm inflation scenario. The mean temperature is smaller than the GUT temperature (i.e., $`T_r<T_{GUT}10^{15}`$ GeV). In this theory the classical matter field lead to the expansion of the universe, while the fluctuations of the matter field generate the inhomogeneities of the metric, and thus local curvature of the spacetime. However the expectation value of the curvature is zero, in consistency with a globally flat perturbed Friedmann - Robertson - Walker metric here considered. In this framework the quantum fluctuations averaged over a scale much bigger than the observable universe are responsible (not only of the radiation energy density fluctuations, thermal fluctuations and matter energy density fluctuations - these topics were studied in another previous works) for the fluctuations of the metric. This calculation was done with the assumption that the quantum fluctuations are very small. Thus, was possible to expand the coarse - grained field as free. In a more appropiate treatment for the coarse - grained field, it must be consider as interactuant with the metric.
In the example here considered I observe that the fluctuations of the metric decreases with time, in a power - law expansion when the expansion is very rapid. Furthermore, the radiation energy fluctuations decreases with time for a sufficiently large rate of expansion of the universe. Hence, when the universe expands very rapidly, the thermal equilibrium holds.
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# Spin transitions in a small Si quantum dot
\[
## Abstract
We have studied the magnetic field dependence of the ground state energies in a small Si quantum dot. At low fields the first five electrons are added in a spin-up – spin-down sequence minimizing the total spin. This sequence does not hold for larger number of electrons in the dot. At high fields the dot undergoes transitions between states with different spins driven entirely by Zeeman energy. We identify some features that can be attributed to transitions between different spin configurations preserving the total spin of the dot. For a few peaks we observed large linear shifts that correspond to the change of the spin of the dot by 3/2. Such a change requires that an electron in the dot flips its spin during every tunneling event.
\]
The spin degree of freedom is an essential part of mesoscopic physics. Quite often the knowledge of the spin advances our understanding of electron-electron interactions, which can reveal themselves in a non-trivial spin configuration of a mesoscopic system. From this perspective, quantum dots can be regarded as model systems for the study of spin–related phenomena because they contain just a few electrons and different coupling parameters can be tuned almost independently. Within the simplest model of non-interacting electrons each additional electron is added into the dot to the next single-particle energy level (there is also a constant energy associated with the charging of the environment). The spin is accounted for by allowing two electrons to fill the same single-particle energy level, thus the total spin of the system should alternate between $`s=0`$ and $`s=1/2`$. There are several ways to determine the spin of the dot. The most direct way is by studying the Kondo effect. If the spin of the ground state $`s>0`$, the Coulomb blockade is lifted in the corresponding conduction valley at low temperatures. Indeed, valleys with Kondo-enhanced conductivity were found to alternate with the regular Coulomb blockade valleys of vanishing conductivity. However, if the dot is weakly coupled to the leads the Kondo temperature can be too small to be achieved experimentally. In such dots individual energy levels are sharp and spin of the tunneling electron can be determined from the Zeeman shift of the energy level. The shift of consecutive peaks has been shown to alternate with $`\pm g^{}\mu _BB`$ in small Al clusters and carbon nanotubes, supporting the alternating spin filling of the dot (here $`g^{}`$ is the effective $`g`$-factor and $`\mu _B`$ is the Bohr magneton). In lateral semiconductor quantum dots with a large number of electrons the Zeeman shift is masked by much larger orbital effects, and direct determination of the spin is a formidable task. Peaks fluctuate as a function of $`B`$, reflecting the orbital shift of the levels. An indirect information about the spin can be obtained from the comparison of such ”magnetic fingerprints”, in order to find whether the two consecutive electrons fill the same energy level. The underlying assumption here is that the addition of an electron does not change the spectrum of the dot significantly. Using this method, strong deviations from the alternating spin filling has been reported. The most dramatic example of a non-alternating spin filling is the polarization of small vertical dots due to exchange interactions, similar to the Hund’s rule in atomic physics.
The application of a magnetic field alters the spin configuration. Energy levels shift differently with magnetic field and cross each other. In the vicinity of such a magnetically-induced level crossing exchange interactions may lift the degeneracy by favoring the formation of a triplet state. A sufficiently large field gradually polarizes the dot by collapsing all electrons into the lowest Landau level.
In this Letter we examine the electron transport through a small Si quantum dot. Unlike the previously studied semiconductor dots, in our dot the field-dependent shift of energy levels is dominated by the Zeeman energy rather than by orbital effects. Thus, we can measure the spin of the dot directly for ground states with different number of electrons, starting from one. We also study the evolution of the total spin as a function of the magnetic field. We find that the field dependence of energy levels consists of several linear segments with different slopes. Comparison with a simple model for non-interacting electrons allows us to identify a set of spin transitions for the first few ground states. In addition to the singlet-triplet and triplet-polarized transitions there are some features in the spectra, which we attribute to transitions between different realizations of the triplet state. A detailed analysis reveals some deviation from the model that hints for the importance of the underlying interactions. As the number of electrons increases, the spectra become more complicated. For example, tunneling of an electron into the dot can change the total spin of the other electrons in the dot. Such a tunneling process is beyond the scope of a model of non-interacting electrons.
The measurements were performed on a small Si quantum dot fabricated from a silicon-on-insulator wafer. The dot resides inside a narrow bridge patterned from the top Si layer (see inset in Fig. 1). A 50 nm thick layer of thermal oxide is grown around the bridge followed by a poly-Si gate. The fabrication steps have been described previously. Gate capacitance is estimated to be 0.8-1.0 aF, the total capacitance $`C15`$ aF and the charging energy $`U_c=e/C10`$ meV. Spacing between excited levels $`\delta 14`$ meV, measured using non-zero bias spectroscopy, is comparable to the charging energy and is consistent with the lithographical size of the dot $`l\sqrt{\mathrm{}/m^{}\delta }100190`$ Å . The gate voltage – to – energy conversion coefficient, measured from both non-zero bias spectroscopy and $`T`$-dependent scaling of the peak width, is $`\alpha 14`$ mV/meV. The sample was studied in three separate cooldowns and the reported phenomena were found to be insensitive to redistribution of background charges, thus reflecting intrinsic properties of the dot.
At high temperatures, $`T>120`$ K, the device exhibits regular metal-oxide-semiconductor field effect transistor (MOSFET) characteristics with a threshold gate voltage of $`V_{th}0.2`$ V. At $`T<100`$ K the Coulomb blockade emerges and the conductance oscillates as a function of the gate voltage $`V_g`$ for $`V_g>V_{th}`$. A representative trace of the conductance $`G`$ as a function of $`V_g`$ at $`T=1.5`$ K is shown in Fig. 1. There is a series of sharp peaks, spaced by 150-200 mV. The peaks, corresponding to the entrance of the first two electrons, cannot be reliably measured, but their positions can be determined from high bias spectroscopy. Commonly for this type of devices, the sample has a parallel conducting channel, which exhibits Coulomb blockade at $`V_g<1.2`$ V with peaks separated by 60 mV. At $`V_g>1.2`$ V the extra channel has finite conductance with some broad features as a function of $`V_g`$. Coulomb blockade peaks, originated from the lithographical dot, are not broadened at high $`V_g`$, thus electrical transport through the dot and the parallel channel are decoupled. Charging of remote impurities (which can be accomplished by wide gate voltage scans $`\mathrm{\Delta }V_g>3`$ V) changes positions of the extra peaks at $`V_g<1.2`$ V and the value of the background conductance at $`V_g>1.2`$ V without altering the position and amplitude of the main peaks.
We studied the peak positions $`V_g^p`$ as a function of magnetic field $`B`$ for the first 30 peaks. For $`V_g<0.4`$ V (the first three peaks), electron density in the contacts is low and the contacts are spin polarized by a moderate magnetic field. For $`V_g>0.4`$ V ($`N4`$) both spin subbands in the contacts are occupied within the experimental range of $`0<B<13`$ T. Thus, the Fermi energy $`E_F`$ is field-independent and the peak shift reflects only the field dependence of the energy levels in the dot (mobility of the two-dimensional gas is low, $`300`$ cm<sup>2</sup>/V$``$s at 4.2 K, and there is no measurable modulation of $`E_F`$ due to Subnikov-de-Haas oscillations for $`B`$ up to 13 T).
The evolution of several peaks with $`B`$ is shown in Fig. 2a (peaks 4-7) and in Fig. 3a (peaks 21-23). Clearly, $`V_g^p`$ and the peak amplitudes $`G^p`$ change non-monotonically with $`B`$. Analysis of ”magnetic fingerprints” reveals that there is no apparent pairing of the neighboring peaks within the first 30 peaks. In fact, we observed an unexpected tripling of the peaks: two bunches of peaks have similar ”magnetic fingerprints” for three consecutive peaks (13,14,15 and 16,17,18, not shown in the figures). The measurements were repeated for two different orientations of $`B`$, defined in the inset in Fig. 1. We found that $`V_g^p`$ is insensitive to the direction of the magnetic field: aligning $`B`$ with the current direction ($`B_{||}`$, in-plane) or perpendicular to the plane of the sample ($`B_{}`$) does not change $`V_g^p`$ significantly. The dot is lithographically asymmetric and the orbital effects are expected to depend on the field direction. Thus, we conclude that in our small dot the $`B`$-dependence of $`V_g^p`$ is dominated by spin effects. This conclusion is also supported by the observation that, in the range of $`B`$ when the contacts are fully spin-polarized, the $`V_g^p`$ for peaks 1-3 does not depend on $`B`$ at all.
Unlike $`V_g^p`$, the peak amplitude $`G^p`$ depends on the direction of the magnetic field. The $`G^p`$ reflects the tunneling probability and depends exponentially on the overlap of wavefunctions in the dot and in the contacts. As such, $`G^p`$ is sensitive to a particular configuration of the wavefunction within the dot, and redistribution of the wavefunction due to small orbital effects can result in a significant change of $`G^p`$.
What physics is behind the field-dependence of the peak position? At zero bias $`V_g^p`$ is determined by the degeneracy condition that the electrochemical potentials for the ground states with $`N1`$ and $`N`$ electrons in the dot are equal. Provided that the Fermi energy in the contacts is independent of the magnetic field, a shift of the $`N`$-th peak with $`B`$ reflects the relative change of the ground state energies $`\mathrm{\Delta }U_N^p(B)=\mathrm{\Delta }U(N,B)\mathrm{\Delta }U(N1,B)`$, where $`U(N,B)`$ is the energy of the ground state of $`N`$ electrons in magnetic field $`B`$ and $`\mathrm{\Delta }U(N,B)=U(N,B)U(N,0)`$. In the absence of spin-orbit interactions (which is the case for the bulk Si) the total energy can be separated into spin and orbital terms. The spin term includes Zeeman energy $`s(N)g^{}\mu _BB`$ where $`s(N)`$ is the total spin of the ground state with $`N`$ electrons and $`g^{}`$ is the effective $`g`$-factor. Thus, the Zeeman-related peak shift is $`[s(N)s(N1)]g^{}\mu _BB=(\pm n\pm 1/2)g^{}\mu _BB`$, where the spin $`\pm 1/2`$ is carried by the tunneling electron and $`n=0,1,2,\mathrm{}`$ is the number of electrons in the dot that flip their spins upon the tunneling event. In the simplest case of no interactions peaks should shift linearly with $`B`$ by $`\pm 1/2g^{}\mu _BB`$.
Experimentally, peaks do not shift linearly with $`B`$. Instead, $`dV_g^p/dB`$ changes both its value and sign as $`B`$ is varied from 0 to 13 T. For a quantitative analysis, peak positions are extracted from $`G`$ vs $`V_g`$ scans, and the peak shifts $`\mathrm{\Delta }U^p(B)=[V_g^p(B)V_g^p(0)]/\alpha `$ are plotted as a function of $`B`$ in Figs. 2b and 3b. The curves are offset for clarity. For a comparison, lines with slopes $`\pm 1/2g^{}\mu _B`$ for $`g^{}=2`$ are also shown (solid lines). First, let us focus on the low-field ($`B<2`$ T) region. Peaks 4 and 5 shift linearly with $`B`$ and the corresponding slopes are + and $`1/2g^{}\mu _B`$. In the same low-field region the preceding peaks 2 and 3 also shift with + and $`1/2g^{}\mu _B`$ slopes correspondingly. Thus, at low fields the ground states with up to 5 electrons in the dot have the lowest spin configuration and the dot is filled in a spin-down – spin-up sequence $``$ (in the order the levels are filled). Such a filling sequence requires that the valley degeneracy is lifted and two electrons with different spins can occupy the same energy level.
This simple picture of alternating filling does not hold for $`N>5`$ even at low fields. At $`B<2`$ T peak 6 consists of three peaks separated by $`0.5`$ meV at zero field, none of which shifts with $`1/2g^{}\mu _BB`$ (the zero-field positions of the three peaks are marked by triangles in Fig. 2a). The slope of the lowest-energy branch is close to $`3/2g^{}\mu _B`$, the other two branches have small negative slopes. The shift of the next, the 7-th, peak has a positive slope, while the lowest-spin arrangement for a dot with 7 electrons should have negative Zeeman energy. We conclude that the ground state with 6 electrons is spontaneously polarized and the total spin $`s(6)>1/2`$. Transitions between ground states that involve a change in spin by $`\mathrm{\Delta }s>1/2`$ have low probability and the corresponding peaks are expected to be suppressed (so-called spin blockade). Indeed, the overall conductance of peak 6 is strongly suppressed and, presumably, the appearance of several branches can be explained by the instability of the polarized state.
The low-field spin configuration is not preserved at high magnetic fields. For peak 4, $`dV_g^p/dB`$ changes sign from positive to negative at $`B=2.5`$ T, back to positive at $`B=9`$ T, and, again, to negative at $`B12`$ T. The spin of the tunneling electron changes from being $`+1/21/2+1/21/2`$. The corresponding spin transitions of the ground state can be understood from a simple model for non-interacting electrons. Let us consider four single-particle levels $`E_i`$, as shown in Fig. 2c. Each level is spin-degenerate at zero field and splits into two levels $`E_i\pm 1/2g^{}\mu _BB`$ for $`B>0`$. In the absence of interactions position of the $`N`$-th peak is determined by $`U(N,B)U(N1,B)=_k^NE(k,B)_k^{N1}E(k,B)=E(N,B)`$, where $`E(k,B)`$ is the energy of the $`k`$-th electron, including the Zeeman contribution. $`E(4,B)`$ is outlined by the thick solid line in Fig. 2c. Qualitatively, $`E(4,B)`$ captures the main features of $`V_g^p`$ vs. $`B`$ for the 4-th peak and the kinks can be attributed to the corresponding level crossings. Each level crossing results in a change of the spin configuration and the ground state of 4 electrons undergoes spin transitions as a function of $`B`$: $``$ (the regions with different spin configurations are separated by dashed vertical lines in Fig. 2c). The first transition is singlet-triplet and the last transition is triplet-spin polarized. There are two intermediate transitions within the triplet state which change the spin configuration within the dot without changing the total spin. At $`B7`$ T the spin configuration of the ground state with 4 electrons changes without reversing the spin of the tunneling electron; such a transition does not change the sign of $`dV_g^p/dB`$. In the absence of interactions there should be no corresponding kink. The second transition flips the spin of the tunneling electron and of an electron in the dot simultaneously, preserving $`s(4)=1`$ but changing the sign of $`dV_g^p/dB`$.
The model, described above, also reproduces the features of peak 5 for $`B<7`$ T (dashed line in Fig. 2c). However, there are some important discrepancies, which cannot be understood within this model of non-interacting electrons. First of all, we cannot describe the evolution of $`N>5`$ peaks within this model. Second, each level crossing should result in a pair of upward–downward kinks in two neighboring peaks at the same value of $`B`$. Clearly, kinks in $`V_g^p(B)`$ for peaks 4 and 5 near 2 T are shifted by $`0.5`$ T. The most notable deviation from this simple model of level crossing is shown in Fig. 3b, where upward kinks at 2.3 T and 5.3 T in $`V_g^p(B)`$ for peak 21 have no downward counterparts in $`V_g^p(B)`$ for peak 22. Third, we have to assume a small single-particle level spacing of $`0.3`$ meV to fit the positions of the observed spin transitions. From non-zero bias spectroscopy, as well as from the statistics of the zero-bias peak spacing, we estimate that excited levels are separated by 1-4 meV.
For most of the peaks $`|dV_g^p/dB|1/2g^{}\mu _B`$. However, there are a few peaks that shift much faster with magnetic field. In Fig. 3b $`\mathrm{\Delta }U^p(B)`$ for peaks 21 and 22 have linear segments with a slope $`3/2g^{}\mu _BB`$. Remarkably, the shift of peak 21 has such a large slope in the whole range $`0<B<13`$ T, although its sign changes four times. We can rule out enhancement of the $`g`$-factor because i) there are segments in the neighboring peak 23 with the slope $`1/2g^{}\mu _B`$ (assuming $`g^{}=2`$), and ii) it is known that interactions renormalize $`g^{}`$ at low electron densities in Si-MOSFETs but $`g^{}`$ approaches the bulk value of 2 as the density increases. Thus, the total spin of the dot changes by $`s(N)s(N1)=3/2`$. A change of the spin by more than 1/2 means that at least one electron in the dot should flip its spin ($`3/2=1+1/2=21/2`$) upon the tunneling of an electron. We want to stress the difference with the spin transitions discussed above: there, the total spin of the dot changes as a function of $`B`$, but it is fixed for any particular $`B`$. In order to change the total spin by 3/2 an electron in the dot has to flip its spin during the tunneling event. In the absence of spin-orbit interactions such a flip is forbidden unless some other spin scattering mechanism is considered. As we mentioned earlier, the absence of an efficient spin scattering should result in a spin blockade with the corresponding suppression of the peak amplitude. Experimentally, there is no apparent suppression of peaks 21 and 22, which have the $`3/2g^{}\mu _B`$ slopes, compared to the amplitude of peak 23, which has the regular slope of $`1/2g^{}\mu _B`$.
To summarize, we have analyzed the field dependence of ground state energies in a small Si quantum dot. The dot is in a new regime where the $`B`$-dependence of the energy levels is dominated by the Zeeman energy. There are distinctive features in the data which we attribute to the transitions between different spin configurations of the dot. For the state with 4 electrons in the dot we identified five different spin configurations, including three with the same total spin $`s=1`$. Some peaks have large shift as a function of magnetic field which requires the total spin of the dot to be changed by $`\mathrm{\Delta }s>1/2`$ upon the tunneling of an electron. Surprisingly, we found that such peaks are not necessarily suppressed.
We gratefully acknowledge discussions with Boris Altshuler and Richard Berkovits. The work was supported by ARO, ONR and DARPA.
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# Detection of Lead in the Carbon-Rich, Very Metal-Poor Star LP 625-44: A Strong Constraint on s-Process Nucleosynthesis at Low Metallicity
## 1 Introduction
The slow neutron-capture process (s-process) is considered one of the major pathways for the creation of nuclei heavier than iron, and the asymptotic giant-branch (AGB) phase of low- and intermediate-mass stars has been studied as its most likely astrophysical site. One important component in understanding s-process nucleosynthesis is the correct identification of the neutron sources involved. Two reactions – <sup>22</sup>Ne($`\alpha ,n`$)<sup>25</sup>Mg and <sup>13</sup>C($`\alpha ,n`$)<sup>16</sup>O – have received most attention. Recent models of AGB stars prefer <sup>13</sup>C as the main source, because the temperature of the He burning shell hardly reaches $`3\times 10^8`$ K required for the <sup>22</sup>Ne($`\alpha ,n`$)<sup>25</sup>Mg reaction (e.g., Gallino et al., 1998). This is supported by the observed metallicity dependence of the abundance ratios of heavier s-process elements (e.g., Ba, Nd) to lighter ones (e.g., Sr, Zr) found in s-process-enhanced objects such as MS- and S-type stars (Smith & Lambert, 1990), barium stars (Luck & Bond, 1991) and CH stars (Vanture, 1992; Norris et al., 1997a). While the seed nuclei for the s-process, such as iron, are secondary (i.e., their abundances are proportional to metallicity), the production of <sup>13</sup>C in AGB stars is primary, contrary to that of <sup>22</sup>Ne. Therefore, higher neutron exposure, and thus larger enhancement of the heavier elements, is expected from <sup>13</sup>C in stars of lower metallicity. Models of nucleosynthesis in AGB stars by Gallino et al. (1998), followed by Busso et al. (1999), successfully reproduced the observed trend for lighter elements (Sr-Zr), as well as for heavier ones (Ba-Gd).
For stars of very low metallicity, according to these models, a large excess of lead (Pb) is expected. For instance, the enhancement of Pb by two or three orders of magnitude relative to that expected for solar-abundance stars is predicted for AGB stars with \[Fe/H\] $`=2.0`$ <sup>1</sup><sup>1</sup>1\[A/B\]$`\mathrm{log}(N_\mathrm{A}/N_\mathrm{B})`$ $`\mathrm{log}(N_\mathrm{A}/N_\mathrm{B})_{}`$, and $`\mathrm{log}ϵ_\mathrm{A}`$ $`\mathrm{log}(N_\mathrm{A}/N_\mathrm{H})+12`$ for elements A and B, while that of Ba is at most one order of magnitude (Busso et al., 1999). Thus Pb in metal-poor, s-process-enhanced, stars should provide an excellent diagnostic for models of s-process nucleosynthesis in AGB stars.
However, the abundance of Pb is difficult to measure in most stars.
Lead abundances for the metal-poor stars HD 115444 and HD 126238 were derived from Hubble Space Telescope ultraviolet spectra (Sneden et al., 1998), but Pb has not yet been detected in the optical spectra of these objects. The Sneden et al. (2000) analysis of a high-S/N Keck HIRES spectrum of the r-process-rich star CS 22892-052 detected Pb I lines in the visual spectrum of this star, and derived its abundance. We note that the Pb observed in all three of these stars is attributed primarily to the r-, rather than the s-process, due to the strong enhancements of other r-process-dominated nuclei, such as Eu. One study of the s-process for a solar metallicity star by Gonzalez et al. (1998) reports the Pb abundance of the post-AGB star FG-Sge, based on an analysis of the Pb I 7229Å line.
In this Letter we report the detection of Pb I 4057.8Å and derive a Pb abundance in the carbon-rich, very metal-poor star LP 625-44. This object was shown by Norris et al. (1997a) to exhibit very large excesses of carbon, nitrogen, and neutron-capture elements. Their interpretation was that the large excesses of heavy elements were likely to have originated from s-process nucleosynthesis in an AGB binary companion which provided LP 625-44 with carbon-rich material by mass transfer. The updated abundance pattern (see §3), and variation of radial velocity (see §2), reported in the present work make this interpretation quite convincing. The determination of a Pb abundance for this star (§3) provides the opportunity, for the first time, to test models of nucleosynthesis in AGB stars for s-process elements between Sr and Pb (§4).
## 2 Observations and Measurements
A high-resolution spectrum of LP 625-44 was obtained with the University College London coudé échelle spectrograph (UCLES) and Tektronix 1024$`\times `$1024 CCD at the Anglo-Australian Telescope on August 5, 1996. The wavelength range 3700–4720 Å was covered with resolving power $`R40,000`$. We also obtained a red spectrum (5015–8500 Å) with the same instrument on June 16, 1994. The numbers of detected photons are 12000 per 0.04Å pixel at 4300Å ($`S/N150`$ per resolution element) and 2000 per 0.06Å pixel at 6000Å ($`S/N60`$ per resolution element).
Data reduction was performed in the standard way within the IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation environment. Equivalent widths were measured by fitting Gaussian profiles to the absorption lines, and will be reported in Aoki et al. (2000). The error for lines weaker than 60 mÅ, determined from the comparison of two measurements of lines which appear on adjacent échelle orders, was about 4 mÅ and 6 mÅ in the blue and red spectra, respectively. There is no systematic difference between the equivalent widths of Fe I lines measured in this work and those by Norris et al. (1996), even though our $`S/N`$ is substantially higher.
Additional spectra were obtained on 1998 August 11 and 2000 January 26, the former using UCLES, and the latter with the Utrecht échelle spectrograph (UES) on the William Herschel Telescope (WHT). Both have lower S/N than is necessary for an abundance analysis, the sole aim being to measure radial velocities. In each case, HD 140283 was also observed to provide a template for cross-correlation. That star has a similar metallicity but is free of the CH blends that affect many lines in LP 625-44. Radial velocities for LP 625-44 were obtained relative to HD 140283 by cross-correlation, and by measuring the radial velocity of HD 140283 from the central wavelengths of 175 (1998) and 122 (2000) unblended lines. Error estimates were based on the variation in velocity from different échelle orders and from the standard error in the measurement of HD 140283. The heliocentric values, which extend those presented by Norris et al. (1997a), are: HJD 2451037.00: $`v_{\mathrm{rad}}`$ = 33.5$`\pm `$0.2 (1$`\sigma `$) km s<sup>-1</sup>; and HJD 2451569.80: $`v_{\mathrm{rad}}`$ = 30.0$`\pm `$0.3 (1$`\sigma `$) km s<sup>-1</sup>. Ryan et al. (1999) estimated external errors of 0.3 km s<sup>-1</sup> for a similar procedure; this has been added to the internal errors for Fig. 1. The data confirm that LP 625-44 is a binary with a period of at least 12 years.
## 3 Abundance Analysis and Results
In the region near the Pb I line at 4058Å , line-blending is so severe that the Pb abundance was derived by spectrum synthesis. This method was also applied to lines which are affected by blending and/or hyperfine splitting. The standard analysis, based on the equivalent widths, was applied to single (unblended) lines.
The abundance analysis used model atmospheres in the ATLAS9 grid of Kurucz (1993a). We adopted an effective temperature $`T_{\mathrm{eff}}=5500`$K, determined by Norris et al. (1997a) from the $`RI`$ color. This color is not severely affected by strong carbon and nitrogen features in stars of this temperature and abundance (Aoki et al., 2000). Surface gravity ($`\mathrm{log}g`$), metallicity, and microturbulent velocity ($`\xi `$), were re-determined in the present work. The surface gravity was obtained from the ionization balance between Fe I and Fe II, the metallicity was estimated from the abundance analysis of those lines, and $`\xi `$ was determined from the Fe I lines by demanding no dependence of the derived abundance on equivalent widths. The results are: $`\mathrm{log}g=2.8`$, $`\xi =1.6`$ km s<sup>-1</sup>, and \[Fe/H\] $`=2.72`$. The agreement with the results of Norris et al. (1997a) is good, with the exception of the microturbulent velocity, for which they derived 1.0 km s<sup>-1</sup>.
Pb lines are difficult to measure in optical stellar spectra. Even for the sun, only four Pb I lines (3639 Å , 3683 Å , 3739 Å , and 4057 Å) have been studied in the visual region. From these, Youssef & Khalil (1989) determined the abundance $`\mathrm{log}ϵ_{\mathrm{Pb}}=2.0`$, which agrees fairly well with meteoritic measurements ($`\mathrm{log}ϵ_{\mathrm{Pb}}=2.06;`$ Grevesse et al., 1996).
Our spectra covered the lines at 3739.9Å and 4057.8Å , listed in Table 1. They are expected to be weak, and no clear absorption appears in the solar spectrum. However, Pb I 4057.8Å was clearly detected in LP 625-44, as shown in the upper panel of Fig. 2, where the synthetic spectra fitted to the observed data are also shown. In spite of the presence of other lines, the contribution of Pb I is remarkable. To check possible contamination of the Pb region by other elements, we examined the spectrum of HD 140283, a very metal-poor subgiant ($`T_{\mathrm{eff}}=5750`$ K, $`\mathrm{log}g=3.4`$ and \[Fe/H\] $`=`$2.54, Ryan et al. 1996). We found no distinct absorption feature at 4057.8Å (see Fig. 1b in Norris et al., 1996). As a further check for contamination due to CH and CN, the observed and synthetic spectra of CS 22957-027 are shown in the lower panel of Fig. 2. Norris et al. (1997b) showed that this very metal-poor giant ($`T_{\mathrm{eff}}=4850`$K, $`\mathrm{log}g=1.9`$ and \[Fe/H\] $`=3.38`$) has very large excesses of <sup>12</sup>C, <sup>13</sup>C and N but no excess of heavy elements. This spectrum indicates that the absorption feature at 4057.8Å in LP 625-44 is not due to CH and CN lines.
To check our procedure for the determination of the Pb abundance, we also analyzed the solar spectrum (Kurucz, 1993b) using a solar photospheric model (Holweger & Müller, 1974). Our result agrees very well with that of Youssef & Khalil (1989) for Pb I 3683Å, which is the clearest Pb I line in the optical range. This demonstrates the reliability of the basic data (e.g., partition functions) and the software used in our analysis. (Line contamination is so severe at 4057.8Å in the solar spectrum that the exact abundance cannot be derived from this line.)
An abundance ratio \[Pb/Fe\] $`=2.65`$ was derived for LP 625-44 from a comparison between the synthetic spectra and the observed one. The other Pb I line covered by our spectrum is at 3739.9Å, but there is no distinct feature at this wavelength. Hence, we derive an upper limit on the abundance ratio \[Pb/Fe\] $`<+3.2`$ from this non-detection, which supports the Pb I 4057.8Å result (\[Pb/Fe\]=2.65). This upper limit is important, as in the next section we show that this Pb abundance is lower than predicted by some models of $`s`$-process nucleosynthesis in very metal-poor AGB stars.
Our derived abundances for the heavy elements are similar to the results presented by Norris et al. (1997a). Abundances of Er, Tm and Hf, not previously known in this star, could also be determined due to the better quality of the new spectra. All new results are given in Table 2. The line data used in the analysis will be compiled in Aoki et al. (2000).
Errors in our estimated abundances were obtained as follows. Errors arising from uncertainties in the atmospheric parameters were evaluated by adding in quadrature the individual errors on the parameters – $`\mathrm{\Delta }T_{\mathrm{eff}}=100`$K, $`\mathrm{\Delta }\mathrm{log}g=0.3`$, and $`\mathrm{\Delta }\xi =0.5`$km s<sup>-1</sup>. The internal errors were estimated by assuming the random error in the equivalent width measurements to be 4 mÅ (and 6 mÅ for Ba II in the red region; see §2), and taking the random error in less-certain $`gf`$ values to be 0.1 dex.
## 4 Discussion and Concluding Remarks
Fig. 3 presents derived abundances as a function of atomic species for LP 625-44. The thick solid line indicates the abundance pattern of the main s-process component determined by Arlandini et al. (1999), while the thin line indicates the r-process component. The dotted line is the total solar abundance adopted by Arlandini et al. (1999). We see good agreement between the observed abundances of elements heavier than Ba with the scaled s-process component. This fact, found by Norris et al. (1997a) for Ba to Dy, is now extended to heavier elements and made even more compelling. The excesses of these elements, and their s-process nature, are interpreted as a result of the transfer of material rich in s-process elements across a binary system including an AGB star. Our new radial velocity measurements confirm the binarity and strengthen this interpretation. Since the excess of heavy elements is very large (e.g., \[Ba/Fe\]=2.7), the material from the AGB star should dominate the surface abundances of LP 625-44. Thus, the relative abundances of the heavy elements in this star should provide an almost pure representation of the nucleosynthesis products of the previously existing very metal-poor (\[Fe/H\] $`=2.7`$) AGB companion.
With the adoption of this interpretation the abundances in LP 625-44 can be compared with theoretical predictions of nucleosynthesis in AGB stars. Gallino et al. (1998) and Busso et al. (1999) showed that, at low abundance, the metallicity effect on s-process yields favors the production of heavier elements. As found in Fig. 3, the Sr-Zr enhancement relative to the solar s-component is much smaller than that of heavy elements (Ba-Hf) in LP 625-44, a result in qualitative agreement with the expected metallicity dependence.
The metallicity effect is essentially due to the level of neutron exposure, which is expected to be higher at lower metallicity (see §1). Higher neutron exposure necessarily requires larger production of the heaviest s-process element, Pb, in very metal-poor stars. Busso et al. (1999) explicitly showed the metallicity effect on the enhancement factors of s-process elements relative to solar abundances for $`3.2<`$ \[Fe/H\] $`<0.4`$ in their Figure 12, where the enhancement factor of Pb is larger by about two orders of magnitude than that of Ba at \[Fe/H\] $`=2.7`$. However, the enhancement of Pb in LP 625-44 (\[Pb/Fe\] = 2.65) is nearly the same as that of Ba (\[Ba/Fe\] = 2.74). If the observed Pb abundance of LP 625-44 generally represents the yields from very metal-poor AGB stars, their models of nucleosynthesis in AGB stars may overestimate its production by about two orders of magnitude at very low metallicity.
This conflict might be resolved by tuning the models of Gallino et al. (1998) or Busso et al. (1999). For instance, the neutron flux can be changed by modifying the extension or chemical profile of the <sup>13</sup>C-pocket, which is a free parameter in their models. Another parameter is the mass of the AGB star, upon which the number of thermal pulses (and hence episodes of neutron exposure) is strongly dependent.
Another possibility is that the s-process production mechanism in very metal-poor (e.g., \[Fe/H\]$`<2.5`$) AGB stars is quite different from that in more metal-rich stars. The calculation of low-mass stellar evolution in metal-deficient stars by Fujimoto et al. (2000) showed that hydrogen mixing occurs during the helium shell flash for 1$``$3.5 M stars with \[Fe/H\]$`<2.5`$ (their case II’), contrary to the situation for stars with higher metallicity (their case IV). Their result suggests that the production of <sup>13</sup>C, and subsequent s-process nucleosynthesis, is possible in the helium convective region during thermal pulses in these very metal-poor stars.
Our detection of Pb in the very metal-poor, carbon- and s-process-rich star, LP 625-44 provides a strong constraint on models of nucleosynthesis in AGB stars. The observed abundance patterns for heavier elements (Ba-Pb) agree well with the solar main s-process component, rather than with nucleosynthesis models for very metal-poor AGB stars. Further observation of s-process elements, including Pb, for objects similar to LP 625-44, and revisions of the theoretical approach to the nucleosynthesis in very metal-poor environments, will impact on our understanding of the evolution of low- and intermediate-mass stars, as well as of the enrichment of heavy elements in the early Galaxy. In this context, we note that we have also measured Pb in the star LP 706-7, which is similar to LP 625-44 in many respects (Norris et al., 1997a). That object, whose s-process abundances nevertheless differ from those of LP 625-44, will be discussed separately in a future paper.
We are grateful to the Director and staff of the Anglo-Australian Observatory, and the Australian Time Allocation Committee for providing the observational facilities used in this study. W.A. would like to acknowledge fruitful discussions with T. Kajino. T.C.B. acknowledges partial support of this work from grant AST 95-29454 awarded by the (US) National Science Foundation.
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# PREDICTIONS FOR ASSOCIATED PRODUCTION OF GAUGINOS AND GLUINOS AT NLO IN SUSY-QCD
## 1 Introduction
The search for supersymmetry (SUSY) is a major goal of the Tevatron Run II and LHC physics programs. If SUSY exists at the electroweak scale, SUSY partners of the Standard Model (SM) particles will either be discovered at these hadron colliders, or a very large region of SUSY parameter space will be excluded, provided reliable theoretical predictions in next-to-leading order (NLO) SUSY-QCD are available $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$. We calculate the NLO contributions for associated production of gauginos and gluinos at hadron colliders $`^{\mathrm{?},\mathrm{?}}`$. This production channel is enhanced by the strong coupling of the gluino and by the mass of the gaugino which is small in many popular models of SUSY breaking. The leptonic decay of the gaugino makes this process a good candidate for a mass determination of the gluino or the discovery or exclusion of a light gluino.
## 2 NLO SUSY-QCD Formalism
The associated production of a gluino and a gaugino proceeds in leading order (LO) through a quark-antiquark initial state and the exchange of an intermediate squark in the $`t`$-channel or $`u`$-channel. At NLO, virtual loop corrections must be considered which involve the exchange of intermediate SM or SUSY particles in self-energy, vertex correction, or box diagrams. Ultraviolet and infrared divergences appear at the upper and lower boundaries of integration over unobserved loop momenta. They are regulated dimensionally and removed through renormalization or cancellation with corresponding divergences in 2 to 3 parton (real emission) diagrams that have an additional gluon radiated into the final state. At NLO there is a second 2 to 3 parton process, one with a $`qg`$ initial state and a light quark emitted into the final state. In addition to soft divergences, real emission contributions have collinear divergences that are factored into the NLO parton densities.
## 3 Tevatron and LHC Cross Sections
To obtain quantitative predictions for the associated production of gauginos and gluinos at Run II of the Tevatron and at the LHC, we convolve LO and NLO partonic cross sections with CTEQ5 parton densities in LO and NLO ($`\overline{\mathrm{MS}}`$) along with 1- and 2-loop expressions for $`\alpha _s`$, the corresponding values of $`\mathrm{\Lambda }`$, and five active quark flavors. To constrain the SUSY parameter space, we choose an illustrative SUGRA model with $`m_0=100`$ GeV, $`A_0`$=300 GeV, tan $`\beta `$ = 4, sgn $`\mu `$ = +, and we vary $`m_{1/2}`$ between 100 and 400 GeV. The resulting masses for $`\stackrel{~}{\chi }_{1\mathrm{}4}^0`$ vary between 31…162, 63…317, 211…665, and 241…679 GeV, and $`\stackrel{~}{\chi }_{1,2}^\pm `$ are almost mass degenerate with $`\stackrel{~}{\chi }_{2,4}^0`$. The mass $`m_{\stackrel{~}{\chi }_3^0}<0`$ inside a polarization sum. Our method is not restricted to the SUGRA case and can be applied to any SUSY breaking model.
We present the total hadronic cross sections in Figure 1
as functions of the gluino mass. The light gaugino channels should be observable at both colliders, while the heavier Higgsino channels are suppressed by about one order of magnitude and might be observable only at the LHC.
The impact of the NLO corrections can be seen more readily in the ratio of NLO to LO cross sections computed at a renormalizaton scale set equal to the average mass of the final state particles. Figure 2
shows that the NLO effects are moderate (of $`𝒪`$ (10%)) at the Tevatron, while at the LHC the NLO contributions can increase the cross sections by as much as a factor of two. Enhancements of this size can shift mass determinations or discovery limits for SUSY particles by tens of GeV and must therefore always be taken into account.
An important measure of the theoretical uncertainty is the variation of the hadronic cross section with the renormalization and factorization scales. At LO, these scales enter only in the strong coupling constant and the parton densities, while at NLO they appear also explicitly in the hard cross section. As a result, the scale dependence is reduced considerably, as can be seen in Figure 3.
The Tevatron (LHC) cross sections vary by $`\pm 23(12)\%`$ at LO, but only by $`\pm 8(4.5)\%`$ in NLO when the scale is varied by a factor of two around the central scale.
For experimental searches, distributions in transverse momentum are important since cuts on $`p_T`$ help to enhance the signal. In Figure 4
we demonstrate that NLO contributions can have a large impact on $`p_T`$ spectra, especially at the LHC, where contributions from the $`gq`$ initial state become important. At the Tevatron the NLO $`p_T`$-distribution is shifted to lower $`p_T`$ with respect to the LO expectation.
## 4 Conclusions
The direct search for SUSY particles is a major goal at hadron colliders. For associated production of gauginos and gluinos, we demonstrate that NLO SUSY-QCD corrections stabilize the LO estimates against variations of the renormalization and factorization scale. The cross sections are increased by as much as a factor of two at the LHC, and the $`p_T`$ spectra are softened.
## Acknowledgments
M.K. thanks the organizers of the XXXVth Rencontres de Moriond for the kind invitation and financial support. Work at ANL is supported by the U.S. Department of Energy, Division of High Energy Physics, under Contract W-31-109-ENG-38. M.K. is supported by the Bundesministerium für Bildung und Forschung under Contract 05 HT9GUA 3, by Deutsche Forschungsgemeinschaft under Contract KL 1266/1-1, and by the European Commission under Contract ERBFMRXCT980194.
## References
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# Variational Approach to Hydrogen Atom in Uniform Magnetic Field of Arbitrary Strength
## I Introduction
The quantum statistical and quantum mechanical properties of a hydrogen atom in an external magnetic field are not exactly calculable. Perturbative approaches yield good results only for weak uniform fields as discussed in detail by Le Guillou and Zinn-Justin , who interpolated with analytic mapping techniques the ground state energy between weak- and strong-field. Other approaches are based on recursive procedures in higher-order perturbation theory . Zero-temperature properties were also investigated with the help of an operator optimization method in a second-quantized variational procedure . The behaviour at high uniform fields was inferred from treatments of the one-dimensional hydrogen atom . Hydrogen in strong magnetic fields is still a problem under investigation, since its solution is necessary to understand the properties of white dwarfs and neutron stars, as emphasized in Refs. .
A compact and detailed presentation of the bound states and highly accurate numerically values for the energy levels is given in Ref. .
Equations for a first-order variational approach to the ground-state energy of hydrogen in a uniform magnetic field based on the Jensen-Peierls inequality were written down a long time ago , but never evaluated. Apparently, they merely served as a preparation for attacking the more complicated problem of a polaron in a magnetic field .
In our approach, we calculate the quantum statistical properties of the system by an extension of variational perturbation theory . The crucial quantity is the effective classical potential. In the zero-temperature limit, it yields the ground state energy. Our calculations in a magnetic field require an extension of the formalism in Ref. which derives the effective classical potential from the phase space representation of the partition function.
Variational perturbation theory has an important advantage over other approaches: The calculation yields a good effective classical potential for all temperatures and coupling strengths. The quantum statistical partition function is obtained from a simple integral over a Boltzmann-factor involving the effective classical potential. The ground state energy is then obtained from its zero-temperature limit. The asymptotic behaviour in the strong-coupling limit is emerging automatically and does not have to be derived from other sources.
## II Effective Classical Representations for the Quantum Statistical Partition Function
A point particle in $`D`$ dimensions with a potential $`V(𝐱)`$ and a vector potential $`𝐀(𝐱)`$ is described by a Hamiltonian
$$H(𝐩,𝐱)=\frac{1}{2M}\left[𝐩\frac{e}{c}𝐀(𝐱)\right]^2+V(𝐱).$$
(1)
The quantum statistical partition function is given by the euclidean phase space path integral
$$Z=𝒟_{}^{}{}_{}{}^{D}x𝒟^Dpe^{𝒜[𝐩,𝐱]/\mathrm{}}$$
(2)
with an action
$$𝒜[𝐩,𝐱]=_0^\mathrm{}\beta 𝑑\tau \left[i𝐩(\tau )\dot{𝐱}(\tau )+H(𝐩(\tau ),𝐱(\tau ))\right],$$
(3)
and the path measure
$$𝒟_{}^{}{}_{}{}^{D}x𝒟^Dp=\underset{N\mathrm{}}{lim}\underset{n=1}{\overset{N+1}{}}\left[\frac{d^Dx_nd^Dp_n}{(2\pi \mathrm{})^D}\right].$$
(4)
The parameter $`\beta =1/k_BT`$ denotes the usual inverse thermal energy at temperature $`T`$, where $`k_B`$ is the Boltzmann constant. From $`Z`$ we obtain the free energy of the system:
$$F=\frac{1}{\beta }\mathrm{ln}Z.$$
(5)
In perturbation theory, one treats the external potential $`V(𝐱)`$ as a small quantity, and expands the partition function into powers of $`V(𝐱)`$. Such a naive expansion is applicable only for extremely weak couplings, and has a vanishing radius of convergence. Convergence is achieved by variational perturbation theory , which yields good approximations for all potential strengths, as we shall see in the sequel.
### A Effective Classical Potential
All quantum-mechanical systems studied so far in variational perturbation theory were governed by a Hamiltonian of the standard form
$$H(𝐩,𝐱)=\frac{𝐩^2}{2M}+V(𝐱).$$
(6)
The simple quadratic dependence on the momenta makes the momentum integrals in the path integral (2) trivial. The remaining configuration space representation of the partition function is used to define an effective classical potential $`V_{\mathrm{eff}}(𝐱_0)`$, from which quantum mechanical partition function is found by a classically looking integral
$$Z=\frac{d^Dx_0}{\lambda _{\mathrm{th}}^D}\mathrm{exp}\left[\beta V_{\mathrm{eff}}(𝐱_0)\right],$$
(7)
where $`\lambda _{\mathrm{th}}=\sqrt{2\pi \mathrm{}^2\beta /M}`$ is the thermal wavelength. The Boltzmann factor plays the role of a local partition function $`Z^{𝐱_0}`$, which is calculated from the restricted path integral
$$e^{\beta V_{\mathrm{eff}}(𝐱_0)}Z^{𝐱_0}=\lambda _{\mathrm{th}}^D𝒟^Dx\delta (𝐱_0\overline{𝐱(\tau )})e^{𝒜[𝐱]/\mathrm{}},$$
(8)
with the action
$$𝒜[𝐱]=_0^\mathrm{}\beta 𝑑\tau \left[\frac{M}{2}\dot{𝐱}^2(\tau )+V(𝐱(\tau ))\right],$$
(9)
and the path measure
$$𝒟^Dx=\underset{N\mathrm{}}{lim}\underset{n=1}{\overset{N+1}{}}\left\{\frac{d^Dx_n}{[2\pi \mathrm{}^2\beta /M(N+1)]^{D/2}}\right\}.$$
(10)
The special treatment of the temporal average of the Fourier path
$$𝐱_0=\overline{𝐱(\tau )}=\frac{1}{\mathrm{}\beta }_0^\mathrm{}\beta 𝑑\tau 𝐱(\tau )$$
(11)
is essential for the quality of the results. It subtracts from the harmonic fluctuation width $`𝐱^2^{\mathrm{cl}}`$ the classical divergence proportional to $`T=1/k_B\beta `$ of the Dulong-Petit law . Such diverging fluctuations cannot be treated perturbatively, and require the final integration in expression (7) to be done numerically.
For the Coulomb potential $`V(𝐱)=e^2/4\pi \epsilon _0|𝐱|`$ in three dimensions, the effective classical potential in Eq. (8) can be approximated well by variational perturbation theory .
### B Effective Classical Hamiltonian
In order to deal with Hamiltonians like (1) which contain a $`𝐩𝐀(𝐱)`$-term, we must generalize the variational procedure. Extending (8), we define an effective classical Hamiltonian by the phase space path integral
$$e^{\beta H_{\mathrm{eff}}(𝐩_0,𝐱_0)}Z^{𝐩_0,𝐱_0}=(2\pi \mathrm{})^D𝒟_{}^{}{}_{}{}^{D}x𝒟^Dp\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})e^{𝒜[𝐩,𝐱]/\mathrm{}},$$
(12)
with the action (3) and the measure (4). This allows us to express the partition function as the classically looking phase space integral
$$Z=\frac{d^Dx_0d^Dp_0}{(2\pi \mathrm{})^D}\mathrm{exp}\left[\beta H_{\mathrm{eff}}(𝐩_0,𝐱_0)\right],$$
(13)
where $`𝐩_0`$ is the temporal average of the momentum:
$$𝐩_0=\overline{𝐩(\tau )}=\frac{1}{\mathrm{}\beta }_0^\mathrm{}\beta 𝑑\tau 𝐩(\tau ).$$
(14)
The fixing of $`𝐩_0`$ is done for the same reason as that for $`𝐱_0`$, since the classical expectation value $`𝐩^2^{\mathrm{cl}}`$ is diverging linearly with $`T`$, just as $`𝐱^2^{\mathrm{cl}}`$.
In the special case of a standard Hamiltonian (6), the effective Hamiltonian in Eq. (13) reduces to the effective classical potential, since the momentum integral in Eq. (12) can then be easily performed, and the resulting restricted partition function becomes
$$Z^{𝐩_0,𝐱_0}=\mathrm{exp}\left(\beta \frac{𝐩_0^2}{2M}\right)Z^{𝐱_0}$$
(15)
with the local partition function $`Z^{𝐱_0}=\mathrm{exp}[\beta V_{\mathrm{eff}}(𝐱_0)]`$ of Eq. (8). Thus the complete quantum statistical partition function is given by (13), with an effective classical Hamilton function
$$H_{\mathrm{eff}}(𝐩_0,𝐱_0)=\frac{𝐩_0^2}{2M}+V_{\mathrm{eff}}(𝐱_0).$$
(16)
As a consequence of the purely quadratic momentum dependence of $`H(𝐩,𝐱)`$ in (6), the $`𝐩_0`$-integral in (13) can be done, thus expressing the quantum statistical partition function as a pure configuration space integral over the Boltzmann factor involving the effective classical potential $`V_{\mathrm{eff}}(𝐱_0)`$, as in Eq. (7).
### C Exact Effective Classical Hamiltonian for an Electron in a Constant Magnetic Field
The effective classical Hamiltonian for the electron moving in a constant magnetic field can be calculated exactly. We consider a magnetic field $`𝐁=B𝐞_z`$ pointing along the positive $`z`$-axis. The only nontrivial motion of the electron is in the $`xy`$-plane. In symmetric gauge the vector potential is given by
$$𝐀(𝐱)=\frac{B}{2}(y,x,0).$$
(17)
The choice of the gauge does not affect the partition function since the periodic path integral (2) is gauge invariant. Ignoring the trivial free particle motion along the $`z`$-direction, we may restrict our attention to the two-dimensional Hamiltonian
$$H(𝐩,𝐱)=\frac{𝐩^2}{2M}\frac{1}{2}\omega _cl_z(𝐩,𝐱)+\frac{1}{8}M\omega _c^2𝐱^2$$
(18)
with $`𝐱=(x,y)`$ and $`𝐩=(p_x,p_y)`$. Here, $`\omega _c=eB/Mc`$ is the Landau frequency, and
$$l_z(𝐩,𝐱)=(𝐱\times 𝐩)_z=xp_yyp_x$$
(19)
the third component of the orbital angular momentum. The partition function of the problem is given by Eq. (13), with $`D=2`$. Being interested in an effective classical formulation, we have to calculate the path integral (12). First we express the $`\delta `$-function for the averaged momentum as a Fourier integral
$`\delta (𝐩_0\overline{𝐩(\tau )})`$ $`=`$ $`{\displaystyle \frac{d^2\xi }{(2\pi \mathrm{})^2}\mathrm{exp}\left(\frac{i}{\mathrm{}}𝝃𝐩_0\right)\mathrm{exp}\left[\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐯_0(𝝃)𝐩(\tau )\right]}`$ (20)
involving an auxiliary source
$$𝐯_0(𝝃)=\frac{i}{\mathrm{}\beta }𝝃$$
(21)
which is constant in time. Substituting the $`\delta `$-function in Eq. (12) by this source representation, the partition function reads
$`Z^{𝐩_0,𝐱_0}`$ $`=`$ $`{\displaystyle d^2\xi \mathrm{exp}\left(\frac{i}{\mathrm{}}𝝃𝐩_0\right)𝒟_{}^{}{}_{}{}^{2}x𝒟^2p\delta (𝐱_0\overline{𝐱(\tau )})}`$ (23)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau \left[i𝐩(\tau )\dot{𝐱}(\tau )+H(𝐩(\tau ),𝐱(\tau ))+𝐯_0(𝝃)𝐩(\tau )\right]\right\}.`$
Evaluating the momentum integrals and utilizing the periodicity property $`𝐱(0)=𝐱(\mathrm{}\beta )`$, we obtain the configuration space path integral
$`Z^{𝐩_0,𝐱_0}=\underset{\mathrm{\Omega }0}{lim}`$ $`{\displaystyle d^2\xi \mathrm{exp}\left(\frac{i}{\mathrm{}}𝝃𝐩_0\frac{M}{2\mathrm{}^2\beta }𝝃^2\right)𝒟^2x\delta (𝐱_0\overline{𝐱(\tau )})}`$ (25)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau \left[{\displaystyle \frac{M}{2}}\dot{𝐱}^2(\tau )+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }^2𝐱^2(\tau ){\displaystyle \frac{i}{2}}M\omega _c(𝐱(\tau )\times \dot{𝐱}(\tau ))_z+𝐱(\tau )𝐣_1(𝝃)\right]\right\},`$
where the source $`𝐯_0`$ coupled to the momentum in (23) has turned to a source $`𝐣_1`$ coupled to the path in configuration space , with components
$$𝐣_1(𝝃)=\frac{M}{2}\omega _c(v_{0y}(𝝃),v_{0x}(𝝃))=\frac{i\omega _cM}{2\mathrm{}\beta }(\xi _y,\xi _x).$$
(26)
We have introduced an additional harmonic oscillator in Eq. (25) which will turn out to be useful at intermediate stages of the development. At the end of the calculation, only the limit $`\mathrm{\Omega }0`$ will be relevant.
Expressing the $`\delta `$-function in the path integral of Eq. (25) by the Fourier integral
$$\delta (𝐱_0\overline{𝐱(\tau )})=\frac{d^2\kappa }{(2\pi )^2}\mathrm{exp}\left(\frac{i}{\mathrm{}}𝜿𝐱_0\right)\mathrm{exp}\left[\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐣_2(𝜿)𝐱(\tau )\right]$$
(27)
with the new source
$$𝐣_2(𝜿)=\frac{i𝜿}{\beta },$$
(28)
the partition function (25) can be written as
$$Z^{𝐩_0,𝐱_0}=\underset{\mathrm{\Omega }0}{lim}d^2\xi \mathrm{exp}\left(\frac{i}{\mathrm{}}𝝃𝐩_0\frac{M}{2\mathrm{}^2\beta }𝝃^2\right)\frac{d^2\kappa }{(2\pi )^2}\mathrm{exp}\left(\frac{i}{\mathrm{}}𝜿𝐱_0\right)Z_\mathrm{\Omega }[𝐉(𝝃,𝜿)].$$
(29)
The functional $`Z_\mathrm{\Omega }[𝐉(𝝃,𝜿)]`$ is defined as the configuration space path integral
$$Z_\mathrm{\Omega }[𝐉(𝝃,𝜿)]=𝒟^2x\mathrm{exp}\left[\frac{1}{2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}𝐱(\tau )𝐆^1(\tau ,\tau ^{})𝐱(\tau ^{})\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐉(𝝃,𝜿)𝐱(\tau )\right],$$
(30)
where we have introduced the combined source $`𝐉(𝝃,𝜿)=𝐣_1(𝝃)+𝐣_2(𝜿)`$. Formally, the solution reads
$$Z_\mathrm{\Omega }[𝐉(𝝃,𝜿)]=Z_\mathrm{\Omega }[0]\mathrm{exp}\left[\frac{1}{2\mathrm{}^2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}𝐉(𝝃,𝜿)𝐆(\tau ,\tau ^{})𝐉(𝝃,𝜿)\right],$$
(31)
where $`𝐆(\tau ,\tau ^{})`$ is the matrix of Green functions obtained by inverting
$$𝐆^1(\tau ,\tau ^{})=\frac{M}{\mathrm{}}\left(\begin{array}{cc}\frac{d^2}{d\tau ^2}+\mathrm{\Omega }^2& i\omega _c\frac{d}{d\tau }\\ i\omega _c\frac{d}{d\tau }& \frac{d^2}{d\tau ^2}+\mathrm{\Omega }^2\end{array}\right)\delta (\tau \tau ^{}).$$
(32)
The inversion is easily done in frequency space after spectrally decomposing the $`\delta `$-function into the Matsubara frequencies $`\omega _m=2\pi m/\mathrm{}\beta `$,
$$\delta (\tau \tau ^{})=\frac{1}{\mathrm{}\beta }\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}e^{i\omega _m(\tau \tau ^{})}.$$
(33)
The result is
$$\stackrel{~}{𝐆}(\omega _m)=\frac{\mathrm{}}{M}\frac{1}{\mathrm{det}\stackrel{~}{𝐆}}\left(\begin{array}{cc}\omega _m^2+\mathrm{\Omega }^2& \omega _c\omega _m\\ \omega _c\omega _m& \omega _m^2+\mathrm{\Omega }^2\end{array}\right).$$
(34)
At this point, the additional oscillator in Eq. (25) proves useful: It ensures that the determinant
$$\mathrm{det}\stackrel{~}{𝐆}(\omega _m)=(\omega _m^2+\mathrm{\Omega }^2)^2+\omega _c^2\omega _m^2$$
(35)
is nonzero for $`m=0`$, thus playing the role of an infrared regulator. The Fourier expansion
$$𝐆(\tau ,\tau ^{})=\frac{1}{\mathrm{}\beta }\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\stackrel{~}{𝐆}(\omega _m)e^{i\omega _m(\tau \tau ^{})}$$
(36)
yields the matrix of Green functions
$$𝐆(\tau ,\tau ^{})=\left(\begin{array}{cc}G_{xx}(\tau ,\tau ^{})& G_{xy}(\tau ,\tau ^{})\\ G_{yx}(\tau ,\tau ^{})& G_{yy}(\tau ,\tau ^{})\end{array}\right)$$
(37)
which inherits the symmetry properties from the kernel (32):
$$G_{xx}(\tau ,\tau ^{})=G_{yy}(\tau ,\tau ^{}),G_{xy}(\tau ,\tau ^{})=G_{yx}(\tau ,\tau ^{}).$$
(38)
A more detailed description of these Green functions is given in Apps. A and B.
Since the current $`𝐉`$ does not depend on the euclidean time, the expression (31) simplifies therefore to
$$Z_\mathrm{\Omega }[𝐉(𝝃,𝜿)]=Z_\mathrm{\Omega }[0]\mathrm{exp}\left[\frac{1}{\mathrm{}^2}𝐉^2(𝝃,𝜿)_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}G_{xx}(\tau ,\tau ^{})\right].$$
(39)
The Green function has the Fourier decomposition
$$G_{xx}(\tau ,\tau ^{})=\frac{1}{M\beta }\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\omega _m^2+\mathrm{\Omega }^2}{(\omega _m^2+\mathrm{\Omega }_+^2)(\omega _m^2+\mathrm{\Omega }_{}^2)}e^{i\omega _m(\tau \tau ^{})},$$
(40)
where $`\mathrm{\Omega }_\pm `$ are the frequencies
$$\mathrm{\Omega }_\pm =\sqrt{\mathrm{\Omega }^2+\frac{1}{2}\omega _c^2\pm \omega _c\sqrt{\mathrm{\Omega }^2+\frac{1}{4}\omega _c^2}}.$$
(41)
The ratios in the sum of (40) can be decomposed into two partial fractions, each of them representing a single harmonic oscillator with frequency $`\mathrm{\Omega }_+`$ and $`\mathrm{\Omega }_{}`$, respectively. The analytic form of the periodic Green function of a single harmonic oscillator is well known (see Chap. 3 in ), and we obtain for the present Green function (41):
$$G_{xx}(\tau ,\tau ^{})=\frac{1}{M\beta }\left(\frac{\mathrm{}\beta }{2\mathrm{\Omega }_+}\frac{\mathrm{\Omega }_+^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\frac{\mathrm{cosh}\mathrm{\Omega }_+(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\frac{\mathrm{}\beta }{2\mathrm{\Omega }_{}}\frac{\mathrm{\Omega }_{}^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\frac{\mathrm{cosh}\mathrm{\Omega }_{}(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}\right).$$
(42)
By writing the determinant (35) as
$$\mathrm{det}\stackrel{~}{𝐆}(\omega _m)=(\omega _m^2+\mathrm{\Omega }_+^2)(\omega _m^2+\mathrm{\Omega }_{}^2)$$
(43)
and summing over the logarithms of this, we calculate the partition function as a product of two single harmonic oscillators:
$$Z_\mathrm{\Omega }=Z_\mathrm{\Omega }[0]=\frac{1}{2\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\frac{1}{2\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}.$$
(44)
The results (42) and (44) determine the generating functional (39). The euclidean time integrations are then easily done, and subsequently the $`𝜿`$\- and $`𝝃`$-integrations in (29). As a result, we obtain the restricted partition function
$$Z^{𝐩_0,𝐱_0}=\underset{\mathrm{\Omega }0}{lim}\mathrm{exp}\left\{\beta \left(\frac{1}{\beta }\mathrm{ln}\frac{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}{\mathrm{}\beta \mathrm{\Omega }_+/2}\frac{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{}\beta \mathrm{\Omega }_{}/2}+\frac{𝐩_0^2}{2M}\frac{1}{2}\omega _cl_z(𝐩_0,𝐱_0)+\frac{1}{8}M\omega _c^2𝐱_0^2+\frac{1}{2}M\mathrm{\Omega }^2𝐱_0^2\right)\right\}.$$
(45)
If we now remove the additional oscillator by taking the limit $`\mathrm{\Omega }0`$, we find from (41): $`\mathrm{\Omega }_+\omega _c`$, $`\mathrm{\Omega }_{}0`$, and therefore
$$\underset{\mathrm{\Omega }0}{lim}\frac{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}{\mathrm{}\beta \mathrm{\Omega }_+/2}=\frac{\mathrm{sinh}\mathrm{}\beta \omega _c/2}{\mathrm{}\beta \omega _c/2},\underset{\mathrm{\Omega }0}{lim}\frac{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{}\beta \mathrm{\Omega }_{}/2}=1.$$
(46)
Recalling the definition (12), we identify the exact effective classical Hamiltonian for an electron in a magnetic field as
$$H_{\mathrm{eff}}(𝐩_0,𝐱_0)=\frac{1}{\beta }\mathrm{ln}\frac{\mathrm{sinh}\mathrm{}\beta \omega _c/2}{\mathrm{}\beta \omega _c/2}+\frac{𝐩_0^2}{2M}\frac{1}{2}\omega _cl_z(𝐩_0,𝐱_0)+\frac{1}{8}M\omega _c^2𝐱_0^2.$$
(47)
Integrating out the momenta in Eq. (13), the configuration space representation (7) for the partition function contains the effective classical potential for a charged particle in the plane perpendicular to the direction of a uniform magnetic field
$$V_{\mathrm{eff}}(𝐱_0)=\frac{1}{\beta }\mathrm{ln}\frac{\mathrm{sinh}\mathrm{}\beta \omega _c/2}{\mathrm{}\beta \omega _c/2}.$$
(48)
Note that this is a constant potential.
Denoting the area $`d^2x_0`$ by $`A`$, we find the exact quantum statistical partition function
$$Z=\frac{A}{\lambda _{\mathrm{th}}^2}\frac{\mathrm{}\beta \omega _c/2}{\mathrm{sinh}\mathrm{}\beta \omega _c/2}.$$
(49)
After these preparations, we can turn our attention to the system we want to study in this paper: the hydrogen atom in a uniform magnetic field, where the additional Coulomb interaction prevents us from finding an exact solution for the effective classical Hamilton function.
## III Hydrogen Atom in Constant Magnetic Field
The zero-temperature properties of the hydrogen atom without external fiels are exactly known. For the quantum statistics at finite temperatures, an analytic expression exists, but it is hard to evaluate. It is easier to find an accurate approximative result with the help of variational perturbation theory . Similar calculations have been performed for the electron-proton pair distribution function which can be interpreted as the unnormalized density matrix .
Here we extend this method of calculation to the hydrogen atom in a constant magnetic field. This extension is quite nontrivial since the weak- and strong-field limits will turn out to exhibit completely different asymptotic behaviours. Let us first generalize variational perturbation theory to an electron in a constant magnetic field and arbitrary potential.
### A Generalized Variational Perturbation Theory
We consider once more the effective classical form (13) of the quantum statistical partition function, which requires the path integration (12) in phase space. Fluctuations parallel and vertical to the magnetic field lines are now both nontrivial, and we must deal with the full three-dimensional system and the components of the electron position and momentum are now denoted by $`𝐱=(x,y,z)`$ and $`𝐩=(p_x,p_y,p_z)`$. For the uniform magnetic field pointing along the $`z`$-axis, the vector potential $`𝐀(𝐱)`$ is used in the gauge (17). Thus the Hamilton function of an electron in a magnetic field and an arbitrary potential $`V(𝐱)`$ is
$$H(𝐩,𝐱)=\frac{𝐩^2}{2M}\frac{1}{2}\omega _cl_z(𝐩,𝐱)+\frac{1}{8}M\omega _c^2𝐱^2+V(𝐱).$$
(50)
The orbital angular momentum $`l_z(𝐩,𝐱)`$ was introduced in Eq. (19), and the Landau frequency $`\omega _c`$ below Eq. (18). The importance of the separation of the zero frequency components $`𝐱_0`$ and $`𝐩_0`$ was discussed in Sect. II. Their divergence with the temperature $`T`$ prevents a perturbative treatment. Thus it is essential to set up the perturbation theory only for the fluctuations around $`𝐱_0`$ and $`𝐩_0`$. For this we rewrite the action functional (3) associated with the Hamiltonian (50) as
$$𝒜[𝐩,𝐱]=𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]+𝒜_{\mathrm{int}}[𝐩,𝐱],$$
(51)
where we have introduced the fluctuation action
$`𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]={\displaystyle _0^\mathrm{}\beta }d\tau \{`$ $`i[𝐩(\tau )𝐩_0]\dot{𝐱}(\tau )+{\displaystyle \frac{1}{2M}}[𝐩(\tau )𝐩_0]^2+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_1l_z(𝐩(\tau )𝐩_0,𝐱(\tau )𝐱_0)`$ (53)
$`+{\displaystyle \frac{1}{8}}M\mathrm{\Omega }_2^2[𝐱^{}(\tau )𝐱_0^{}]^2+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_{}^2[z(\tau )z_0]^2\},`$
in which $`𝐱^{}=(x,y)`$ denotes the transverse part of $`𝐱`$. The interaction is now
$$𝒜_{\mathrm{int}}[𝐩,𝐱]=_0^\mathrm{}\beta 𝑑\tau V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))=𝒜[𝐩,𝐱]𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]$$
(54)
with the interaction potential
$`V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))`$ $`=`$ $`{\displaystyle \frac{1}{2M}}\left\{𝐩^2(\tau )\left[𝐩(\tau )𝐩_0\right]^2\right\}+{\displaystyle \frac{1}{2}}\omega _c𝐩^{}(\tau )\times 𝐱^{}(\tau )`$ (57)
$`{\displaystyle \frac{1}{2}}\mathrm{\Omega }_1(𝐩^{}(\tau )𝐩_0^{})\times (𝐱^{}(\tau )𝐱_0^{})+{\displaystyle \frac{1}{8}}M\omega _c^2𝐱_{}^{}{}_{}{}^{2}(\tau )`$
$`{\displaystyle \frac{1}{8}}M\mathrm{\Omega }_2^2\left[𝐱^{}(\tau )𝐱_0^{}\right]^2{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_{}^2[z(\tau )z_0]^2+V(𝐱(\tau )),`$
where $`𝐩^{}=(p_x,p_y)`$. The frequencies $`𝛀=(\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{\Omega }_{})`$ are for the moment arbitrary. The decomposition (51) forms the basis for the variational approach, where the first term in the action (51) allows an exact treatment. The transverse part was given in Sec. II C and the longitudinal part is trivial, since it is harmonic with frequency $`\mathrm{\Omega }_{}`$. The associated partition function is given by the path integral
$$Z_𝛀^{𝐩_0,𝐱_0}=𝒟_{}^{}{}_{}{}^{3}x𝒟^3p\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})e^{𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]/\mathrm{}},$$
(58)
which can be performed. Details are given in Appendix C. The result is
$$Z_𝛀^{𝐩_0,𝐱_0}=\frac{\mathrm{}\beta \mathrm{\Omega }_+/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\frac{\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}\frac{\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2},$$
(59)
where auxiliary frequencies are composed of the frequencies $`\mathrm{\Omega }_1,\mathrm{\Omega }_2`$ in the action (53) as
$$\mathrm{\Omega }_\pm (\mathrm{\Omega }_1,\mathrm{\Omega }_2)=\frac{1}{2}|\mathrm{\Omega }_1\pm \mathrm{\Omega }_2|.$$
(60)
This partition function serves in the subsequent pertubation expansion as trial system which depends explicitly on the frequencies $`𝛀`$. The correlation functions are a straightforward generalization of (37) to three dimensions:
$$𝐆^{𝐱_0}(\tau ,\tau ^{})=\left(\begin{array}{ccc}G_{xx}^{𝐱_0}(\tau ,\tau ^{})& G_{xy}^{𝐱_0}(\tau ,\tau ^{})& 0\\ G_{yx}^{𝐱_0}(\tau ,\tau ^{})& G_{yy}^{𝐱_0}(\tau ,\tau ^{})& 0\\ 0& 0& G_{zz}^{𝐱_0}(\tau ,\tau ^{})\end{array}\right),$$
(61)
whose explicit form is derived in App. C.
The $`𝛀`$-dependent action in Eq. (51) is treated perturbatively. Writing the partition function (12) as
$$Z^{𝐩_0,𝐱_0}=(2\pi \mathrm{})^3𝒟_{}^{}{}_{}{}^{3}x𝒟^3p\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})\mathrm{exp}\left\{\frac{1}{\mathrm{}}𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]\right\}\mathrm{exp}\left\{\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))\right\},$$
(62)
the second exponential is expanded into a Taylor series, yielding
$`Z^{𝐩_0,𝐱_0}`$ $`=`$ $`(2\pi \mathrm{})^3{\displaystyle 𝒟_{}^{}{}_{}{}^{3}x𝒟^3p\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})\mathrm{exp}\left\{\frac{1}{\mathrm{}}𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]\right\}}`$ (64)
$`\times \left[1{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))+{\displaystyle \frac{1}{2!\mathrm{}^2}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau _1{\displaystyle _0^\mathrm{}\beta }𝑑\tau _2V_{\mathrm{int}}(𝐩(\tau _1),𝐱(\tau _1))V_{\mathrm{int}}(𝐩(\tau _2),𝐱(\tau _2))\mathrm{}\right].`$
Defining harmonic expectation values with respect to the restricted path integral as
$$\mathrm{}_𝛀^{𝐩_0,𝐱_0}=\frac{(2\pi \mathrm{})^3}{Z_𝛀^{𝐩_0,𝐱_0}}𝒟_{}^{}{}_{}{}^{3}x𝒟^3p\mathrm{}\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})\mathrm{exp}\left\{\frac{1}{\mathrm{}}𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]\right\},$$
(65)
the perturbation expansion for the partition function (64) reads
$$Z^{𝐩_0,𝐱_0}=Z_𝛀^{𝐩_0,𝐱_0}\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{\mathrm{}^nn!}\left(_0^\mathrm{}\beta 𝑑\tau V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))\right)^n_𝛀^{𝐩_0,𝐱_0}.$$
(66)
This power series expansion can be rewritten in the exponential form
$$Z^{𝐩_0,𝐱_0}=Z_𝛀^{𝐩_0,𝐱_0}\mathrm{exp}\left\{\underset{n=0}{\overset{\mathrm{}}{}}\frac{(1)^n}{\mathrm{}^nn!}\left(_0^\mathrm{}\beta 𝑑\tau V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))\right)^n_{𝛀,c}^{𝐩_0,𝐱_0}\right\},$$
(67)
where the subscript $`c`$ on the expectation values indicates cumulants. The lowest cumulants are related to the full expectation values as follows:
$`O_1(𝐩(\tau _1),𝐱(\tau _1))_{𝛀,c}^{𝐩_0,𝐱_0}`$ $`=`$ $`O_1(𝐩(\tau _1),𝐱(\tau _1))_𝛀^{𝐩_0,𝐱_0},`$ (68)
$`O_1(𝐩(\tau _1),𝐱(\tau _1))O_2(𝐩(\tau _2),𝐱(\tau _2))_{𝛀,c}^{𝐩_0,𝐱_0}`$ $`=`$ $`O_1(𝐩(\tau _1),𝐱(\tau _1))O_2(𝐩(\tau _2),𝐱(\tau _2))_𝛀^{𝐩_0,𝐱_0}`$ (70)
$`O_1(𝐩(\tau _1),𝐱(\tau _1))_𝛀^{𝐩_0,𝐱_0}O_2(𝐩(\tau _2),𝐱(\tau _2))_𝛀^{𝐩_0,𝐱_0},`$
$`\mathrm{}`$ $`,`$ (71)
where $`O_i(𝐩(\tau _j),𝐱(\tau _j))`$ denotes any observable depending on position and momentum. Recalling the relation (12) between partition function (67) and effective classical Hamiltonian $`H_{\mathrm{eff}}(𝐩_0,𝐱_0)`$, we obtain from (67) the effective classical Hamiltonian as a cumulant expansion:
$$H_{\mathrm{eff}}(𝐩_0,𝐱_0)=\frac{1}{\beta }\mathrm{ln}Z_𝛀^{𝐩_0,𝐱_0}+\frac{1}{\beta }\underset{n=1}{\overset{\mathrm{}}{}}\frac{(1)^n}{\mathrm{}^nn!}\left(_0^\mathrm{}\beta 𝑑\tau V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))\right)^n_{𝛀,c}^{𝐩_0,𝐱_0}.$$
(72)
Up to now, we did not make any approximation. The expansion on the right-hand side is an exact expression for the effective classical Hamiltonian for any $`𝛀`$.
For systems with a nontrivial interaction, we are capable of calculating only some initial truncated part of the series (72), say up to the $`N`$th order, leading to the approximate effective classical Hamiltonian
$$_𝛀^{(N)}(𝐩_0,𝐱_0)=\frac{1}{\beta }\mathrm{ln}Z_𝛀^{𝐩_0,𝐱_0}+\frac{1}{\beta }\underset{n=1}{\overset{N}{}}\frac{(1)^n}{\mathrm{}^nn!}\left(_0^\mathrm{}\beta 𝑑\tau V_{\mathrm{int}}(𝐩(\tau ),𝐱(\tau ))\right)^n_{𝛀,c}^{𝐩_0,𝐱_0}.$$
(73)
This depends explicitly on the three parameters $`𝛀`$. Since the exact expression (72) is independent of $`𝛀`$, the best approximation for $`_𝛀^{(N)}(𝐩_0,𝐱_0)`$ should depend on $`𝛀`$ minimally. Thus the optimal solution will be found by determining the parameters from the conditions
$$_𝛀_𝛀^{(N)}(𝐩_0,𝐱_0)\stackrel{!}{=}0.$$
(74)
Let us denote the optimal variational parameters to $`N`$th order by
$$𝛀^{(N)}=(\mathrm{\Omega }_1^{(N)}(𝐩_0,𝐱_0),\mathrm{\Omega }_2^{(N)}(𝐩_0,𝐱_0),\mathrm{\Omega }_{}^{(N)}(𝐩_0,𝐱_0)).$$
(75)
Inserting these into Eq. (73) yields the optimal effective classical Hamiltonian $`^{(N)}(𝐩_0,𝐱_0)`$.
### B First-Order Effective Classical Potential
The first-order approximation of the effective classical Hamiltonian (73) reads
$$_𝛀^{(1)}(𝐩_0,𝐱_0)=\frac{1}{\beta }\mathrm{ln}Z_𝛀^{𝐩_0,𝐱_0}V_{\mathrm{int}}(𝐩,𝐱)_𝛀^{𝐩_0,𝐱_0}.$$
(76)
In writing the last term we have used the fact that, as a consequence of the time translation invariance of the system, the first-order expectation value of $`V_{\mathrm{int}}(𝐱)`$ is independent of the euclidean time $`\tau `$.
In order to calculate $`_𝛀^{(1)}(𝐩_0,𝐱_0)`$, we use the two-point correlation functions derived in App. C, and the vanishing of the linear expectations, e.g.
$$p_x(\tau )p_{0}^{}{}_{x}{}^{}_𝛀^{𝐩_0,𝐱_0}=0$$
(77)
to find
$$_𝛀^{(1)}(𝐩_0,𝐱_0)=\frac{𝐩_0^2}{2M}\frac{1}{2}\omega _cl_z(𝐩_0,𝐱_0)+\frac{1}{8}M\omega _c^2(x_0^2+y_0^2)+W_𝛀^{(1)}(𝐱_0),$$
(78)
where we have collected all terms depending on the variational parameters $`𝛀`$ in the potential
$$W_𝛀^{(1)}(𝐱_0)=\frac{1}{\beta }\mathrm{ln}Z_𝛀^{𝐩_0,𝐱_0}+(\omega _c\mathrm{\Omega }_1)b_{}^2(𝐱_0)\frac{1}{4}\left(\mathrm{\Omega }_2^2\omega _c^2\right)a_{}^2(𝐱_0)\frac{1}{2}M\mathrm{\Omega }_{}^2a_{}^2(𝐱_0)+V(𝐱)_𝛀^{𝐩_0,𝐱_0}.$$
(79)
The quantities $`a_{}^2(𝐱_0)`$ and $`a_{}^2(𝐱_0)`$ are the transverse and longitudinal fluctuation widths
$$a_{}^2(𝐱_0)=G_{xx}^{𝐩_0,𝐱_0}(0),a_{}^2(𝐱_0)=G_{zz}^{𝐩_0,𝐱_0}(0),b_{}^2(𝐱_0)=G_{xp_y}^{𝐩_0,𝐱_0}(0).$$
(80)
Note that the potential (79) is independent of $`𝐩_0`$. This means that the approximation (78) to the effective classical Hamiltonian contains no coupling of the momentum $`𝐩_0`$ to a variational parameter $`𝛀`$, such that the optimal $`𝛀^{(1)}`$ determined by minimizing $`_𝛀^{(1)}(𝐩_0,𝐱_0)`$ is independent of $`𝐩_0`$. We may therefore integrate out $`𝐩_0`$ in the phase space representation of the first-order approximation for the partition function
$$Z^{(1)}=\frac{d^3x_0d^3p_0}{(2\pi \mathrm{}^3)}e^{\beta _𝛀^{(1)}(𝐩_0,𝐱_0)}$$
(81)
to find the pure configuration space integral
$$Z^{(1)}=\frac{d^3x_0}{\lambda _{\mathrm{th}}^3}e^{\beta W_𝛀^{(1)}(𝐱_0)},$$
(82)
in which $`W_𝛀^{(1)}(𝐱_0)`$ is the first-order approximation to the effective classical potential of an electron in a potential $`V(𝐱)`$ and a uniform magnetic field.
### C Application to the Hydrogen Atom in a Magnetic Field
We now apply the formulas of the preceding section to the Hamiltonian (50) with an attracting Coulomb potential
$$V(𝐱)=\frac{e^2}{4\pi \epsilon _0|𝐱|},$$
(83)
where $`|𝐱|`$ is the distance between the electron and the proton. The only nontrivial problem is the calculation of the expectation value $`V(𝐱(\tau ))_𝛀^{𝐩_0,𝐱_0}`$ in Eq. (79). This is done using the so-called smearing formula, which is a Gaussian convolution of $`V(𝐱)`$. This formula was first derived by Feynman and Kleinert , and exists now also in an extension to arbitrary order . The generalization to position and momentum dependent observables was given in the phase space formulation . We briefly rederive the first-order smearing formula. The expectation value is defined by
$$V(𝐱(\tau ^{}))_𝛀^{𝐩_0,𝐱_0}=\frac{(2\pi \mathrm{})^3}{Z_𝛀^{𝐩_0,𝐱_0}}𝒟_{}^{}{}_{}{}^{3}x𝒟^3pV(𝐱(\tau ^{}))\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})e^{𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]/\mathrm{}},$$
(84)
Now we substitute the potential by the expression
$$V(𝐱(\tau ^{}))=d^3xV(𝐱)\delta (𝐱𝐱(\tau ^{}))=d^3xV(𝐱)\frac{d^3\kappa }{(2\pi )^3}\mathrm{exp}\left[i𝜿(𝐱𝐱_0)\right]\mathrm{exp}\left\{\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐣(\tau )[𝐱(\tau )𝐱_0]\right\},$$
(85)
where we have introduced the source
$$𝐣(\tau )=i\mathrm{}𝜿\delta (\tau \tau ^{}).$$
(86)
Inserting the expression (85) into Eq. (84) we obtain
$$V(𝐱(\tau ^{}))_𝛀^{𝐩_0,𝐱_0}=\frac{1}{Z_𝛀^{𝐩_0,𝐱_0}}d^3xV(𝐱)\frac{d^3\kappa }{(2\pi )^3}\mathrm{exp}\left[i𝜿(𝐱𝐱_0)\right]Z_𝛀^{𝐩_0,𝐱_0}[𝐣],$$
(87)
with the harmonic generating functional
$$Z_𝛀^{𝐩_0,𝐱_0}[𝐣]=(2\pi \mathrm{})^3𝒟_{}^{}{}_{}{}^{3}x𝒟^3p\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})\mathrm{exp}\left\{\frac{1}{\mathrm{}}𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐣(\tau )[𝐱(\tau )𝐱_0]\right\}.$$
(88)
The solution is
$$Z_𝛀^{𝐩_0,𝐱_0}[𝐣]=Z_𝛀^{𝐩_0,𝐱_0}\mathrm{exp}\left[\frac{1}{2\mathrm{}^2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}𝐣(\tau )𝐆^{𝐱_0}(\tau ,\tau ^{})𝐣(\tau ^{})\right]$$
(89)
with the $`3\times 3`$-matrix of Green functions of Eq. (61). The properties of the Green functions are discussed in the Appendices A and B. Expressing the source $`𝐣(\tau )`$ in terms of $`𝜿`$ via Eq. (86) and performing the $`\tau `$-integrations, we arrive at
$$V(𝐫(\tau ^{}))_𝛀^{𝐩_0,𝐱_0}=d^3xV(𝐱)\frac{d^3\kappa }{(2\pi )^3}\mathrm{exp}\left\{i𝜿[𝐱𝐱_0]\right\}\mathrm{exp}\left[\frac{1}{2}𝜿𝐆^{𝐱_0}(0)𝜿\right].$$
(90)
Recognizing that $`G_{yx}^{𝐱_0}(0)=G_{xy}^{𝐱_0}(0)`$ vanish, the $`𝜿`$-integral is easily calculated and leads to the first-order smearing formula for an arbitrary position dependent potential
$$V(𝐱(\tau ^{}))_𝛀^{𝐩_0,𝐱_0}=\frac{1}{(2\pi )^{3/2}a_{}^2(𝐱_0)\sqrt{a_{}^2(𝐱_0)}}d^3xV(𝐱)\mathrm{exp}\left[\frac{(xx_0)^2+(yy_0)^2}{2a_{}^2(𝐱_0)}\frac{(zz_0)^2}{2a_{}^2(𝐱_0)}\right],$$
(91)
the right-hand side containing the Gaussian fluctuation widths (80).
For the Coulomb potential (83) that we are interested in, the integral in the smearing formula (91) can not be done exactly. An integral representation for a simple numerical treatment is
$`{\displaystyle \frac{e^2}{4\pi \epsilon _0|𝐱|}}_𝛀^{𝐩_0,𝐱_0}`$ $`=`$ $`{\displaystyle \frac{e^2}{4\pi \epsilon _0}}\sqrt{{\displaystyle \frac{2}{\pi }}a_{}^2(𝐱_0)}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{d\xi }{a_{}^2(𝐱_0)+\xi ^2[a_{}^2(𝐱_0)a_{}^2(𝐱_0)]}}`$ (93)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\xi ^2}{2}}\left({\displaystyle \frac{x_0^2+y_0^2}{a_{}^2(𝐱_0)+\xi ^2[a_{}^2(𝐱_0)a_{}^2(𝐱_0)]}}+{\displaystyle \frac{z_0^2}{a_{}^2(𝐱_0)}}\right)\right\}.`$
With this expression we know the entire first-order effective classical potential (79) for an electron in a Coulomb potential and a uniform magnetic field which has to be optimized in the variational parameters $`𝛀`$.
## IV Results
We are now going to optimize the effective classical potential by extremizing it in $`𝛀`$ at different temperatures and magnetic field strengths. In the zero-temperature limit this will produce the ground state energy.
### A Effective Classical Potential for Different Temperatures and Magnetic Field Strengths
The optimization of $`W_𝛀^{(1)}(𝐱_0)`$ proceeds by minimization in $`𝛀`$ and must be done for each value of $`𝐱_0`$. Reinserting the optimal parameters $`𝛀^{(1)}(𝐱_0)`$ into the expressions (79) and (93), we obtain the optimal first-order effective classical potential $`W^{(1)}(𝐱_0)`$. The calculations are done numerically, where we used natural units $`\mathrm{}=e^2/4\pi \epsilon _0=k_B=c=M=1`$. This means that energies are measured in units of $`ϵ_0=Me^4/(4\pi \epsilon _0)^2\mathrm{}^22\mathrm{Ryd}27.21\mathrm{eV}`$, temperatures in $`ϵ_0/k_B3.16\times 10^5\mathrm{K}`$, distances in Bohr radii $`a_B=(4\pi \epsilon _0)^2\mathrm{}^2/Me^20.53\times 10^{10}\mathrm{m}`$, and magnetic field strengths in $`B_0=e^3M^2/\mathrm{}^3(4\pi \epsilon _0)^22.35\times 10^5\mathrm{T}=2.35\times 10^9\mathrm{G}`$. Figure 1 shows the resulting curves for various magnetic field strengths $`B`$ and an inverse tempature $`\beta =1/T=1`$. Examples of the lower temperature behaviour are shown in Fig. 2 for $`\beta =100`$. To see the expected anisotropy of the curves in the magnetic field direction and in the plane perpendicular to it, we plot simultanously the curves for $`W^{(1)}(𝐱_0)`$ transversal to the magnetic field as a function of $`\rho _0=\sqrt{x_0^2+y_0^2}`$ at $`z=0`$ (solid curves) and parallel as a function of $`z_0`$ at $`\rho _0=0`$ (dashed curves). The curves become strongly anisotropic for low temperatures and increasing field strengths (Fig. 2). At a given field strength $`B`$, the two curves converge for large distances from the origin, where the proton resides, to the same constant depending on $`B`$. This is due to the decreasing influence of the Coulomb interaction which shows the classical $`1/r`$-behaviour in each direction. When approaching the classical high-temperature limit, the effect of anisotropy becomes less important since the violent thermal fluctuations do not have a preferred direction (see Fig. 1). For $`\rho _0\mathrm{}`$ or $`z_0\mathrm{}`$, the expectation value of the Coulomb potential (93) tends to zero. The remaining effective classical potential
$$W_𝛀^{(1)}(𝐱_0)\frac{1}{\beta }\mathrm{ln}Z_𝛀^{𝐩_0,𝐱_0}+(\omega _c\mathrm{\Omega }_1)b_{}^2\frac{1}{4}\left(\mathrm{\Omega }_2^2\omega _c^2\right)a_{}^2\frac{1}{2}M\mathrm{\Omega }_{}^2a_{}^2$$
(94)
is a constant with regard to the position $`𝐱_0`$, and the optimization yields $`\mathrm{\Omega }_1^{(1)}=\mathrm{\Omega }_2^{(1)}=\omega _c`$ and $`\mathrm{\Omega }_{}^{(1)}=0`$, leading to the asymptotic constant value
$$W^{(1)}(𝐱_0)\frac{1}{\beta }\mathrm{ln}\frac{\mathrm{}\beta \omega _c/2}{\mathrm{sinh}\mathrm{}\beta \omega _c/2}.$$
(95)
The $`B=0`$ -curves are of course identical with those obtained from variational perturbation theory for the hydrogen atom .
### B Ground State Energy of the Hydrogen Atom in Uniform Magnetic Field
In what follows we investigate the zero-temperature behaviour of the theory. Figures 1 and 2 show that the minimum of each potential curve lies at the origin. This means that the first-order approximation to the ground state energy for a fixed magnitude of the magnetic field $`B`$ is found by considering the zero-temperature limit of the first-order effective classical potential in the origin
$$E^{(1)}=\underset{\beta \mathrm{}}{lim}W^{(1)}(0).$$
(96)
Thus we obtain from Eq. (79) the variational expression for the ground state energy:
$$E_𝛀^{(1)}(B)=\frac{\mathrm{}}{4\mathrm{\Omega }_2}\left(\mathrm{\Omega }_2^2+\omega _c^2\right)+\frac{\mathrm{}\mathrm{\Omega }_{}}{4}\frac{e^4}{4\pi \epsilon _0}\frac{1}{|𝐱|}_𝛀^\mathrm{𝟎},$$
(97)
where the expectation value for the Coulomb potential (93) can now be calculated exactly since the exponential in the integral simplifies to unity:
$$\frac{1}{|𝐱|}_𝛀^\mathrm{𝟎}=2\sqrt{\frac{M}{\pi \mathrm{}}}\times \{\begin{array}{cc}\sqrt{\frac{\mathrm{\Omega }_{}\mathrm{\Omega }_2}{\mathrm{\Omega }_{}\mathrm{\Omega }_2}}\mathrm{arctan}\sqrt{\frac{2\mathrm{\Omega }_{}}{\mathrm{\Omega }_2}1},& 2\mathrm{\Omega }_{}>\mathrm{\Omega }_2,\hfill \\ \sqrt{\mathrm{\Omega }_{}},& 2\mathrm{\Omega }_{}=\mathrm{\Omega }_2,\hfill \\ \frac{1}{2i}\sqrt{\frac{\mathrm{\Omega }_{}\mathrm{\Omega }_2}{\mathrm{\Omega }_{}\mathrm{\Omega }_2}}\mathrm{ln}\frac{1+i\sqrt{2\mathrm{\Omega }_{}/\mathrm{\Omega }_21}}{1i\sqrt{2\mathrm{\Omega }_{}/\mathrm{\Omega }_21}},& 2\mathrm{\Omega }_{}<\mathrm{\Omega }_2.\hfill \end{array}$$
(98)
The equations (97) and (98) are independent of the frequency parameter $`\mathrm{\Omega }_1`$ such that the optimization of the first-order expression for the ground state energy (97) requires the satisfying of the equations
$$\frac{E_𝛀^{(1)}(B)}{\mathrm{\Omega }_2}\stackrel{!}{=}0,\frac{E_𝛀^{(1)}(B)}{\mathrm{\Omega }_{}}\stackrel{!}{=}0.$$
(99)
Reinserting the resulting values $`\mathrm{\Omega }_2^{(1)}`$ and $`\mathrm{\Omega }_{}^{(1)}`$ into Eq. (97) yields the first-order approximation for the ground state energy $`E^{(1)}(B)`$. In the absence of the Coulomb interaction the optimization with respect to $`\mathrm{\Omega }_2`$ yields $`\mathrm{\Omega }_2^{(1)}=\omega _c`$, rendering the ground state energy $`E^{(1)}(B)=\omega _c/2`$, which is the zeroth Landau level. An optimal value for $`\mathrm{\Omega }_{}`$ does not exist since the dependece of the ground state energy of this parameter is linear in Eq. (97) in this special case. To obtain the lowest energy, this parameter can be set to zero (all optimal frequency parameters used in the optimization procedure turn out to be nonnegative). For a vanishing magnetic field, $`B=0`$, Eq. (97) exactly reproduces the first-order variational result for the ground state energy of the hydrogen atom, $`E^{(1)}(B=0)0.42[2\mathrm{Ryd}]`$, obtained in Ref. .
To investigate the asymptotics in the strong-field limit $`B\mathrm{}`$, it is useful to extract the leading term $`\omega _c/2`$. Thus we define the binding energy
$$\epsilon (B)\frac{\omega _c}{2}E(B)$$
(100)
which possesses an characteristic strong-field behaviour to be discussed in detail subsequently. The result is shown in Fig. 3 as a function of the magnitude of the magnetic field $`B`$, where it is compared with the high-accuracy results of Ref. . As a first-order approximation, this result is satisfactory. It is of the same quality like other first-order results, for example those from the operator optimization method in first order of Ref. . The advantage of variational perturbation theory is that it yields good results over the complete range of the coupling strength, here the magnetic field. Moreover, as a consequence of the exponential convergence \[16, Chap. 5\], higher orders of variational perturbation theory push the approximative result of any quantity very rapidly towards the exact value.
#### 1 The Weak-Field Case
We investigate now the weak-field behaviour of our theory starting from the expression (100) and the expectation value of the Coulomb potential (98) in natural units:
$$\epsilon _{\eta ,\mathrm{\Omega }}^{(1)}(B)=\frac{B}{2}\frac{\mathrm{\Omega }}{4}\left(1+\frac{\eta }{2}\right)\frac{B^2}{4\mathrm{\Omega }}\sqrt{\frac{\eta \mathrm{\Omega }}{2\pi }}h(\eta )$$
(101)
with
$$h(\eta )=\frac{1}{\sqrt{1\eta }}\mathrm{ln}\frac{1\sqrt{1\eta }}{1+\sqrt{1\eta }}.$$
(102)
In comparison with Eq. (97) we introduced new variational parameters
$$\eta \frac{2\mathrm{\Omega }_{}}{\mathrm{\Omega }_2},\mathrm{\Omega }\mathrm{\Omega }_2$$
(103)
and utilized, as the calculations for the binding energy showed, that always $`\eta 1`$. Performing the derivatives with respect to these variational parameters and setting them zero yields conditional equations which can be written after some manipulations as
$`{\displaystyle \frac{\mathrm{\Omega }}{8}}+\sqrt{{\displaystyle \frac{\mathrm{\Omega }}{2\pi \eta }}}{\displaystyle \frac{1}{1\eta }}\left(1+{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\sqrt{1\eta }}}\mathrm{ln}{\displaystyle \frac{1\sqrt{1\eta }}{1+\sqrt{1\eta }}}\right)`$ $`\stackrel{!}{=}`$ $`0,`$ (104)
$`{\displaystyle \frac{1}{4}}+{\displaystyle \frac{\eta }{8}}{\displaystyle \frac{B^2}{4\mathrm{\Omega }^2}}+{\displaystyle \frac{1}{2}}\sqrt{{\displaystyle \frac{\eta }{2\pi \mathrm{\Omega }}}}{\displaystyle \frac{1}{\sqrt{1\eta }}}\mathrm{ln}{\displaystyle \frac{1\sqrt{1\eta }}{1+\sqrt{1\eta }}}`$ $`\stackrel{!}{=}`$ $`0.`$ (105)
Expanding the variational parameters into perturbation series of the square magnetic field $`B^2`$,
$$\eta (B)=\underset{n=0}{\overset{\mathrm{}}{}}\eta _nB^{2n},\mathrm{\Omega }(B)=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Omega }_nB^{2n}$$
(106)
and inserting these expansions into the self-consistency conditions (104) and (105) we obtain order by order the coefficients given in Table I. Inserting these values into the expression for the binding energy (101) and expand with respect to $`B^2`$, we obtain the perturbation series
$$\epsilon ^{(1)}(B)=\frac{B}{2}\underset{n=0}{\overset{\mathrm{}}{}}\epsilon _nB^{2n}.$$
(107)
The first coefficients are also given in Table I. We find thus the important result that the first-order variational perturbation solution possesses a perturbative behaviour with respect to the square magnetic field strength $`B^2`$ in the weak-field limit thus yielding the correct asymptotics. The coefficients differ in higher order from the exact ones but are improved in higher orders of the variational perturbation theory \[16, Chap. 5\].
#### 2 Asymptotical Behaviour in the Strong-Field Regime
In the discussion of the pure magnetic field below Eq. (99) we have mentioned that the variational calculation for the ground state energy which is thus associated with the zeroth Landau level yields a frequency $`\mathrm{\Omega }_2B`$ while $`\mathrm{\Omega }_{}=0`$. Therefore we use the assumption (with $`\mathrm{\Omega }_{}\mathrm{\Omega }_2`$)
$$\mathrm{\Omega }_{}2\mathrm{\Omega }_{},\mathrm{\Omega }_{}B$$
(108)
for the consideration of the ground state energy (97) of the hydrogen atom in a strong magnetic field. In a first step we expand the last expression of the expectation value (98) which corresponds to the condition (108) in terms of $`2\mathrm{\Omega }_{}/\mathrm{\Omega }_{}`$ and reinsert this expansion in the equation of the ground state energy (97). Then we omit all terms proportional to $`C/\mathrm{\Omega }_{}`$ where $`C`$ stands for any expression with a value much smaller than the field strength $`B`$. In natural units, we thus obtain the strong-field approximation for the first-order binding energy (100)
$$\epsilon _{\mathrm{\Omega }_{},\mathrm{\Omega }_{}}^{(1)}=\frac{B}{2}\left(\frac{\mathrm{\Omega }_{}}{4}+\frac{B^2}{4\mathrm{\Omega }_{}}+\frac{\mathrm{\Omega }_{}}{4}+\sqrt{\frac{\mathrm{\Omega }_{}}{\pi }}\mathrm{ln}\frac{\mathrm{\Omega }_{}}{2\mathrm{\Omega }_{}}\right).$$
(109)
As usual, we consider the zeros of the derivatives with respect to the variational parameters
$$\frac{\epsilon _{\mathrm{\Omega }_{},\mathrm{\Omega }_{}}^{(1)}}{\mathrm{\Omega }_{}}\stackrel{!}{=}0,\frac{\epsilon _{\mathrm{\Omega }_{},\mathrm{\Omega }_{}}^{(1)}}{\mathrm{\Omega }_{}}\stackrel{!}{=}0,$$
(110)
which lead to the self-consistence equations
$`\sqrt{\mathrm{\Omega }_{}}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}\left(\mathrm{ln}\mathrm{\Omega }_{}\mathrm{ln}\mathrm{\Omega }_{}+2\mathrm{ln}\mathrm{\hspace{0.17em}2}\right),`$ (111)
$`\mathrm{\Omega }_{}`$ $`=`$ $`2\sqrt{{\displaystyle \frac{\mathrm{\Omega }_{}}{\pi }}}+B\sqrt{1+4{\displaystyle \frac{\mathrm{\Omega }_{}}{\pi B^2}}}`$ (112)
Let us first consider the last equation. Utilizing the second of the conditions (108) we expand the second root around unity yielding the expression
$$\mathrm{\Omega }_{}=B+2\sqrt{\frac{\mathrm{\Omega }_{}}{\pi }}+2\frac{\mathrm{\Omega }_{}}{\pi B}4\frac{\mathrm{\Omega }_{}^2}{\pi ^2B^3}+\mathrm{},$$
(113)
where the terms are sorted with regard to their contribution starting with the biggest. Since we are interested in the strong $`B`$ limit, we can obviously neglect terms suppressed by powers of $`1/B`$. Thus we only consider the following terms for the moment:
$$\mathrm{\Omega }_{}B+2\sqrt{\frac{\mathrm{\Omega }_{}}{\pi }}.$$
(114)
Inserting this into the other condition (111), expanding the corresponding logarithm, and, once more, neglecting terms of order $`1/B`$, we find
$$\sqrt{\mathrm{\Omega }_{}}\frac{2}{\sqrt{\pi }}\left(\mathrm{ln}B\mathrm{ln}\mathrm{\Omega }_{}^{(1)}+\mathrm{ln}\mathrm{\hspace{0.17em}2}2\right).$$
(115)
To obtain a tractable approximation for $`\mathrm{\Omega }_{}`$, we perform some iterations starting from
$$\sqrt{\mathrm{\Omega }_{}^{(1)}}=\frac{2}{\sqrt{\pi }}\mathrm{ln}\mathrm{\hspace{0.17em}2}Be^2$$
(116)
Reinserting this on the right-hand side of Eq. (115), one obtains the second iteration $`\sqrt{\mathrm{\Omega }_{}^{(2)}}`$. We stop this procedure after an additional reinsertion which yields
$$\sqrt{\mathrm{\Omega }_{}^{(3)}}=\frac{2}{\sqrt{\pi }}\left(\mathrm{ln}\mathrm{\hspace{0.17em}2}Be^22\mathrm{l}\mathrm{n}\left[\frac{2}{\sqrt{\pi }}\left\{\mathrm{ln}\mathrm{\hspace{0.17em}2}Be^22\mathrm{l}\mathrm{n}\left(\frac{2}{\sqrt{\pi }}\mathrm{ln}\mathrm{\hspace{0.17em}2}Be^2\right)\right\}\right]\right).$$
(117)
The reader may convince himself that this iteration procedure indeed converges. For a subsequent systematical extraction of terms essentially contributing to the binding energy, the expression (117) is not satisfactory. Therefore it is better to separate the leading term in the curly brackets and expand the logarithm of the remainder. Then this proceeding is applied to the expression in the angular brackets and so on. Neglecting terms of order $`\mathrm{ln}^3B`$, we obtain
$$\sqrt{\mathrm{\Omega }_{}^{(3)}}\frac{2}{\sqrt{\pi }}\left(\mathrm{ln}\mathrm{\hspace{0.17em}2}Be^2+\mathrm{ln}\frac{\pi }{4}2\mathrm{l}\mathrm{n}\mathrm{l}\mathrm{n}\mathrm{\hspace{0.17em}2}Be^2\right).$$
(118)
The double-logarithmic term can be expanded in a similar way as described above:
$$\mathrm{lnln}\mathrm{\hspace{0.17em}2}Be^2=\mathrm{ln}\left[\mathrm{ln}B\left(1+\frac{\mathrm{ln}\mathrm{\hspace{0.17em}2}2}{\mathrm{ln}B}\right)\right]=\mathrm{lnln}B+\frac{\mathrm{ln}\mathrm{\hspace{0.17em}2}2}{\mathrm{ln}B}\frac{1}{2}\frac{(\mathrm{ln}\mathrm{\hspace{0.17em}2}2)^2}{\mathrm{ln}^2B}+𝒪(\mathrm{ln}^3B).$$
(119)
Thus the expression (118) may be rewritten as
$$\sqrt{\mathrm{\Omega }_{}^{(3)}}=\frac{2}{\sqrt{\pi }}\left(\mathrm{ln}B2\mathrm{l}\mathrm{n}\mathrm{l}\mathrm{n}B+\frac{2a}{\mathrm{ln}B}+\frac{a^2}{\mathrm{ln}^2B}+b\right)+𝒪(\mathrm{ln}^3B)$$
(120)
with abbreviations
$$a=2\mathrm{ln}\mathrm{\hspace{0.17em}2}1.307,b=\mathrm{ln}\frac{\pi }{2}21.548.$$
(121)
The first observation is that the variational parameter $`\mathrm{\Omega }_{}`$ is always much smaller than $`\mathrm{\Omega }_{}`$ in the high $`B`$-field limit. Thus we can further simplify the approximation (114) by replacing
$$\mathrm{\Omega }_{}B\left(1+\frac{2}{B}\sqrt{\frac{\mathrm{\Omega }_{}}{\pi }}\right)B$$
(122)
without affecting the following expression for the binding energy. Inserting the solutions (120) and (122) into the equation for the binding energy (109) and expanding the logarithmic term once more as described, we find up to the order $`\mathrm{ln}^2B`$:
$`\epsilon ^{(1)}(B)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\left(\mathrm{ln}^2B4\mathrm{ln}B\mathrm{lnln}B+4\mathrm{ln}^2\mathrm{ln}B4b\mathrm{lnln}B+2(b+2)\mathrm{ln}B+b^2{\displaystyle \frac{1}{\mathrm{ln}B}}\left[8\mathrm{ln}^2\mathrm{ln}B8b\mathrm{lnln}B+2b^2\right]\right)`$ (124)
$`+𝒪(\mathrm{ln}^2B)`$
Note that the prefactor $`1/\pi `$ of the leading $`\mathrm{ln}^2B`$-term differs from a value $`1/2`$ obtained by Landau and Lifschitz . Our different value is a consequence of using a harmonic trial system. The calculation of higher orders in variational perturbation theory would improve the value of the prefactor.
At a magnetic field strength $`B=10^5B_0`$, which corresponds to $`2.35\times 10^{10}\mathrm{T}=2.35\times 10^{14}\mathrm{G}`$, the contribution from the first six terms is $`22.87[2\mathrm{Ryd}]`$. The next three terms suppressed by a factor $`\mathrm{ln}^1B`$ contribute $`2.29[2\mathrm{Ryd}]`$, while an estimate for the $`\mathrm{ln}^2B`$-terms yields nearly $`0.3[2\mathrm{Ryd}]`$. Thus we find
$$\epsilon ^{(1)}(10^5)=20.58\pm 0.3[2\mathrm{Ryd}].$$
(125)
This is in very good agreement with the value $`20.60[2\mathrm{Ryd}]`$ obtained from the full treatment described in Sec. IV B.
Table II lists the values of the first six terms of Eq. (124). This shows in particular the significance of the second-leading term $`(4/\pi )\mathrm{ln}B\mathrm{lnln}B`$, which is of the same order of the leading term $`(1/\pi )\mathrm{ln}^2B`$ but with an opposite sign. In Fig. 3, we have plotted the expression
$$\epsilon _L(B)=\frac{1}{2}\mathrm{ln}^2B$$
(126)
from Landau and Lifschitz to illustrate that it gives far too large binding energies even at very large magnetic fields, e.g. at $`2000B_010^{12}\mathrm{G}`$.
This strength of magnetic field appears on surfaces of neutron stars ($`10^{10}10^{12}\mathrm{G}`$). A recently discovered new type of neuton star is the so-called magnetar. In these, charged particles such as protons and electrons produced by decaying neutrons give rise to the giant magnetic field of $`10^{15}\mathrm{G}`$. Magnetic fields of white dwarfs reach only up to $`10^610^8\mathrm{G}`$. All these magnetic field strengths are far from realization in experiments. The strongest magnetic fields ever produced in a laboratory were only of the order $`10^5\mathrm{G}`$, an order of magnitude larger than the fields in sun spots which reach about $`0.4\times 10^4\mathrm{G}`$. Recall, for comparison, that the earth’s magnetic field has the small value of $`0.6\mathrm{G}`$.
As we see in Fig. 3, the nonleading terms in Eq. (124) give important contributions to the asymptotic behaviour even at such large magnetic fields. It is an unusual property of the asymptotic behaviour that the absolute value of the difference between the Landau-expression (126) and our approximation (124) diverges with increasing magnetic field strengths $`B`$, only the relative difference decreases.
## V Summary
We have calculated the effective classical potential for the hydrogen atom in a magnetic field. For this we have generalized variational perturbation theory to make it applicable to physical systems with uniform external magnetic field.
The effective classical potential containing the complete quantum statistical information of the system was determined in first-order variational perturbation theory. For zero-temperature, it gave the energy of the system. Our result consists of a single analytic expression which is quite accurate at all temperatures and magnetic field strengths.
## Acknowledgments
We thank Prof. J. Čížek and Dr. J. Weniger for useful hints and references of a perturbative treatment of the ground state properties of the hydrogen atom in a magnetic field. The authors also thank Dr. J. Ortner and M. Steinberg for discussions of the finite temperature behaviour of this system. For interesting discussions we would also like to thank Prof. J.T. Devreese and Prof. G. Wunner.
One of us (M.B.) is grateful for support by the Studienstiftung des deutschen Volkes.
## A Generating Functional for Particle in Magnetic Field and Harmonic Oscillator Potential
For the determination of the correlation functions of a system, we need to know the solution of the two-dimensional generating functional in the presence of an external source $`𝐣=(j_x,j_y)`$:
$$Z^{𝐱_0}[𝐣]=\lambda _{\mathrm{th}}^2𝒟^2x\delta (𝐱_0\overline{𝐱(\tau )})e^{𝒜^{𝐱_0}[𝐱;𝐣]/\mathrm{}}.$$
(A1)
The action of a particle in a magnetic field in $`z`$-direction and a harmonic oscillator reads
$$𝒜^{𝐱_0}[𝐱;𝐣]=_0^\mathrm{}\beta 𝑑\tau \left[\frac{M}{2}\dot{𝐱}^2(\tau )\frac{i}{2}M\omega ([𝐱(\tau )𝐱_0]\times \dot{𝐱}(\tau ))_z+\frac{1}{2}M\mathrm{\Omega }^2[𝐱(\tau )𝐱_0]^2+𝐣(\tau )(𝐱(\tau )𝐱_0)\right].$$
(A2)
The position dependent terms are centered around $`𝐱_0=(x_0,y_0)`$, which is the temporal average of the path $`𝐱(\tau )`$, and thus equal to the zero frequency component of the Fourier path, is
$$𝐱(\tau )=𝐱_0+\underset{m=1}{\overset{\mathrm{}}{}}\left(𝐱_me^{i\omega _m\tau }+𝐱_m^{}e^{i\omega _m\tau }\right)$$
(A3)
with the Matsubara frequencies $`\omega _m=2\pi m/\mathrm{}\beta `$ and complex Fourier coefficients $`𝐱_m=𝐱_m^{\mathrm{re}}+i𝐱_m^{\mathrm{im}}`$. Introducing a similar Fourier decomposition for the current $`𝐣(\tau )`$ with Fourier components $`𝐣_m`$ and using the orthonormality relation
$$\frac{1}{\mathrm{}\beta }_0^\mathrm{}\beta 𝑑\tau e^{i(\omega _m\omega _n)\tau }=\delta _{mn},$$
(A4)
the generating functional can be written as
$$Z^{𝐱_0}[𝐣]=\underset{m=1}{\overset{\mathrm{}}{}}\left[\frac{dx_m^{\mathrm{re}}dx_m^{\mathrm{im}}dy_m^{\mathrm{re}}dy_m^{\mathrm{im}}}{(\pi /M\beta \omega _m^2)^2}e^{𝒜_m(𝐱_m,𝐱_m^{};𝐣_m,𝐣_m^{})/\mathrm{}}\right]$$
(A5)
with
$`𝒜_m(𝐱_m,𝐱_m^{};𝐣_m,𝐣_m^{})`$ $`=`$ $`\mathrm{}\beta M(\omega _m^2+\mathrm{\Omega }^2)([x_m^{\mathrm{re}}]^2+[x_m^{\mathrm{im}}]^2+[y_m^{\mathrm{re}}]^2+[y_m^{\mathrm{im}}]^2)+2i\mathrm{}\beta M\omega \omega _m(x_m^{\mathrm{re}}y_m^{\mathrm{im}}x_m^{\mathrm{im}}y_m^{\mathrm{re}})`$ (A7)
$`+2\mathrm{}\beta (x_m^{\mathrm{re}}j_{x}^{}{}_{m}{}^{\mathrm{re}}+x_m^{\mathrm{im}}j_{x}^{}{}_{m}{}^{\mathrm{im}}+y_m^{\mathrm{re}}j_{y}^{}{}_{m}{}^{\mathrm{re}}+y_m^{\mathrm{im}}j_{y}^{}{}_{m}{}^{\mathrm{im}}).`$
Expression (A5) is equivalent to the path integral (A1) and we obtain after performing the integrations and retransforming the currents
$$𝐣_m=\frac{1}{\mathrm{}\beta }_0^\mathrm{}\beta 𝑑\tau 𝐣(\tau )e^{i\omega _m\tau }$$
(A8)
the resulting generating functional
$$Z^{𝐱_0}[𝐣]=Z^{𝐱_0}\mathrm{exp}\left\{\frac{1}{2\mathrm{}^2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}𝐣(\tau )𝐆^{𝐱_0}(\tau ,\tau ^{})𝐣(\tau ^{})\right\}$$
(A9)
with the partition function
$$Z^{𝐱_0}Z^{𝐱_0}[0]=\underset{m=1}{\overset{\mathrm{}}{}}\frac{\omega _m^4}{\omega ^2\omega _m^2+(\omega _m^2+\mathrm{\Omega }^2)^2}$$
(A10)
and the $`2\times 2`$-matrix of Green functions
$$𝐆^{𝐱_0}(\tau ,\tau ^{})=\left(\begin{array}{cc}G_{xx}^{𝐱_0}(\tau ,\tau ^{})& G_{xy}^{𝐱_0}(\tau ,\tau ^{})\\ G_{yx}^{𝐱_0}(\tau ,\tau ^{})& G_{yy}^{𝐱_0}(\tau ,\tau ^{})\end{array}\right).$$
(A11)
The elements of this matrix are position-position correlation functions what can be easily proved by applying two functional derivatives with respect to the desired component of the current to the functional (A1), for example
$$G_{xx}^{𝐱_0}(\tau ,\tau ^{})=(x(\tau )x_0)(x(\tau ^{})x_0)^{𝐱_0}=\left[\mathrm{}^2\frac{1}{Z^{𝐱_0}[𝐣]}\frac{\delta ^2}{\delta j_x(\tau )\delta j_x(\tau ^{})}Z^{𝐱_0}[𝐣]\right]_{𝐣=0},$$
(A12)
where we have defined expectation values by
$$\mathrm{}^{𝐱_0}=\frac{\lambda _{\mathrm{th}}^2}{Z^{𝐱_0}}𝒟^2x\mathrm{}\delta (𝐱_0\overline{𝐱(\tau )})e^{𝒜^{𝐱_0}[𝐱;0]/\mathrm{}}.$$
(A13)
From the above calculation we find the following expressions for the Green functions in Fourier space ($`0\tau ,\tau ^{}\mathrm{}\beta `$):
$`G_{xx}^{𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{x}(\tau )\stackrel{~}{x}(\tau ^{})^{𝐱_0}=G_{yy}^{𝐱_0}(\tau ,\tau ^{})=\stackrel{~}{y}(\tau )\stackrel{~}{y}(\tau ^{})^{𝐱_0}`$ (A14)
$`=`$ $`{\displaystyle \frac{2}{M\beta }}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\omega _m^2+\mathrm{\Omega }^2}{\omega ^2\omega _m^2+(\omega _m^2+\mathrm{\Omega }^2)^2}}e^{i\omega _m(\tau \tau ^{})},`$ (A15)
$`G_{xy}^{𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{x}(\tau )\stackrel{~}{y}(\tau ^{})^{𝐱_0}=G_{yx}^{𝐱_0}(\tau ,\tau ^{})=\stackrel{~}{y}(\tau )\stackrel{~}{x}(\tau ^{})^{𝐱_0}`$ (A16)
$`=`$ $`{\displaystyle \frac{2\omega }{M\beta }}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\omega _m}{\omega ^2\omega _m^2+(\omega _m^2+\mathrm{\Omega }^2)^2}}e^{i\omega _m(\tau \tau ^{})},`$ (A17)
where, for simplicity, $`\stackrel{~}{𝐱}(\tau )=𝐱(\tau )𝐱_0`$. It is desirable to find analytical expressions for the Green functions and the partition function (A10). All these quantities possess the same dominator which can be decomposed as
$$\omega ^2\omega _m^2+(\omega _m^2+\mathrm{\Omega }^2)^2=(\omega _m^2+\mathrm{\Omega }_+^2)(\omega _m^2+\mathrm{\Omega }_{}^2)$$
(A18)
with frequencies
$$\mathrm{\Omega }_\pm (\omega ,\mathrm{\Omega })=\sqrt{\mathrm{\Omega }^2+\frac{1}{2}\omega ^2\pm \omega \sqrt{\mathrm{\Omega }^2+\frac{1}{4}\omega ^2}}.$$
(A19)
Therefore the partition function (A10) can be split into two products, each of which known from the harmonic oscillator \[16, Chap. 5\]:
$$Z^{𝐱_0}=\underset{m=1}{\overset{\mathrm{}}{}}\left[\frac{\omega _m^2}{\omega _m^2+\mathrm{\Omega }_+^2}\right]\underset{m=1}{\overset{\mathrm{}}{}}\left[\frac{\omega _m^2}{\omega _m^2+\mathrm{\Omega }_{}^2}\right]=\frac{\mathrm{}\beta \mathrm{\Omega }_+/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\frac{\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}.$$
(A20)
Now we apply the property (A18) to decompose the Green functions (A14) into partial fractions, yielding
$$G_{xx}^{𝐱_0}(\tau ,\tau ^{})=G_{yy}^{𝐱_0}(\tau ,\tau ^{})=\frac{1}{M\beta }\left(\alpha _1\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{\omega _m^2+\mathrm{\Omega }_+^2}e^{i\omega _m(\tau \tau ^{})}+\alpha _2\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{\omega _m^2+\mathrm{\Omega }_{}^2}e^{i\omega _m(\tau \tau ^{})}\frac{1}{\mathrm{\Omega }^2}\right)$$
(A21)
with coefficients
$$\alpha _1=\frac{\mathrm{\Omega }_+^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2},\alpha _2=\frac{\mathrm{\Omega }_{}^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}.$$
(A22)
Following Ref. \[16, Chap. 3\], sums of the kind occuring in expression (A21) are spectral decompositions of the correlation function for the harmonic oscillator and can be summed up:
$$\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{\omega _m^2+\mathrm{\Omega }_\pm ^2}e^{i\omega _m(\tau \tau ^{})}=\frac{\mathrm{}\beta }{2\mathrm{\Omega }_\pm }\frac{\mathrm{cosh}\mathrm{\Omega }_\pm (|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_\pm /2}.$$
(A23)
Thus, the $`xx`$\- and $`yy`$-correlation functions can be expressed by
$`G_{xx}^{𝐱_0}(\tau ,\tau ^{})=G_{yy}^{𝐱_0}(\tau ,\tau ^{})`$ (A24)
$`={\displaystyle \frac{1}{M\beta }}\left({\displaystyle \frac{\mathrm{}\beta }{2\mathrm{\Omega }_+}}{\displaystyle \frac{\mathrm{\Omega }_+^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}}{\displaystyle \frac{\mathrm{cosh}\mathrm{\Omega }_+(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}}{\displaystyle \frac{\mathrm{}\beta }{2\mathrm{\Omega }_{}}}{\displaystyle \frac{\mathrm{\Omega }_{}^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}}{\displaystyle \frac{\mathrm{cosh}\mathrm{\Omega }_{}(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}}{\displaystyle \frac{1}{\mathrm{\Omega }^2}}\right),`$ (A25)
where, from Eq. (A19), $`\mathrm{\Omega }_\pm =\mathrm{\Omega }_\pm (\omega ,\mathrm{\Omega })`$ are functions of the original frequencies $`\omega `$ from the magnetic field and $`\mathrm{\Omega }`$ from the additional harmonic oscillator (A2). It is obvious that expression (A24) reduces to the Green function of the harmonic oscillator for $`\omega 0`$:
$$\underset{\omega 0}{lim}G_{ii}^{𝐱_0}(\tau ,\tau ^{})=\frac{1}{M\beta \mathrm{\Omega }^2}\left(\frac{\mathrm{}\beta \mathrm{\Omega }}{2}\frac{\mathrm{cosh}\mathrm{\Omega }(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }/2}1\right)$$
(A26)
with $`i\{x,y\}`$. In this limit, the partition function (A20) turns out to be the usual one \[16, Chap. 5\] for such a harmonic oscillator
$$\underset{\omega 0}{lim}Z^{𝐱_0}=\frac{\mathrm{}\beta \mathrm{\Omega }/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }/2}.$$
(A27)
It is worth mentioning that with the last term in Green function (A24) the classical harmonic fluctuation width
$$G_{xx}^{\mathrm{cl}}=x^2^{\mathrm{cl}}=\frac{1}{M\beta \mathrm{\Omega }^2}$$
(A28)
is subtracted. This is the consequence of the exclusion of the zero frequency mode of the Fourier path (A3) in the generating functional (A1). The necessity to do this has already been discussed in Sect. II. The other terms in Eq. (A24) are those which we would have obtained without separation of the $`x_0`$-component. Thus these terms represent the quantum mechanical Green function containing all quantum as well as thermal fluctuations. It is a nice property of all Green functions discussed in this paper that
$$G_{xx}^{𝐱_0}(\tau ,\tau ^{})=G_{xx}^{\mathrm{qm}}(\tau ,\tau ^{})G_{xx}^{\mathrm{cl}}.$$
(A29)
Such a relation exists for all other Green functions appropriately, including momentum-position correlations which we consider subsequently.
The knowledge of relation (A23) makes it quite easy to determine the algebraic expression for the mixed $`xy`$-correlation functions. Rewriting Eq. (A16) as
$$G_{xy}^{𝐱_0}(\tau ,\tau ^{})=G_{yx}^{𝐱_0}(\tau ,\tau ^{})=\frac{i\omega }{M\beta }\frac{1}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\frac{}{\tau }\left(\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{\omega _m^2+\mathrm{\Omega }_+^2}e^{i\omega _m(\tau \tau ^{})}+\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\frac{1}{\omega _m^2+\mathrm{\Omega }_{}^2}e^{i\omega _m(\tau \tau ^{})}\right)$$
(A30)
and applying the derivative with respect to $`\tau `$ to relation (A23), we obtain the following expression for the mixed Green function:
$`G_{xy}^{𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`G_{yx}^{𝐱_0}(\tau ,\tau ^{})`$ (A31)
$`=`$ $`{\displaystyle \frac{\mathrm{}\omega }{2iM}}{\displaystyle \frac{1}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}}\left\{\mathrm{\Theta }(\tau \tau ^{})[g_{\mathrm{\Omega }_+}(\tau ,\tau ^{})g_\mathrm{\Omega }_{}(\tau ,\tau ^{})]\mathrm{\Theta }(\tau ^{}\tau )[g_{\mathrm{\Omega }_+}(\tau ^{},\tau )g_\mathrm{\Omega }_{}(\tau ^{},\tau )]\right\},`$ (A32)
where we have used the abbreviation
$$g_{\mathrm{\Omega }_\pm }(\tau ,\tau ^{})=\frac{\mathrm{sinh}\mathrm{\Omega }_\pm (\tau \tau ^{}\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_\pm /2},\tau ,\tau ^{}(0,\mathrm{}\beta ).$$
(A33)
Note that classically $`xy^{\mathrm{cl}}=0`$ such that Eq. (A29) reduces to
$$G_{xy}^{𝐱_0}(\tau ,\tau ^{})=G_{xy}^{\mathrm{qm}}(\tau ,\tau ^{}).$$
(A34)
The Heaviside function in Eq. (A31) is defined symmetrically:
$$\mathrm{\Theta }(\tau \tau ^{})=\{\begin{array}{cc}1& \tau >\tau ^{},\\ 1/2& \tau =\tau ^{},\\ 0& \tau <\tau ^{}.\end{array}$$
(A35)
In the quantum mechanical limit of zero-temperature ($`\beta \mathrm{}`$), the Green function (A24) simplifies to
$$\underset{\beta \mathrm{}}{lim}G_{xx}^{𝐱_0}(\tau ,\tau ^{})=\underset{\beta \mathrm{}}{lim}G_{yy}^{𝐱_0}(\tau ,\tau ^{})=\frac{\mathrm{}}{2M}\left(\frac{1}{\mathrm{\Omega }_+}\frac{\mathrm{\Omega }_+^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}e^{\mathrm{\Omega }_+|\tau \tau ^{}|}\frac{1}{\mathrm{\Omega }_{}}\frac{\mathrm{\Omega }_{}^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}e^{\mathrm{\Omega }_{}|\tau \tau ^{}|}\right),$$
(A36)
while in Eq. (A31) only $`g_{\mathrm{\Omega }_\pm }(\tau ,\tau ^{})`$ changes:
$$\underset{\beta \mathrm{}}{lim}g_{\mathrm{\Omega }_\pm }(\tau ,\tau ^{})=e^{\mathrm{\Omega }_\pm (\tau \tau ^{})}.$$
(A37)
## B Properties of Green Functions
In this section we list properties of the Green functions (A24) and (A31) which are important for the forthcoming consideration of the generating functional with sources coupling linearily to position or momentum in Appendix C. For all relations we suppose that $`0\tau ,\tau ^{}\mathrm{}\beta `$.
### 1 General Properties
A first observation is the temporal translational invariance of the Green functions:
$$G_{ij}^{𝐱_0}(\tau ,\tau ^{})=G_{ij}^{𝐱_0}(\tau \tau ^{}),$$
(B1)
where each of the indices $`i,j`$ stands for $`x`$ or $`y`$, respectively. For equal times we find
$$G_{ij}^{𝐱_0}(\tau ,\tau )=\frac{1}{M\beta }\left(\frac{\mathrm{}\beta }{2\mathrm{\Omega }_+}\frac{\mathrm{\Omega }_+^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\mathrm{coth}\mathrm{}\beta \mathrm{\Omega }_+/2\frac{\mathrm{}\beta }{2\mathrm{\Omega }_{}}\frac{\mathrm{\Omega }_{}^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\mathrm{coth}\mathrm{}\beta \mathrm{\Omega }_{}/2\frac{1}{\mathrm{\Omega }^2}\right)\times \{\begin{array}{cc}1& i=j,\\ 0& ij.\end{array}$$
(B2)
Moreover we read off the following symmetries from the expressions (A24) and (A31):
$$G_{ij}^{𝐱_0}(\tau ,\tau ^{})=G_{ij}^{𝐱_0}(\tau ^{},\tau )\times \{\begin{array}{cc}1& i=j,\\ 1& ij.\end{array}$$
(B3)
Otherwise,
$$G_{ij}^{𝐱_0}(\tau ,\tau ^{})=G_{ji}^{𝐱_0}(\tau ^{},\tau ).$$
(B4)
Throughout the paper we always use periodic paths. Hence it is obvious that all Green functions are periodic, too:
$$G_{ij}^{𝐱_0}(0,\tau ^{})=G_{ij}^{𝐱_0}(\mathrm{}\beta ,\tau ^{}),G_{ij}^{𝐱_0}(\tau ,0)=G_{ij}^{𝐱_0}(\tau ,\mathrm{}\beta ).$$
(B5)
### 2 Derivatives of Green Functions
We now proceed with derivatives of the Green functions (A24) and (A31), since these are essential for the derivation of the generating functional of position and momentum dependent correlations in the forthcoming Appendix C.
Before considering the concrete expressions we introduce a new symbol indicating uniquely to which argument the derivative is applied. A dot on the left-hand side means to perform the derivative with respect to the first argument and the dot on the right-hand side indicates that to differentiate with respect to the other argument. Having a dot on both sides the Green function is derived with respect to both arguments:
$${}_{}{}^{}G_{ij}^{𝐱_0}(\tau ,\tau ^{})=\frac{G_{ij}^{𝐱_0}(\tau ,\tau ^{})}{\tau },G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,\tau ^{})=\frac{G_{ij}^{𝐱_0}(\tau ,\tau ^{})}{\tau ^{}},{}_{}{}^{}G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,\tau ^{})=\frac{^2G_{ij}^{𝐱_0}(\tau ,\tau ^{})}{\tau \tau ^{}}.$$
(B6)
Applying such derivatives to the Green functions (A24), we obtain ($`i\{x,y\}`$):
$${}_{}{}^{}G_{ii}^{𝐱_0}(\tau ,\tau ^{})=\frac{\mathrm{}}{2M}\frac{1}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\left[\mathrm{\Theta }(\tau \tau ^{})g(\tau ,\tau ^{})\mathrm{\Theta }(\tau ^{}\tau )g(\tau ^{},\tau )\right]=G_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\tau ^{})$$
(B7)
with
$$g(\tau ,\tau ^{})=(\mathrm{\Omega }_+^2\mathrm{\Omega }^2)g_{\mathrm{\Omega }_+}(\tau ,\tau ^{})(\mathrm{\Omega }_{}^2\mathrm{\Omega }^2)g_\mathrm{\Omega }_{}(\tau ,\tau ^{}),$$
(B8)
where $`g_\pm (\tau ,\tau )`$ was defined in Eq. (A33). Performing the derivatives to both arguments leads to the expression
$${}_{}{}^{}G_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\tau ^{})={}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\tau ^{})+\frac{\mathrm{}}{M}\delta (\tau \tau ^{}),$$
(B9)
where we have introduced the partial function
$${}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\tau ^{})=\frac{\mathrm{}}{2M}\left[\mathrm{\Omega }_+\frac{\mathrm{\Omega }_+^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\frac{\mathrm{sinh}\mathrm{\Omega }_+(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\mathrm{\Omega }_{}\frac{\mathrm{\Omega }_{}^2\mathrm{\Omega }^2}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\frac{\mathrm{sinh}\mathrm{\Omega }_{}(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}\right]$$
(B10)
which is finite for equal times.
Applying derivatives with respect to the first respective second argument to the mixed correlation function (A31), we find:
$${}_{}{}^{}G_{xy}^{𝐱_0}(\tau ,\tau ^{})=\frac{\mathrm{}\omega }{2iM}\frac{1}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\left[\mathrm{\Omega }_+\frac{\mathrm{cosh}\mathrm{\Omega }_+(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\mathrm{\Omega }_{}\frac{\mathrm{cosh}\mathrm{\Omega }_{}(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}\right]=G_{}^{}{}_{xy}{}^{𝐱_0}(\tau ,\tau ^{})$$
(B11)
and
$${}_{}{}^{}G_{yx}^{𝐱_0}(\tau ,\tau ^{})={}_{}{}^{}G_{xy}^{𝐱_0}(\tau ,\tau ^{}).$$
(B12)
Differentiating each argument of the mixed Green function results in
$${}_{}{}^{}G_{}^{}{}_{xy}{}^{𝐱_0}(\tau ,\tau ^{})=\frac{i\mathrm{}\omega }{2M}\frac{1}{\mathrm{\Omega }_+^2\mathrm{\Omega }_{}^2}\left[\mathrm{\Theta }(\tau \tau ^{})h(\tau ,\tau ^{})\mathrm{\Theta }(\tau ^{}\tau )h(\tau ^{},\tau )\right]={}_{}{}^{}G_{}^{}{}_{yx}{}^{𝐱_0}(\tau ,\tau ^{})$$
(B13)
with
$$h(\tau ,\tau ^{})=\mathrm{\Omega }_+^2g_{\mathrm{\Omega }_+}(\tau ,\tau ^{})\mathrm{\Omega }_{}^2g_\mathrm{\Omega }_{}(\tau ,\tau ^{}).$$
(B14)
An additional property we read off from Eqs. (B7) and (B11) is ($`i,j\{x,y\}`$):
$`{}_{}{}^{}G_{ij}^{𝐱_0}(\tau ,\tau ^{})={}_{}{}^{}G_{ij}^{𝐱_0}(\tau ^{},\tau )\times \{\begin{array}{cc}1& i=j,\\ 1& ij,\end{array}`$ (B17)
$`G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,\tau ^{})=G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ^{},\tau )\times \{\begin{array}{cc}1& i=j,\\ 1& ij.\end{array}`$ (B20)
The double-sided derivatives (B9), (B10), and (B13) imply
$${}_{}{}^{}G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,\tau ^{})={}_{}{}^{}G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ^{},\tau )\times \{\begin{array}{cc}1& i=j,\\ 1& ij.\end{array}$$
(B21)
The derivatives (B7), (B10), (B11), and (B13) are periodic:
$`{}_{}{}^{}G_{ij}^{𝐱_0}(\tau ,0)={}_{}{}^{}G_{ij}^{𝐱_0}(\tau ,\mathrm{}\beta ),{}_{}{}^{}G_{ij}^{𝐱_0}(0,\tau ^{})={}_{}{}^{}G_{ij}^{𝐱_0}(\mathrm{}\beta ,\tau ^{}),`$ (B22)
$`G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,0)=G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,\mathrm{}\beta ),G_{}^{}{}_{ij}{}^{𝐱_0}(0,\tau ^{})=G_{}^{}{}_{ij}{}^{𝐱_0}(\mathrm{}\beta ,\tau ^{}),`$ (B23)
$`{}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,0)={}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\mathrm{}\beta ),{}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{ii}{}^{𝐱_0}(0,\tau ^{})={}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{ii}{}^{𝐱_0}(\mathrm{}\beta ,\tau ^{}),`$ (B24)
$`{}_{}{}^{}G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,0)={}_{}{}^{}G_{}^{}{}_{ij}{}^{𝐱_0}(\tau ,\mathrm{}\beta ),{}_{}{}^{}G_{}^{}{}_{ij}{}^{𝐱_0}(0,\tau ^{})={}_{}{}^{}G_{}^{}{}_{ij}{}^{𝐱_0}(\mathrm{}\beta ,\tau ^{}),(ij).`$ (B25)
## C Generating Functional for Position- and Momentum-Dependent Correlation Functions
With the discussion of the generating functional for position-dependent correlation functions and, in particular, the Green functions in Appendix A and their properties in Appendix B, we have layed the foundation to derive the generating functional for correlation functions depending on both, position and momentum. Following the framework presented in an earlier work , such a functional involving sources coupled to the momentum can always be reduced to one containing position-coupled sources only.
We start from the three-dimensional effective classical representation for the generating functional
$$Z_𝛀[𝐣,𝐯]=\frac{d^3x_0d^3p_0}{(2\pi \mathrm{})^3}Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]$$
(C1)
with zero frequency components $`𝐱_0=(x_0,y_0,z_0)=\mathrm{const}.`$ and $`𝐩_0=(p_{x}^{}{}_{0}{}^{},p_{y}^{}{}_{0}{}^{},p_{y}^{}{}_{0}{}^{})=\mathrm{const}.`$ of the Fourier path separated. The reduced functional is
$$Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]=(2\pi \mathrm{})^3𝒟_{}^{}{}_{}{}^{3}x𝒟^3p\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})\mathrm{exp}\left\{\frac{1}{\mathrm{}}𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱;𝐣,𝐯]\right\},$$
(C2)
where the path integral measure is that defined in Eq. (4). Extending the action (3) by source terms, considering a more general Hamilton function than (17), and introducing an additional harmonic oscillator in $`z`$-direction, the action functional in Eq. (C2) shall read
$`𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱;𝐣,𝐯]={\displaystyle _0^\mathrm{}\beta }d\tau \{`$ $`i\stackrel{~}{𝐩}(\tau )\dot{𝐱}(\tau )+{\displaystyle \frac{1}{2M}}\stackrel{~}{𝐩}^2(\tau ){\displaystyle \frac{1}{2}}\mathrm{\Omega }_1l_z(\stackrel{~}{𝐩},\stackrel{~}{𝐱})+{\displaystyle \frac{1}{8}}M\mathrm{\Omega }_2^2\left[\stackrel{~}{x}^2(\tau )+\stackrel{~}{y}^2(\tau )\right]+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_{}^2\stackrel{~}{z}^2(\tau )`$ (C4)
$`+𝐣(\tau )\stackrel{~}{𝐱}(\tau )+𝐯(\tau )\stackrel{~}{𝐩}(\tau )\}`$
with shifted positions and momenta
$$\stackrel{~}{𝐱}=𝐱(\tau )𝐱_0,\stackrel{~}{𝐩}=𝐩(\tau )𝐩_0.$$
(C5)
The orbital angular momentum $`l_z(𝐩,𝐱)`$ is defined in Eq. (19) and is used in Eq. (C4) with the shifted phase space coordinates (C5). We have introduced three different frequencies in (C4), $`𝛀=(\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{\Omega }_{})`$, where the first both components are used in regard to the oscillations in the plane perpendicular to the direction of the magnetic field which shall be considered here to point into $`z`$-direction. The last component, $`\mathrm{\Omega }_{}`$, is the frequency of a trial oscillator parallel to the field lines.
Due to the periodicity of the paths, we suppose that the sources might also be periodic:
$$𝐣(0)=𝐣(\mathrm{}\beta ),𝐯(0)=𝐯(\mathrm{}\beta ).$$
(C6)
Since we want to simplify expression (C2) such that we can use the results obtained in Appendix A, the momentum path integral is solved in the following. In a first step we reexpress the momentum $`\delta `$-function in (C2) by
$$\delta (𝐩_0\overline{𝐩(\tau )})=\frac{d^3\xi }{(2\pi \mathrm{})^3}\mathrm{exp}\left\{\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐯_0[𝐩(\tau )𝐩_0]\right\},$$
(C7)
where
$$𝐯_0(𝝃)=\frac{i}{\mathrm{}\beta }𝝃$$
(C8)
is an additional current which is coupled to the momentum and is constant in time. Defining the sum of all sources coupled to the momentum by
$$𝐕(𝝃,\tau )=𝐯(\tau )+𝐯_0(𝝃),$$
(C9)
the functional (C2) can be written as
$`Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]`$ $`=`$ $`{\displaystyle }d^3\xi {\displaystyle }𝒟_{}^{}{}_{}{}^{3}x𝒟^3p\delta (𝐱_0\overline{𝐱(\tau )})\mathrm{exp}\{{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle _0^\mathrm{}\beta }d\tau [i𝐩(\tau )\dot{𝐱}(\tau )+{\displaystyle \frac{𝐩^2(\tau )}{2M}}{\displaystyle \frac{1}{2}}\mathrm{\Omega }_1l_z(𝐩(\tau ),\stackrel{~}{𝐱}(\tau ))`$ (C11)
$`+{\displaystyle \frac{1}{8}}M\mathrm{\Omega }_2^2\{\stackrel{~}{x}^2(\tau )+\stackrel{~}{y}^2(\tau )\}+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_{}^2\stackrel{~}{z}^2(\tau )+𝐣(\tau )\stackrel{~}{𝐱}(\tau )+𝐕(𝝃,\tau )𝐩(\tau )]\},`$
where we have used the translation invariance $`\stackrel{~}{𝐩}𝐩`$ of the path integral. To solve the momentum path integral, it is useful to express it in its discretized form. Performing quadratic completions such that the momentum path integral separates into an infinite product of simple Gaussian integrals which are easily calculated, the remaining functional is reduced to the configuration space path integral
$`Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]={\displaystyle d^3\xi \mathrm{exp}\left[\frac{M}{2\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐕^2(𝝃,\tau )\right]𝒟^3x\delta (𝐱_0\overline{𝐱(\tau )})\mathrm{exp}\left\{\frac{1}{\mathrm{}}𝒜_𝛀^{𝐩_0,𝐱_0}[𝐱;𝐣,𝐕]\right\}}`$ (C12)
with the measure (10) for $`D=3`$. The action functional is
$`𝒜_𝛀^{𝐩_0,𝐱_0}[𝐱;𝐣,𝐕]`$ $`=`$ $`{\displaystyle _0^\mathrm{}\beta }d\tau [{\displaystyle \frac{M}{2}}\dot{𝐱}^2(\tau )+{\displaystyle \frac{1}{2}}iM\mathrm{\Omega }_1\{\dot{x}(\tau )\stackrel{~}{y}(\tau )\dot{y}(\tau )\stackrel{~}{x}(\tau )\}+{\displaystyle \frac{1}{8}}M(\mathrm{\Omega }_2^2\mathrm{\Omega }_1^2)\{\stackrel{~}{x}^2(\tau )+\stackrel{~}{y}^2(\tau )\}`$ (C15)
$`+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_{}^2\stackrel{~}{z}^2(\tau )+\stackrel{~}{x}(\tau )\left[j_x(\tau )+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_1V_y(𝝃,\tau )\right]`$
$`+\stackrel{~}{y}(\tau )[j_y(\tau ){\displaystyle \frac{1}{2}}M\mathrm{\Omega }_1V_x(𝝃,\tau )]+\stackrel{~}{z}(\tau )j_z(\tau )]{\displaystyle \frac{iM}{\mathrm{}}}{\displaystyle }_0^\mathrm{}\beta d\tau \dot{𝐱}(\tau )𝐕(𝝃,\tau ),`$
where the last term simplifies by the following consideration. A partial integration of this term yields
$$_0^\mathrm{}\beta 𝑑\tau \dot{𝐱}(\tau )𝐕(𝝃,\tau )=_0^\mathrm{}\beta 𝑑\tau (𝐱(\tau )𝐱_0)\dot{𝐕}(𝝃,\tau ).$$
(C16)
The surface term vanishes as a consequence of the periodicity of the path and the source. This periodicity is also the reason why we could shift $`𝐱(\tau )`$ by the constant $`𝐱_0`$ on the right-hand side of Eq. (C16). Obviously, the importance of this expression lies in the coupling of the time derivative of $`𝐕(𝝃,\tau )`$ to the path $`𝐱(\tau )`$. Thus, $`\dot{𝐕}(𝝃,\tau )`$ can be handled like a $`𝐣(\tau )`$-current and the action (C15) can be written as
$$𝒜_𝛀^{𝐩_0,𝐱_0}[𝐱;𝐣,𝐕]=𝒜_𝛀^{𝐩_0,𝐱_0}[𝐱;𝐉,0]=𝒜_𝛀^{𝐩_0,𝐱_0}[𝐱;0,0]\frac{1}{\mathrm{}}_0^\mathrm{}\beta 𝑑\tau \stackrel{~}{𝐱}(\tau )𝐉(𝝃,\tau )$$
(C17)
with the new current vector $`𝐉(𝝃,\tau )`$ which has the components
$`J_x(𝝃,\tau )`$ $`=`$ $`j_x(\tau )+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_1V_y(𝝃,\tau )iM\dot{V}_x(𝝃,\tau ),`$ (C18)
$`J_y(𝝃,\tau )`$ $`=`$ $`j_y(\tau ){\displaystyle \frac{1}{2}}M\mathrm{\Omega }_1V_x(𝝃,\tau )iM\dot{V}_y(𝝃,\tau ),`$ (C19)
$`J_z(𝝃,\tau )`$ $`=`$ $`j_z(\tau ){\displaystyle \frac{1}{2}}M\mathrm{\Omega }_{}V_z(𝝃,\tau ).`$ (C20)
and couples to the path $`𝐱(\tau )`$ only. With the expression (C12) for the generating functional and the action (C17), we have derived a representation similar to Eq. (A1) with the action (A2), extended by an additional oscillator in $`z`$-direction. We identify
$$\omega \mathrm{\Omega }_1,\mathrm{\Omega }\frac{1}{4}\left(\mathrm{\Omega }_2^2\mathrm{\Omega }_1^2\right),j_xJ_x,j_yJ_y.$$
(C21)
Thus the auxiliary frequencies $`\mathrm{\Omega }_\pm `$ (A19) become
$$\mathrm{\Omega }_\pm (\mathrm{\Omega }_1,\mathrm{\Omega }_2)=\frac{1}{2}|\mathrm{\Omega }_1\pm \mathrm{\Omega }_2|.$$
(C22)
Inserting the substitutions (C21) into the solution (A9) for the generating functional in two dimensions and performing the usual calculation for a harmonic oscillator with external source \[16, Chaps. 3,5\] in $`z`$-direction, we obtain an intermediate result for the generating functional in three dimensions (C2):
$$Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]=\lambda _{\mathrm{th}}^3Z_𝛀^{𝐩_0,𝐱_0}d^3\xi \mathrm{exp}\left\{\frac{M}{2\mathrm{}}_0^\mathrm{}\beta 𝑑\tau 𝐕^2(𝝃,\tau )\right\}\mathrm{exp}\left\{\frac{1}{2\mathrm{}^2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}𝐉(𝝃,\tau )𝐆^{𝐱_0}(\tau ,\tau ^{})𝐉(𝝃,\tau ^{})\right\}.$$
(C23)
The partition function follows from Eqs. (A20) and (A27)
$$Z_𝛀^{𝐩_0,𝐱_0}=Z_𝛀^{𝐩_0,𝐱_0}[0,0]=\frac{\mathrm{}\beta \mathrm{\Omega }_+/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\frac{\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}\frac{\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}$$
(C24)
and $`𝐆^{𝐱_0}(\tau ,\tau ^{})`$ is the $`3\times 3`$-matrix of Green functions
$$𝐆^{𝐱_0}(\tau ,\tau ^{})=\left(\begin{array}{ccc}G_{xx}^{𝐱_0}(\tau ,\tau ^{})& G_{xy}^{𝐱_0}(\tau ,\tau ^{})& 0\\ G_{yx}^{𝐱_0}(\tau ,\tau ^{})& G_{yy}^{𝐱_0}(\tau ,\tau ^{})& 0\\ 0& 0& G_{zz}^{𝐱_0}(\tau ,\tau ^{})\end{array}\right).$$
(C25)
Except $`G_{zz}^{𝐱_0}(\tau ,\tau ^{})`$, the Green functions are given by the expressions in Eqs. (A24) and (A31) with frequencies (C22). The Green function of the pure harmonic oscillator in $`z`$-direction
$$G_{zz}^{𝐱_0}(\tau ,\tau ^{})=\frac{1}{M\beta \mathrm{\Omega }_{}^2}\left(\frac{\mathrm{}\beta \mathrm{\Omega }_{}}{2}\frac{\mathrm{cosh}\mathrm{\Omega }_{}(|\tau \tau ^{}|\mathrm{}\beta /2)}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}1\right)$$
(C26)
follows directly from the limit (A26). Since the current $`𝐉`$ (C18) still depends on time derivatives of $`𝐕`$, we have to perform some partial integrations in the functional (C23). This is a very extensive but straightforward work and thus we only present an instructive example. For that we apply the properties and the time derivatives of the Green functions which we presented in Appendix B. Consider the integral
$$I=\frac{M^2}{2\mathrm{}^2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}\dot{V}_i(𝝃,\tau )G_{ii}^{𝐱_0}(\tau ,\tau ^{})\dot{V}_i(𝝃,\tau ^{})$$
(C27)
occuring in the second exponential of Eq. (C23) with $`i\{x,y,z\}`$. A partial integration in the $`\tau ^{}`$-integral leads to
$`I`$ $`=`$ $`{\displaystyle \frac{M^2}{2\mathrm{}^2}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau \dot{V}_i(𝝃,\tau )\left(G_{ii}^{𝐱_0}(\tau ,\tau ^{})V_i(𝝃,\tau ^{})|_{\tau ^{}=0}^{\tau ^{}=\mathrm{}\beta }{\displaystyle _0^\mathrm{}\beta }𝑑\tau ^{}{\displaystyle \frac{G_{ii}^{𝐱_0}(\tau ,\tau ^{})}{\tau ^{}}}V_i(𝝃,\tau ^{})\right)`$ (C28)
$`=`$ $`{\displaystyle \frac{M^2}{2\mathrm{}^2}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau {\displaystyle _0^\mathrm{}\beta }𝑑\tau ^{}\dot{V}_i(𝝃,\tau )G_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\tau ^{})V_i(𝝃,\tau ^{}).`$ (C29)
The surface term in the first line vanishes as a consequence of the periodicity of the current (C6) and the Green function (B5). A second partial integration, now in the $`\tau `$-integral, results in
$`I`$ $`=`$ $`{\displaystyle \frac{M^2}{2\mathrm{}^2}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau {\displaystyle _0^\mathrm{}\beta }𝑑\tau ^{}V_i(𝝃,\tau ){}_{}{}^{}G_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\tau ^{})V_i(𝝃,\tau ^{})`$ (C30)
$`=`$ $`{\displaystyle \frac{M^2}{2\mathrm{}^2}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau {\displaystyle _0^\mathrm{}\beta }𝑑\tau ^{}V_i(𝝃,\tau ){}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{ii}{}^{𝐱_0}(\tau ,\tau ^{})V_i(𝝃,\tau ^{}){\displaystyle \frac{M}{2\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑\tau V_i^2(𝝃,\tau ).`$ (C31)
Here we have applied the periodicity property of the right-hand derivative of the Green function (B23), leading to a vanishing surface term in this case, too. In the second line, we have used the decomposition (B9) of the double-sided differentiated Green function. Note that the last term just cancels the appropriate term in the first exponential of the right-hand side of Eq. (C23). Eventually, after performing all such partial integrations, we reexpress Eq. (C23) by
$$Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]=\lambda _{\mathrm{th}}^3Z_𝛀^{𝐩_0,𝐱_0}d^3\xi \mathrm{exp}\left\{\frac{1}{2\mathrm{}^2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}\stackrel{~}{𝐬}(𝝃,\tau )𝐇^{𝐱_0}(\tau ,\tau ^{})\stackrel{~}{𝐬}(𝝃,\tau ^{})\right\}$$
(C32)
with six-dimensional sources
$$\stackrel{~}{𝐬}(𝝃,\tau )=(𝐣(\tau ),𝐕(𝝃,\tau )).$$
(C33)
and the $`6\times 6`$-matrix $`𝐇^{𝐱_0}(\tau ,\tau ^{})`$ which has no significance as long as we have not done the $`𝝃`$-integration. We explicitly insert the decomposition (C9) into expression (C33) of the source vector $`\stackrel{~}{𝐬}`$. Since $`𝐯_0(𝝃)`$ from Eq. (C8) is constant in time, some temporal integrals in the exponential of Eq. (C32) can be calculated and we obtain
$`Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]`$ $`=`$ $`\lambda _{\mathrm{th}}^3Z_𝛀^{𝐩_0,𝐱_0}\mathrm{exp}\left\{{\displaystyle \frac{1}{2\mathrm{}^2}}{\displaystyle _0^\mathrm{}\beta }𝑑\tau {\displaystyle _0^\mathrm{}\beta }𝑑\tau ^{}𝐬(\tau )𝐇^{𝐱_0}(\tau ,\tau ^{})𝐬(\tau ^{})\right\}`$ (C35)
$`\times {\displaystyle }d^3\xi \mathrm{exp}\{{\displaystyle \frac{M}{2\mathrm{}^2\beta }}𝝃^2+i{\displaystyle \frac{M}{\mathrm{}^2\beta }}𝝃{\displaystyle _0^\mathrm{}\beta }d\tau 𝐯(\tau )\}`$
with the new $`6`$-vector
$$𝐬(\tau )=(𝐣(\tau ),𝐯(\tau ))$$
(C36)
consisting of the original sources $`𝐣`$ and $`𝐯`$ only. The Gaussian $`\xi `$-integral in Eq. (C35) can be easily solved and the terms appearing from quadratic completion modify the above matrix $`𝐇^{𝐱_0}(\tau ,\tau ^{})`$. The final result for the generating functional of all position and momentum dependent correlations is given by
$$Z_𝛀^{𝐩_0,𝐱_0}[𝐣,𝐯]=Z_𝛀^{𝐩_0,𝐱_0}\mathrm{exp}\left\{\frac{1}{2\mathrm{}^2}_0^\mathrm{}\beta 𝑑\tau _0^\mathrm{}\beta 𝑑\tau ^{}𝐬(\tau )𝐆^{𝐩_0,𝐱_0}(\tau ,\tau ^{})𝐬(\tau ^{})\right\}.$$
(C37)
The complete $`6\times 6`$-matrix $`𝐆^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ contains all possible Green functions describing position-position, position-momentum, and momentum-momentum correlations. As a consequence of separating the fluctuations into those perpendicular and parallel to the direction of the magnetic field, all correlations between $`x,y`$ on the one and $`z`$ on the other hand vanish as well as those for the appropriate momenta. The symmetries for the Green functions and their derivatives were investigated in detail in Appendix B and lead to a further reduction of the number of significant matrix elements. It turns out that only 9 elements are independent of each other. Therefore we can write the matrix
$$𝐆^{𝐱_0,𝐩_0}(\tau ,\tau ^{})=\left(\begin{array}{cccccc}G_{xx}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& G_{xy}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& 0& G_{xp_x}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& G_{xp_y}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& 0\\ G_{xy}^{𝐩_0,𝐱_0}(\tau ^{},\tau )& G_{xx}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& 0& G_{xp_y}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& G_{xp_x}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& 0\\ 0& 0& G_{zz}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& 0& 0& G_{zp_z}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})\\ G_{xp_x}^{𝐩_0,𝐱_0}(\tau ^{},\tau )& G_{xp_y}^{𝐩_0,𝐱_0}(\tau ^{},\tau )& 0& G_{p_xp_x}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& G_{p_xp_y}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& 0\\ G_{xp_y}^{𝐩_0,𝐱_0}(\tau ^{},\tau )& G_{xp_x}^{𝐩_0,𝐱_0}(\tau ^{},\tau )& 0& G_{p_xp_y}^{𝐩_0,𝐱_0}(\tau ^{},\tau )& G_{p_xp_x}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})& 0\\ 0& 0& G_{zp_z}^{𝐩_0,𝐱_0}(\tau ^{},\tau )& 0& 0& G_{p_zp_z}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})\end{array}\right).$$
(C38)
The matrix decays into four $`3\times 3`$-blocks, each of the which describing another type of correlation: the upper left position-position, the upper right position-momentum (as well as the lower left one), and the lower right momentum-momentum correlations. The different elements of the matrix are
$`G_{xx}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{x}(\tau )\stackrel{~}{x}(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=G_{xx}^{𝐱_0}(\tau ,\tau ^{}),`$ (C39)
$`G_{xy}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{x}(\tau )\stackrel{~}{y}(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=G_{xy}^{𝐱_0}(\tau ,\tau ^{}),`$ (C40)
$`G_{zz}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{z}(\tau )\stackrel{~}{z}(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=G_{zz}^{𝐱_0}(\tau ,\tau ^{}),`$ (C41)
$`G_{xp_x}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{x}(\tau )\stackrel{~}{p}_x(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=iMG_{}^{}{}_{xx}{}^{𝐱_0}(\tau ,\tau ^{}){\displaystyle \frac{1}{2}}M\mathrm{\Omega }_1G_{xy}^{𝐱_0}(\tau ,\tau ^{}),`$ (C42)
$`G_{xp_y}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{x}(\tau )\stackrel{~}{p}_y(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=iMG_{}^{}{}_{xy}{}^{𝐱_0}(\tau ,\tau ^{})+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_1G_{xx}^{𝐱_0}(\tau ,\tau ^{}),`$ (C43)
$`G_{zp_z}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{z}(\tau )\stackrel{~}{p}_z(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=iMG_{}^{}{}_{zz}{}^{𝐱_0}(\tau ,\tau ^{}),`$ (C44)
$`G_{p_xp_x}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{p}_x(\tau )\stackrel{~}{p}_x(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=M^2{}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{xx}{}^{𝐱_0}(\tau ,\tau ^{})iM^2\mathrm{\Omega }_1{}_{}{}^{}G_{xy}^{𝐱_0}(\tau ,\tau ^{})+{\displaystyle \frac{1}{4}}M^2\mathrm{\Omega }_1^2G_{xx}^{𝐱_0}(\tau ,\tau ^{}){\displaystyle \frac{M}{\beta }},`$ (C45)
$`G_{p_xp_y}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{p}_x(\tau )\stackrel{~}{p}_y(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=iM^2{}_{}{}^{}G_{xx}^{𝐱_0}(\tau ,\tau ^{})M^2{}_{}{}^{}G_{}^{}{}_{xy}{}^{𝐱_0}(\tau ,\tau ^{})+{\displaystyle \frac{1}{4}}M^2\mathrm{\Omega }_1^2G_{xy}^{𝐱_0}(\tau ,\tau ^{}),`$ (C46)
$`G_{p_zp_z}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ $`=`$ $`\stackrel{~}{p}_z(\tau )\stackrel{~}{p}_z(\tau ^{})_𝛀^{𝐩_0,𝐱_0}=M^2{}_{}{}^{}\stackrel{~}{G}_{}^{}{}_{zz}{}^{𝐱_0}(\tau ,\tau ^{}){\displaystyle \frac{M}{\beta }},`$ (C47)
where the expectation values are defined by Eq. (65). Note that all these Green functions are invariant under time translations such that
$$G_{\mu \nu }^{𝐩_0,𝐱_0}(\tau ,\tau ^{})=G_{\mu \nu }^{𝐩_0,𝐱_0}(\tau \tau ^{})$$
(C48)
with $`\mu ,\nu \{x,y,z,p_x,p_y,p_z\}`$.
It is quite instructive to prove that all these Green functions can be decomposed into a quantum statistical and a classical part as we did it in Eq. (A24). Since we know that the classical correlation functions do not depend on the euclidean time, all derivative terms in Eqs. (C39)–(C47) do not contribute a classical term. We can write each Green function
$$G_{\mu \nu }^{𝐩_0,𝐱_0}(\tau ,\tau ^{})=G_{\mu \nu }^{\mathrm{qm}}(\tau ,\tau ^{})G_{\mu \nu }^{\mathrm{cl}},$$
(C49)
This relation has been already checked for Eqs. (C39)-(C41) in Appendix A. The classical contribution is zero in Eqs. (C42), (C44), and (C46) following from the absence of classical terms in derivatives of the Green functions and mixed correlations like (A34). It seems surprising that the correlation (C43) contains a classical term while (C42) possesses none. This is, however, a consequence of the cross product of the orbital angular momentum appearing in the action (C4) and the explicit classical calculation entails
$$G_{xp_x}^{\mathrm{cl}}=xp_x^{\mathrm{cl}}=0,G_{xp_y}^{\mathrm{cl}}=xp_y^{\mathrm{cl}}=\frac{2}{\beta }\frac{\mathrm{\Omega }_1}{\mathrm{\Omega }_2^2\mathrm{\Omega }_1^2},$$
(C50)
where the latter is the subtracted classical term in Eq. (A24) when considering the first two substitutions in (C21). In Eq. (C47), the second term is obviously the classical one since
$$G_{p_zp_z}^{\mathrm{cl}}=p_zp_z^{\mathrm{cl}}=\frac{M}{\beta }.$$
(C51)
The extraction of the classical terms
$$G_{p_xp_x}^{\mathrm{cl}}=p_xp_x^{\mathrm{cl}}=\frac{M}{\beta }\left(1+\frac{\mathrm{\Omega }_1^2}{\mathrm{\Omega }_2^2\mathrm{\Omega }_1^2}\right)$$
(C52)
in the case of the Green function $`G_{p_xp_x}^{𝐩_0,𝐱_0}(\tau ,\tau ^{})`$ requires the consideration of the last two terms in Eq. (C45). Thus we have shown that the decomposition (C49) holds for each of the Green functions (C39)–(C47). Note the necessity of subtracting the classical terms since they all diverge in the classical limit of high temperatures ($`\beta 0`$).
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# Monte Carlo Predictions of Far-Infrared Emission from Spiral Galaxies
## 1 Introduction
Dark lanes in edge-on spirals clearly reveal the presence of dust in galactic disks. Depending on the dust amount and distribution, the extinction of starlight may affect to different degrees our knowledge of the objects harboring dust and of the distant universe in the background. However, deriving the amount of dust from the extinction features in spiral galaxies is not easy, because of the uncertainty in the relative distributions of dust and stars. Indeed, realistic models are needed to avoid equivocal results (Disney et al. 1989). Recently, Xilouris et al. (1999) have analised a sample of seven edge-on spirals by fitting optical and NIR images with a complex radiative transfer model inclusive of scattering. They derived a small opacity, with their exponential dust disks having a mean face-on optical depth in the B-band $`\tau _\mathrm{B}=0.8`$. Unfortunately, the method can only be applied successfully to a limited numer of cases, i.e. edge-on galaxies, where it is possible to analyze the vertical surface-brightness distribution and the dust extinction is maximised by the inclination.
Further constrains on the structure and opacity of dust disks can be obtained by studying dust emission. Assuming that all the stellar radiation absorbed by dust is reemitted in the MIR and (mainly) in the FIR, the amount of extinction in a galaxy can be simply derived by comparing the total dust emission with the stellar one (the energy balance technique). The extinction is then used, together with a radiative transfer model, to study the dust disk properties (Evans 1992; Xu & Buat 1995; Trewhella 1998a). However, the dust emission is treated as a bulk and generally only the disk optical depth is derived. The parameters for the geometrical distributions of dust and stars have to be assumed a priori.
Three-dimensional models of FIR emission from spiral galaxies can help to retrieve information about the relative star-dust geometry from the SED of dust emission and the spatial distribution of the radiation, when available. However, an additional factor is introduced in such models. Not only it is necessary to assume a dust distribution, but also a knowledge of the variation of the ISRF heating throughout the disk is needed. The ISRF could be in principle derived from the observed profiles of stellar emission, but a complicate de-projection is needed involving assumption on the dust distribution. This approach is adopted by Walterbos & Greenawalt (1996). They compute fluxes and surface brightness gradients at 60$`\mu `$m and 100$`\mu `$m for a sample of 20 spirals. The diffuse ISRF is derived from B band profiles and the dust distribution from the HI column density. The problem of defining an ISRF can be avoided by assuming a temperature gradient for dust grains. Davies et al. (1997) modelled the 140$`\mu `$m and 240$`\mu `$m Galactic emission observed by DIRBE, by adopting a gradient that matches the longitudinal variation of temperature observed in the data themselves.
Otherwise, the dust heating by the ISRF can be derived in a self-consistent way from a radiative transfer model of the stellar radiation in the dusty environment. Sauty et al. (1998) produced a Monte Carlo radiative transfer model for the clumpy ISM of the spiral galaxy NGC 6946, although a self-consistent ISRF is derived only for the UV radiation, relying on optical profiles for radiation at longer wavelelngth. Dust heating in a clumpy galactic ISM is also included in the photometric evolution model of Silva et al. (1998). They carry out the radiative transfer in a proper way for stellar emission embedded in the clumps and use an approximation for the smooth medium. In the emission model of Wolf et al. (1998) the radiative transfer is computed with the Monte Carlo technique for the whole spectral range, both in the smooth medium and clumps. However, the adopted geometry is typical of star-formation environment.
In this paper, we present a self-consistent galactic model for the FIR emission. The model is based on the Bianchi, Ferrara & Giovanardi (1996, hereafter, BFG) Monte Carlo code for the radiative transfer (complete with scattering) in dusty spiral galaxies. A map of the energy absorbed by dust is derived from the output of the radiative transfer for stellar radiation at several wavelength. The grain temperature along the dust distribution is therefore computed from the radiative transfer itself. Luminosity, spectrum and surface brightness distribution of the FIR radiation can be retrieved, for any wavelength. The main advantage with respect to other models of FIR emission lies in the relative simplicity. We have limited the number of parameters involved to a minimum, adopting, whenever possible, empyrical data. This will allow an isolation of the effect of the dust distribution on the FIR emission. We will try to explain the characteristic of FIR emission with diffuse dust heated by a diffuse ISRF. Clumping of dust and embedded stellar emission (Bianchi et al. 2000b) are not yet included in this work.
Recent FIR observations suggest large scalelengths for dust emission. In a sample of seven spirals resolved by the ISOPHOT instrument aboard the ISO satellite (Kessler et al. 1996; Lemke et al. 1996), Alton et al. (1998c) found that the 200 $`\mu `$m scalelengths are 70% larger than those in the B-band. The large ratio may be a result of a dust disk larger than the stellar distribution. However, the result may also be affected by the transient effects of the ISOPHOT detectors. In fact, the P32 mapping mode with which the Alton et al. (1998c) images were obtained is not yet scientifically validated (Klaas & Richards 2000). We will use the model presented here to test if indeed the observations are compatible with large disks. Other works seems to support the same hypothesis. Davies et al. (1997) could model the Galactic FIR surface brightness only with a dust disk with radial and vertical scalelengths larger than the stellar by a factor 1.5 and 2.0, respectively. Nelson, Zaritsky & Cutri (1998) combined 100 $`\mu `$m IRAS images of galaxies with similar angular size and concluded that dust is present well beyond the observed stellar disk. Xilouris et al. (1999) derived intrinsic dust scalelengths larger than the stellar one, by a mean factor 1.4 in the seven edge-on spiral sample. If confirmed, such disks may have a large impact on the observations of the distant universe, because of their larger cross section to radiation from backgroud objects.
The paper is structured as follows: Sect. 2 describes the radiative transfer code, with a particular attention to the chosen dust and stellar distributions and to the modification with respect to the original BFG code. Sect. 3 and 4 show how the map of absorbed energy derived from the radiative transfer model can be converted into a map of the dust temperature distribution and hence, of FIR emission. We then show an application of the model to the spiral galaxy NGC 6946 with the intent of analysing the behaviour of our model and test the influence of the various parameters on the FIR emission. The observed properties of NGC 6946 used in the model are described in Sect. 5. The results are presented in Sect. 6 and discussed in Sect. 7. A summary is given in the last section.
## 2 The radiative transfer model
The original BFG radiative transfer code carries out the radiative transfer for typical galactic geometries. Simulation of optical surface brightness, as well as light polarisation, are produced. Dust properties are computed from the Draine & Lee (1984) dust model, using Mie’s theory for spherical grains. In this paper we use a simplified version of the code: the treatment of polarization has been omitted; empirical dust properties and scattering phase functions are used. The same approach has recently been adopted in Ferrara et al. (1999) and Bianchi et al. (2000b).
In this Section, we describe the stellar and dust distributions adopted in the galaxy model, and the chosen dust extinction and scattering properties. A brief description is then given of the radiative transfer code.
### 2.1 The stellar disk
As in BFG, the new version of the code allows the use of both disks and spheroidal bulges to describe the spatial distribution of stars (Bianchi 1999). However, to limit the number of free parameters in the modelling, we adopt here a single stellar disk. The chosen disk is exponential along both the radial and vertical directions. The luminosity density distribution is thus described by
$$\rho =\rho _0\text{exp}(r/\alpha _{}z/\beta _{}).$$
(1)
where r and z are the galactocentric distance and the height above the galactic plane, respectively, and $`\alpha _{}`$ and $`\beta _{}`$ the relative scalelengths. While there is a consensus for the radial exponential behaviour (de Vaucouleurs 1959; Freeman 1970), a number of expressions for the vertical distribution have been used in the literature. van der Kruit & Searle (1981) use a sech<sup>2</sup>, the solution for a self gravitating isothermal sheet, while van der Kruit (1988) propose a sech, consistent with measures of stellar velocity dispersion out of the Galactic plane. The exponential adopted here (Wainscoat, Freeman & Hyland 1989) has the advantage of mathematical simplicity, but has no firm physical justification. Analysing the vertical profiles of a sample of 24 edge-on galaxies in the relatively dust-free K-band, de Grijs, Peletier & van der Kruit (1997) find that a distribution with a peak intermediate between the exp and the sech fits better the central peak (see also de Grijs & van der Kruit 1996). Since a small inclination from the pure edge-on case can produce a less sharp profile, these results are consistent with an exponential distribution.
The radiative transfer model of the galaxy (Sect. 2.4) is run for several wavelength bands. The radial scalelength of the stellar distribution, $`\alpha _{}`$, is assumed to be the same for all of these bands. Therefore, any observed colour gradient along the disk of a spiral is interpreted as due to differential extinction with $`\lambda `$, rather than being a reflection of different stellar components. The wavelength dependence of the intrinsic stellar radial scalelength is uncertain, because a knowledge of the dust distribution is needed for its determination. Peletier et al. (1995) analyzed the variation of the ratio of B and K-band scalelengths with inclination in a sample of 37 Sb-Sc galaxies and concluded that the observed gradient is due to extinction. On the other hand, de Jong (1996) compared colour-colour plots of a sample of 86 face-on galaxies with a Monte Carlo radiative transfer model and finds that reasonable dust models cannot be responsible for the observed colour gradients. However, such a conclusion does not seem to include models with dust disks more extended that the stellar, that is one of the cases studied here. The fits of edge-on galaxies of Xilouris et al. (1999) show a slight variation of the intrinsic scalelength with wavelength.
Observations on the Galaxy show that different stellar populations have different vertical scalelengths. For the sake of simplicity, a single mean value for the vertical scalelength has been adopted here, for each wavelength. Wainscoat et al. (1992) provide a table of vertical scalelengths, B, V, J, H, K-band absolute magnitudes and local number density for the main Galactic stellar sources. A mean value has been computed averaging over the total disk luminosity integrated from the B to the K band. Assuming $`\alpha _{}=3`$ kpc (Kent, Dame & Fazio 1991; Fux & Martinet 1994), we derive $`\alpha _{}/\beta _{}=14.4`$. Similar values for $`\alpha _{}/\beta _{}`$ can be found in the literature (van der Kruit & Searle 1982; de Grijs & van der Kruit 1996; Xilouris et al. 1999).
The stellar disk is truncated at 6$`\alpha _{}`$. Fits of star counts for faint Galactic sources suggest a similar truncation (Wainscoat et al. 1992; Robin, Crézé & Mohan 1992; Ruphy et al. 1996). The truncation along the vertical direction is at 6$`\beta _{}`$.
### 2.2 The dust disk
The parameters for the dust disk are far more uncertain than those for the stellar disk. Usually the same functional form is used as for the stellar distribution, with independent scalelengths (Kylafis & Bahcall 1987; Byun et al. 1994; Bianchi et al. 1996; Xilouris et al. 1999).
A very good correlation is usually found between dust and gas in spirals, especially with the molecular component which is dominant in the inner galaxy (Devereux & Young 1990; Xu et al. 1997; Bianchi et al. 1998; Alton et al. 1998a; Bianchi et al. 2000a). It is therefore reasonable to use the gas distribution as a tracer of dust.
In luminous, face-on, late-type spirals, H<sub>2</sub> peaks in the centre and falls off monotonically with increasing galactocentric distance, while the HI distribution shows a central depression and a nearly constant surface density across the rest of the optical disk (Young & Scoville 1991). Some early type galaxies show the same behaviour, although a good fraction of them presents a central depression and a flatter ring distribution for the molecular gas. Indeed, flat or ring-like dust distributions have been derived from FIR and sub-mm observations (Xu & Helou 1996b; Sodroski et al. 1997; Bianchi et al. 1998).
Nevertheless, the gas distribution of most galaxies is dominated by a centrally peaked molecular gas component. This is the case of the late-type spiral NGC 6946 (Tacconi & Young 1986), whose dust distribution is the main concern of this paper (Sect. 5). Therefore, in this paper dust is assumed to be distributed in a smooth radial and vertical exponential disk, similar to the stellar one (Sect. 2.1). The number density of dust grains can be written as
$$n(r,z)=n_0\mathrm{exp}(r/\alpha _dz/\beta _d),$$
(2)
with $`n_0`$ the central number density. The dust scalelengths $`\alpha _d`$ and $`\beta _d`$ can be selected independently from the analogous stellar parameters. It is usually assumed in radiative transfer models that $`\alpha _d\alpha _{}`$ and $`\beta _d0.5\beta _{}`$ (Byun et al. 1994; Bianchi et al. 1996; Davies et al. 1997, and references therein). We will refer to this as the standard model. The choice of the vertical scalelengths is mainly dictated by the impossibility of simulating the extinction lanes in edge-on galaxies with a dust distribution thicker than the stars (Xilouris et al. 1999). The assumption of a similar radial scalelength for stars and dust has no firm justification. As already said in the introduction, recent analysis of extinction of starlight and FIR emission suggests more extended dust distributions.
The central number density of dust $`n_0`$ is normalised from the chosen optical depth of the dust disk. For a central face-on optical depth in the V-band, $`\tau _V`$, $`n_0`$ is given by
$$n_0=\frac{\tau _V}{2\beta _d\sigma _{\mathrm{ext}}(V)},$$
(3)
where $`\sigma _{\mathrm{ext}}(V)`$ is the extinction cross section. The optical depth at any wavelength can then be computed by integrating the absorption coefficient
$$k_\lambda (r,z)=n(r,z)\sigma _{\mathrm{ext}}(\lambda )$$
(4)
along any path through the dust distribution. The integral involves knowledge of the ratio $`\sigma _{\mathrm{ext}}(\lambda )/\sigma _{\mathrm{ext}}(V)`$, which is found from the assumed extinction law. As for the stellar disk, the dust disk has been truncated at 6 scalelengths both in the vertical and radial direction.
In section 6 various values of $`\alpha _d/\alpha _{}`$, $`\beta _d/\beta _{}`$ and $`\tau _V`$ will be explored.
### 2.3 Dust extinction
The extinction law adopted here is that typical of diffuse Galactic dust, with a strong bump at 2175Å. Regions of high star formation and starburst galaxies usually have a weaker bump and a steeper far-UV rise (Whittet 1992; Calzetti, Kinney & Storchi-Bergmann 1994; Gordon, Calzetti & Witt 1997). Extinction laws in the optical show smaller differences. From radiative transfer models of seven edge-on galaxies, Xilouris et al. (1999) derive an extinction law similar to the Galactic one longward of the U-band.
The extinction law used is tabulated in Tab. (1), for 17 bands in the spectral range of stellar emission. The bands from UV1 to K have essentially the same spectral coverage as the homonymous described by Gordon et al. (1997), from which the extinction law is taken. Two bands (namely EUV and LMN) have been added to extend the spectral coverage in the ultraviolet up to the ionization limit and in the near infrared. The extinction law for these two bands has been taken from Whittet (1992) and Rieke & Lebofsky (1985), respectively.
We have used the Henyey & Greenstein (1941) phase function for the scattering. This is given by
$$\varphi (\theta )=\frac{1}{2}\frac{1g^2}{(1+g^22g\mathrm{cos}\theta )^{3/2}},$$
(5)
with $`\theta `$ the angle of scattering. The asymmetry parameter $`g`$, as well as the dust albedo $`\omega `$, are derived from models of dust scattering by reflection nebulae in the Milky Way. The adopted $`g`$ and $`\omega `$ are shown in Tab. (1). Again, the values for the bands from UV1 to K are from Gordon et al. (1997). For the EUV band we have used Witt et al. (1993) data at 1000Å while the values for the LMN band have been assumed equal to those in the K band.
### 2.4 The Monte Carlo code
We give here a brief description of the radiative transfer code, referring the interested reader to the BFG paper for more details.
Adhering to the Monte Carlo method, the code follows the life of each energy unit (a *photon*) through scattering and absorption, until the radiation is able to escape the dusty medium. The code is monochromatic, since the optical properties of dust must be specified for a particular wavelength. In principle, the geometrical distributions of stars may depend on $`\lambda `$ too. For a given star-dust geometry and central face-on optical depth $`\tau _V`$, a radiative transfer simulation is produced for each of the 17 bands of Sect. 2.3. Typically, a simulation consists of 10<sup>7</sup> photons.
The main steps of the computation scheme are:
the position of a photon in the 3-D space is derived according to the stellar distributions described in Sect. 2.1. Photons are emitted isotropically and with unit energy.
The absorption coefficient $`k_\lambda (r,z)`$ (Eqn. 4) is integrated from the emission position along the photon travelling direction to derive the total optical depth $`\tau _T`$ through the dust distribution. A fraction $`e^{\tau _T}`$ of all the energy travelling in that direction propagates through the dust. With the Monte Carlo method it is then possible to extract the optical depth $`\tau `$ at which the photon impinges on a dust grain. This optical depth can be computed inverting
$$_0^\tau e^\xi 𝑑\xi =R,$$
(6)
where $`R`$ is a random number in the range . If the derived $`\tau `$ is smaller than $`\tau _T`$, the photon suffer extinction, otherwise it escapes the dusty medium. This process is quite inefficient when the optical depth of the dust distribution is small, most of the photons leaving the dust distribution unaffected. To overcome this problem, the forced scattering method is used (Cashwell & Everett 1959; Witt 1977): essentially, a fraction $`e^{\tau _T}`$ of the photon energy is unextinguished and the remaining $`1e^{\tau _T}`$ is forced to scatter. When the optical depth is small ($`\tau _T<10^4`$) or the photon path is free of dust, the photon escapes the cycle. Once $`\tau `$ is known, the integral of $`k_\lambda (r,z)`$ is inverted to derive the geometrical position of the interaction between the photon and the dust grain.
a fraction of the photon energy, given by the albedo $`\omega `$, is scattered, while the remaining $`(1\omega )`$ is absorbed. The scattering polar angle $`\theta `$, i.e. the angle between the original photon path and the scattered direction, is computed using the Henyey & Greenstein (1941) scattering phase function (Eqn. 5), inverting
$$_0^\theta \varphi (\theta ^{})\mathrm{sin}\theta ^{}=R,$$
(7)
with $`R`$ a random number. The inversion of Eqn. (7) is given by the analytical formula
$$\theta =\mathrm{arccos}\left[\frac{1}{2g}\left(1+g^2\frac{(1g^2)^2}{(1+g(12R))^2}\right)\right].$$
(8)
No preferential direction perpendicular to the original photon path is assumed, as for scattering by spherical grains. The azimuthal angle is thus extracted randomly in the range $`[0,2\pi ]`$.
The amount of energy absorbed by the dust grain is stored as a function of the galactocentric distance and the height above the plane, using the two model symmetries, i.e. the symmetry around the vertical axis and the symmetry about the galactic plane, to improve the signal-to-noise. Maps of absorbed energy were not produced by the original BFG code.
the last two steps are then repeated, using the new direction of the scattered photon, the coordinates of the scattering point and the energy reduced by absorption, until the energy of the photon falls below a threshold value ($`10^4`$ of the initial energy) or until the exit conditions on $`\tau `$ are verified.
After the exit conditions are satisfied, the photon is characterised by the last scattering point, its travelling direction and its energy. To reduce the computational time, the two model symmetries are exploited to produce a total of 4 photons from each one. Photons are then classified according to the angle between the last direction and the symmetry axis. An image is produced collecting all the photons that fall in a solid angle band of width $`4\pi /15`$ (BFG) and mean polar angle equal to the chosen model inclination.
Finally, all the photons in an angle band are projected into the plane of the sky according to their point of last scattering. Maps of 201x201 pixels are produced, covering an area of 12x12 stellar radial scalelengths around the centre of the galaxy. Maps of absorbed energy cover 6 dust scalelengths in the radial direction and in the positive vertical direction in 101x101 pixels.
### 2.5 Normalisation of the radiative transfer output
Each set of 17 radiative transfer simulations is then normalised according to a chosen SED for the stellar radiation. Two kinds of normalization are available: in the first, the *input* normalization, a SED is chosen for the intrinsic, unextinguished, stellar energy emission. The SED is integrated over the bands limits and the total energy emitted in a band is used to scale each monochromatic simulation. This is suitable, for instance, to predict the FIR emission associated with a synthetic galactic SED.
The second normalisation mode, the *output* normalisation, is instead based on the observed SED of a galaxy. An observed SED is constructed from observations in several wavelength and the model is scaled to produce images that have the same SED. The intrinsic, unextinguished, SED is then inferred from the radiative transfer model. This second mode, suitable for fitting the FIR emission, is the one used in this paper (Sect. 5).
Together with the optical images, the absorbed energy maps are also normalised. Thus, the output of a model consists of a set of 17 images in the wavelength range of stellar emission and a set of 17 maps of the energy absorbed by dust from stellar radiation emitted in each band. The images are produced in units of surface brightness, while the maps of absorbed energy are in units of energy per unit time per unit volume. The absorbed energy maps can be added together to produce a map of the total energy absorbed by dust. However, for the purpose of modelling the FIR emission, a further step is needed before coadding. This is explained in the next Section.
## 3 The MIR correction
The energy absorbed from photons heats up dust grains and is re-emitted in the infrared, preferentially at $`\lambda >10\mu `$m. We do not consider here the heating of dust grains by collisions with the interstellar gas: this process is normally negligible and the dust temperature is almost entirely determined by radiative processes (Spitzer 1978).
When the energy of the absorbed photon is small compared to the internal energy of the grain, radiation is emitted at the thermal equilibrium. This is the case for larger grains, responsible for most of the FIR emission. For smaller grains, the absorption of a single high-energy photon can substantially alter the internal energy. The grain undergoes temperature fluctuations of several degrees and cools by emission of radiation, mainly in the MIR range (Whittet 1992). The models of this paper are limited to thermal FIR emission only. Therefore, it is necessary to exclude from the total absorbed energy the fraction that goes into non-equilibrium heating.
Désert, Boulanger & Puget (1990) derived an empirical dust model, from an analysis of the features in the Galactic extinction law and in the infrared emission. Three dust components were needed in the model: i) big grains ($`0.015\mu \text{m}<a<0.11\mu \text{m}`$), responsible for the FIR emission; ii) very small grains ($`0.0012\mu \text{m}<a<0.015\mu \text{m}`$) and iii) PAHs, responsible for the MIR emission at $`\lambda <80\mu \text{m}`$. From the parameters and functional forms of Désert et al. (1990), we have derived the mean absorption cross-section of the model, $`\sigma _{\mathrm{abs}}(\lambda )`$, and for each of its three dust components<sup>1</sup><sup>1</sup>1Note that $`\sigma _{\mathrm{ext}}=\sigma _{\mathrm{abs}}+\sigma _{\mathrm{sca}}`$. In the Désert et al. (1990) model, very small grains and PAHs are pure absorbers (i.e. $`\sigma _{\mathrm{sca}}=0`$).. Of the light impinging on the grain mixture, a fraction of energy proportional to $`\sigma _{\mathrm{abs}}(\lambda )`$ is absorbed; therefore, the contribution of very small grains and PAHs to the absorption is given by
$$F_{\mathrm{MIR}}(\lambda )=\frac{\sigma _{\mathrm{abs}}^{\mathrm{VSG}}(\lambda )+\sigma _{\mathrm{abs}}^{\mathrm{PAHs}}(\lambda )}{\sigma _{\mathrm{abs}}(\lambda )}.$$
(9)
Values of $`F_{\mathrm{MIR}}(\lambda )`$, the MIR correction, are given in Tab. (1) for each of the model bands and plotted in Fig. 1 as a function of $`1/\lambda `$. $`F_{\mathrm{MIR}}(\lambda )`$ has a local maximum in the position of the 2175Å bump of the extinction curve: absorption by very small grains in the Désert et al. (1990) model is responsible for this feature. The rise in the Far-UV is due to PAHs.
After multiplying by $`(1F_{\mathrm{MIR}}(\lambda )`$), the 17 absorbed energy maps store the energy that is absorbed by large grains only. The final coadded map, $`W_{\mathrm{abs}}(r,z)`$, contains the energy absorbed by dust per unit time and unit volume, that goes into thermal equilibrium emission only. The derivation of the dust temperature and FIR emission from $`W_{\mathrm{abs}}(r,z)`$ is shown in the next Section.
## 4 The dust emission model
The power absorbed by a single dust grain can be derived by dividing $`W_{\mathrm{abs}}(r,z)`$ by the grain number density $`n(r,z)`$. Because of the conservation of energy, this is equal to the power each grain radiates. Assuming that all dust grains in our model are spherical and with the same mean radius $`a`$, the thermal equilibrium imposes that
$$\frac{W_{\mathrm{abs}}(r,z)}{n(r,z)}=4\pi a^2_0^{\mathrm{}}Q_{\mathrm{em}}(\lambda )\pi B_\lambda (T_\mathrm{d}(r,z))𝑑\lambda ,$$
(10)
where $`T_\mathrm{d}(r,z)`$ is the dust temperature distribution and $`Q_{\mathrm{em}}(\lambda )`$ the emission efficiency. By Kirchhof’s law, $`Q_{\mathrm{em}}=Q_{\mathrm{abs}}=\sigma _{\mathrm{abs}}/\pi a^2`$ (Whittet 1992).
For the exponential dust distribution of Eqn. (2) and the central number density of Eqn. (3), Eqn. (10) can be rewritten as
$`{\displaystyle \frac{\beta _dW_{\mathrm{abs}}(r,z)}{2\tau _\mathrm{V}\mathrm{exp}(r/\alpha _dz/\beta _d)}}=`$
$`\text{ }{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{Q_{\mathrm{em}}(\lambda )}{Q_{\mathrm{ext}}(V)}}\pi B_\lambda (T_\mathrm{d}(r,z))𝑑\lambda ,`$ (11)
with $`Q_{\mathrm{ext}}(V)=\sigma _{\mathrm{ext}}(V)/\pi a^2`$. The map of dust temperature $`T_\mathrm{d}(r,z)`$ is derived by inverting Eqn. (11).
For the ratio $`Q_{\mathrm{em}}(\lambda )/Q_{\mathrm{ext}}(V)`$, we use a value derived for Galactic dust (Bianchi, Davies & Alton 1999),
$$\frac{Q_{\mathrm{em}}(\lambda )}{Q_{\mathrm{ext}}(V)}=\frac{2.005}{2390}\frac{\left(\text{100 }\mu \text{m}/\lambda \right)^2}{\left[1+\left(\text{200 }\mu \text{m}/\lambda \right)^6\right]^{1/6}}.$$
(12)
Eqn.(12) has also been derived assuming that all grains have the same radius. The wavelength dependence of the ratio changes smoothly from $`\lambda ^1`$ to $`\lambda ^2`$ at 200 $`\mu `$m, as observed by Reach et al. (1995) using high S/N FIR spectra of the Galactic plane. The absolute value of the ratio has been derived using dust column density maps calibrated to Galactic extinction (Schlegel, Finkbeiner & Davis 1998).
Finally, FIR images are created integrating the emission coefficient,
$`j_\lambda (r,z)=n(r,z)\sigma _{\mathrm{abs}}(\lambda )B_\lambda (T_\mathrm{d}(r,z))`$
$`={\displaystyle \frac{\tau _\mathrm{V}}{2\beta _\mathrm{d}}}\mathrm{exp}(r/\alpha _dz/\beta _d){\displaystyle \frac{Q_{\mathrm{em}}(\lambda )}{Q_{\mathrm{ext}}(V)}}B_\lambda (T_\mathrm{d}(r,z)),`$ (13)
along a given line of sight through the dust distribution, under the assumption that dust is optically thin to FIR radiation. Using the emissivity of Eqn. (12) this is justified for any model with reasonable values of $`\tau _\mathrm{V}`$.
As for the optical images, the far infrared images are produced in units of surface brightness. FIR images have the same extent and resolution of the optical images, i.e. a region of 12x12 stellar radial scalelengths around the centre of the galaxy is mapped in 201x201 pixels.
## 5 A worked example: NGC 6946
In the following we show an application of the radiative transfer and dust emission model to the spiral galaxy NGC 6946. The final goal is to find the parameters of the dust distribution compatible with the available observations in the optical and FIR range.
NGC 6946 is a large nearby Sc galaxy (A few basic properties are presented in Table 2). The optical appearance of the galaxy is characterised by six prominent spiral arms of which the three arms originating from the NE quadrant are brighter and more developed than those in the SW (Tacconi & Young 1990). Trewhella (1998a, b) finds that the NE interarm region suffers high extinction. The region appear fainter because of the dust effect, rather than being intrinsically less luminous. Indeed, polarised light is observed in the interarm regions as well as in the spiral arms (Fendt et al. 1998).
NGC 6946 shows a centrally peaked molecular gas distribution (Tacconi & Young 1989) dominant in the inner 10$`\mathrm{}`$ over the flatter atomic gas component (Tacconi & Young 1986). The galaxy is marginally resolved in FIR observations with KAO and IRAS (Engargiola 1991; Alton et al. 1998c), with the dust emission tending to follow the spiral arms and the bright central emission. The 117$`\mathrm{}`$ resolution 200$`\mu `$m ISO image shows a morphology similar to the 100$`\mu `$m IRAS observation (Alton et al. 1998c). Higher resolution SCUBA images show a tight correlation between the dust emission at 850$`\mu `$m and the dominant gas phase (Bianchi et al. 2000a).
We now describe the observed SED for the stellar emission, used in the normalization of the radiative transfer, as shown in Sect. 2.5, and the SED of dust emission, that will be compared to the model results. The procedure adopted in the modelling are presented in Sect. 5.3.
### 5.1 The stellar Spectral Energy Distribution
We have derived the observed stellar SED from literature data. Values for the flux inside an aperture of 5$`\mathrm{}`$ (corresponding to the B-band half light radius; Engargiola 1991) are presented in Table 3. All the data have been corrected for Galactic extinction (Table 2) using the assumed extinction law.
Optical and Near Infrared data in the bands U, B, g, V, r, I, J, H, K are from Engargiola (1991). The observed NIR emission is considered as purely stellar. Using the IRAS 12 $`\mu `$m flux as a template of small grain emission and the Désert et al. (1990) model, we derived a dust contribution to the K-band emission of only 0.3%. The stellar emission at 5 $`\mu `$m has been extrapolated from the K-band, using the Rayleigh-Jeans spectral region in the synthetic galactic SEDs of Fioc & Rocca-Volmerange (1997).
The SED in the non-ionising UV is taken from Rifatto, Longo & Capaccioli (1995b). The authors derive fluxes for a large sample of galaxies, by homogenising observation from several satellites (notably IUE), balloon and rocket-borne experiments, with different apertures and sensitivities. An aperture correction depending on the morphologycal type is applied to each flux, and total magnitudes are given for three photometric bands centred at 1650Å (*short*-UV), 2500Å (*medium*-UV), 3150Å (*long*-UV). Using an appropriate aperture correction and calibrating as in Rifatto, Longo & Capaccioli (1995a), we have derived fluxes for our selected aperture. The quoted errors come mainly from the aperture correction. The flux at the Lyman limit (912Å) has been extrapolated after observing that in a $`\nu f_\nu `$ versus $`\mathrm{log}\lambda `$ plot the observed SED is flat for the *short*\- and *medium*-UV. This trend is assumed to be valid down to the ionization limit. The ionising UV is not included in the model (Sect. 7).
### 5.2 The dust Spectral Energy Distribution
The SED of dust emission is also shown in Table 3. Surface brightnesses over the half-light aperture shortward of 100$`\mu `$m have been derived from IRAS High Resolution (HiRes) images (Rice 1993; Alton et al. 1998b). Derived values have been colour corrected (Rice et al. 1988) and are consistent (within a 20% error; Alton et al. 1998c) with the analogous data provided by Engargiola (1991), derived on previous enhanced resolution IRAS images. The value at 160$`\mu `$m has been derived by Engargiola (1991) from the air-borne KAO telescope.
Data at 200$`\mu `$m are derived from P32 images taken from the ISOPHOT instrument aboard the ISO satellite (Alton et al. 1998c). The P32 mapping mode is still not scientifically validated and its photometric calibration is highly uncertain (Klaas & Richards 2000). Alton et al. (1998c) compared both the integrated flux with the value derived by Engargiola (1991) on 200$`\mu `$m KAO images and the measured background with an extrapolation from the 100$`\mu `$m value. They conclude that the calibration may overestimate the flux by about 30%. A 30% error is shown in Table 3.
Finally, lower and upper limit of the flux at 450 and 850 $`\mu `$m are derived from SCUBA images (Bianchi et al. 2000a). The lower limit correspond to the signal coming from regions 3$`\sigma `$ brighter than the sky, inside the selected aperture. The upper limit is derived assuming a 1$`\sigma `$ emission for the regions without detected signal. Sub-mm fluxes are only given for completeness, the FIR models being constrained mainly by the fluxes at 100, 160 and 200 $`\mu `$m. Emission at 12, 25 and 60 $`\mu `$m, dominated by small grains (Désert et al. 1990), is used in Sect. 7 to validate the MIR correction.
### 5.3 Modelling procedure
After choosing the star-dust geometry, a radiative transfer simulation is run for each of 17 bands of the model (Sect. 2.5) and images for the inclination of NGC 6946 are produced (Table 2). An intrinsic scalelength $`\alpha _{}=2.5`$ kpc has been used for the stellar emission, derived from a K-band image (Trewhella 1998b), where the effects of extinction are smaller. For the $`\alpha _{}/\beta _{}`$ ratio of Sect. 2.1, this leads to $`\beta _{}=`$170 pc. The simulations do not depend on the absolute values of the scalelengths, as long as they are shown as functions of scaled galactocentric and vertical distances. However, absolute values are needed if we want to derive correct values for the emitted and absorbed energies in each band.
The simulated images are then integrated inside the half light radius aperture (derived from the B-band simulation). Mean fluxes over each band of Table 1 are obtained by integrating a continuous SED derived from the data in Table 3. These fluxes are then used to scale the results of the aperture photometry on the simulation. Because of this normalisation, the stellar SED derived from the simulated images is the same for each model. The intrinsic *unextinguished* energy emitted by stars is then derived from the Monte Carlo code. The half-light radius aperture was chosen in the earlier stages of our analysis because of the availability of one more UV data point in Engargiola (1991). However, if the distribution we have adopted is a correct description of the galaxy luminosity density, the normalization is independent of the aperture dimension, being equivalent to assigning a value for $`\rho _0`$ in Eqn. (1). After the normalization, FIR images for the selected inclination are produced and the fluxes derived, as for the images of stellar emission. The simulated SED in the FIR is then compared to the observed one.
For a sample of seven spirals, Alton et al. (1998c) measured B-band and FIR scalelengths in a galactocentric distance range from 1.5$`\mathrm{}`$ to 3.5$`\mathrm{}`$ after smoothing the optical and the IRAS images to the poorest resolution of the 200$`\mu `$m ISO map (FWHM=117$`\mathrm{}`$). For NGC 6946, they found that the B band scalelength is slightly smaller than the 200$`\mu \text{m}`$ scalelength, by a factor 0.9, while it is larger than the 100$`\mu `$m scalelength, by a factor 1.8. The absolute values of the scalelengths are presented in Table 4. The large 200$`\mu \text{m}`$ scalelengths, compared to the optical and IRAS data, are a general property of the sample. The FIR simulation of this work are compared with the observations of Alton et al. (1998c). FIR images are smoothed to the ISO resolution and scalelengths are derived on the same distance range as the observations<sup>2</sup><sup>2</sup>2As described in Sect. 2.4, a region of 12$`\alpha _{}`$ is covered by 201 pixels, thus giving a pixel size of 5.7$`\mathrm{}`$ (150 pc), for the adopted $`\alpha _{}`$=95$`\mathrm{}`$. Therefore the ISO beam can be modelled by a gaussian of FWHM$``$20 pixels..
Keeping fixed the stellar distribution, several models are produced with different parameters for the dust distribution. The final goal is to obtain a single model able to describe both the observed FIR SED and the spatial distribution of the emission.
## 6 Results
The first models presented here have the geometrical parameters of the standard model, i.e. a star-dust geometry with $`\alpha _\mathrm{d}=\alpha _{}`$ and $`\beta _\mathrm{d}=0.5\beta _{}`$. As already said in Sect. 2.2, the choice is mainly motivated by the presence of extinction lanes along the major axis of edge-on galaxies. The H<sub>2</sub> column density in NGC 6946 has an exponential radial profile, with a scalelength of 90$`\mathrm{}`$ (Tacconi & Young 1986). Incidentally, this is very close to $`\alpha _{}`$. Thus, the dust disk discussed here can be thought of as a dust component associated with the molecular phase. Four values for the V-band face-on optical depth have been chosen, $`\tau _\mathrm{V}=`$0.5, 1, 5 and 10.
The SED of the four models in presented in Fig. 2. The thick solid line represents the stellar emission of the galaxy as derived from the simulations, normalised to the observed data points as described in the previous section. Each different line shows the intrinsic unextinguished SED in the short-wavelength part of the spectrum and the FIR emission in the long-wavelength side. For the dust-free emission, the SED is derived from the total intrinsic energy, assuming isotropic emission and measuring the half light radius in a transparent model. All the models presented in this paper show a spike at $`\lambda `$2000Å in the unextinguished SED. This is due to the extinction feature at 2175Å present in the assumed extinction law (Sect. 2.3), while the observed SED is flat. A possible dip in the stellar SED caused by this extinction feature could be masked by the broad band of the observation or the large errors in the UV photometry. The extinction feature could also be absent, in which case our assumption of Galactic dust for NGC 6946 is not correct.
In all the models, the wavebands from EUV to U contributes tho 30% of the total absorbed energy, while the wavebands from B to LMN to 60%. Therefore, most of the radiation is absorbed from the Optical-NIR light. A similar result was obtained for the same galaxy by Trewhella (1998b).
The MIR corrections are quite similar for all the models: approximately 32% of the total absorbed energy is estimated to go into MIR emission, the remaining 68% being available for thermal equilibrium processes and FIR emission. Since small-grains and PAHs responsible for no-equilibrium processes have a higher absorption efficiency at shorter wavelength, the contribution of absorption from Optical-NIR wavebands is higher after the MIR correction. So the radiation originally emitted at $`\lambda >4000`$Å contributes $``$ 70% of the FIR emission. The same result applies to all the models presented in this paper, therefore MIR corrections are not discussed separately in each case. A comparison between the estimated and observed MIR emission is given in Sect. 7.
The temperature distributions for each model are shown in Fig. 3, as a function of the galactocentric radius and height above the plane. Apart from the central region, the distributions are very similar. For a dust distribution narrower than the stellar one, the stellar radiation field is expected to increase with height above the plane in an optically thick model (Draine & Lee 1984; Rowan-Robinson 1986), because the stars closer to the plane are shielded. This is evident in the central regions ($`R<1.5\alpha _{}`$) of the models. When the optical depth increases, dust at higher temperature is found at higher positions above the plane. In the models with higher extinction, the effect can still be seen at larger galactocentric distances, the region of higher temperature approaching the galactic plane at large distances. Vertical gradients are very shallow, because of the greater extent of the galaxy in the radial direction with respect to the vertical and because the stellar distribution is smooth.
The FIR spectrum is shown in Fig. 2. Only models with central optical depth between $`\tau _\mathrm{V}`$=5 and 10 produce enough energy to match the observational data. This is a general property of all the models we are going to discuss: a substantial extinction is necessary to produce the observed SED in the FIR. The total amount of energy absorbed and re-emitted (in both MIR and FIR) is between 1.5 and $`310^{10}`$ L, for the two high $`\tau _\mathrm{V}`$ models, respectively. This corresponds to a fraction 0.3-0.4 of the intrinsic energy produced by the stars. Therefore, $`1/3`$ of the bolometric luminosity of NGC 6946 is absorbed by dust.
The optical (B) and FIR (100$`\mu `$m and 200 $`\mu `$m) scalelengths measured on standard models are shown in Table 4. The B-band scalelength increases with $`\tau _\mathrm{V}`$ and it is close to that observed if $`\tau _\mathrm{V}5`$. The FIR scalelengths increase with the wavelength. A slight increase with $`\tau _\mathrm{V}`$ is also observed, because of the smaller temperature in the centre. However, FIR scalelengths are never as large as observed, especially at 200$`\mu `$m. While the 100$`\mu `$m scalelength for the $`\tau _\mathrm{V}=5`$ model is underestimated by $``$15%, the 200$`\mu `$m scalelength is about one half of that measured on ISO images. Therefore, the ratio between the B-band and 200$`\mu `$m scalelengths is smaller than that measured by Alton et al. (1998c). It is interesting to note that, because of the more rapid increase of the B-band scalelength with $`\tau _\mathrm{V}`$, the lowest values of the ratio are obtained for optically thin cases. Nevertheless, optically thin cases are unable to explain the large 200$`\mu `$m scalelength and the energy output. We note here that, by adopting a single $`\alpha _{}`$ measured in the K-band, we have assumed that the larger optical scalelength with respect to that measured in the NIR is entirely due to dust extinction. If indeed an intrinsic difference exist, a smaller optical depth may be able to reproduce the observed B band scalelength. We will discuss the effect of our assumption in Sect. 7.
Alton et al. (1998c) suggest that the large observed scalelengths at 200$`\mu `$m is due to a dust distribution more extended than the stellar disk. To test this hypothesis, we run models with $`\alpha _\mathrm{d}=1.5\alpha _{}`$ (Davies et al. 1997; Xilouris et al. 1999), keeping the other parameters as for the standard model. The $`\tau _\mathrm{V}=5`$ case is shown in Fig. 4. As for standard models, only optically thick cases are able to match the observed SED. For the same optical depth the extended model has a higher extinction (e.g. 44% of the energy is absorbed in the V-band against the 34% of the standard model).
The temperature distribution for the extended model is shown in the central panel of Fig. 5. For ease of comparison, the temperature distribution of the standard model is shown again (in the right panel), with the same scale as for the new model. Within a radius of 6 $`\alpha _{}`$ (the extent of the stellar disk) the temperature pattern of the extended model is quite similar, apart from a small difference due to the normalisation. This is reflected in the peak of the FIR SED, that is essentially the same in both the models. Outside of 6$`\alpha _{}`$, the truncation of the stellar distribution, dust is colder and it does not modify the shape of the SED. The steeper temperature gradient at 6$`\alpha _{}`$ is an artifact due to the truncation of the stellar disk. A truncation is indeed suggested by counts of faint stellar sources in the Galaxy (Sect. 2.1). A few tests have been conducted with stellar distributions truncated at the same distance as the dust disks, to avoid having dust in regions without local stellar emission. The steeper gradient disappears and a larger distance is needed to reach the same temperature. However, changes are small, the general trend in the temperature distribution and in the FIR emission distribution are essentially the same. It is interesting to note that for $`R>6\alpha _{}`$, dust is colder on the plane than above, because starlight is seen through higher optical depths along the plane.
The extended dust distribution causes an increase of the FIR scalelengths with respect to the standard model (Table 4). Scalelengths at 100$`\mu `$m are quite close to those observed (within a 3% for the $`\tau _\mathrm{V}=5`$ case). However, the 200$`\mu `$m scalelengths are still underestimated, by at least 30%. B band scalelengths are similar to those for the standard model (the effect of an extended dust distribution being appreciable only for high optical depths). As in the standard model, despite the increase of the FIR scalelength with $`\tau _\mathrm{V}`$, B/200$`\mu `$m scalelength ratios are close to those of Alton et al. (1998c) only in the optically thin case.
The analysis of Galactic FIR emission of Davies et al. (1997) suggests that the vertical scalelength of dust as well is larger than the stellar one, by a factor of two. Models that included a thicker dust disk, however, do not produce better results than those presented here (Bianchi 1999).
Larger FIR scalelengths can be produced by extending further the dust distribution. Davies et al. (1999b) note that the IRAS emission tends to follow the molecular gas, while the $`200\mu `$m profile is closer to the much broader atomic gas component. Could dust associated with the HI distribution be responsible of the emission at $`200\mu `$m? A model was produced with two dust disks. As said earlier, a disk with $`\alpha _\mathrm{d}=\alpha _{}`$ can describe dust associated with the molecular gas. As for the dust associated to the atomic gas, we chose a disk with a flat radial distribution up to 2.5$`\alpha _{}`$, then falling off exponentially with a scalelength of 3 $`\alpha _{}`$. This mimics the observed column density of HI (Tacconi & Young 1986).
Using the relation between extinction and gas column density of Bohlin et al. (1978) and the column density measured by Tacconi & Young (1986), we derived an optical depth $`\tau _\mathrm{V}0.5`$ for the broader dust disk. For the dust disk associated with the molecular phase we adopted $`\tau _\mathrm{V}=5`$, as in the previous standard disks. Using the observed H<sub>2</sub> column density, a value larger at least by a factor of two would have been derived. However, too much energy would have been absorbed from dust (see, for instance, the SED of the $`\tau _\mathrm{V}=10`$ standard model in Fig. 2). The discrepancy between the value derived from the gas column density and that needed to produce the right amount of absorption may be due to the smaller effective absorption in a clumpy structure of dust. Clumping is discussed in Sect. 7.
The SED of this new model is shown in Fig. 4. The temperature distribution is presented in the bottom panel of Fig. 5. Despite the colder temperature of the dust associated to HI, the behaviour of the FIR radiation in the region where scalelengths are measured is dominated by the dust disk associated with H<sub>2</sub><sup>3</sup><sup>3</sup>3By extending the dust disk further, the FIR scalelength can became as large as that observed. This happens, for example, for a single HI-like dust distribution. However, optically thick models ($`\tau _\mathrm{V}=2`$ for the HI-like distribution) are still needed. Such models would have a dust mass larger than what implied assuming the canonical gas-to-dust mass ratio (Sect. 7).. The derived scalelengths are similar to those of an extended model ($`\alpha _\mathrm{d}=1.5\alpha _{}`$) with a single disk.
In conclusion, the optically thick regime is required to match the observed FIR SED, both in models with standard and extended dust distribution. FIR scalelengths are larger in extended models with respect to standard models. The 100$`\mu `$m scalelength is very similar to that derived on IRAS images, but the 200$`\mu `$m scalelength is always smaller than that measured by ISO. It is not possible to produce B/200$`\mu `$m scalelength ratios as small as in Alton et al. (1998c), for any of the models presented here.
## 7 Discussion
As shown in the previous sections, models with a central face-on optical depth $`\tau _V5`$ are necessary to explain the SED observed in the FIR for NGC 6946. A high optical depth through the central regions of the galaxy has also been found by Engargiola (1991) and Devereux & Young (1993). Evans (1992) and Trewhella (1998b) apply the energy balance method to the stellar and dust emission of NGC 6946, using a TRIPLEX model, i.e. an analytical approximation for the radiative transfer, neglecting scattering, in a standard mode (Disney et al. 1989). They both derived high optical depths for the disk, using the data inside the half light radius. Evans (1992) measured $`\tau _V=67`$, while Trewhella (1998b) $`\tau _V=4\pm 1`$. A high optical depth is also suggested by high-resolution sub-mm SCUBA images: the diffuse inter-arm emission in the NE spiral arms at a distance of 2$`\mathrm{}`$ ($`\alpha _{}`$) is compatible with $`\tau _V=2.2`$ (Bianchi et al. 2000a). Xu & Buat (1995) carried out an energy balance on a sample of 134 nearby spirals with available UV, B and IRAS fluxes. They derived a mean optical depth $`\tau _B=0.60`$. A direct comparison between their result and the optical depth needed by our model to explain the FIR output may lead to the wrong conclusion that we have overestimated the amount of absorbed energy. However, the mean ratio of bolometric luminosity absorbed by dust for their sample ($``$ 1/3) is similar to that predicted by our models for NGC 6946. The difference depend on their adopted geometry, a plane-parallel homogeneous model for dust and stars. If the dust is associated with the dominant gas component, like observations suggest, our choice of an exponential distribution seems more appropriate. For the same optical depth, their model results much more effective in extinguishing radiation. For example, in the UV, Xu & Buat adopt dust and stellar distributions of the same thickness (a slab model Disney et al. 1989) and isotropic scattering with albedo $`\omega `$=0.18. By opportunely setting scalelengths and truncations, the BFG model can produce results for a slab. At their UV reference wavelength $`\lambda =2030`$Å, 60% of the radiation is absorbed by dust, for their adopted UV extinction law. Because of their low albedo, scattering does not play a relevant role in the radiative transfer. In our corresponding UV4 band, the $`\tau _V=5`$ standard model absorbs only 50% of the radiation. In the optical, they adopt a sandwich model, with the thickness of the stellar distribution twice that for dust. At $`\lambda =4400`$Å, their model absorbs 20% of the radiation, while the $`\tau _V=5`$ standard model 40%.
In our simulations most of the radiation is absorbed in the Optical-NIR spectral range (60% for light at $`\lambda >4000`$Å). This was alredy noted for NGC 6946 by Trewhella (1998b). Xu & Buat (1995) reach opposite conclusions, with the non-ionising UV (912Å$`<\lambda <`$3650Å) contributing to 60$`\pm 9`$% of the absorbed radiation. This is in part due to the model differences we have already discussed. Xu & Buat adopted a different thickness for the UV stellar distribution to simulate the thinner distribution of younger stars. We tried a similar approach by adopting the same scaleheight for dust and star at smaller $`\lambda `$ (Bianchi 1999), but the fraction of radiation absorbed in the UV increased only by a 5%. Furthermore, the SED observed in NGC 6946, with a peak in the NIR, may be different from the mean characteristic of the Xu & Buat sample. The UV selected sample may be biased toward bright UV galaxies, although this hypothesis is dismissed in Buat & Xu (1996).
Another major difference between our model and that of Xu & Buat (1995) is in the ionising UV longward of 912Å. In their work the absorption of Lyman continuum photons contributes as much as 20$`\pm 1`$% to the total FIR emission in a sample of 23 late type galaxies. Our omission of this spectral range from the radiative transfer, due to the lack of empyrical data, may result in an overestimate of the optical depth. To test our approximation, we derived the ionising UV from H$`\alpha `$ observations of NGC 6946 (Kennicutt & Kent 1983), assuming standard ionisation conditions as in HII regions (Lequeux 1980). After correcting for the \[NII\] contamination (25% for spirals), Galactic extinction ($`A_{\mathrm{H}\alpha }1`$) and internal extinction (30% for the R-band, where the H$`\alpha `$ line is located, for a standard $`\tau _V=5`$ model), we derived an intrinsic unextinguished H$`\alpha `$ flux,
$$f(\text{H}\alpha )=\mathrm{9.1\hspace{0.33em}10}^{11}\text{erg cm}\text{-2}\text{ s}\text{-1}.$$
(14)
Following Xu & Buat (1995), the Lyman continuum flux can be derived as
$$f(\text{Lyc})=33.9f(\text{H}\alpha )=\mathrm{3.1\hspace{0.33em}10}^9\text{erg cm}\text{-2}\text{ s}\text{-1},$$
(15)
of which 75% is assumed to be absorbed by gas and converted into emission lines at larger wavelengths (see also Mezger 1978; DeGioia-Eastwood 1992). Allowing the remaining 25% to be entirely absorbed by dust and adopting a lower limit MIR correction of $`70`$% (Sect. 9), the total FIR luminosity originating from the absorption of ionising photons is
$$L^{\mathrm{FIR}}(\text{Lyc})=\mathrm{2.2\hspace{0.33em}10}^8\text{L}_{},$$
(16)
for the assumed distance of 5.5 Mpc. This correspond to only $``$ 2% of the total FIR luminosity emitted by the standard $`\tau _V`$=5 model. As a comparison, the contribution to the FIR from the EUV band for the same model is 3.6%. The ionising UV contribution is similar for all the models that roughly provide the same amount of FIR energy as that observed. Therefore, it is justified to omit the ionising radiation in the model. The approach of Xu & Buat (1995) is different from the one presented here. Their UV contribution includes direct absorption of Lyc photons and indirect (via emission lines), while in the present model the absorption of emission line photons is taken care of in the spectral band where the emission occurs (e.g. in the R-band for the H$`\alpha `$ line), the total contribution of re-combination is summed up to the stellar SED for each band. Nevertheless, their ratio between Lyc emission and total absorbed energy is similar to the one derived here.
A wrong MIR correction could also be responsible for an underestimation of the FIR emission in our model. The fraction of absorbed energy that goes into MIR emission depends essentially on the absorption of light from the short wavelength spectrum, since the absorption efficiency of small grains responsible for non-equilibrium processes is higher in the UV (Sect. 9). For the model presented here, with a dust scaleheight smaller than the stellar one, the amount of energy absorbed from UV bands does not increase very much with the optical depth (The saturation effect, Bianchi et al. 1996). The MIR corrections for these models are therefore quite constant, $``$32% of the total absorbed energy being re-emitted in the MIR. Even for models with higher efficiency in extinguishing radiation, like the optically thick model with a thicker dust distribution, the fraction of the total energy emitted in the MIR is not very different from this value (Bianchi 1999).
For the local ISRF the contribution of small grain emission to the 60 $`\mu `$m IRAS band is $``$62%, while at 100 $`\mu `$m it is only 14% and at 200 $`\mu `$m 4% (Désert et al. 1990). Therefore, the fraction of energy emitted in non equilibrium heating can be roughly estimated by measuring the MIR emission shortward of 60 $`\mu `$m. After integrating a continuous SED interpolated from the data points in Table 3, the MIR energy is derived to be 34% of the total infrared energy emitted by dust. The value derived from observation is very close to the model one. This justifies the use of the Désert et al. (1990) dust model as described in Sect. 9. It is interesting to note that the infrared galactic spectrum used in the Désert et al. (1990) model is different from the one of NGC 6946. As an example, the ratio between fluxes at 60 $`\mu `$m and 100 $`\mu `$m is 0.2, while it is 0.5 from the NGC 6946 data. This does not necessarily mean that the dust model of Désert et al. (1990) cannot be applied to NGC 6946. The different ratio could be due to different heating conditions in the local interstellar radiation field, with respect to the mean radiation field of NGC 6946. Larger ratios between 60 $`\mu `$m and 100 $`\mu `$m can be derived from the Désert et al. (1990) model when the ISRF is larger than the local. As already said, emission from small grains can contribute to part of the observed FIR flux. However, for a wide range of heating conditions, the contribution close to the peak of FIR emission is minimal Désert et al. (1990). Since modelled and observed fluxes are compared in this spectral range, the assumption that all FIR emission occurs at thermal equilibrium does not affect sensibly our results.
The position of the peak emission in the modelled FIR SED clearly shows that the dust temperature in the simulations is not severely different from the real one. Recently, Stickel et al. (2000) extracted high signal-to-noise sources from the ISOPHOT Serendipity Survey. Fluxes at 170 $`\mu `$m were measured for 115 sources having a galaxy association, most of which are spirals, and integrated with 100$`\mu `$m IRAS data. The distribution of the ratio of 170$`\mu `$m and 100$`\mu `$ fluxes is quite well confined, with half of the galaxies having a ratio between 1 and 1.5. For the emissivity adopted here, this translates in a colour temperature in the range 23-28 K. Similar colour temperatures (24-26K) can be derived from the ratios of the total fluxes at 100 and 200$`\mu `$m, in any of our models. Using emissivities derived from Galactic emission and exinction, Bianchi et al. (1999) found that the value of the dust mass for a small sample of galaxies does not depend dramatically on the assumed spectral behaviour of the emissivity. This can be easily interpred if dust temperatures in external galaxies are similar to those in the Milky Way. Because of these similarities, and of the small dependence of the temperature on the assumed model, it is not too audacious to compare values measured on the present models with the high quality observations of the Galaxy.
Dust temperatures at high latitude are quite constant in the Galaxy (Reach et al. 1995; Lagache et al. 1998; Schlegel et al. 1998). For the emissivity adopted here, a temperature of $``$21K can be derived (Bianchi et al. 1999). Remarkably, nearly all the models present a similar temperature at a galactocentric distance of 3$`\alpha _{}`$ (the Sun position, for $`R_{}`$=8.5kpc and a Galactic scalelength of 3kpc). Sodroski et al. (1997) decompose the Galactic FIR emission observed by DIRBE into three components, associated with the atomic, molecular and the ionised gas phases, and derive the temperature for four annulii at different galactocentric distance. Fig. 6 shows the temperature of the dust associated with the atomic gas (supposedly heated mainly by the diffuse ISRF) for each annulus. Data have been scaled to a Galactic scalelength $`\alpha _{}=3`$ kpc and corrected for the emissivity law used in this work. In Fig. 6 we also plot the temperature radial profile for four representative models: the $`\tau _\mathrm{V}=0.5`$ and $`\tau _\mathrm{V}=5`$ standard models; the $`\tau _\mathrm{V}=5`$ $`\alpha _\mathrm{d}=1.5\alpha _{}`$ model; and the model with a dust disk associated with each of the two gas phases. Clearly, all the simulations presented here have a temperature gradient compatible with the Galactic one.
Because of the lack of high resolution observations in the FIR, the derivation of temperature gradients in external galaxies is more difficult. Davies et al. (1999b) derived the temperatures from 100$`\mu `$m IRAS and 200$`\mu `$m ISO fluxes at two different position on NGC 6946, in the centre and on the disk at 3 arcmin from the centre. After smoothing for the ISO resolution, our models are compatible with those observations (Bianchi 1999).
For any model, the FIR scalelengths increase with the wavelength of emission, the $`\alpha _\mathrm{d}/\alpha _{}`$ ratio and, although slightly, with the optical depth. Therefore, larger FIR scalelengths can be found in optically thick extended models. For the $`\tau _\mathrm{V}=5`$ extended model, that provides a good fit to the SED, the 100$`\mu `$m scalelenght is very close to the value derived from IRAS images. The $`200\mu \mathrm{m}`$ scalelength is larger than scalelength at 100$`\mu `$m. However, it is 30% smaller than that observed by ISO. It has not been possible to find a model able to reproduce the large 200$`\mu `$m scalelength measured by Alton et al. (1998c) in a sample of seven galaxies including NGC 6946. The B-band ratio is close to that observed for the $`\tau _\mathrm{V}=5`$ extended and standard models. Alton et al. derived a ratio 0.9 between the B and 200$`\mu `$m scalelengths. Because of the dependence of the B-band scalelength on $`\tau _\mathrm{V}`$, similar ratios can be obtained only for optically thin extended models. Such models fail to predict both the SED and the absolute values of the FIR scalelengths. A possible reason for the discrepancy between observations and models may reside in the transition effects of the ISOPHOT detectors. Memory effects during the scanning of a bright source affects the P32 mapping mode. Alton et al. checked for this problem by conducting the scalelength analysis on both in-scan and cross-scan directions and concluded that the effects were negligible. However, a proper analysis requires a knowledge of the underlying light distribution and a description of the transient effects, which are still poorly understood. The P32 mapping mode is still not scientifically validated (Klaas & Richards 2000).
We have carried out a few test to check for the influence of different stellar distributions, like including a small bulge appropriate for a late type galaxy or having a smaller vertical scalelength for younger stars (Bianchi 1999). The result were not significantly different from those presented here. A basic assumption of all our models is that the intrinsic (not extinguished) radial stellar scalelength is the same at any wavelength. According to this view, the larger scalelengths observed in the optical are due to extinction, rather than to intrinsic color gradients. As outlined in Sect. 2.1, the problem is still very debated. A simple test can be conducted here to assess the influence of scalelength variations with $`\lambda `$. Let’s assume that the measured B-band scalelength (1.3$`\alpha _{}`$) is more representative of the intrinsic radial variation of the stellar distribution. By substituting $`\alpha _{}`$ with 1.3$`\alpha _{}`$, the results of Table 4show that standard or extended models with $`\tau _\mathrm{V}=0.5`$ or 1 have an observed (after smoothing) B-band scalelength quite similar to the one of Alton et al. (1998c), the difference between observed and intrinsic scalelength due to the smoothing process, rather than to extinction. Since the modelled B/200$`\mu `$m scalelength ratio remains the same, these optically thin cases may have a 200$`\mu `$m scalelength close to the observed. However, the fraction of energy absorbed in these models will remain the same as well. Therefore, the optically thin models will still be unable to produce the required energy output in the FIR. Furthermore, most of the energy is absorbed in the NIR, where the peak of stellar emission occurs. In this spectral range our K-band based $`\alpha _{}`$ is surely a better description for the galaxy radial scalelength.
We also produced a model with two dust distributions, a $`\tau _\mathrm{V}=5`$ standard disk associated with the molecular gas dominant in the central part of NGC 6946, and a $`\tau _\mathrm{V}=0.5`$ disk associated with the broader distribution of atomic gas. Even with the presence of such an extended distribution, the double disk model is dominated by the optically thick disk associated to the H<sub>2</sub> necessary for the FIR emission. For what concerns FIR scalelengths, the same results as for the $`\tau _\mathrm{V}=5`$ extended model are obtained. It is interesting to note that for the parameters adopted here for the two dust disks and assuming the atomic + molecular gas mass of Table 2, the gas-to-dust mass ratio is 180, close to the Galactic value of 160 (Sodroski et al. 1994). Similar values can be retrieved from all the optically thick models.
The high optical depth of the models contrasts with the recent determination of optical depth in edge-on spirals by Xilouris et al. (1999). Using a sample of seven edge-on galaxies, they find a mean central face-on optical depth $`\tau _\mathrm{V}0.5`$. The higher opacity of NGC 6946 may be a result of the galaxy being very gas-rich (Tacconi & Young 1990); or it may be due to clumping of the ISM, affecting in a different way FIR and optical determinations of the optical depth. While FIR observations detect all of the dust (at least when the temperature of the clump and inter-clump medium are similar), optical observations may be affected preferentially by the extinction of the smoother, lower density (and optical depth) inter-clump medium.
To our knowledge, only two works in literature discuss the effects of clumping within the framework of a radiative transfer model for spiral galaxies. Kuchinski et al. (1998) distribute clumps in the dust disk assuming a constant filling factor and a value for the ratio between densities in clumps and in the nearby smooth medium. The other model is by Bianchi et al. (2000b), based on the BFG radiative transfer code, as the present work. Bianchi et al. (2000b) modelled the clump distribution in a similar way to the distribution of molecular gas in the Galaxy. The clump properties were derived from those observed in Giant Molecular Clouds. A comparison with the Kuchinski et al. (1998) results indicates a strong dependence of the observed brightness profiles on the detailed internal and spatial distribution properties of clumps. This makes the interpretation of the data very difficult. For the same dust mass, it is found that a clumpy dust medium has lower extinction. However, when a fraction of the stellar radiation is allowed to be emitted from inside the clumps, as for Giant Molecular Clouds hosting star-formation, extinction increases and can reach higher values than those for homogeneous models.
Predictions of the influence of clumping on the FIR models are not easy. If the dust component associated with the H<sub>2</sub> is clumped, it could be responsible for most of the FIR emission, when embedded stellar emission is considered. A diffuse component associated with the HI would be responsible for a reduced $`\tau _\mathrm{V}`$ derived from the analysis of edge-on galaxies. For the models of Bianchi et al. (2000b), it is unlikely that a clumpy distribution with a large dust mass (corresponding to an optically thick smooth distribution) has an optically thin apparent $`\tau _\mathrm{V}`$ as measured by Xilouris et al. (1999). However, the models of Bianchi et al. (2000b) are based on the Galaxy, while clumping based on the different distributions of atomic and molecular gas in NGC 6946 may have a different behaviour. Clumping will also affect the spatial distribution of the FIR emission. Dust in clumps is shielded from the ISRF and heated to lower temperature. Depending on the distribution of cold dust clouds, a broader 200$`\mu `$m emission could be produced.
Two recent models include clumping of dust and embedded stellar emission to describe the radiative transfer and FIR emission of NGC 6946. However, the complexity of the models prevents an isolation of the effect of the dust distribution on the FIR heating. Silva et al. (1998) fit the observed optical and FIR SED of NGC 6946 with their galactic photometric evolution model, adopting a simplified treatment for the radiative transfer in the diffuse medium and a separate model for dust in clumps. They find that nearly half of the dust emission comes from within molecular clouds. The adopted distributions for smooth dust and stars are different from those in Sect. 6: the radial scalelengths are the same for both dust and stars, as well as the vertical scalelengths; the radial scalelength is twice the one adopted here; the disk is thicker by a factor of two. Since we use the same SED and their model is optically thin to B-band radiation (only 10% of radiation absorbed), most of the FIR emission must come from absorption of UV photons in the molecular cloud distribution. A distribution of molecular clouds compatible with the galactic gravitational potential, together with a smooth phase modelled on the atomic gas, is used by Sauty et al. (1998) to study the FIR emission in NGC 6946. The radiative transfer is carried out using the Monte Carlo technique only for the UV radiation emitted in star-forming regions within the molecular clouds. The ISRF at $`\lambda >2000`$Å is derived from a R-band map of the galaxy, scaled on the Galactic local ISRF. As in our case, the observed FIR emission can be explained only with substantial extinction. However, they find that emission at $`\lambda <2000`$Å contributes 72% of the total FIR, while in the present work energy absorbed from optical radiation is dominant.
Therefore, according to the models of Silva et al. (1998) and Sauty et al. (1998), the dominant contribution to the dust heating comes from young stars. If this is true, the recent star-formation rate of a galaxy could be measured using FIR fluxes. These are more readily available than other tracers of star-formation (Devereux & Young 1991). In a series of papers, Devereux & Young (1993, and references therein) compared FIR and H$`\alpha `$ fluxes and suggested that the FIR luminosity is dominated by warm dust absorbing radiation from OB stars. However, this result is debated. Walterbos & Greenawalt (1996) modelled the FIR emission in spiral galaxies, deriving the diffuse ISRF from optical profiles and estimating the dust column density from the atomic gas. For a sample of 20 galaxies, they found that dust heated by a diffuse ISRF can account for, on average, half of the observed IRAS fluxes. Using large nearby objects it is possible to decompose the FIR emission into the different sources of heating. Xu & Helou (1996a) measure the ratio between the IRAS 60$`\mu `$m and H$`\alpha `$ fluxes from bright FIR-resolved sources in M31. Using the total H$`\alpha `$ luminosity and the Désert et al. (1990) dust model, they extrapolate the fraction of the total luminosity emitted by dust that is associated with HII regions and star-formation. A value of 30$`\pm `$14% only is found. A decomposition of the Galactic flux observed by COBE at 140$`\mu `$m and 240$`\mu `$m shows that most of the dust emission ($``$70%) arises from dust associated with the atomic gas (Sodroski et al. 1994), with temperature gradients compatible with a diffuse ISRF heating. Only 20% of the FIR is emitted by dust associated with the molecular component and 10% is due to hot dust associated with the HII phase.
The smooth models presented here attempt to explain the FIR emission without invoking dust heated in star-forming regions. Under the assumption of a diffuse stellar emission and dust, we find that radiation at $`\lambda <4000`$Å contributes $``$ 30% to the total FIR, this including the ionising UV. This percentage will surely increase when a clumpy dust structure is adopted. The high optical depths necessary to produce the observed energy may be an overestimation because we neglect clumpy hot dust. However, without a proper model for the radiative transfer and dust emission in a clumpy medium, it is not easy to understand if the results shown here are severely biased by the assumption of smooth distributions.
Finally, we tried to model the FIR emission with an extra dust component, a spherical halo. The presence of dust in the halo of normal, non-active, spirals is predicted as a result of the imbalance between the radiation pressure and the galactic gravitational force (Davies et al. 1999a). A halo of cold gas clouds, able to explain part of the dark matter in spiral galaxies, may be stabilized by the presence of dust grains (Gerhard & Silk 1996). Unfortunately, it is difficult to obtain information about a putative dusty halo. Because such a distribution would act as a screen for the galactic disk, the radiative transfer fits of Xilouris et al. (1999) are unable to detect it. Analysing the difference in colours of background objects between fields at different distances from the centre of two nearby galaxies, Zaritsky (1994) find a B-I colour excess in the inner fields. He derived a dust halo scalelength of 31$`\pm `$8 kpc, although the measured colour difference is only 2$`\sigma `$ of the statistical noise.
Because of the high uncertainty of the halo parameters, we simply used a spherical homogeneous dust halo, together with a standard disk. The halo has the same radial dimension as the dust and stellar disks, i.e. 6$`\alpha _{}`$. Among the models we tried, the case in which $`\tau _\mathrm{V}=1`$ for both halo and disk provides a reasonable fit to the FIR emitted energy and scalelengths (Bianchi 1999). However, when the galaxy is seen edge-on, such a dust distribution would be easily observable, within the resolution and sensitivity of available FIR observations. Alton et al. (1998b) studied the FIR emission in 24 edge-on galaxies, including starburst and quiescent objects, using HiRes IRAS images. None of the object was found to be resolved along the minor axis.
## 8 Summary and conclusions
We have described in this paper a model for dust emission in spiral galaxies, based on the Monte Carlo radiative transfer code of Bianchi, Ferrara & Giovanardi (1996). For each relative star-dust geometry and dust optical depth, the radiative transfer code is carried out for 17 different photometric bands, covering the spectral range of stellar emission. In the application to NGC 6946 presented here, we have chosen a model stellar SED and dust distribution that produces the same fluxes as those observed in the UV-Optical-NIR images. The code also produces a map of the total energy absorbed by dust. For each position inside the galaxy, dust is heated by an ISRF consistent with the radiative transfer itself, without any other assumption. The dust temperature is computed from the absorbed energy, using the emissivity derived by Bianchi et al. (1999) for Galactic dust and correcting for the contribution of non-thermal equilibrium processes to the emission. Hence, FIR maps can be easily obtained for any wavelength, integrating along a specific line of sight.
The model optical and FIR scalelengths and the SED have been compared to NGC 6946 observations. Several models have been explored. It is found that optically thick dust disks ($`\tau _\mathrm{V}`$5) are needed to match the observed FIR output. Approximately one third of the total stellar radiation is absorbed by dust in this case. The temperature distributions are quite similar, for any of the dust disk models. Temperature values in the models are comparable with those observed in our Galaxy and other spirals.
We have compared the modelled FIR scalelengths with the observations of Alton et al. (1998c). Extended dust disk model with $`\alpha _\mathrm{d}=1.5\alpha _{}`$ (Davies et al. 1997; Xilouris et al. 1999) have larger FIR scalelengths than standard models with $`\alpha _\mathrm{d}=\alpha _{}`$. For optically thick cases, the 100$`\mu `$m scalelength is close to the value derived on IRAS images. Alton et al. (1998c) found that the B-band scalelength of NGC 6946 is smaller than the 200$`\mu `$m one by a factor 0.9 (0.8 for a sample of seven galaxies). We have not been able to reproduce the large FIR scalelengths measured on 200$`\mu `$m ISO data. In the required optically thick regime, the scalelength ratio B/200$`\mu `$m is always larger than observed. Smaller ratios can be obtained only in optically thin cases, but the absolute values for the scalelengths are smaller. The results are not improved if two dust disks modelled on the gas distribution are used. The behaviour of the model is dominated by the standard optically thick disk associated with the molecular component, rather than the very extended dust distribution associated with HI. A spherical dust halo could produce results closed to those observed, but would also be easily detected in currently available FIR observations, which is not.
The high optical depth found for NGC 6946 contrasts with recent determinations on edge-on spiral galaxies Xilouris et al. (1999). This may be a result of our assumption of a smooth distribution for stars and dust. The inclusion of clumping in a proper model of radiative transfer and FIR emission is therefore desirable. However, the heavy dependence of clumping on the assumed model makes the modelling more complicate. Future high resolution and sensitivity instrumentation will therefore be essential to define the dust distribution and limit the number of parameters in the model. Clumping may also be responsible for the discrepancy between the observed and modelled scalelengths at 200 $`\mu `$m, if a diffuse component of cold dust clumps shielded from the ISRF is present. On the other hand, it is necessary to remind here that the results of Alton et al. (1998c) at 200$`\mu `$m are based on data which is not yet scientifically validated. Again, future FIR instrumentation or a set of validated data will help to asses if the large scalelengths are an artefact of the ISOPHOT detector transients or if they are the genuine results of a temperature distribution different from the one we have derived here.
The main purpose of the paper was the presentation of the model itself. A test has been conducted on a well studied galaxy. However, NGC 6946 is very gas rich (Tacconi & Young 1990) and its characteristics may be different from those of a ’mean’ galaxy. Optical and FIR data are available for several galaxies. A future paper will be devoted to their analysis and more general conclusions about the dust distribution and extinction will be drawn.
###### Acknowledgements.
The work of this paper has benefitted from comments and discussion with many people. Among them we wish to remember M. Trewhella, Z. Morshidi, R. Smith, A. Kambas, J. Haynes, Rh. Evans, Rh. Morris, A. Whitworth, A. Ferrara, S. Kitsionas, P. Gladwin, N. Francis and an anonymous referee. S.B. acknowledges a PhD studentship from the Department of Physics and Astronomy at Cardiff University.
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# I Introduction
## I Introduction
The parity-violating (PV) nucleon-nucleon interaction has been a subject of interest in nuclear and particle physics for some time. To date, PV observables generated by this interaction remain the only experimental windows on the $`\mathrm{\Delta }S=0`$, nonleptonic weak interaction. Since the 1970’s, the PV NN interaction has been studied in a variety of processes, including $`\stackrel{}{p}p`$ and $`\stackrel{}{p}`$nucleus scattering, $`\gamma `$-decays of light nuclei, the scattering of epithermal neutrons from heavy nuclei, and atomic PV (for a review, see Refs. . The on-going interest in the subject has spawned new PV experiments in few-body systems, including high-energy $`\stackrel{}{p}p`$ scattering at COSY, $`\stackrel{}{n}+pd+\gamma `$ at LANSCE, $`\gamma +dn+p`$ at JLab, and the rotation of polarized neutrons in helium at NIST.
The theoretical analysis of these PV observables is complicated by the short range of the low-energy weak interaction. The Compton wavelength of the weak gauge bosons ($`0.002`$ fm) implies that direct $`W^\pm `$ and $`Z`$ exchange between nucleons is highly suppressed by the short-range repulsive core of the strong NN interaction. In the conventional framework, longer range PV effects arise from the exchange of light mesons between nucleons. One requires the exchange of the $`\pi `$, $`\rho `$ and $`\omega `$ in order to saturate the seven spin-isospin channels associated with the quantum numbers of the underlying four-quark strangeness-conserving PV interaction, $`_W^{PV}(\mathrm{\Delta }S=0)`$ (henceforth, the $`\mathrm{\Delta }S=0`$ will be understood). These exchanges are parameterized by PV meson-nucleon couplings, $`h_M`$, whose values may be extracted from experiment. At present, there appear to be discrepancies between the values extracted from different experiments. In particular, the values of the isovector $`\pi NN`$ coupling, $`h_\pi `$, and the isoscalar $`\rho NN`$ coupling, $`h_\rho ^0`$ – as extracted from $`\stackrel{}{p}p`$ scattering and the $`\gamma `$-decay of <sup>18</sup>F, do not appear to agree with the corresponding values implied by the anapole moment of <sup>133</sup>Cs as measured in atomic PV .
The origin of this discrepancy is not understood. One possibility is that the use of $`\rho `$ and $`\omega `$ exchange to describe the short-range part of the PV NN interaction is inadequate. An alternate approach, using effective field theory (EFT), involves an expansion of the short-range PV NN interaction in a series of four-nucleon contact interactions whose coefficients are a priori unknown but in principle could be determined from experiment. The use of $`\rho `$ and $`\omega `$ exchange amounts to adoption of a model – rather than the use of experiment – to determine the coefficients of the higher-derivative operators in this expansion. Whether or not the application of EFT to nuclear PV can yield a more self-consistent set of PV low-energy constants than the meson-exchange approach remains to be seen. A comprehensive analysis of nuclear PV observables using EFT has yet to be performed.
The least ambiguous element – shared by both approaches – involves the long-range $`\pi `$-exchange interaction. At leading order in the derivative expansion, the PV $`\pi NN`$ interaction is a purely isovector, Yukawa interaction. The strength of this interaction is characterized by the same constant – $`h_\pi `$ – in both the EFT and meson-exchange approaches. At the level of the Standard Model (SM), $`h_\pi `$ is particularly sensitive to the neutral current component of $`_W^{PV}`$. In this respect, the result of <sup>18</sup>F PV $`\gamma `$-decay measurement is puzzling:
$$h_\pi =(0.73\pm 2.3)g_\pi ,$$
(1)
where $`g_\pi =3.8\times 10^8`$ gives the scale of the $`h_M`$ in the absence of neutral currents. This result is especially significant, since the relevant two-body nuclear parity-mixing matrix element can be obtained by isospin symmetry from the $`\beta `$-decay of <sup>18</sup>Ne . The result in Eq. (1) is, thus, relatively insensitive to the nuclear model. Theoretical calculations of $`h_\pi `$ starting from $`_W^{PV}`$ have been performed using SU(6)<sub>w</sub> symmetry and the quark model , the Skyrme model , and QCD sum rules . As a benchmark for comparison with experiment, we refer the SU(6)<sub>w</sub>/quark model analysis of Refs. —hereafter referred to as DDH,FCDH <sup>*</sup><sup>*</sup>*Note that although the DDH analysis used the symmetry group SU(6)<sub>w</sub> in order to connect weak vector meson and pion couplings the predictions relating pion couplings alone to hyperon decay data rely only SU(3).. These authors quote a “best value” and “reasonable range” for the $`h_M`$:
$`h_\pi (\text{best})`$ $`=`$ $`7g_\pi `$ (2)
$`h_\pi (\text{range})`$ $`:`$ $`(030)g_\pi .`$ (3)
where here the “best value” is more aptly described as an educated guess, while the “reasonable range” indicates a set of numbers such that theory would be very hard-pressed to explain were the experimental value not found to be within this band. Neverthess, the difference between the “best value” of Eq. (2) and the <sup>18</sup>F result would appear to call for an explanation, and in the following note we comment on a possible source of the discrepancy.
In general, the problem of relating the fundamental weak quark-quark interaction to the low-energy constants which parameterize hadronic matrix elements of that interaction is non-trivial. In the framework of EFT, one may define these constants at tree-level in the hadronic effective theory. The quantities extracted from experiment in the conventional analysis, however, are not the tree-level parameters, but rather renormalized couplings. Denoting the latter as $`h_\pi ^{EFF}`$, one has
$$h_\pi ^{EFF}=Z_N\sqrt{Z_\pi }h_\pi ^{BARE}+\mathrm{\Delta }h_\pi ,$$
(4)
where $`h_\pi ^{BARE}`$ is the coefficient of the leading-order, PV Yukawa interaction in the effective theory, $`\sqrt{Z_N}`$ and $`\sqrt{Z_\pi }`$ denotes chiral loop renormalizations of the nucleon and pion wavefunctions, respectively, and $`\mathrm{\Delta }h_\pi `$ denotes contributions from chiral loops and higher-dimension operators to the Yukawa interactions (only the finite parts of these couplings are implied; loop divergences are cancelled by the corresponding pole terms in $`h_\pi ^{BARE}`$ and the $`Z_{N,\pi }`$). At leading order in $`1/\mathrm{\Lambda }_\chi `$, one has $`Z_{N,\pi }=1`$, $`\mathrm{\Delta }h_\pi =0`$, and $`h_\pi ^{EFF}=h_\pi ^{BARE}`$. The renormalized coupling appears as the coefficient in the one-pion-exchange (OPE) PV NN potential
$$\widehat{H}_{PV}^{OPE}=i\frac{g_{NN\pi }h_\pi ^{EFF}}{\sqrt{2}}\left(\frac{\tau _1\times \tau _2}{2}\right)_z(\stackrel{}{\sigma }_1+\stackrel{}{\sigma }_2)[\frac{\stackrel{}{p}_1\stackrel{}{p}_2}{2m_N},f_\pi (r)],$$
(5)
where $`g_{NN\pi }`$ is the strong $`\pi NN`$ coupling and $`f_\pi (r)=\mathrm{exp}(m_\pi r)/4\pi r`$. Neglecting the effects of three-body PV forces and $`2\pi `$-exchange interactions, it is $`h_\pi ^{EFF}`$ to which the result in Eq. (1) corresponds.
The relationship between $`h_\pi ^{EFF}`$ and the coupling obtained by computing $`N\pi |_W^{PV}|N`$ in a microscopic model is not immediately transparent. In what follows, we make several observations about this relationship. We first show that $`Z_N\sqrt{Z_\pi }`$ and $`\mathrm{\Delta }h_\pi `$ are substantial, so that $`h_\pi ^{EFF}`$ differs significantly from $`h_\pi ^{BARE}`$. To that end, we compute all of the chiral corrections to the PV Yukawa interaction through $`𝒪(1/\mathrm{\Lambda }_\chi ^3)`$, where $`\mathrm{\Lambda }_\chi =4\pi F_\pi `$. We work to leading order in $`1/m_N`$ in heavy baryon chiral perturbation theory (HBChPT). Of particular significance is the dependence of $`\mathrm{\Delta }h_\pi `$ on other low-energy constants parameterizing PV $`2\pi `$ production and the PV $`N\pi \pi \mathrm{\Delta }`$ transition. We subsequently reexamine the SU(6)<sub>w</sub>/quark model calculation of Refs. and argue that most – if not all – of the chiral loop effects which renormalize $`h_\pi `$ are not included in the microscopic calculation. Thus, the relationship between $`h_\pi ^{EFF}`$ and microscopic calculations remains ambiguous at best. This ambiguity is unlikely to be resolved until an unquenched lattice QCD calculation of $`h_\pi `$ using light quarks becomes tenable. In the meantime, one should not necessarily view a discrepancy between the experimental value of $`h_\pi ^{EFF}`$ and microscopic model calculations as disturbing.
Our discussion of these observations is organized as follows. In Section 2 we summarize our conventions and notation, including the PV chiral Lagrangians relevant to $`h_\pi `$ renormalization. Section 3 gives a discussion of the loop calculations. In Section 4 we comment on the scale of the loop corrections and provide simple estimates of some of the new PV low-energy constants appearing in the analysis. Section 5 gives our discussion of the relationship between $`h_\pi ^{EFF}`$ and the calculation of Refs. . Section 6 summarizes our conclusions. Some technical details are relegated to the Appendices.
## II Notations and Conventions
We follow standard HBChPT conventions and introduce
$$\mathrm{\Sigma }=\xi ^2,\xi =\mathrm{exp}(\frac{i\pi }{F_\pi }),\pi =\frac{1}{2}\pi ^a\tau ^a$$
(6)
with $`F_\pi =92.4`$ MeV being the pion decay constant. The chiral vector and axial vector currents are given by
$`𝒟_\mu `$ $`=`$ $`D_\mu +V_\mu `$ (7)
$`A_\mu `$ $`=`$ $`{\displaystyle \frac{i}{2}}(\xi D_\mu \xi ^{}\xi ^{}D_\mu \xi )={\displaystyle \frac{D_\mu \pi }{F_\pi }}+O(\pi ^3)`$ (8)
$`V_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}(\xi D_\mu \xi ^{}+\xi ^{}D_\mu \xi ).`$ (9)
For the $`\mathrm{\Delta }`$, we use the isospurion formalism, treating the $`\mathrm{\Delta }`$ field $`T_\mu ^i(x)`$ as a vector spinor in both spin and isospin space with the constraint $`\tau ^iT_\mu ^i(x)=0`$. The components of this field are
$$T_\mu ^3=\sqrt{\frac{2}{3}}\left(\begin{array}{c}\mathrm{\Delta }^+\hfill \\ \mathrm{\Delta }^0\hfill \end{array}\right)_\mu ,T_\mu ^+=\left(\begin{array}{c}\mathrm{\Delta }^{++}\hfill \\ \mathrm{\Delta }^+/\sqrt{3}\hfill \end{array}\right)_\mu ,T_\mu ^{}=\left(\begin{array}{c}\mathrm{\Delta }^0/\sqrt{3}\hfill \\ \mathrm{\Delta }^{}\hfill \end{array}\right)_\mu .$$
(10)
The field $`T_\mu ^i`$ also satisfies the constraints for the ordinary Schwinger-Rarita spin-$`\frac{3}{2}`$ field,
$$\gamma ^\mu T_\mu ^i=0\text{and}p^\mu T_\mu ^i=0.$$
(11)
We eventually convert to the heavy baryon expansion, in which case the latter constraint becomes $`v^\mu T_\mu ^i=0`$ with $`v_\mu `$ the heavy baryon velocity.
The relativistic parity-conserving (PC) Lagrangian for $`\pi `$, $`N`$, $`\mathrm{\Delta }`$ interactions needed here is
$`^{PC}`$ $`=`$ $`{\displaystyle \frac{F_\pi ^2}{4}}TrD^\mu \mathrm{\Sigma }D_\mu \mathrm{\Sigma }^{}+\overline{N}(i𝒟_\mu \gamma ^\mu m_N)N+g_A\overline{N}A_\mu \gamma ^\mu \gamma _5N`$ (15)
$`T_i^\mu [(i𝒟_\alpha ^{ij}\gamma ^\alpha m_\mathrm{\Delta }\delta ^{ij})g_{\mu \nu }{\displaystyle \frac{1}{4}}\gamma _\mu \gamma ^\lambda (i𝒟_\alpha ^{ij}\gamma ^\alpha m_\mathrm{\Delta }\delta ^{ij})\gamma _\lambda \gamma ^\nu `$
$`+{\displaystyle \frac{g_1}{2}}g_{\mu \nu }A_\alpha ^{ij}\gamma ^\alpha \gamma _5+{\displaystyle \frac{g_2}{2}}(\gamma _\mu A_\nu ^{ij}+A_\mu ^{ij}\gamma _\nu )\gamma _5+{\displaystyle \frac{g_3}{2}}\gamma _\mu A_\alpha ^{ij}\gamma ^\alpha \gamma _5\gamma _\nu ]T_j^\nu `$
$`+g_{\pi N\mathrm{\Delta }}[\overline{T}_i^\mu (g_{\mu \nu }+z_0\gamma _\mu \gamma _\nu )\omega _i^\nu N+h.c.],`$
where $`\omega _\mu ^i=\mathrm{tr}[\tau ^iA_\mu ]/2`$ while $`D_\mu `$ and $`𝒟_\mu `$ are the gauge and chiral covariant derivatives, respectively. Explicit expressions for the fields and the transformation properties can be found in . Here, $`z_0`$ is an off-shell parameter, which is not relevant in the present work .
In order to obtain proper chiral counting for the nucleon, we employ the conventional heavy baryon expansion of $`^{PC}`$, and in order to cosistently include the $`\mathrm{\Delta }`$ we follow the small scale expansion proposed in . In this approach energy-momenta and the delta and nucleon mass difference $`\delta `$ are both treated as small expansion parameters in chiral power counting. The leading order vertices in this framework can be obtained via $`P_+\mathrm{\Gamma }P_+`$ where $`\mathrm{\Gamma }`$ is the original vertex in the relativistic Lagrangian and
$$P_\pm =\frac{1\pm \overline{)}v}{2}.$$
(16)
are projection operators for the large, small components of the Dirac wavefunction respectively. We collect some of the relevant terms below:
$`_v^{PC}`$ $`=`$ $`\overline{N}[ivD+2g_ASA]Ni\overline{T}_i^\mu [ivD^{ij}\delta ^{ij}\delta +g_1SA^{ij}]T_\mu ^j`$ (18)
$`+g_{\pi N\mathrm{\Delta }}[\overline{T}_i^\mu \omega _\mu ^iN+\overline{N}\omega _\mu ^iT_i^\mu ]`$
where $`S_\mu `$ is the Pauli-Lubanski spin operator and $`\delta m_\mathrm{\Delta }m_N`$.
The PV analog of Eq. (15) can be constructed using the chiral fields $`X_{L,R}^a`$ defined as :
$$X_L^a=\xi ^{}\tau ^a\xi ,X_R^a=\xi \tau ^a\xi ^{},X_\pm ^a=X_L^a\pm X_R^a.$$
(19)
We find it convenient to follow the convention in Ref. and separate the PV Lagrangian into its various isospin components.
The hadronic weak interaction has the form
$$_W=\frac{G_\mu }{\sqrt{2}}J_\lambda J^\lambda +\text{h.c.},$$
(20)
where $`J_\lambda `$ denotes either a charged or neutral weak current built out of quarks. In the Standard Model, the strangeness conserving charged currents are pure isovector, whereas the neutral currents contain both isovector and isoscalar components. Consequently, $`_W`$ contains $`\mathrm{\Delta }I=0,1,2`$ pieces and these channels must all be accounted for in any realistic hadronic effective theory.
We quote here the relativistic Lagrangians, but employ the heavy baryon projections, as described above, in computing loops. It is straightforward to obtain the corresponding heavy baryon Lagrangians from those listed below, so we do not list the PV heavy baryon terms below. For the $`\pi N`$ sector we have
$`_{\mathrm{\Delta }I=0}^{\pi N}`$ $`=`$ $`h_V^0\overline{N}A_\mu \gamma ^\mu N`$ (21)
$`_{\mathrm{\Delta }I=1}^{\pi N}`$ $`=`$ $`{\displaystyle \frac{h_V^1}{2}}\overline{N}\gamma ^\mu NTr(A_\mu X_+^3){\displaystyle \frac{h_A^1}{2}}\overline{N}\gamma ^\mu \gamma _5NTr(A_\mu X_{}^3)`$ (24)
$`{\displaystyle \frac{h_\pi }{2\sqrt{2}}}F_\pi \overline{N}X_{}^3N`$
$`_{\mathrm{\Delta }I=2}^{\pi N}`$ $`=`$ $`h_V^2^{ab}\overline{N}[X_R^aA_\mu X_R^b+X_L^aA_\mu X_L^b]\gamma ^\mu N`$ (27)
$`{\displaystyle \frac{h_A^2}{2}}^{ab}\overline{N}[X_R^aA_\mu X_R^bX_L^aA_\mu X_L^b]\gamma ^\mu \gamma _5N,`$
where $`^{ab}`$ is a matrix coupling the $`X^{a,b}`$ to $`I=2,I_3=0`$. The above Lagrangian was first given by Kaplan and Savage (KS). However, the coefficients used in our work are slightly different from those of Ref. since our definition of $`A_\mu `$ differs by an overall phase.
The term proportional to $`h_\pi `$ contains no derivatives. At leading-order in $`1/F_\pi `$, it yields the PV $`NN\pi `$ Yukawa coupling traditionally used in meson-exchange models for the PV NN interaction . Unlike the PV Yukawa interaction, the vector and axial vector terms in Eqs. (21-27) contain derivative interactions. The terms containing $`h_A^1`$ and $`h_A^2`$ start off with $`NN\pi \pi `$ interactions, while all the other terms start off as $`NN\pi `$. Such derivative interactions have not been included in conventional analyses of nuclear and hadronic PV experiments. Consequently, the experimental constraints on the low-energy constants $`h_V^i`$, $`h_A^i`$ are unknown.
It is useful to list the first few terms obtained by expanding the Lagrangians in Eqs. (21-27) in $`1/F_\pi `$. For the present purposes, the following terms are needed:
$`_{\text{Yukawa}}^{\pi NN}=ih_\pi (\overline{p}n\pi ^+\overline{n}p\pi ^{})[1{\displaystyle \frac{1}{3F_\pi ^2}}(\pi ^+\pi ^{}+{\displaystyle \frac{1}{2}}\pi ^0\pi ^0)]`$ (28)
$`_V^{\pi NN}={\displaystyle \frac{h_V^0+4/3h_V^2}{\sqrt{2}F_\pi }}[\overline{p}\gamma ^\mu nD_\mu \pi ^++\overline{n}\gamma ^\mu pD_\mu \pi ^{}]`$ (29)
$`_A^{\pi NN}=i{\displaystyle \frac{h_A^1+h_A^2}{F_\pi ^2}}\overline{p}\gamma ^\mu \gamma _5p(\pi ^+D_\mu \pi ^{}\pi ^{}D_\mu \pi ^+)`$ (30)
$`+i{\displaystyle \frac{h_A^1h_A^2}{F_\pi ^2}}\overline{n}\gamma ^\mu \gamma _5n(\pi ^+D_\mu \pi ^{}\pi ^{}D_\mu \pi ^+)`$ (31)
$`+i{\displaystyle \frac{\sqrt{2}h_A^2}{F_\pi ^2}}\overline{p}\gamma ^\mu \gamma _5n\pi ^+D_\mu \pi ^0i{\displaystyle \frac{\sqrt{2}h_A^2}{F_\pi ^2}}\overline{p}\gamma ^\mu \gamma _5n\pi ^+D_\mu \pi ^0`$ $`.`$ (32)
For the PV $`\pi NN`$ Yukawa coupling we have also kept terms with three pions.
The corresponding PV Lagrangians involving a $`N\mathrm{\Delta }`$ transition are somewhat more complicated. We relegate the complete expressions to Appendix A, and give here only the leading terms required for our calculation. As noted in Ref. , the one-pion $`\pi N\mathrm{\Delta }`$ PV Lagrangian vanishes at leading order in the heavy baryon expansion. The two-pion terms are
$`_A^{\pi N\mathrm{\Delta }}={\displaystyle \frac{ih_A^{p\mathrm{\Delta }^{++}\pi ^{}\pi ^0}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^{++}D^\mu \pi ^{}\pi ^0{\displaystyle \frac{ih_A^{p\mathrm{\Delta }^{++}\pi ^0\pi ^{}}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^{++}D^\mu \pi ^0\pi ^{}`$ (33)
$`{\displaystyle \frac{ih_A^{p\mathrm{\Delta }^+\pi ^0\pi ^0}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^+D^\mu \pi ^0\pi ^0{\displaystyle \frac{ih_A^{p\mathrm{\Delta }^+\pi ^+\pi ^{}}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^+D^\mu \pi ^+\pi ^{}`$ (34)
$`{\displaystyle \frac{ih_A^{p\mathrm{\Delta }^+\pi ^{}\pi ^+}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^+D^\mu \pi ^{}\pi ^+{\displaystyle \frac{ih_A^{p\mathrm{\Delta }^0\pi ^+\pi ^0}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^0D^\mu \pi ^+\pi ^0`$ (35)
$`{\displaystyle \frac{ih_A^{p\mathrm{\Delta }^0\pi ^0\pi ^+}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^0D^\mu \pi ^0\pi ^+{\displaystyle \frac{ih_A^{p\mathrm{\Delta }^{}\pi ^+\pi ^+}}{F_\pi ^2}}\overline{p}\mathrm{\Delta }_\mu ^{}D^\mu \pi ^+\pi ^+`$ (36)
$`{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^{++}\pi ^{}\pi ^{}}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^{++}D^\mu \pi ^{}\pi ^{}{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^+\pi ^{}\pi ^0}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^+D^\mu \pi ^{}\pi ^0`$ (37)
$`{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^+\pi ^0\pi ^{}}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^+D^\mu \pi ^0\pi ^{}{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^0\pi ^0\pi ^0}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^0D^\mu \pi ^0\pi ^0`$ (38)
$`{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^0\pi ^+\pi ^{}}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^0D^\mu \pi ^+\pi ^{}{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^0\pi ^{}\pi ^+}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^0D^\mu \pi ^{}\pi ^+`$ (39)
$`{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^{}\pi ^+\pi ^0}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^{}D^\mu \pi ^+\pi ^0{\displaystyle \frac{ih_A^{n\mathrm{\Delta }^{}\pi ^0\pi ^+}}{F_\pi ^2}}\overline{n}\mathrm{\Delta }_\mu ^{}D^\mu \pi ^0\pi ^++\text{h.c.}`$ $`,`$ (40)
where the couplings $`h_A^{p\mathrm{\Delta }^{++}\pi ^{}\pi ^0}`$ are defined in terms of the various SU(2) PV low-energy constants in Appendix A.
The PV $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ Lagrangians, also listed in Appendix A, contain terms analogous to the Yukawa, $`V`$, and $`A`$ terms in Eqs. (21-27). Since we compute corrections up to one-loop order only, and since the initial and final states are nucleons, the PV $`\pi \pi \mathrm{\Delta }\mathrm{\Delta }`$ terms ($`A`$-type) are not relevant here. The leading, single-$`\pi `$ Yukawa and $`V`$-type interactions are
$`_{\text{Yukawa}}^{\pi \mathrm{\Delta }\mathrm{\Delta }}=i{\displaystyle \frac{h_\mathrm{\Delta }}{\sqrt{3}}}(\overline{\mathrm{\Delta }}^{++}\mathrm{\Delta }^+\pi ^+\overline{\mathrm{\Delta }}^+\mathrm{\Delta }^{++}\pi ^{})`$ (41)
$`i{\displaystyle \frac{h_\mathrm{\Delta }}{\sqrt{3}}}(\overline{\mathrm{\Delta }}^0\mathrm{\Delta }^{}\pi ^+\overline{\mathrm{\Delta }}^{}\mathrm{\Delta }^0\pi ^{})`$ (42)
$`i{\displaystyle \frac{2h_\mathrm{\Delta }}{3}}(\overline{\mathrm{\Delta }}^+\mathrm{\Delta }^0\pi ^+\overline{\mathrm{\Delta }}^0\mathrm{\Delta }^+\pi ^{})`$ (43)
$`_V^{\pi \mathrm{\Delta }\mathrm{\Delta }}={\displaystyle \frac{h_V^{\mathrm{\Delta }^{++}\mathrm{\Delta }^+}}{F_\pi }}(\overline{\mathrm{\Delta }}^{++}\gamma _\mu \mathrm{\Delta }^+D^\mu \pi ^++\overline{\mathrm{\Delta }}^+\gamma _\mu \mathrm{\Delta }^{++}D^\mu \pi ^{})`$ (44)
$`{\displaystyle \frac{h_V^{\mathrm{\Delta }^+\mathrm{\Delta }^0}}{F_\pi }}(\overline{\mathrm{\Delta }}^+\gamma _\mu \mathrm{\Delta }^0D^\mu \pi ^++\overline{\mathrm{\Delta }}^0\gamma _\mu \mathrm{\Delta }^+D^\mu \pi ^{})`$ (45)
$`{\displaystyle \frac{h_V^{\mathrm{\Delta }^0\mathrm{\Delta }^{}}}{F_\pi }}(\overline{\mathrm{\Delta }}^0\gamma _\mu \mathrm{\Delta }^{}D^\mu \pi ^++\overline{\mathrm{\Delta }}^{}\gamma _\mu \mathrm{\Delta }^0D^\mu \pi ^{})`$ (46)
where the coefficients are given in Appendix A.
One may ask whether there exist additional PV effective interactions that could contribute at the order to which we work. In the pionic sector there exists one CP-conserving, PV Lagrangian:
$$_\pi ^{PV}=ϵ_{ijk}\omega _\mu ^i\omega _\nu ^j(D^\mu \omega _k^\nu D^\nu \omega _k^\mu ).$$
(47)
At leading order in $`1/F_\pi `$, $`_\pi `$ contains five pions. Its lowest order contribution appears at two-loop order at best, so we do not consider it here.
Similarly, one may consider possible contributions from two-derivative operators. There exists one CP-conserving, PV operator:
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}\sigma ^{\mu \nu }[D_\mu A_\nu D_\nu A_\mu ]N.$$
(48)
There exist three independent PC, two-derivative operators . For example, one may choose the following three:
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}i\gamma _5D_\mu A^\mu N,$$
(49)
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}A^\mu A_\mu N,$$
(50)
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}\sigma ^{\mu \nu }[A_\mu ,A_\nu ]N.$$
(51)
As we discuss in Appendix B, none of the two-derivative operators in Eqs. (48) - (51) contribute to the renormalization of $`h_\pi `$ at the order to which we work in the present analysis.
## III The loop corrections
The leading order loop corrections to the Yukawa interaction of Eq. (28) are generated by the diagrams of Figs. 1-2. As we discuss in Appendix B, the contributions from many of the diagrams which nominally renormalize $`h_\pi `$ vanish at the order at which we truncate. In particular, none of the vector ($`V`$-type) $`\pi NN`$ and $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ terms contribute to this order. In what follows, we discuss only the non-vanishing Yukawa and $`A`$-type contributions. Details regarding the vanishing of the other contributions appear in Appendix B. Following the conventional practice, we regulate the loop integrals using dimensional regularization. The pole terms proportional to $`1/D4`$ are cancelled by appropriate counterterms. We identify only the terms nonanalytic in quark masses with the loops. All other analytic terms are indistinguishable from finite parts of the corresponding counterterms.
The nonvanishing contribution from Fig. 1(a) arises from the insertion of the $`3\pi `$ part of the Yukawa interaction of Eq. (24). The nonanalytic term is
$$iM_{(a)}=\frac{5}{6}\frac{m_\pi ^2}{\mathrm{\Lambda }_\chi ^2}\mathrm{ln}(\frac{\mu }{m_\pi })^2h_\pi \tau ^+,$$
(52)
where $`\mathrm{\Lambda }_\chi =4\pi F_\pi `$ and $`\mu `$ is the subtraction scale introduced in dimensional regularization. For simplicity, we show here only the contributions for $`np\pi ^{}`$. The terms for $`pn\pi ^+`$ are equal in magnitude and opposite in sign since it is the hermitian conjugate of the $`np\pi ^{}`$ piece. This property holds to all orders of chiral expansion.
The nonvanishing contribution from Fig. 1 (b) arises from strong vertex correction to the leading order $`\pi NN`$ Yukawa interaction:
$$iM_{(b)}=\frac{3}{4}g_A^2\frac{m_\pi ^2}{\mathrm{\Lambda }_\chi ^2}\mathrm{ln}(\frac{\mu }{m_\pi })^2h_\pi \tau ^+.$$
(53)
The terms in Figs. 1(c1)-(c2) are generated by the PV axial $`\pi \pi NN`$ couplings proportional to the $`h_A^i`$. We have
$$iM_{(c1)+(c2)}=2\sqrt{2}\pi g_A\frac{m_\pi ^3}{F_\pi \mathrm{\Lambda }_\chi ^2}h_A^1\tau ^+.$$
(54)
The contribution from $`h_A^2`$ to these two diagrams cancels out, leaving only the dependence on $`h_A^1`$. We note that although this term is propotional to $`m_q^{3/2}`$ and, thus, nominally suppressed, the coefficient of $`h_A^1`$ is fortuitously large ($`1/4`$). The two pion vertex in FIG. 1 (d1)-(d2) comes from the chiral connection $`V_\mu `$:
$$iM_{(d1)+(d2)}=\frac{1}{2}\frac{m_\pi ^2}{\mathrm{\Lambda }_\chi ^2}\mathrm{ln}(\frac{\mu }{m_\pi })^2h_\pi \tau ^+.$$
(55)
The leading contribution involving $`\mathrm{\Delta }`$ intermediate states arises from Fig. 2(a). The corresponding amplitude receives contributions from three different isospin combinations for the $`\mathrm{\Delta }`$ intermediate states. Their sum reads
$$iM_{2(a)}=\frac{20}{9}\frac{g_{\pi N\mathrm{\Delta }}^2h_\mathrm{\Delta }}{\mathrm{\Lambda }_\chi ^2}[(2\delta ^2m_\pi ^2)\mathrm{ln}(\frac{\mu }{m_\pi })^24\delta \sqrt{\delta ^2m_\pi ^2}\mathrm{ln}\frac{\delta +\sqrt{\delta ^2m_\pi ^2}}{m_\pi }]\tau ^+.$$
(56)
The corrections generated by the PV $`\pi \pi N\mathrm{\Delta }`$ vertices are
$$iM_{2(b1)+2(b2)}=\frac{2}{3}\frac{g_{\pi N\mathrm{\Delta }}}{F_\pi \mathrm{\Lambda }_\chi ^2}[(\delta ^2\frac{3}{2}m_\pi ^2)\delta \mathrm{ln}(\frac{\mu }{m_\pi })^22(\delta ^2m_\pi ^2)^{3/2}\mathrm{ln}\frac{\delta +\sqrt{\delta ^2m_\pi ^2}}{m_\pi }]h_A^\mathrm{\Delta }\tau ^+,$$
(57)
where $`h_A^\mathrm{\Delta }`$ is defined as
$$h_A^\mathrm{\Delta }=\frac{1}{\sqrt{3}}(h_A^{n\mathrm{\Delta }^0\pi ^+\pi ^{}}+h_A^{p\mathrm{\Delta }^+\pi ^{}\pi ^+})+\sqrt{\frac{2}{3}}(h_A^{n\mathrm{\Delta }^+\pi ^0\pi ^{}}h_A^{p\mathrm{\Delta }^0\pi ^0\pi ^+})h_A^{n\mathrm{\Delta }^{++}\pi ^{}\pi ^{}}h_A^{p\mathrm{\Delta }^{}\pi ^+\pi ^+}.$$
(58)
Summing all the non-vanishing loop contributions yields the following expression for $`\mathrm{\Delta }h_\pi `$:
$`\mathrm{\Delta }h_\pi `$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle \frac{m_\pi ^2}{\mathrm{\Lambda }_\chi ^2}}\mathrm{ln}({\displaystyle \frac{\mu }{m_\pi }})^2h_\pi +{\displaystyle \frac{3}{4}}g_A^2{\displaystyle \frac{m_\pi ^2}{\mathrm{\Lambda }_\chi ^2}}\mathrm{ln}({\displaystyle \frac{\mu }{m_\pi }})^2h_\pi +2\sqrt{2}\pi g_A{\displaystyle \frac{m_\pi ^3}{F_\pi \mathrm{\Lambda }_\chi ^2}}h_A^1`$ (61)
$`{\displaystyle \frac{20}{9}}{\displaystyle \frac{g_{\pi N\mathrm{\Delta }}^2h_\mathrm{\Delta }}{\mathrm{\Lambda }_\chi ^2}}[(2\delta ^2m_\pi ^2)\mathrm{ln}({\displaystyle \frac{\mu }{m_\pi }})^24\delta \sqrt{\delta ^2m_\pi ^2}\mathrm{ln}{\displaystyle \frac{\delta +\sqrt{\delta ^2m_\pi ^2}}{m_\pi }}]`$
$`+{\displaystyle \frac{2}{3}}{\displaystyle \frac{g_{\pi N\mathrm{\Delta }}}{F_\pi \mathrm{\Lambda }_\chi ^2}}[(\delta ^2{\displaystyle \frac{3}{2}}m_\pi ^2)\delta \mathrm{ln}({\displaystyle \frac{\mu }{m_\pi }})^22(\delta ^2m_\pi ^2)^{3/2}\mathrm{ln}{\displaystyle \frac{\delta +\sqrt{\delta ^2m_\pi ^2}}{m_\pi }}]h_A^\mathrm{\Delta }`$
The final nonvanishing corrections arise from $`N`$ and $`\pi `$ wavefunction renormalization. These corrections, which have been computed previously , generate deviations from unity of $`Z_N`$ and $`\sqrt{Z_\pi }`$ appearing in the expression for $`h_\pi ^{EFF}`$ in Eq. (4). In the case of $`Z_N`$, the nonvanishing contributions arise from Figs. 1(e1)-(e2) and 2(c1)-(c2):
$`Z_N1={\displaystyle \frac{9}{4}}g_A^2{\displaystyle \frac{m_\pi ^2}{\mathrm{\Lambda }_\chi ^2}}\mathrm{ln}({\displaystyle \frac{\mu }{m_\pi }})^24g_{\pi N\mathrm{\Delta }}^2[{\displaystyle \frac{2\delta ^2m_\pi ^2}{\mathrm{\Lambda }_\chi ^2}}\mathrm{ln}({\displaystyle \frac{\mu }{m_\pi }})^2`$ (62)
$`4{\displaystyle \frac{\delta \sqrt{\delta ^2m_\pi ^2}}{\mathrm{\Lambda }_\chi ^2}}\mathrm{ln}{\displaystyle \frac{\delta +\sqrt{\delta ^2m_\pi ^2}}{m_\pi }}].`$ (63)
The pion’s wavefunction renormalization arises from Fig. 2 (k) :
$$\sqrt{Z_\pi }1=\frac{1}{3}\left(\frac{m_\pi }{\mathrm{\Lambda }_\chi }\right)^2\mathrm{ln}\left(\frac{\mu }{m_\pi }\right)^2.$$
(64)
Numercially, the loop contributions to $`\sqrt{Z_\pi }`$ are small compared to those entering $`Z_N`$.
Note the one loop renormalization of $`h_\pi `$ from the PV Yukawa $`\pi NN`$ and $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ vertices is already at the order $`1/\mathrm{\Lambda }_\chi ^2`$. An additional loop will introduce a factor of $`1/\mathrm{\Lambda }_\chi ^2`$. Loops containing the axial vector $`NN\pi \pi `$ and $`N\mathrm{\Delta }\pi \pi `$ vertices and one strong $`NN\pi `$ or $`N\mathrm{\Delta }\pi `$ vertex are of $`𝒪(1/\mathrm{\Lambda }_\chi ^2F_\pi )`$. To obtain contributions of $`𝒪(1/\mathrm{\Lambda }_\chi ^3)`$, one would require the insertion of operators carrying explicit factors of $`1/\mathrm{\Lambda }_\chi `$ into one loop graphs. We find no such contributions.
## IV The scale of loop corrections
We may estimate the numerical importance of the loop corrections to $`h_\pi ^{BARE}`$ by taking $`\delta =0.3`$ GeV, $`g_A=1.267`$ and $`g_{\pi N\mathrm{\Delta }}=1.05`$ and by choosing $`\mu =\mathrm{\Lambda }_\chi =1.16`$ GeVSince the dependence on $`\mu `$ is logarithmic, one may choose other values, such as $`\mu =m_\rho `$, without affecting the numerical results significantly. With these inputs, the value of $`Z_N\sqrt{Z_\pi }`$ is completely determined. The vertex corrections, which appear as $`\mathrm{\Delta }h_\pi `$ in Eq. (4), depend on the PV couplings $`h_\pi `$, $`h_A^1`$, $`h_\mathrm{\Delta }`$, and $`h_A^\mathrm{\Delta }`$. We obtain
$$h_\pi ^{EFF}=0.5h_\pi +0.25h_A^10.24h_\mathrm{\Delta }+0.079h_A^\mathrm{\Delta }.$$
(65)
Note that the effect of the wavefunction renormalization corrections is to reduce the dependence on $`h_\pi ^{BARE}`$ by roughly 50%. In addition, the dependence of $`h_\pi ^{EFF}`$ on $`h_A^1`$ and $`h_\mathrm{\Delta }`$ is non-negligible. Their coefficients are only a factor of two smaller than that of $`h_\pi ^{BARE}`$. Although these contributions arise at $`𝒪(p^2,p^3)`$, they are fortuitously enhanced numerically. Thus, in a complete anaysis of the OPE PV interaction one should not ignore these constants.
At present, one has no direct experimental constraints on the parameters $`h_A^1`$, $`h_\mathrm{\Delta }`$, and $`h_A^\mathrm{\Delta }`$, as a comprehensive analysis of hadronic PV data including the full chiral structure of the PV hadronic interaction has yet to be performed. Consequently, one must rely on theoretical input for guidance regarding the scale of the unknown constants. Estimates of $`h_A^1`$ are given by the authors of Ref. . These authors observe that the usual pole dominance approximation for P-wave non-leptonic hyperon decays typically underpredicts the experimental amplitudes by a factor of two. The difference may be resolved by the inclusion of local, parity-conserving operators having structures analogous to the $`A`$-type terms in Eq. (24). The requisite size of the $`\mathrm{\Delta }S=1`$ contact terms may imply a scale for the analogous $`\mathrm{\Delta }I=1`$ PV terms. If so, one might conclude that $`h_A^1`$ should be on the order of $`10g_\pi `$. On the other hand, a simple factorization estimate leads to $`h_A^10.2g_\pi `$. While the sign of $`h_A^1`$ is fixed in the factorization approximation, the sign of the larger value is undetermined. Thus, it is reasonable to conclude that $`h_A^1`$ may be large enough to significantly impact $`h_\pi ^{EFF}`$, though considerably more analysis is needed to yield a firm conclusion.
The $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ Yukawa coupling $`h_\mathrm{\Delta }`$ has been estimated in Ref. using methods similar to those of Ref. . The authors quote a “best value” of $`h_\mathrm{\Delta }=20g_\pi `$, with a “reasonable range” of $`(510)\times g_\pi `$.This coupling is denoted $`f_{\mathrm{\Delta }\mathrm{\Delta }\pi }`$ in Ref. . Naïvely, subsitution of the best value into Eq. (65) would increase the value of $`h_\pi ^{EFF}`$, whereas the <sup>18</sup>F result would seem to require a reduction. As we argue below, however, the relationship between the couplings computed in Refs. and the parameters appearing in Eq. (65) is somewhat ambiguous. Direct substitution of the theoretical value into $`h_\pi ^{EFF}`$ may not be entirely appropriate.
To date, no theoretical estimate of the $`A`$-type $`\pi \pi N\mathrm{\Delta }`$ coupling has been performed. A simple estimate of the scale is readily obtained using the factorization approximation. To that end, we work with tree-level form of $`_W^{PV}`$. Neglecting short-distance QCD corrections and terms containing strange quarks, one has
$`_W^{PV}(\mathrm{\Delta }S=0)`$ $`=`$ $`{\displaystyle \frac{G_F}{\sqrt{2}}}\{\mathrm{cos}^2\theta _c\overline{u}\gamma _\lambda (1\gamma _5)d\overline{d}\gamma ^\lambda (1\gamma _5)u`$ (67)
$`2(12\mathrm{sin}^2\theta _W)V_\lambda ^{(3)}A^{(3)\lambda }+{\displaystyle \frac{4}{3}}\mathrm{sin}^2\theta _WV_\lambda ^{(0)}A^{(3)\lambda }\},`$
where $`V_\lambda ^{(3)}`$ and $`A_\lambda ^{(3)}`$ denote the third components of the octet of vector and axial vector currents, respectively, and
$$V_\lambda ^{(0)}=\frac{1}{2}(\overline{u}\gamma _\lambda u+\overline{d}\gamma _\lambda d).$$
(68)
Consider now the first term in the expression for $`h_A^\mathrm{\Delta }`$ given in Eq. (58). In the factorization approximation, $`_W^{PV}`$ contributes only to the antisymmetric combination
$$\frac{1}{2}(h_A^{n\mathrm{\Delta }^0\pi ^+\pi ^{}}h_A^{n\mathrm{\Delta }^0\pi ^{}\pi ^+}).$$
(69)
The neutral current contribution to this combination, which arises only from the term containing $`V_\lambda ^{(3)}`$, is
$$\sqrt{2}G_FF_\pi ^2(12\mathrm{sin}^2\theta _W)C_5^A(n\mathrm{\Delta }^0)2g_\pi C_5^A(n\mathrm{\Delta }^0),$$
(70)
where $`C_5^A(n\mathrm{\Delta }^0)𝒪(1)`$ is the axial vector $`n\mathrm{\Delta }^0`$ form factor at the photon point. After Fierz re-ordering, the charged current component of $`_W^{PV}`$ contributes roughly
$$(4g_\pi /3)C_5^A(n\mathrm{\Delta }^0),$$
(71)
yielding a total factorzation contribution of about $`(2g_\pi /3)C_5^A(n\mathrm{\Delta }^0)`$. Thus, one would expect the scale of the axial vector $`\pi \pi N\mathrm{\Delta }`$ couplings to be on the order of a few $`\times g_\pi `$.
In the particular case of the combination appearing in $`h_A^\mathrm{\Delta }`$, however, the sum of factorization contributions cancels identically. As one sees from the expressions for the $`h_A^{N\mathrm{\Delta }\pi \pi }`$ given in Appendix A, isospin requires
$$h_A^{n\mathrm{\Delta }^0\pi ^+\pi ^{}}+h_A^{p\mathrm{\Delta }^+\pi ^{}\pi ^+}=0.$$
(72)
The factorization contributions independently satisfy this sum rule. The second combination of constants appearing in Eq. (58),
$$h_A^{n\mathrm{\Delta }^+\pi ^0\pi ^{}}h_A^{p\mathrm{\Delta }^0\pi ^0\pi ^+},$$
(73)
also vanishes in the factorization approximation, even though the individual couplings do not. The third pair of couplings received no factorization contributions. Thus, one has $`h_A^\mathrm{\Delta }=0`$ in this approximation. In principle, non-factorization contributions yield a non-zero value for $`h_A^\mathrm{\Delta }`$. Although we have not evaluated these contributions, we do not expect the scale to be significantly larger than the factorization value for the individual $`h_A^{N\mathrm{\Delta }\pi \pi }`$ couplings. Consequently, we estimate a reasonable range for $`h_A^\mathrm{\Delta }`$ of $`(0\text{few})\times g_\pi `$.
These theoretical estimates suggest considerable ambiguity in the prediction for $`h_\pi ^{EFF}`$. In principle, some of this ambiguity might be removed by performing the comprehensive analysis of hadronic PV suggested above, in which the various constants would be determined entirely by experiment. The viability of such a program remains to be seen.
## V Comparing with Microscopic Calculations
The results in Eqs. (61-64) embody the full SU(2) chiral structure at $`𝒪(p^3)`$ of $`N\pi |_W^{PV}|N`$ at leading order in the pion momentum. Any microscopic calculation of this matrix element which respects the symmetries of QCD should display the dependence on light quark masses appearing in $`h_\pi ^{EFF}`$. In principle, an unquenched lattice QCD calculation with light quarks would manifest this chiral structure. In practice, however, unquenched calculations remain difficult, and even quenched calculations require the use of heavy quarks. For a lattice determination of $`N\pi |_W^{PV}|N`$, the expressions in Eqs. (61-64) could be used to extrapolate to the light quark limit, much as the chiral structure of baryon mass and magnetic moment can be used for similar extrapolations .
In the absence of a first principles QCD calculation, one must rely on symmetries and/or models to obtain the PV $`NN\pi `$ coupling. A variety of such approaches have been undertaken, including the SU(6)<sub>w</sub>/quark model calculation of Refs. , the Skyrme model , and QCD sum rules . To date, the DDH/FCDH analysis remains the most comprehensive and has become the benchmark for comparison between experiment and theory. Consequently, we focus on this work as a “case study” in the problem of matching microscopic calculations onto hadronic effective theory.
The DDH/FCDH approach relies heavily on symmetry methods in order to relate the PV $`\mathrm{\Delta }S=0`$ matrix elements to experimental $`\mathrm{\Delta }S=1`$ nonleptonic hyperon decay amplitudes. All the charged current (CC) contributions to the $`\mathrm{\Delta }S=0,1`$ $`BB^{}M`$ amplitudes, where $`M`$ is a pseudoscalar meson, can be related using SU(3) arguments. Likewise, the neutral current (NC) component of the effective weak Hamiltonian belonging to the same multiplets as the CC components (i.e. those arising from a product of purely left-handed currents) can also be related via SU(3). The remaining NC contributions to the $`\mathrm{\Delta }S=0`$ PV amplitudes are computed using factorization and the MIT bag model. The DDH approach also employs SU(6)<sub>w</sub> symmetry arguments in order to calculate parity-violating vector meson couplings. Although one requires only SU(3) to determine the pseudoscalar couplings, we refer below to the general SU(6)<sub>w</sub> formalism used in Refs. .
The general SU(6)<sub>w</sub> analysis employed by DDH/FCDH introduces five reduced matrix elements: $`a_{t,v}`$, $`b_{t,v}`$, and $`c_v`$. These constants correspond to SU(6)<sub>w</sub> components of the weak Hamiltonian:
$$[(\overline{B}B)_{35}M_{35}]_{35}c_v$$
(74)
$$[(\overline{B}B)_{405}M_{35}]_{280,\overline{280}}b_t,b_v$$
(75)
$$[(\overline{B}B)_{405}M_{35}]_{280,\overline{280}}a_t,a_v$$
(76)
One may represent these different components of $`_W^{PV}`$ diagramatically as in Fig. 3. The components shown in Fig. 3a,b correspond to $`b_{t,v}`$ and $`c_v`$, respectively. In practice, these contributions are determined entirely from empirical hyperon decay data. The term in Fig. 3a corresponds to $`a_{t,v}`$ and is computed in Refs. using factorization.
The PV $`NN\pi `$ Yukawa coupling can be expressed in terms of these SU(6)<sub>w</sub> reduced matrix elements plus an additional factorization/quark model term. Temporarily neglecting short-distance QCD corrections to $`_W^{PV}`$, one has
$`p\pi ^{}|_W^{PV}|n`$ $`=`$ $`{\displaystyle \frac{1}{3\sqrt{2}}}\mathrm{tan}\theta _cc_v`$ (77)
$``$ $`{\displaystyle \frac{2}{9\sqrt{2}}}\mathrm{csc}2\theta _c\mathrm{sin}^2\theta _W(2c_vb_t)+{\displaystyle \frac{1}{3}}\mathrm{sin}^2\theta _Wy,`$ (78)
where $`\theta _c`$ and $`\theta _W`$ are the Cabibbo and Weinberg angles, respectively, and $`y`$ denotes a Fierz/factorization contribution. The first term on the RHS of Eq. (77) gives the CC contribution, while the remaining terms arise from weak NC’s. Including short-distance QCD renormalization of $`_W^{PV}`$ leads to a modification of Eq. (77):
$`p\pi ^{}|_W^{PV}|n`$ $`=`$ $`\left\{[12\mathrm{sin}^2\theta _W]\gamma (K)+\mathrm{sin}^2\theta _c\right\}{\displaystyle \frac{\rho }{\mathrm{sin}^2\theta _c}}g_\pi `$ (80)
$`+\mathrm{sin}^2\theta _c(B_1+B_2),`$
where
$`g_\pi `$ $`=`$ $`{\displaystyle \frac{1}{3\sqrt{2}}}\mathrm{tan}\theta _cc_v`$ (81)
$`B_1`$ $`=`$ $`{\displaystyle \frac{4}{9\sqrt{2}}}\eta E(K)\left({\displaystyle \frac{1}{\mathrm{sin}\theta _c\mathrm{cos}\theta _c}}\right)(b_v/6b_t/12c_v/2)`$ (82)
$`B_2`$ $`=`$ $`{\displaystyle \frac{1}{3}}F(K)y,`$ (83)
and $`\gamma (K)`$, $`E(K)`$ and $`F(K)`$ are summed leading log (renormalization group) factors dependent on
$$K=1\frac{\alpha _s(\mu )}{\pi }[11\frac{2}{3}N_f]\mathrm{ln}\frac{M_W^2}{\mu ^2}.$$
(84)
The overall scale factor $`\rho `$ appearing in Eq. (80) was introduced in Ref. in order to account for various theoretical uncertainties entering the analysis.
The appearance of $`c_v`$, $`b_t`$, and $`b_v`$ in $`g_\pi `$ and $`B_1`$ relies on tree-level SU(6)<sub>w</sub> symmetry—long-distance chiral corrections of the types shown in Fig. 4 have not been explicitly included. Inclusion of such corrections would necessitate a reanalysis of the $`\mathrm{\Delta }S=1`$ amplitudes in much the same way that one treats the octet of baryon axial vector currents or magnetic moments . For example, letting $`A(\mathrm{\Lambda }_{}^0)`$ denote the amplitude for $`\mathrm{\Lambda }p\pi ^{}`$ one has at tree-level
$$A(\mathrm{\Lambda }_{}^0)=\frac{1}{\sqrt{3}}(b_v/6b_t/12c_v/2).$$
(85)
Including the leading chiral corrections would yield the modification
$$A(\mathrm{\Lambda }_{}^0)=\frac{1}{\sqrt{3}}\sqrt{Z_\mathrm{\Lambda }Z_pZ_\pi }(b_v/6b_t/12c_v/2)+\mathrm{\Delta }A(\mathrm{\Lambda }_{}^0),$$
(86)
where $`\mathrm{\Delta }A(\mathrm{\Lambda }_{}^0)`$ denotes vertex corrections and possible contributions from higher-dimension operators. Similar corrections would appear in the SU(6)<sub>w</sub> symmetry terms in Eqs. (77, 80). Given the absence of these corrections from the DDH/FCDH analysis, the symmetry components $`p\pi ^{}|_W^{PV}|n`$ do not formally embody the subleading chiral structure of $`h_\pi ^{EFF}`$. The numerical impact of applying chiral corrections to the DDH/FCDH tree-level SU(6)<sub>w</sub> analysis is much less clear, since some of the chiral modifications can be absorbed into renormalized values of the chiral couplings, which are determined empirically. Nevertheless, the potentially sizeable effect of the SU(2) chiral corrections on $`h_\pi ^{EFF}`$ should give one pause.
A related issue is the degree to which ambiguities introduced by kaon and $`\eta `$ loops in SU(3) HBChPT could plague an analysis of the $`\mathrm{\Delta }S=1`$ amplitudes. Here recent work by Donoghue and Holstein argues that finite nucleon size call for long-distance regularization of such heavy meson loops, which substantially reduces their effects. Results are then similar to what arises from use of a cloudy bag approach to such matrix elements. A comprehensive study of such issues – and their impact on the DDH/FCDH calculation of $`h_\pi `$ – goes beyond the scope of the present work. Nevertheless, the sizeable impact of the chiral corrections in $`h_\pi ^{EFF}`$ and the use of tree-level symmetry arguments in Refs. points to a possibly significant mismatch between $`h_\pi ^{EFF}`$ and $`h_\pi ^{DDH}`$.
The remaining terms in the DDH/FCDH analysis – involving the parameters $`\eta `$ and $`y`$ – are determined by explicit MIT bag model calculations. One may ask whether the latter effectively includes any part of the subleading chiral structure of $`h_\pi ^{EFF}`$. In order to address this question, we make three observations:
1. Sea quarks and gluons generate $`c_v`$. The parameter $`c_v`$ vanishes identically in any quark model in which baryons consist solely of three constituent quarks. The $`\mathrm{\Delta }S=1`$ hyperon decay data, however, clearly implies that $`c_v0`$. In order to obtain a nonzero value in a quark model, one requires the presence of sea quarks and gluons. It is shown in for example, that $`c_v0`$ when gluons are added to the MIT bag model. Similarly, one would expect contributions from the $`q\overline{q}`$ pairs in the sea. Since relativistic quark models already contain $`q\overline{q}`$ pairs in the form of “Z-graphs” , it is likely that disconnected $`q\overline{q}`$ insertions (see Fig. 5b) give the dominant sea quark contribution to $`c_v`$. In a chirally corrected analysis of nonleptonic decays, the long-distance part of the disconnected $`q\overline{q}`$ insertions appear explicitly in the guise of pseudoscalar loops, while the short-distance contributions are subsumed into the value of $`c_v`$ and possible higher dimension operators. “Quenched” quark models without explicit pionic degress of freedom generally do not include the long-distance physics of disconnected insertions.
2. The $`m_q`$-dependence is different. In conventional HBChPT analyses of hadronic observables, one only retains the loop contributions non-analytic in the light quark mass. The constituent quark model (without explicit pions) has a difficult time producing these non-analytic contributions. The simplest, illustrative example is the nucleon isovector charge radius, $`r^2_{T=1}`$, which is singular in the chiral limit . This chiral singularity, of the form $`\mathrm{ln}m_\pi ^2\mathrm{ln}m_q`$, is produced by $`\pi `$ loops. Relativistic quark models, such as the MIT bag model, yield a finite value for $`r^2_{T=1}`$ as $`m_q0`$. One cannot produce the chiral singularity in a quark model without including disconnected $`q\overline{q}`$ insertions dressed as mesons.
The corresponding argument in the case of $`h_\pi ^{EFF}`$ is less direct, but still straightforward. In the limit of a degenerate $`N`$ and $`\mathrm{\Delta }`$, the non-analytic terms in $`h_\pi ^{EFF}`$ have quark mass-dependences of the form $`m_q\mathrm{ln}m_q`$ or $`m_q^{3/2}`$. As we show in Appendix C, bag model matrix elements of the four quark operators appearing in $`_W^{PV}`$ have a Taylor series expansion about $`m_q=0`$. Thus, the parameters $`\eta `$ and $`y`$ cannot contain the non-analytic structures generated by the diagrams in Figs. 1-2.
3. Graphs are missing. This observation is simply a diagrammatic summary of the previous two observations. For simplicity, consider a subset of the quark-level diagrams associated with the appearance of $`h_A^i`$ in $`h_\pi ^{EFF}`$. Typical contributions to the axial $`NN\pi \pi `$ PV vertex are shown in Fig. 5a. The corresponding loop contributions to $`h_\pi ^{EFF}`$ appear in Fig. 5b,c. Those in Fig. 5b involve disconnected $`q\overline{q}`$ insertions, which do not occur in the constituent quark model. The contribution of Fig. 5c involves Z-graphs, which are produced in a relativistic quark model<sup>§</sup><sup>§</sup>§e.g., as a correction to the $`b_{t,v}`$ terms of Fig. 3b. In principle, the $`3q+q\overline{q}`$ intermediate state could contain an $`N\pi `$ pair. As argued previously, however, the Z-graphs implicit in the MIT bag model calculation of $`h_\pi `$ do not produce the nonanalytic structure of the corresponding $`\pi `$ loop. Apparently, only an unquenched quark model, which generates the disconnected insertions of Fig. 5b, could produce the requisite nonanalytic terms.
From this “case study” of the DDH/FCDH calculation of $`h_\pi `$, we conclude that the SU(6)<sub>w</sub>/quark model approach used in Refs. does not incorporate the chiral structure of $`h_\pi ^{EFF}`$. Were the numerical impact of the chiral corrections negligible, this observation would not be bothersome. The actual impact of the chiral corrections, however, is significant. Thus, it is perhaps not surprising that the <sup>18</sup>F result and the DDH/FCDH “best value” do not agree.
## VI Conclusions
With the confirmation of the electroweak sector of the Standard Model at the 1% level or better in a variety of leptonic and semi-leptonic processes, one has little reason to doubt its validity in the purely hadronic domain. Similarly, the predictions of QCD in the perturbative regime have been confirmed with a high degree of confidence. Thus, one may justifiably consider $`_W^{PV}`$, the effective Hamiltonian including its perturbative strong interaction correction, to be well understood. Moreover, the precision available with present and future hadronic PV experiments is unlikely to match the levels achieved in leptonic and semileptonic processes. Consequently, one has little hope of detecting small deviations in $`_W^{PV}`$ from its SM structure due to “new physics”. On the other hand, much about QCD in the non-perturbative regime remains mysterious: the mechanism of confinement, the dynamics of chiral symmetry breaking, the role or sea quarks in the low-energy structure of the nucleon, and so forth. Each of these issues bears on one’s understanding of matrix elements of $`_W^{PV}`$. In this sense, the low-energy, PV hadronic weak interaction constitutes a probe of the dynamics of low-energy QCD, in a manner analogous to the probe provided by the electromagnetic interaction.
From a phenomenological standpoint, the matrix element one may hope to extract from hadronic PV observables with the least ambiguity is $`N\pi |_W^{PV}|N`$. In this study, we have argued that any theoretical interpretation of this matrix element must take into account the consequences of chiral symmetry. Indeed the chiral corrections to the tree-level, PV $`\pi NN`$ Yukawa coupling are not small. At $`𝒪(1/\mathrm{\Lambda }_\chi ^3)`$, the effective coupling measured in experiments, $`h_\pi ^{EFF}`$, depends not only on the bare coupling, $`h_\pi ^{BARE}`$, but also on new (and experimentally undetermined) PV low-energy constants, $`h_A^1`$, $`h_A^\mathrm{\Delta }`$, and $`h_\mathrm{\Delta }`$, as well. Furthermore, the coefficients of $`h_\pi ^{BARE}`$, $`h_A^1`$, and $`h_\mathrm{\Delta }`$ are comparable in magnitude. At present, one has only simple theoretical estimates of the magnitudes of the $`h_A^1`$ and $`h_A^\mathrm{\Delta }`$ in addition to the FCDH calculation of $`h_\mathrm{\Delta }`$. These estimates suggest that the new PV couplings appearing in $`h_\pi ^{EFF}`$ could be as large as $`h_\pi ^{BARE}`$. Since no experimental constraints have been obtained for the new couplings, there exists considerable latitude in the theoretical expectation for $`h_\pi ^{EFF}`$.
For two decades now, the benchmark theoretical calculation of $`N\pi |_W^{PV}|N`$ has been the SU(6)<sub>w</sub>/quark model approach of Ref. , updated in Ref. . We have argued, however, that the DDH/FCDH calculation does not manifest the general strictures of chiral invariance obtained in the present analysis. At the quark level, this chiral structure reflects the importance of the “disconnected” $`q\overline{q}`$ components of the sea. While relativistic quark models contain $`q\overline{q}`$ sea quark effects in the guise of Z-graphs or lower-component wavefunctions, the most common “quenched” versions do not include explicit disconnected pairsSome effects of disconnected $`q\overline{q}`$ pairs may, however, hide in the effective parameters of the quark model, such as the string tension .. Given the size of the chiral corrections associated in part with the disconnected insertions, it may then be not so surprising to find a possible discrepancy between the experimental value for $`h_\pi ^{EFF}`$ and the DDH/FCDH “best value”.
Applying chiral corrections to the SU(3) analysis of $`\mathrm{\Delta }S=1`$ hyperon decays may help to close the gap between $`h_\pi ^{EFF}`$ and $`h_\pi ^{DDH}`$. Presumably, similar corrections should be applied in other approaches not containing explicit pionic degress of freedom. In the longer run, one may be able to use the chiral structure of $`h_\pi ^{EFF}`$ to extrapolate an unquenched lattice calculation with heavy quarks into the physical regime.
## Acknowledgement
It is a pleasure to thank J.L. Goity and N. Isgur for useful discussions. This work was supported in part under U.S. Department of Energy contract #DE-AC05-84ER40150, the National Science Foundation, and a National Science Foundation Young Investigator Award.
## A PV Lagrangians
Here we present the full expressions for some of the PV Lagrangians not included in the main body of the paper. The analogues of Eqs. (21-27) are
$`_{\mathrm{\Delta }I=0}^{\pi \mathrm{\Delta }N}`$ $`=`$ $`f_1ϵ^{abc}\overline{N}i\gamma _5[X_L^aA_\mu X_L^b+X_R^aA_\mu X_R^b]T_c^\mu `$ (A2)
$`+g_1\overline{N}[A_\mu ,X_{}^a]_+T_a^\mu +g_2\overline{N}[A_\mu ,X_{}^a]_{}T_a^\mu +\text{h.c.}`$
$`_{\mathrm{\Delta }I=1}^{\pi \mathrm{\Delta }N}`$ $`=`$ $`f_2ϵ^{ab3}\overline{N}i\gamma _5[A_\mu ,X_+^a]_+T_b^\mu +f_3ϵ^{ab3}\overline{N}i\gamma _5[A_\mu ,X_+^a]_{}T_b^\mu `$ (A7)
$`+{\displaystyle \frac{g_3}{2}}\overline{N}[(X_L^aA_\mu X_L^3X_L^3A_\mu X_L^a)(X_R^aA_\mu X_R^3X_R^3A_\mu X_R^a)]T_a^\mu `$
$`+{\displaystyle \frac{g_4}{2}}\{\overline{N}[3X_L^3A^\mu (X_L^1T_\mu ^1+X_L^2T_\mu ^2)+3(X_L^1A^\mu X_L^3T_\mu ^1+X_L^2A^\mu X_L^3T_\mu ^2)`$
$`2(X_L^1A^\mu X_L^1+X_L^2A^\mu X_L^22X_L^3A^\mu X_L^3)T_\mu ^3](LR)\}+\text{h.c.}`$
$`_{\mathrm{\Delta }I=2}^{\pi \mathrm{\Delta }N}`$ $`=`$ $`f_4ϵ^{abd}^{cd}\overline{N}i\gamma _5[X_L^aA_\mu X_L^b+X_R^aA_\mu X_R^b]T_c^\mu `$ (A11)
$`+f_5ϵ^{ab3}\overline{N}i\gamma _5[X_L^aA_\mu X_L^3+X_L^3A_\mu X_L^a+(LR)]T_b^\mu `$
$`+g_5^{ab}\overline{N}[A_\mu ,X_{}^a]_+T_b^\mu +g_6^{ab}\overline{N}[A_\mu ,X_{}^a]_{}T_b^\mu +\text{h.c.},`$
where the terms containing $`f_i`$ and $`g_i`$ start off with one- and two-pion vertices, respectively. In the heavy baryon expansion, the terms containing the $`f_i`$ start to contribute at $`𝒪(1/m_N)`$. The leading order term vanishes since $`P_+i\gamma _5P_+=0`$. Since we work only to lowest order in the $`1/m_N`$ expansion, we obtain no contribution from the terms containing the $`f_i`$.
For the pv $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ effective Lagrangians we have
$$_{\mathrm{\Delta }I=0}^{\pi \mathrm{\Delta }}=j_0\overline{T}^iA_\mu \gamma ^\mu T_i,$$
(A12)
$`_{\mathrm{\Delta }I=1}^{\pi \mathrm{\Delta }}={\displaystyle \frac{j_1}{2}}\overline{T}^i\gamma ^\mu T_iTr(A_\mu X_+^3){\displaystyle \frac{k_1}{2}}\overline{T}^i\gamma ^\mu \gamma _5T_iTr(A_\mu X_{}^3)`$ (A13)
$`{\displaystyle \frac{h_{\pi \mathrm{\Delta }}^1}{2\sqrt{2}}}f_\pi \overline{T}^iX_{}^3T_i{\displaystyle \frac{h_{\pi \mathrm{\Delta }}^2}{2\sqrt{2}}}f_\pi \{3T^3(X_{}^1T^1+X_{}^2T^2)+3(\overline{T}^1X_{}^1+\overline{T}^2X_{}^2)T^3`$ (A14)
$`2(\overline{T}^1X_{}^3T^1+\overline{T}^2X_{}^3T^22\overline{T}^3X_{}^3T^3)\}+j_2\{3[(\overline{T}^3\gamma ^\mu T^1+\overline{T}^1\gamma ^\mu T^3)Tr(A_\mu X_+^1)`$ (A15)
$`+(\overline{T}^3\gamma ^\mu T^2+\overline{T}^2\gamma ^\mu T^3)Tr(A_\mu X_+^2)]2(\overline{T}^1\gamma ^\mu T^1+\overline{T}^2\gamma ^\mu T^22\overline{T}^3\gamma ^\mu T^3)Tr(A_\mu X_+^3)\}`$ (A16)
$`+k_2\{3[(\overline{T}^3\gamma ^\mu \gamma _5T^1+\overline{T}^1\gamma ^\mu \gamma _5T^3)Tr(A_\mu X_{}^1)+(\overline{T}^3\gamma ^\mu \gamma _5T^2+\overline{T}^2\gamma ^\mu \gamma _5T^3)Tr(A_\mu X_{}^2)]`$ (A17)
$`2(\overline{T}^1\gamma ^\mu \gamma _5T^1+\overline{T}^2\gamma ^\mu \gamma _5T^22\overline{T}^3\gamma ^\mu \gamma _5T^3)Tr(A_\mu X_{}^3)\}`$ (A18)
$`+j_3\{\overline{T}^a\gamma ^\mu [A_\mu ,X_+^a]_+T^3+\overline{T}^3\gamma ^\mu [A_\mu ,X_+^a]_+T^a\}`$ (A19)
$`+j_4\{\overline{T}^a\gamma ^\mu [A_\mu ,X_+^a]_{}T^3\overline{T}^3\gamma ^\mu [A_\mu ,X_+^a]_{}T^a\}`$ (A20)
$`+k_3\{\overline{T}^a\gamma ^\mu \gamma _5[A_\mu ,X_{}^a]_+T^3+\overline{T}^3\gamma ^\mu \gamma _5[A_\mu ,X_+^a]_+T^a\}`$ (A21)
$`+k_4\{\overline{T}^a\gamma ^\mu \gamma _5[A_\mu ,X_{}^a]_{}T^3\overline{T}^3\gamma ^\mu \gamma _5[A_\mu ,X_+^a]_{}T^a\}`$ $`,`$ (A22)
$`_{\mathrm{\Delta }I=2}^{\pi \mathrm{\Delta }}=j_5^{ab}\overline{T}^a\gamma ^\mu A_\mu T^b+j_6^{ab}\overline{T}^i[X_R^aA_\mu X_R^b+X_L^aA_\mu X_L^b]\gamma ^\mu T_i`$ (A23)
$`+k_5^{ab}\overline{T}^i[X_R^aA_\mu X_R^bX_L^aA_\mu X_L^b]\gamma ^\mu \gamma _5T_i`$ (A24)
$`+k_6ϵ^{ab3}[\overline{T}^3i\gamma _5X_+^bT^a+\overline{T}^ai\gamma _5X_+^bT^3]`$ $`,`$ (A25)
where we have suppressed the Lorentz indices of the $`\mathrm{\Delta }`$ field, i.e., $`\overline{T}^\nu \mathrm{}T_\nu `$. The vertices with $`k_i,h_\mathrm{\Delta }`$ contain two pions. All other vertices contain one pion when expanded to the leading order. At first sight the leading order term with $`k_6`$ in (A23) has no pions. However such a term cancels its hermitian conjugate exactly. The constants $`h_{\pi \mathrm{\Delta }}^i`$ are the PV $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ Yukawa coupling constants.
In Section 2, the leading terms in the above Lagrangians were expressed in terms of effective $`\pi \pi N\mathrm{\Delta }`$ and $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ coupling constants. These constants may be expressed in terms of the $`f_i`$, $`g_i`$, $`k_i`$, $`j_i`$ and $`h_{\pi \mathrm{\Delta }}^i`$ as follows:
$`h_A^{p\mathrm{\Delta }^{++}\pi ^{}\pi ^0}=2g_1+2g_2g_33g_4{\displaystyle \frac{2}{3}}g_5+{\displaystyle \frac{2}{3}}g_6`$ (A26)
$`h_A^{p\mathrm{\Delta }^{++}\pi ^0\pi ^{}}=2g_1+g_3+6g_4+{\displaystyle \frac{2}{3}}g_5`$ (A27)
$`h_A^{p\mathrm{\Delta }^+\pi ^0\pi ^0}={\displaystyle \frac{\sqrt{6}}{9}}(6g_2+9g_4+2g_6)`$ (A28)
$`h_A^{p\mathrm{\Delta }^+\pi ^+\pi ^{}}={\displaystyle \frac{\sqrt{6}}{9}}(6g_19g_4+4g_5+6g_6)`$ (A29)
$`h_A^{p\mathrm{\Delta }^+\pi ^{}\pi ^+}={\displaystyle \frac{\sqrt{6}}{9}}(6g_16g_24g_5+4g_6)`$ (A30)
$`h_A^{p\mathrm{\Delta }^0\pi ^+\pi ^0}={\displaystyle \frac{\sqrt{3}}{9}}(6g_1+6g_23g_3+9g_4+2g_5+2g_6)`$ (A31)
$`h_A^{p\mathrm{\Delta }^0\pi ^0\pi ^+}={\displaystyle \frac{\sqrt{3}}{9}}(6g_1+12g_2+3g_3+18g_42g_58g_6)`$ (A32)
$`h_A^{p\mathrm{\Delta }^{}\pi ^+\pi ^+}={\displaystyle \frac{\sqrt{2}}{3}}(6g_29g_4+2g_6)`$ (A33)
$`h_A^{n\mathrm{\Delta }^{++}\pi ^{}\pi ^{}}={\displaystyle \frac{\sqrt{2}}{3}}(6g_29g_4+2g_6)`$ (A34)
$`h_A^{n\mathrm{\Delta }^+\pi ^{}\pi ^0}={\displaystyle \frac{\sqrt{3}}{9}}(6g_1+6g_2+3g_39g_4+2g_5+2g_6)`$ (A35)
$`h_A^{n\mathrm{\Delta }^+\pi ^0\pi ^{}}={\displaystyle \frac{\sqrt{3}}{9}}(6g_1+12g_23g_318g_42g_58g_6)`$ (A36)
$`h_A^{n\mathrm{\Delta }^0\pi ^0\pi ^0}={\displaystyle \frac{\sqrt{6}}{9}}(6g_2+9g_42g_6)`$ (A37)
$`h_A^{n\mathrm{\Delta }^0\pi ^+\pi ^{}}={\displaystyle \frac{\sqrt{6}}{9}}(6g_1+6g_2+4g_54g_6)`$ (A38)
$`h_A^{n\mathrm{\Delta }^0\pi ^{}\pi ^+}={\displaystyle \frac{\sqrt{6}}{9}}(6g_19g_44g_56g_6)`$ (A39)
$`h_A^{n\mathrm{\Delta }^{}\pi ^+\pi ^0}=2g_12g_2+g_3+3g_4+{\displaystyle \frac{2}{3}}g_5{\displaystyle \frac{2}{3}}g_6`$ (A40)
$`h_A^{n\mathrm{\Delta }^{}\pi ^0\pi ^+}=2g_1g_36g_4{\displaystyle \frac{2}{3}}g_5`$ (A41)
$`h_\mathrm{\Delta }=h_{\pi \mathrm{\Delta }}^1+h_{\pi \mathrm{\Delta }}^2`$ (A42)
(A43)
$`h_V^{\mathrm{\Delta }^{++}\mathrm{\Delta }^+}={\displaystyle \frac{1}{\sqrt{6}}}(j_0+{\displaystyle \frac{4}{3}}j_6)2\sqrt{6}j_2{\displaystyle \frac{2\sqrt{6}}{3}}(j_3+j_4)+{\displaystyle \frac{j_5}{3\sqrt{6}}}`$ (A44)
$`h_V^{\mathrm{\Delta }^+\mathrm{\Delta }^0}={\displaystyle \frac{\sqrt{2}}{3}}(j_0+{\displaystyle \frac{4}{3}}j_6){\displaystyle \frac{2\sqrt{2}}{9}}j_5`$ (A45)
$`h_V^{\mathrm{\Delta }^0\mathrm{\Delta }^{}}={\displaystyle \frac{1}{\sqrt{6}}}(j_0+{\displaystyle \frac{4}{3}}j_6)+2\sqrt{6}j_2+{\displaystyle \frac{2\sqrt{6}}{3}}(j_3+j_4)+{\displaystyle \frac{j_5}{3\sqrt{6}}}`$ (A46)
It is interesting to note there is only one independent PV Yukawa coupling constant $`h_\mathrm{\Delta }`$ for $`\pi \mathrm{\Delta }\mathrm{\Delta }`$ interactions.
## B Vanishing Loop Contributions
As noted in Section 3, a large number of graphs which nominally contribute to $`h_\pi ^{EFF}`$ actually vanish up to $`𝒪(1/\mathrm{\Lambda }_\chi ^3)`$. Here, we summarize the the reasons why.
Consider first the corrections due to the PV vector $`\pi NN`$ vertices. For FIG. 1 (b) we have
$$iM_{(b)}=i\frac{g_A^2}{\sqrt{2}F_\pi ^3}\tau ^+(h_v^0+\frac{4}{3}h_V^2)(vq)\frac{d^Dk}{(2\pi )^D}\frac{(Sk)^2}{vkv(k+q)(k^2m_\pi ^2)}𝒪(1/m_N\mathrm{\Lambda }_\chi ^3),$$
(B1)
where we have used $`vq𝒪(1/m_N)`$. Since we are working to leading order in the $`1/m_N`$ expansion, this amplitude does not contribute. The PV vector interactions also appear in Figs. 1(j1,j2). The corresponding amplitude is
$$iM_{(j1)+(j2)}=i\frac{g_A^2}{\sqrt{2}F_\pi ^3}\tau ^+(h_v^0+2h_V^1\frac{8}{3}h_V^2)\frac{d^Dk}{(2\pi )^D}\frac{[(Sk),(Sq)]_+}{vkv(k+q)(k^2m_\pi ^2)}=0$$
(B2)
This integral vanishes because it is proportional to $`[(Sv),(Sq)]_+`$, which vanishes because $`Sv=0`$. All other possible insertions of PV vector $`\pi NN`$ vertices vanish for similar reasons as either (B1) or (B2). In what follows, we refer only to insertions involving the PV $`\pi NN`$ Yukawa and $`\pi \pi NN`$ axial couplings.
The propagator corrections in FIG. 1 (g1)-(h2) vanish after integration since their amplitude of (g1)-g2) goes as
$$h_\pi \frac{d^Dk}{(2\pi )^D}\frac{vk}{k^2m_\pi ^2}=0.$$
(B3)
while the amplitude of (h1)-(h2) goes as
$$h_A^i\frac{d^Dk}{(2\pi )^D}\frac{Sk}{k^2m_\pi ^2}=0.$$
(B4)
The amplitude of FIG. 1 (i1)-(i4) contains a vanishing integral
$$h_\pi \frac{d^Dk}{(2\pi )^D}\frac{Sk}{vk(k^2m_\pi ^2)}=0.$$
(B5)
Figs. 1 (j1)-(j2) do not contribute for the PV Yukawa coupling $`h_\pi `$ due to charge conservation. The remaining non-zero diagrams are Figs. 1 (a)-(f2) where the insertions in loops are of the Yukawa or axial interations. Figs. (f1), (f2) arises from the insertion of the counter terms of mass and wave function renormalization. Fig. 1 (e1)-(e2) and Fig. 2 (c1)-(c2) contribute to the wave function renormalization in Eq. (62).
Due to the heavy baryon projection $`P_+i\gamma _5P_+=0`$ the one pion PV $`\pi N\mathrm{\Delta }`$ vertex does not contribute in the leading order of heavy baryon expansion. Hence, the chiral loop corrections from FIG. 2 (d1)-(g4) are of higher order. Fig. 2 (h1)-(j2) vanishes after integration for reasons similar to (B2). The remaining, non-vanishing diagrams are discussed explicitly in Section 3.
As pointed out in Section 2, both PC and PV two-derivative operators which conserve CP do not contribute to $`h_\pi `$ renormalization. For example, there exists one CP-conserving, PV such operator:
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}\sigma ^{\mu \nu }[D_\mu A_\nu D_\nu A_\mu ]N.$$
(B6)
After expansion, the leading term starts with three pions. It contributes via Figure 1 (a), at the order of $`1/\mathrm{\Lambda }_\chi F_\pi ^3`$. Moreover the loop integration yields a factor $`g_{\mu \nu }`$ and leads to zero after contraction with $`\sigma ^{\mu \nu }`$.
Another possibility comes from insertions of PC two-derivative nucleon pion operators. There are three PC operators which conserve CP:
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}i\gamma _5D_\mu A^\mu N,$$
(B7)
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}A^\mu A_\mu N,$$
(B8)
$$\frac{1}{\mathrm{\Lambda }_\chi }\overline{N}\sigma ^{\mu \nu }[A_\mu ,A_\nu ]N.$$
(B9)
Note the first two operators are symmetric in the Lorentz indices. Only the last one arises from the antisymmetric operators listed in Eq. (51). The first one starts off with one pion. The relevant Feynman diagrams are Figure 1 (c1)-(c2), where the PV vertex is associated with $`h_A^i`$. Note these diagrams do not contribute at leading order of HBChPT due to the presence of the $`i\gamma _5`$. The remaining two operators start off with two pions. The relevant diagrams are Figure 1 (d1)-(d2). After integration the contribution of the third operator reads
$$h_\pi ϵ^{\mu \nu \alpha \beta }v_\alpha S_\beta v^\mu q^\nu m_\pi ^2\mathrm{ln}m_\pi /\mathrm{\Lambda }_\chi F_\pi ^2.$$
(B10)
So its contribution is zero. In contrast the second operator yields
$$h_\pi (vq)m_\pi ^2\mathrm{ln}m_\pi /(\mathrm{\Lambda }_\chi F_\pi ^2).$$
(B11)
Note $`vq1/m_N`$. So its contribution is of order $`1/(\mathrm{\Lambda }_\chi ^3m_N)`$. In short, none of the two-derivative operators contribute to the renormalization of $`h_\pi `$ at the order to which we work.
## C Bag Model Integrals
Here, we show that the four-quark bag model integrals relevant to the calculation of the DDH/FCDH parameters $`\eta `$ and $`y`$ have a Taylor expansion in light quark mass around $`m_q=0`$. We write a bag model quark wavefunction as
$$\psi (x)=\left(\begin{array}{c}iu(r)\chi \hfill \\ \mathrm{}(r)\stackrel{}{\sigma }\stackrel{}{r}\chi \hfill \end{array}\right)\mathrm{exp}(iEt),$$
(C1)
where $`\chi `$ denotes a two-component Pauli spinor and where wave function normalization yields
$$d^3r(u(r)^2+\mathrm{}(r)^2)=1,$$
(C2)
where the the radial integration runs from 0 to the bag radius, $`R`$. The four quark matrix elements of interest here can depend three different integrals:
$$d^3ru(r)^4,d^3r\mathrm{}(r)^4,d^3ru(r)^2\mathrm{}(r)^2.$$
(C3)
The quark radial wave functions are
$`u(r)`$ $`=`$ $`Nj_0({\displaystyle \frac{p_nr}{R}})`$ (C4)
$`\mathrm{}(r)`$ $`=`$ $`N\left({\displaystyle \frac{\omega _nm_qR}{\omega _n+m_qR}}\right)^{1/2}j_1({\displaystyle \frac{p_nr}{R}}),`$ (C5)
where
$$\mathrm{tan}p_n=\frac{p_n}{\omega _n+m_qR1}(n=1,2,\mathrm{})$$
(C6)
$$p_n=\sqrt{\omega _n^2m_q^2R^2}$$
(C7)
$$N=\sqrt{\frac{p_n^4}{R^3(2\omega _n^22\omega _n+m_qR)\mathrm{sin}^2p_n}}$$
(C8)
$$R^4=\frac{N\omega _nZ_0}{4\pi B}$$
(C9)
B is the bag constant and $`Z_0`$ is a phenomenological parameter involved with the center of mass motion of the bag.
For light quarks and lowest eigenmode
$`\omega _0`$ $`(2.043+0.493m_qR)`$ (C10)
$`N`$ $`2.27/\sqrt{4\pi R^3}.`$ (C11)
It is straightforward to show that the bag model integrals in Eq. (C3) have a Taylor expansion about $`m_q=0`$. The argument proceeds by noting that the quantities $`N`$, $`R`$, $`p_n`$, $`\omega _n`$ and the argument of the spherical Bessel functions all have Taylor series in $`m_q`$ about $`m_q=0`$. The existence of this expansion can be seen be an explicit, iterative construction. First, expand $`\omega _n`$ and $`R`$:
$`\omega _n`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\omega _{n,k}m_q^k`$ (C12)
$`R`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}R_km_q^k.`$ (C13)
Now let $`m_q=0`$ in Eqs. (C6,C7). Doing so eliminates all dependence on $`R`$ and determines $`\omega _{n,0}`$. Next, set $`m_q=0`$ in Eqs. (C8, C9) with $`\omega _n\omega _{n,0}`$. Doing so determines $`R_0`$. Now expand Eqs. (C6,C7) to first order in $`m_q`$. This step yields $`\omega _{n,1}`$ in terms of $`\omega _{n,0}`$ and $`R_0`$. Expanding Eqs. (C8, C9) to first order in $`m_q`$ then determines $`R_1`$ in terms of $`\omega _{n,0}`$, $`\omega _{n,1}`$, and $`R_0`$ and so forth. Note that at any step of the recursion, the argument of any transcendental function is $`\omega _{n,0}`$. Hence, at any order, a solution for the $`\omega _{n,k}`$ and $`R_k`$ exists.
The expansion of the bag model integrals continues by computing their derivatives with respect to $`m_q`$ and using the expansions of $`N`$, $`R`$, etc. in terms of $`m_q`$ as constructed above. Taking $`n`$ derivatives of one of the integrals in Eq. (C3) yields new intregrals involving powers of $`r/R`$ times products of the Bessel functions and their derivatives. Using the standard Bessel function recursion relations, the derivatives of the $`j_k`$ can always be expressed in terms of other spherical Bessel functions. Since the $`j_k`$ and their derivatives are finite at the origin, and since the radial bag integration is bounded above by $`R`$, the $`n`$th derivative of any of the integrals in Eq. (C3) is finite. Thus, each of the integrals in Eq. (C3) can be expanded in a Taylor series about $`m_q=0`$.
Figure Captions
Figure 1. Meson-nucleon intermediate state contributions to the PV $`\pi NN`$ vertex $`h_\pi `$. The shaded circle denotes the PV vertex. The solid and dashed lines correspond to the nucleon and pion respectively.
Figure 2. The chiral corrections from $`\mathrm{\Delta }`$ intermediate state, which is denoted by the double line.
Figure 3.Diagrammatic representation of the SU(6)<sub>w</sub> components of $`B^{}M|_W^{PV}(\mathrm{\Delta }S=0,1)|B`$. Figs. 3a-c correspond, respectively, to $`b_{t,v}`$, $`c_v`$, and $`a_{t,v}`$. The wavy line denotes the action of $`_W^{PV}`$.
Figure 4. Chiral corrections to the $`BB^{}M`$ nonleptonic weak decay.
Figure 5. Quark line diagrams for the renormalization of $`h_\pi `$ due to the axial PV $`\pi \pi NN`$ interaction. As in Fig. 3, the wavey line denotes the action of $`_W^{PV}`$. Fig. 5a shows a typical contribution to $`h_A^i`$. Figs 5b,c denote the corresponding loop corrections to $`h_\pi `$. Fig. 5(b) contains the disconnected $`q\overline{q}`$ insertions, while Fig. 5(c) gives a Z-graph contribution.
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# Functions of Baire Class One
## 1. Preliminaries
Let $`K`$ be a compact metric space. A function $`f:K`$ is said to be of *Baire class one*, or simply, *Baire-1*, if there exists a sequence $`\left(f_n\right)`$ of real-valued continuous functions that converges pointwise to $`f.`$ Let $`𝔅_1\left(K\right)`$ (respectively, $`_1\left(K\right)`$) be the set of all real-valued (respectively, bounded real-valued) Baire-1 functions on $`K.`$ Several authors have studied Baire-1 functions in terms of ordinal ranks associated to each function. (See, e.g., , and ). In this paper, we study the relationship between two of these ordinal ranks, namely the oscillation rank $`\beta `$ and the convergence rank $`\gamma .`$
We begin by recalling the definitions of the indices $`\beta `$ and $`\gamma `$. Suppose that $`H`$ is a compact metric space, and $`f`$ is a real-valued function whose domain contains $`H`$. For any $`\epsilon >0`$, let $`H^0(f,\epsilon )=H`$. If $`H^\alpha (f,\epsilon )`$ is defined for some countable ordinal $`\alpha `$, let $`H^{\alpha +1}(f,\epsilon )`$ be the set of all those $`xH^\alpha (f,\epsilon )`$ such that for every open set $`U`$ containing $`x,`$ there are two points $`x_1`$ and $`x_2`$ in $`UH^\alpha (f,\epsilon )`$ with $`\left|f\left(x_1\right)f\left(x_2\right)\right|\epsilon .`$ For a countable limit ordinal $`\alpha `$, we let
$$H^\alpha (f,\epsilon )=\underset{\alpha ^{}<\alpha }{}H^\alpha ^{}(f,\epsilon ).$$
The index $`\beta _H(f,\epsilon )`$ is taken to be $`\text{the least }\alpha \text{ with }H^\alpha (f,\epsilon )=\mathrm{}`$ if such $`\alpha \text{ exists,}`$ and $`\omega _1`$ otherwise. The oscillation index of $`f`$ is
$$\beta _H\left(f\right)=sup\{\beta _H(f,\epsilon ):\epsilon >0\}.$$
If the ambient space $`H`$ is clear from the context, we write $`\beta (f,\epsilon )`$ and $`\beta (f)`$ in place of $`\beta _H(f,\epsilon )`$ and $`\beta _H(f)`$ respectively.
The $`\gamma `$ index is defined analogously. If $`\left(f_n\right)`$ is a sequence of real-valued functions such that $`H_ndom\left(f_n\right),`$ let $`H^0(\left(f_n\right),\epsilon )=H`$ for any $`\epsilon >0`$. If $`H^\alpha (\left(f_n\right),\epsilon )`$ has been defined for some countable ordinal $`\alpha `$, let $`H^{\alpha +1}(\left(f_n\right),\epsilon )`$ be the set of all those $`xH^\alpha (\left(f_n\right),\epsilon )`$ such that for every open set $`U`$ containing $`x`$ and any $`m,`$ there are two integers $`n_1`$, $`n_2`$ with $`n_1>n_2>m`$ and $`x^{}UH^\alpha (\left(f_n\right),\epsilon )`$ such that $`\left|f_{n_1}\left(x^{}\right)f_{n_2}\left(x^{}\right)\right|\epsilon .`$ Define
$$H^\alpha (\left(f_n\right),\epsilon )=\underset{\alpha ^{}<\alpha }{}H^\alpha ^{}(\left(f_n\right),\epsilon )$$
if $`\alpha `$ is a countable limit ordinal. Let $`\gamma _H(\left(f_n\right),\epsilon )`$ be the least $`\alpha `$ with $`H^\alpha (\left(f_n\right),\epsilon )=\mathrm{}`$ if such $`\alpha `$ exists, and $`\omega _1`$ otherwise. Finally, the convergence index of $`\left(f_n\right)`$ is the ordinal
$$\gamma _H\left(\left(f_n\right)\right)=sup\{\gamma _H(\left(f_n\right),\epsilon ):\epsilon >0\}.$$
Again, if there is no ambiguity about the space $`H`$, we write $`\gamma ((f_n),\epsilon )`$ and $`\gamma ((f_n))`$ for $`\gamma _H((f_n),\epsilon )`$ and $`\gamma _H((f_n))`$ respectively.
It is known that a function $`f:K`$ is Baire-1 if and only if $`\beta \left(f\right)<\omega _1.`$ (See \[3, Proposition 1.2\].) Following , we define the set of functions of small Baire class $`\xi `$ and the set of bounded functions of small Baire class $`\xi `$ for each countable ordinal $`\xi `$ as
$$𝔅_1^\xi \left(K\right)=\{f𝔅_1\left(K\right):\beta \left(f\right)\omega ^\xi \}$$
and
$$_1^\xi \left(K\right)=\{f_1\left(K\right):\beta \left(f\right)\omega ^\xi \}$$
respectively. In , the following results are shown.
###### Theorem 1.1.
Let $`K`$ be a compact metric space.
1. *\[4, Theorem 7\]* If $`\xi `$ is a finite ordinal, then a function $`f_1^{\xi +1}\left(K\right)`$ if and only if there exists a sequence $`(f_n)`$ in $`_1^1\left(K\right)`$ converging pointwise to $`f`$ such that $`\gamma ((f_n))\omega ^\xi `$.
2. *\[4, Corollary 9\]* If $`\xi `$ is an infinite countable ordinal, and $`f_1\left(K\right)`$ is the pointwise limit of a sequence $`(f_n)`$ in $`_1^1\left(K\right)`$ such that $`\gamma ((f_n))\omega ^\xi `$, then $`\beta (f)\omega ^\xi `$.
One of our main results generalizes and unifies the two parts of Theorem 1.1.
###### Theorem 1.2.
Let $`K`$ be a compact metric space and let $`\xi _1`$, $`\xi _2`$ be countable ordinals. A function $`f𝔅_1^{\xi _1+\xi _2}\left(K\right),`$ respectively, $`_1^{\xi _1+\xi _2}\left(K\right),`$ if and only if there exists a sequence $`(f_n)`$ in $`𝔅_1^{\xi _1}\left(K\right),`$ respectively, $`_1^{\xi _1}\left(K\right),`$ converging pointwise to $`f`$ such that $`\gamma ((f_n))\omega ^{\xi _2}`$.
In the course of proving Theorem 1.2, we show that any Baire-$`1`$ function $`f`$ on a closed subspace $`H`$ of a compact metric space $`K`$ can be extended to a Baire-$`1`$ function $`g`$ on $`K`$ such that $`\beta _H(f)=\beta _K(g)`$ (Theorem 3.6). When $`\beta _H(f)=1`$, this is the familiar Tietze Extension Theorem. Proposition 2.1 and Theorem 2.3 in yield that for a bounded Baire-$`1`$ function $`f`$, $`\beta (f)`$ is the smallest ordinal $`\xi `$ such that there exists a sequence of continuous functions $`(f_n)`$ converging pointwise to $`f`$ and having $`\gamma ((f_n))=\xi `$. Theorem 5.5 below shows that the same result holds without the boundedness assumption on the function $`f`$. In the last section, we consider the product of Baire-$`1`$ functions. In contrast to the class $`_1^\xi (K)`$, the class $`𝔅_1^\xi \left(K\right)`$ is not closed under multiplication. Theorem 6.5 shows that if $`f𝔅_1^{\xi _1}\left(K\right)`$ and $`g𝔅_1^{\xi _2}\left(K\right),`$ then $`fg𝔅_1^\xi \left(K\right),`$ where $`\xi =\mathrm{max}\{\xi _1+\xi _2,\xi _2+\xi _1\}.`$ It is also shown that this result is the best possible.
Our notation is standard. In the sequel, $`K`$ will always denote a compact metric space. If $`H`$ is a closed subset of $`K,`$ the derived set $`H^{}`$ is the set of all limit points of $`H.`$ A transfinite sequence of derived sets is defined in the usual manner. Let $`H^{\left(0\right)}=H`$ and $`H^{\left(\alpha +1\right)}=\left(H^{\left(\alpha \right)}\right)^{}`$ for any ordinal $`\alpha .`$ If $`\alpha `$ is a limit ordinal, let
$$H^{\left(\alpha \right)}=\underset{\alpha ^{}<\alpha }{}H^{\left(\alpha ^{}\right)}.$$
Given real-valued functions $`f`$ and $`g`$ defined on a set $`S`$, we let
$$fg_S=sup\{|f(s)g(s)|:sS\}.$$
When there is no cause for confusion, we write $`fg`$ for $`fg_S`$. Since we shall be dealing with unbounded functions in general, this functional can take the value $`\mathrm{}`$ and is not a “norm”. However, it is compatible with the topology of uniform convergence on the set $`^S`$ of all real-valued functions on $`S`$ in the sense that the sets
$$U(f,\epsilon )=\{g:gf_S<\epsilon \}$$
form a basis for the said topology.
## 2. Oscillation and convergence of Baire-1 functions
We begin by proving a result that yields an upper bound of the oscillation index of a Baire-1 function $`f`$ as the product of the convergence index of a sequence of functions $`\left(f_n\right)`$ converging pointwise to $`f`$, and the supremum of the oscillation indices of $`f_n`$’s.
###### Lemma 2.1.
Let $`U`$ and $`L`$ be sets such that $`ULK,`$ where $`U`$ is open in $`K`$ and $`L`$ is closed in $`K.`$ Suppose $`f`$,$`f_n`$ $`\left(n1\right)`$ are Baire-1 functions on $`K,`$ $`\alpha <\omega _1`$, and $`\epsilon >0.`$ Then
(a) $`L^\alpha (f,\epsilon )K^\alpha (f,\epsilon )L,`$
(b) $`L^\alpha (\left(f_n\right),\epsilon )K^\alpha (\left(f_n\right),\epsilon )L,`$
(c) $`K^\alpha (f,\epsilon )UL^\alpha (f,\epsilon ),`$
(d) $`K^\alpha (\left(f_n\right),\epsilon )UL^\alpha (\left(f_n\right),\epsilon ).`$
###### Proof.
We only prove (c). The proof is by induction on $`\alpha .`$ The statement is trivial if $`\alpha =0`$ or a limit ordinal. Suppose the statement is true for all ordinals not greater than $`\alpha .`$ Let $`xK^{\alpha +1}(f,\epsilon )U.`$ If $`N`$ is a neighborhood of $`x`$ in $`K`$, then $`NU`$ is open in $`K.`$ Thus there exist $`x_1,x_2\left(NU\right)K^\alpha (f,\epsilon )=N\left(UK^\alpha (f,\epsilon )\right)NL^\alpha (f,\epsilon )`$ such that $`\left|f\left(x_1\right)f\left(x_2\right)\right|\epsilon .`$ Hence $`xL^{\alpha +1}(f,\epsilon ).`$
###### Proposition 2.2.
Let $`\left(f_n\right)`$ be a sequence in $`𝔅_1\left(K\right)`$ and let $`\epsilon >0.`$ Suppose that $`\beta (f_n,\epsilon )\beta _0`$ for all $`n`$, and $`\gamma (\left(f_n\right),\epsilon )\gamma _0.`$ If $`\left(f_n\right)`$ converges pointwise to a function $`f,`$ then $`\beta (f,3\epsilon )\beta _0\gamma _0.`$
###### Proof.
We first consider the case $`\gamma _0=1.`$ Then $`K^1(\left(f_n\right),\epsilon )=\mathrm{}.`$ For each $`xK,`$ there exist an open neighborhood $`U_x`$ of $`x`$ and $`p_x`$ such that whenever $`n>m>p_x,`$
$$\left|f_n\left(x^{}\right)f_m\left(x^{}\right)\right|<\epsilon $$
for all $`x^{}U_x.`$ By the compactness of $`K,`$ there exist $`x_1,x_2,\mathrm{},x_k`$ such that
$$K\underset{i=1}{\overset{k}{}}U_{x_i}.$$
Let $`p_0=\mathrm{max}\{p_{x_1},p_{x_2},\mathrm{},p_{x_k}\}.`$ Then for all $`n>m>p_0`$ and $`yK,`$ we have $`yU_{x_i}`$ for some $`i,\mathrm{\hspace{0.17em}1}ik.`$ Since $`n>m>p_{x_i},`$
$$\left|f_n\left(y\right)f_m\left(y\right)\right|<\epsilon .$$
Taking limit as $`n\mathrm{},`$ we have
(2.1)
$$ff_m\epsilon \text{ for all }m>p_0.$$
Using (2.1), it is easy to verify by induction that
$$K^\alpha (f,3\epsilon )K^\alpha (f_{p_0+1},\epsilon )$$
for all $`\alpha <\omega _1.`$ In particular,
$$K^{\beta _0}(f,3\epsilon )K^{\beta _0}(f_{p_0+1},\epsilon )=\mathrm{}.$$
Hence $`\beta (f,3\epsilon )\beta _0=\beta _0\gamma _0.`$
Suppose the assertion is true for some $`\gamma _0.`$ Let $`\left(f_n\right)`$ be a sequence in $`𝔅_1\left(K\right)`$ that converges pointwise to a function $`f.`$ Suppose there exists $`\epsilon >0`$ such that $`\beta (f_n,\epsilon )\beta _0`$ for all $`n`$ and $`\gamma (\left(f_n\right),\epsilon )\gamma _0+1.`$ We need to show $`\beta (f,3\epsilon )\beta _0\left(\gamma _0+1\right).`$ Since $`\gamma (\left(f_n\right),\epsilon )\gamma _0+1,`$ we have $`K^{\gamma _0+1}(\left(f_n\right),\epsilon )=\mathrm{}.`$ For each $`m`$, let $`U_m`$ denote the $`\frac{1}{m}`$neighborhood of $`K^{\gamma _0}(\left(f_n\right),\epsilon ).`$ Denote $`KU_m`$ by $`\stackrel{~}{K}_m.`$ From Lemma 2.1(a) and 2.1(b), for each $`n,`$ $`\beta _{\stackrel{~}{K}_m}(f_n,\epsilon )\beta _0`$ and $`\gamma _{\stackrel{~}{K}_m}(\left(f_n\right),\epsilon )\gamma _0`$. By the inductive hypothesis, we see that
$$\beta _{\stackrel{~}{K}_m}(f,3\epsilon )\beta _0\gamma _0.$$
From this and applying Lemma 2.1(c) with $`U=K\overline{U}_m,`$ $`L=\stackrel{~}{K}_m`$ for all $`m`$, we see that $`K^{\beta _0\gamma _0}(f,3\epsilon )K^{\gamma _0}(\left(f_n\right),\epsilon ).`$ Let
$$\stackrel{~}{K}=K^{\beta _0\gamma _0}(f,3\epsilon )K^{\gamma _0}(\left(f_n\right),\epsilon ).$$
Then
$$\beta _{\stackrel{~}{K}}(f_n,\epsilon )\beta _0\text{ and }\gamma _{\stackrel{~}{K}}(\left(f_n\right),\epsilon )=1.$$
Thus
$$\beta _{\stackrel{~}{K}}(f,3\epsilon )\beta _0\text{ by the case when }\gamma _0=1.$$
Therefore
$$K^{\beta _0\left(\gamma _0+1\right)}(f,3\epsilon )=K^{\beta _0\gamma _0+\beta _0}(f,3\epsilon )=\stackrel{~}{K}^{\beta _0}(f,3\epsilon )=\mathrm{}.$$
Hence
$$\beta (f,3\epsilon )\beta _0\left(\gamma _0+1\right).$$
Suppose $`\gamma _0<\omega _1`$ is a limit ordinal and the statement holds for all ordinals $`\gamma <\gamma _0.`$ Let $`\left(f_n\right)𝔅_1\left(K\right)`$ be a sequence that converges pointwise to a function $`f`$ and let $`\epsilon >0`$ be given. Suppose that $`\beta (f_n,\epsilon )\beta _0`$ for all $`n`$, and $`\gamma (\left(f_n\right),\epsilon )\gamma _0.`$ Then $`\gamma (\left(f_n\right),\epsilon )<\gamma _0`$ and $`\beta (f,3\epsilon )\beta _0`$ $`\gamma (\left(f_n\right),\epsilon )<\beta _0\gamma _0.`$
###### Theorem 2.3.
Let $`\left(f_n\right)`$ be a sequence $`𝔅_1\left(K\right)`$ converging pointwise to a function $`f.`$ Suppose $`sup\{\beta \left(f_n\right):n\}\beta _0`$ and $`\gamma \left(\left(f_n\right)\right)\gamma _0.`$ Then $`f`$ is Baire-1 and $`\beta \left(f\right)\beta _0\gamma _0.`$
For the next corollary, recall that $`DBSC\left(K\right)`$ is the space of all differences of semicontinuous functions on $`K.`$ It is known that $`_1^1\left(K\right)`$ is the closure of $`DBSC\left(K\right)`$ in the topology of uniform convergence (\[3, Theorem 3.1\]).
###### Corollary 2.4 (\[4, Corollary 9\]).
Let $`f_1\left(K\right)`$ be the pointwise limit of a sequence $`\left(f_n\right)DBSC\left(K\right)`$. If $`\gamma \left(\left(f_n\right)\right)\omega ^\xi ,`$ $`\omega \xi <\omega _1,`$ then $`\beta \left(f\right)\omega ^\xi .`$
## 3. Extension of Baire-1 functions
In this section, we establish several results regarding the extension of Baire-1 functions. They are analogs of the Tietze Extension Theorem for continuous functions. These results are applied in the next section in proving the converse of Theorem 2.3.
###### Lemma 3.1.
Suppose that $`F`$ is a closed subspace of $`K`$ and that $`f`$ is a Baire-1 function on $`F`$. For any $`\epsilon >0`$, there exists a continuous function $`g:KF^1(f,\epsilon )`$ such that
$$gf_{FF^1(f,\epsilon )}\epsilon .$$
###### Proof.
For any $`xFF^1(f,\epsilon ),`$ choose an open neighborhood $`U_x`$ of $`x`$ in $`K`$ such that $`U_xF^1(f,\epsilon )=\mathrm{}`$ and $`\left|f\left(x_1\right)f\left(x_2\right)\right|<\epsilon \text{ for all }x_1,x_2U_xF.`$ The collection $`𝒰=\{U_x:xFF^1(f,\epsilon )\}\left\{KF\right\}`$ is an open cover of $`KF^1(f,\epsilon )`$. By , Theorems IX.5.3 and VIII.4.2, there exists a partition of unity $`\left(\phi _U\right)_{U𝒰}`$ subordinated to $`𝒰.`$ If $`U=U_x𝒰`$ for some $`xFF^1(f,\epsilon )`$, let $`a_U=f(x);`$ if $`U=KF`$, let $`a_U=0`$. Define $`g:KF^1(f,\epsilon )`$ by $`g=_{U𝒰}a_U\phi _U.`$ The sum is well-defined since $`\{\text{supp }\phi _U:U𝒰\}`$ is locally finite. Let $`xFF^1(f,\epsilon ).`$ Then $`𝒱=\{U𝒰:\phi _U\left(x\right)0\}`$ is a finite set, $`\phi _U\left(x\right)>0`$ for all $`U𝒱`$ and $`_{U𝒱}\phi _U\left(x\right)=1.`$ If $`U𝒱,`$ then $`xUF;`$ hence $`UKF`$. Therefore, $`U=U_y`$ for some $`yFF^1(f,\epsilon ).`$ But then $`x,yU_yF`$ implies that $`\left|a_Uf\left(x\right)\right|=\left|f\left(y\right)f\left(x\right)\right|<\epsilon .`$ It follows that
$`\left|g\left(x\right)f\left(x\right)\right|`$ $`=\left|{\displaystyle \underset{U𝒰}{}}a_U\phi _U\left(x\right)f\left(x\right)\right|=\left|{\displaystyle \underset{U𝒱}{}}a_U\phi _U\left(x\right){\displaystyle \underset{U𝒱}{}}f\left(x\right)\phi _U\left(x\right)\right|`$
$`{\displaystyle \underset{U𝒱}{}}\left|a_Uf\left(x\right)\right|\phi _U\left(x\right)<\epsilon .`$
This shows that
$$gf_{FF^1(f,\epsilon )}\epsilon .$$
Finally, if $`x`$ is a point in $`KF^1(f,\epsilon )`$, there exists an open neighborhood $`V`$ of $`x`$ in $`K`$ such that $`VF^1(f,\epsilon )=\mathrm{}`$ and $`𝒲=\{U𝒰:\text{supp }\phi _UV\mathrm{}\}`$ is finite. Now
$$g_{|V}=\underset{U𝒰}{}a_U\phi _{U|V}=\underset{U𝒲}{}a_U\phi _{U|V}.$$
Hence $`g_{|V}`$ is continuous on $`V`$, since it is a finite linear combination of continuous functions. In particular, $`g`$ is continuous at $`x`$. As $`xKF^1(f,\epsilon )`$ is arbitrary, $`g`$ is continuous on $`KF^1(f,\epsilon )`$. ∎
###### Theorem 3.2.
Suppose that $`F`$ is a closed subspace of $`K`$ and that $`f`$ is a Baire-1 function on $`F`$. For any $`1\beta _0<\omega _1,`$ and any $`\epsilon >0,`$ there exists $`g:KF^{\beta _0}(f,\epsilon )`$ such that
$$gf_{FF^{\beta _0}(f,\epsilon )}\epsilon $$
and
$$\beta _H\left(g\right)\beta _0\text{ for all compact subsets }H\text{ of }KF^{\beta _0}(f,\epsilon ).$$
###### Proof.
Let $`h:KF^1(f,\epsilon )`$ be the function obtained from Lemma 3.1. If $`1\alpha <\beta _0`$, let $`\stackrel{~}{K}=\stackrel{~}{F}=F^\alpha (f,\epsilon )`$. Applying Lemma 3.1 with $`\stackrel{~}{K}`$, $`\stackrel{~}{F}`$, and the function $`f`$ yields a continuous function $`g_\alpha :F^\alpha (f,\epsilon )F^{\alpha +1}(f,\epsilon )`$ such that
$$g_\alpha f_{F^\alpha (f,\epsilon )F^{\alpha +1}(f,\epsilon )}\epsilon .$$
Let $`g=h\left(_{\alpha <\beta _0}g_\alpha \right):KF^{\beta _0}(f,\epsilon )`$. Then $`gf_{FF^{\beta _0}(f,\epsilon )}\epsilon .`$
Suppose that $`\delta >0`$ and $`H`$ is a compact subset of $`KF^{\beta _0}(f,\epsilon ).`$ If $`xF^1(f,\epsilon ),`$ then there exists an open neighborhood $`U`$ of $`x`$ such that
$$\overline{U}F^1(f,\epsilon )=\mathrm{}.$$
Note that $`g_{|\overline{U}}=h_{|\overline{U}}.`$ By Lemma 2.1(c),
$$H^1(g,\delta )U\left(H\overline{U}\right)^1(g,\delta )=\left(H\overline{U}\right)^1(h,\delta )=\mathrm{}$$
by the continuity of $`h.`$ In particular, $`xH^1(g,\delta ).`$ It follows that
$$H^1(g,\delta )HF^1(f,\epsilon ).$$
Repeating the argument inductively yields that
$$H^{\beta _0}(g,\delta )HF^{\beta _0}(f,\epsilon )=\mathrm{}.$$
Hence $`\beta _H\left(g\right)\beta _0`$, as required. ∎
We obtain the following corollaries by taking $`F=K`$ and $`\beta _0=\beta _F\left(f\right)`$ respectively.
###### Corollary 3.3.
Let $`f`$ be a Baire-1 function on $`K`$ such that $`\beta (f,\epsilon )\beta _0`$ for some $`1\beta _0<\omega _1`$ and $`\epsilon >0.`$ Then there exists $`g:K`$ such that
$$gf\epsilon \text{ and }\beta \left(g\right)\beta _0.$$
###### Corollary 3.4.
Let $`F`$ be a closed subspace of $`K.`$ If $`f`$ is a Baire-1 function on $`F`$, then for every $`\epsilon >0`$ there exists a Baire-1 function $`g`$ on $`K`$ such that
$$gf_F\epsilon \text{ and }\beta _K\left(g\right)\beta _F\left(f\right).$$
Next we show that Corollary 3.4 can be improved to an exact extension theorem (i.e., the case $`\epsilon =0`$). In the statement of Lemma 3.5, the vacuous sum $`_{j=1}^0g_j`$ is taken to be the zero function.
###### Lemma 3.5.
Let $`F`$ be a closed subspace of $`K`$ and let $`f`$ be a Baire-1 function on $`F.`$ Then there exists a sequence of Baire-1 functions $`\left(g_n\right)`$ on $`K`$ such that
(a) $`g_n`$ is continuous on $`KF^1(f_{j=1}^{n1}g_j,\frac{1}{2^{n1}})`$ for all $`n,`$
(b) $`f_{j=1}^ng_j_{FF^1(f,\frac{1}{4^{n1}})}{\displaystyle \frac{1}{2^{n1}}},`$ $`n,`$
(c) $`g_n_K{\displaystyle \frac{1}{2^{n2}}}`$ if $`n2,`$ and
(d) $`F^1(f_{j=1}^ng_j,\delta )F^1(f,\frac{\delta }{2^n})`$ if $`0<\delta \frac{1}{2^{n2}},`$ $`n.`$
###### Proof.
The functions $`\left(g_n\right)`$ are constructed inductively. By Lemma 3.1, there exists a continuous function $`g_1:KF^1(f,1)`$ such that $`fg_1_{FF^1(f,1)}1`$. Extend $`g_1`$ to a function on $`K`$ by defining $`g_1`$ to be $`0`$ on $`F^1(f,1)`$. Then (a) and (b) hold. Condition (c) holds vacuously. Moreover, if $`xFF^1(f,\frac{\delta }{2}),\mathrm{\hspace{0.17em}0}<\delta 2,`$ then there exists a neighborhood $`U_1`$ of $`x`$ in $`F`$ such that $`\left|f\left(x_1\right)f\left(x_2\right)\right|<\frac{\delta }{2}`$ for all $`x_1,`$ $`x_2U_1.`$ Note that since $`xFF^1(f,\frac{\delta }{2}),`$ $`g_1`$ is continuous at $`x.`$ Hence there exists a neighborhood $`U_2`$ of $`x`$ in $`F`$ such that $`\left|g_1\left(x_1\right)g_1\left(x_2\right)\right|<\frac{\delta }{2}`$ for all $`x_1,`$ $`x_2U_2.`$ Let $`U=U_1U_2.`$ Then $`U`$ is a neighborhood of $`x`$ in $`F.`$ For all $`x_1,x_2U,`$
$$\left|\left(fg_1\right)\left(x_1\right)\left(fg_1\right)\left(x_2\right)\right|<\delta .$$
Hence $`xF(fg_1,\delta ).`$ This proves (d).
Suppose that $`g_1,g_2,\mathrm{},g_n`$ have been chosen. By Lemma 3.1, there exists a continuous function $`h:KF^1(f_{j=1}^ng_j,\frac{1}{2^n})`$ such that
$$f\underset{j=1}{\overset{n}{}}g_jh_{FF^1(f_{j=1}^ng_j,\frac{1}{2^n})}\frac{1}{2^n}.$$
Define $`\stackrel{~}{h}`$ on $`KF^1(f_{j=1}^ng_j,\frac{1}{2^n})`$ by $`\stackrel{~}{h}=\left(h{\displaystyle \frac{1}{2^{n1}}}\right){\displaystyle \frac{1}{2^{n1}}}.`$ Then $`\stackrel{~}{h}`$ is continuous on $`KF^1(f_{j=1}^ng_j,\frac{1}{2^n}).`$ By (d), $`F^1(f_{j=1}^ng_j,\frac{1}{2^n})F^1(f,{\displaystyle \frac{1}{4^n}}).`$ Hence $`\stackrel{~}{h}`$ is defined and continuous on $`KF^1(f,{\displaystyle \frac{1}{4^n}}).`$ Moreover, it follows from (b) that
(3.1)
$$f\underset{j=1}{\overset{n}{}}g_j_{FF^1(f,\frac{1}{4^n})}\frac{1}{2^{n1}}.$$
From inequality (3.1) and the definition of $`\stackrel{~}{h},`$ we have
$$f\underset{j=1}{\overset{n}{}}g_j\stackrel{~}{h}_{FF^1(f,\frac{1}{4^n})}f\underset{j=1}{\overset{n}{}}g_jh_{FF^1(f,\frac{1}{4^n})}.$$
Therefore, $`f_{j=1}^ng_j\stackrel{~}{h}_{FF^1(f,\frac{1}{4^n})}{\displaystyle \frac{1}{2^n}}.`$ Now define
$$g_{n+1}=\{\begin{array}{cc}\stackrel{~}{h}& \text{on }KF^1(f_{j=1}^ng_j,\frac{1}{2^n})\\ 0& \text{otherwise}\end{array}.$$
Then $`g_{n+1}`$ is continuous on $`KF^1(f_{j=1}^ng_j,\frac{1}{2^n}).`$ This proves (a). Furthermore,
$$f\underset{j=1}{\overset{n+1}{}}g_j_{FF^1(f,\frac{1}{4^n})}=f\underset{j=1}{\overset{n}{}}g_j\stackrel{~}{h}_{FF^1(f,\frac{1}{4^n})}\frac{1}{2^n}.$$
This proves (b). Also,
$$g_{n+1}_K\stackrel{~}{h}_{KF^1(f_{j=1}^ng_j,\frac{1}{2^n})}\frac{1}{2^{n1}}$$
by the definition of $`\stackrel{~}{h}.`$ This proves (c). Finally, suppose $`0<\delta {\displaystyle \frac{1}{2^{n1}}}.`$ Assume that $`xFF^1(f,{\displaystyle \frac{\delta }{2^{n+1}}}).`$ Then $`xF^1(f_{j=1}^ng_j,{\displaystyle \frac{\delta }{2}}).`$ Thus there exists a neighborhood $`U_1`$ of $`x`$ in $`F`$ such that
$$\left|\left(f\underset{j=1}{\overset{n}{}}g_j\right)\left(x_1\right)\left(f\underset{j=1}{\overset{n}{}}g_j\right)\left(x_2\right)\right|<\frac{\delta }{2}$$
whenever $`x_1,x_2U_1.`$ Note that since $`xFF^1(f_{j=1}^ng_j,{\displaystyle \frac{\delta }{2}}),`$ $`g_{n+1}`$ is continuous at $`x.`$ Therefore, there exists a neighborhood $`U_2`$ of $`x`$ in $`F`$ such that $`\left|g_{n+1}\left(x_1\right)g_{n+1}\left(x_2\right)\right|<{\displaystyle \frac{\delta }{2}}`$ for all $`x_1,x_2U_2.`$ Let $`U=U_1U_2.`$ Then $`U`$ is a neighborhood of $`x`$ in $`F`$ such that
$$\left|\left(f\underset{j=1}{\overset{n+1}{}}g_j\right)\left(x_1\right)\left(f\underset{j=1}{\overset{n+1}{}}g_j\right)\left(x_2\right)\right|<\delta $$
whenever $`x_1,x_2U.`$ Hence $`xF^1(f_{j=1}^{n+1}g_j,\delta ).`$ This proves (d). ∎
###### Theorem 3.6.
Let $`F`$ be a closed subspace of $`K`$ and let $`f`$ be a Baire-1 function on $`F.`$ Then there exists a Baire-1 function $`g`$ on $`K`$ such that
$$g_{|F}=f\text{ and }\beta \left(g\right)=\beta _F\left(f\right).$$
###### Proof.
Let $`\left(g_n\right)`$ be the sequence given by Lemma 3.5. Define $`g`$ on $`K`$ by
$$g=\{\begin{array}{cc}_{j=1}^{\mathrm{}}g_j& \text{on }KF\\ f& \text{on }F\end{array}.$$
Note that by (c) of Lemma 3.5, $`_{j=1}^{\mathrm{}}g_j`$ converges uniformly on $`K.`$ Hence $`g`$ is well defined. Obviously, $`g_{|F}=f.`$
Claim. $`K^1(g,\frac{1}{2^{n3}})F^1(f,\frac{1}{4^n})`$ for all $`n.`$
Proof of Claim. Let $`xKF^1(f,\frac{1}{4^n}).`$ We consider two cases. Suppose $`xF.`$ By Lemma 3.5(a), $`g_j`$ is continuous on $`KF`$ for all $`j.`$ Since $`_{j=1}^ng_j`$ converges uniformly to $`g`$ on $`KF,`$ and $`KF`$ is an open subset of $`K,`$ $`g`$ is continuous at $`x.`$ Hence $`xK^1(g,\frac{1}{2^{n3}}).`$ Now suppose $`xF.`$ Then $`xFF^1(f,\frac{1}{4^n}).`$ There is a neighborhood $`U_1`$ of $`x`$ in $`K`$ such that $`\left|f\left(x\right)f\left(x^{}\right)\right|<{\displaystyle \frac{1}{4^n}}`$ for all $`x^{}U_1F.`$ Also, for $`1kn,`$
$`F^1(f{\displaystyle \underset{j=1}{\overset{k}{}}}g_j,{\displaystyle \frac{1}{2^k}})`$ $`F^1(f,{\displaystyle \frac{1}{4^k}})\text{ by Lemma }\text{3.5}\text{(d),}`$
$`F^1(f,{\displaystyle \frac{1}{4^n}}).`$
Since $`g_{k+1}`$ is continuous on $`KF^1(f_{j=1}^kg_j,\frac{1}{2^k}),`$ $`g_{k+1}`$ is continuous on $`KF^1(f,\frac{1}{4^n})`$ for all $`k,`$ $`1kn.`$ Similarly, $`F^1(f,1)F^1(f,\frac{1}{4^n})`$ and $`g_1`$ is continuous on $`KF^1(f,1)`$ by Lemma 3.5(a); thus, $`g_1`$ is continuous on $`KF^1(f,\frac{1}{4^n}).`$ Hence there exists a neighborhood $`U_2`$ of $`x`$ in $`K`$ such that $`U_2KF^1(f,\frac{1}{4^n})`$ and
$$\left|\underset{j=1}{\overset{n+1}{}}g_j\left(x^{}\right)\underset{j=1}{\overset{n+1}{}}g_j\left(x\right)\right|<\frac{1}{2^n}\text{ for all }x^{}U_2.$$
Let $`U=U_1U_2.`$ Then $`U`$ is a neighborhood of $`x`$ in $`K`$. If $`x^{}UF,`$ then $`x^{}U_1F.`$ Thus $`\left|g\left(x^{}\right)g\left(x\right)\right|=\left|f\left(x^{}\right)f\left(x\right)\right|<{\displaystyle \frac{1}{4^n}}<{\displaystyle \frac{1}{2^{n2}}}.`$ If $`x^{}UF,`$ then
$`\left|g\left(x^{}\right)g\left(x\right)\right|`$ $`=\left|{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}g_j\left(x^{}\right)f\left(x\right)\right|`$
$`\left|{\displaystyle \underset{j=1}{\overset{n+1}{}}}g_j\left(x^{}\right){\displaystyle \underset{j=1}{\overset{n+1}{}}}g_j\left(x\right)\right|+\left|{\displaystyle \underset{j=1}{\overset{n+1}{}}}g_j\left(x\right)f\left(x\right)\right|+\left|{\displaystyle \underset{j=n+2}{\overset{\mathrm{}}{}}}g_j\left(x^{}\right)\right|`$
$`<{\displaystyle \frac{1}{2^n}}+|{\displaystyle \underset{j=1}{\overset{n+1}{}}}g_j\left(x\right)f\left(x\right)|+{\displaystyle \underset{j=n+2}{\overset{\mathrm{}}{}}}g_j\text{ since }x^{}U_2,`$
$`{\displaystyle \frac{1}{2^n}}+\left|{\displaystyle \underset{j=1}{\overset{n+1}{}}}g_j\left(x\right)f\left(x\right)\right|+{\displaystyle \underset{j=n+2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{2^{j2}}},\text{ by Lemma }\text{3.5}\text{(c),}`$
$`{\displaystyle \frac{1}{2^n}}+{\displaystyle \frac{1}{2^n}}+{\displaystyle \frac{1}{2^{n1}}},\text{ by Lemma }\text{3.5}\text{(b), since }xFF^1(f,{\displaystyle \frac{1}{4^n}}),`$
$`={\displaystyle \frac{1}{2^{n2}}}.`$
Thus $`\left|g\left(x^{}\right)g\left(x\right)\right|<{\displaystyle \frac{1}{2^{n2}}}`$ if $`x^{}U.`$ Hence $`\left|g\left(x_1\right)g\left(x_2\right)\right|<{\displaystyle \frac{1}{2^{n3}}}`$ whenever $`x_1,`$ $`x_2U.`$ Therefore $`xK^1(g,{\displaystyle \frac{1}{2^{n3}}}).`$ This proves the claim.
It follows by induction that $`K^\alpha (g,{\displaystyle \frac{1}{2^{n3}}})F^\alpha (f,{\displaystyle \frac{1}{4^n}})`$ for $`1\alpha <\omega _1.`$ Indeed, the Claim yields the assertion for $`\alpha =1.`$ If the inclusion holds for some $`\alpha ,`$ $`1\alpha <\omega _1,`$ let $`\stackrel{~}{F}=F^\alpha (f,{\displaystyle \frac{1}{4^n}}).`$ Then $`K^{\alpha +1}(g,{\displaystyle \frac{1}{2^{n3}}})\stackrel{~}{F}^1(g,{\displaystyle \frac{1}{2^{n3}}})=\stackrel{~}{F}^1(f,{\displaystyle \frac{1}{2^{n3}}})\stackrel{~}{F}^1(f,{\displaystyle \frac{1}{4^n}})=F^{\alpha +1}(f,{\displaystyle \frac{1}{4^n}}).`$ Hence the inclusion holds for $`\alpha +1.`$ If the inclusion holds for all $`1\alpha ^{}<\alpha ,`$ where $`\alpha <\omega _1`$ is a limit ordinal, then
$$K^\alpha (g,\frac{1}{2^{n3}})=\underset{1\alpha ^{}<\alpha }{}K^\alpha ^{}(g,\frac{1}{2^{n3}})\underset{1\alpha ^{}<\alpha }{}F^\alpha ^{}(f,\frac{1}{4^n})=F^\alpha (f,\frac{1}{4^n}).$$
This proves the inclusion for $`1\alpha <\omega _1.`$ In particular, if $`\beta _F\left(f\right)=\beta _0,`$ then $`K^{\beta _0}(g,{\displaystyle \frac{1}{2^{n3}}})F^{\beta _0}(f,{\displaystyle \frac{1}{4^n}})=\mathrm{}.`$ Thus $`\beta _K(g,{\displaystyle \frac{1}{2^{n3}}})\beta _0`$ for all $`n.`$ Hence $`\beta _K\left(g\right)\beta _0.`$ Of course, since $`g_{|F}=f,`$ $`\beta _K\left(g\right)\beta _F\left(f\right)\beta _0.`$ Therefore $`\beta _K\left(g\right)=\beta _0=\beta _F\left(f\right).`$
###### Remark 3.7.
*If* $`\beta _F\left(f\right)=1,`$ *Theorem 3.6 is the familiar Tietze Extension Theorem. If* $`\beta _F\left(f\right)`$ *is transfinite, the conclusion of Theorem 3.6 can be obtained easily by defining the extension* $`g`$ *to be* $`0`$ *on* $`KF.`$ *However, we do not see a simple proof for finite* $`\beta _F\left(f\right).`$
## 4. Decomposition of Baire-1 functions
In this section, we give a proof of Theorem 1.2. The extension results in §3 are employed in the course of the proof.
###### Theorem 4.1.
Let $`f`$ be a Baire-1 function on $`K,`$ $`1\beta _0,\gamma _0<\omega _1`$ and $`\epsilon >0.`$ Then there exist
$$\stackrel{~}{f}:KK^{\beta _0\gamma _0}(f,\epsilon )$$
and
$$f_n:KK^{\beta _0\gamma _0}(f,\epsilon )$$
such that $`(f_n)`$ converges to $`f`$ pointwise, $`\stackrel{~}{f}f_{KK^{\beta _0\gamma _0}(f,\epsilon )}\epsilon `$ and $`\beta _H\left(f_n\right)\beta _0,`$ $`\gamma _H\left(\left(f_n\right)\right)\gamma _0`$ for all compact subsets $`H`$ of $`KK^{\beta _0\gamma _0}(f,\epsilon ).`$
###### Proof.
For $`\alpha \gamma _0,`$ let $`K_\alpha =K^{\beta _0\alpha }(f,\epsilon ).`$ If $`n,`$ let $`U_n^\alpha `$ be the $`{\displaystyle \frac{1}{n}}`$neighborhood of $`K_\alpha `$ in $`K.`$ For $`\alpha <\gamma _0,`$ it follows from Theorem 3.2 that there exists $`g_\alpha :K_\alpha K_{\alpha +1}`$ such that $`g_\alpha f_{K_\alpha K_{\alpha +1}}\epsilon `$ and $`\beta _H\left(g_\alpha \right)\beta _0`$ for all compact subsets $`H`$ of $`K_\alpha K_{\alpha +1}.`$ List the ordinals in $`[0,\gamma _0)`$ in a (possibly finite) sequence $`\left(\alpha _n\right)_{n=1}^p.`$ Here $`p`$ or $`p=\mathrm{}.`$ For each $`n,`$ let $`F_n=_{j=1}^{np}\left(K_{\alpha _j}U_n^{\alpha _j+1}\right).`$ Then $`F_n`$ is a closed subset of $`K.`$ It is also easy to see that $`K_\alpha U_n^{\alpha +1}`$ and $`K_\alpha ^{}U_n^{\alpha ^{}+1}`$ are disjoint if $`\alpha \alpha ^{}.`$ Thus $`\left(K_{\alpha _j}U_n^{\alpha _j+1}\right)_{j=1}^{np}`$ is a partition of $`F_n`$ into clopen (in $`F_n`$) subsets. Now define $`\stackrel{~}{g}_n:F_nK`$ to be $`_{j=1}^{np}g_{\alpha _j|K_{\alpha _j}U_n^{\alpha _j+1}}.`$ Since $`H=K_{\alpha _j}U_n^{\alpha _j+1}`$ is a compact subset of $`K_{\alpha _j}K_{\alpha _j+1},`$ $`\beta _H\left(g_{\alpha _j}\right)\beta _0.`$ From the clopeness of the partition $`\left(K_{\alpha _j}U_n^{\alpha _j+1}\right)_{j=1}^{np},`$ it follows readily that $`\beta _{F_n}\left(\stackrel{~}{g}_n\right)\beta _0.`$ By Theorem 3.6, there exists a function $`f_n^{}`$ on $`K`$ such that $`f_{n|F_n}^{}=\stackrel{~}{g}_n`$ and $`\beta _K\left(f_n^{}\right)\beta _0.`$ Finally, define $`f_n`$ to be $`f_{n|KK_{\gamma _0}}^{}`$ and $`\stackrel{~}{f}`$ to be $`_{\alpha <\gamma _0}g_{\alpha |K_\alpha K_{\alpha +1}}.`$ It follows from the choices of the $`g_\alpha `$’s that $`f\stackrel{~}{f}_{KK_{\gamma _0}}\epsilon .`$ Since $`_{n=1}^{\mathrm{}}F_n=KK_{\gamma _0},`$ $`limf_n=\stackrel{~}{f}`$ pointwise on $`KK_{\gamma _0}.`$ Suppose $`H`$ is a compact subset of $`KK_{\gamma _0}.`$ Then $`\beta _H\left(f_n\right)\beta _K\left(f_n^{}\right)\beta _0.`$ To complete the proof, we claim that for any $`\delta >0`$ and any $`\gamma \gamma _0,`$ $`H^\gamma (\left(f_n\right),\delta )K_\gamma .`$ The proof of this is by induction on $`\gamma .`$ The case $`\gamma =0`$ and the limit case is trivial. Now assume that the claim holds for some $`\gamma <\gamma _0.`$ Let $`xH^\gamma (\left(f_n\right),\delta )K_{\gamma +1}.`$ Choose $`j_1,j_2`$ such that $`\alpha _{j_1}=\gamma `$ and $`d(x,K_{\gamma +1}){\displaystyle \frac{1}{j_2}},`$ where $`d`$ is the metric on $`K.`$ Denote $`H^\gamma (\left(f_n\right),\delta )`$ by $`L`$ and the $`{\displaystyle \frac{1}{2j_0}}`$-ball in $`K`$ centered at $`x`$ by $`U,`$ where $`j_0=\mathrm{max}\{j_1,2j_2\}.`$ Note that $`LK_\gamma `$ by the inductive hypothesis: For all $`nj_0=\mathrm{max}\{j_1,2j_2\},`$
$$LUL\overline{U}K_{\alpha _{j_1}}U_n^{\alpha _{j_1}+1}F_n.$$
This implies that $`f_{n|L\overline{U}}=\stackrel{~}{g}_{n|L\overline{U}}=g_{\alpha _{j_1}|L\overline{U}}=g_{\gamma |L\overline{U}}`$ for all $`nj_0.`$ Thus $`\left(L\overline{U}\right)^1(\left(f_n\right),\delta )=\mathrm{}.`$ By Lemma 2.1(d),
$$L^1(\left(f_n\right),\delta )\left(LU\right)=\mathrm{}.$$
In particular,
$$xL^1(\left(f_n\right),\delta )=H^{\gamma +1}(\left(f_n\right),\delta ).$$
Since $`xH^\gamma (\left(f_n\right),\delta )K_{\gamma +1}`$ is arbitrary, this shows that $`H^{\gamma +1}(\left(f_n\right),\delta )K_{\gamma +1}.`$
In particular, if $`\beta _K\left(f\right)\beta _0\gamma _0,`$ we have the following.
###### Theorem 4.2.
Let $`f`$ be a Baire-1 function on $`K,`$ $`1\beta _0,\gamma _0<\omega _1,`$ and $`\beta \left(f\right)\beta _0\gamma _0.`$ For any $`\epsilon >0,`$ there exist $`\stackrel{~}{f}:K`$ and a sequence of functions $`f_n:K`$ such that $`\left(f_n\right)`$ converges to $`\stackrel{~}{f}`$ pointwise, $`\stackrel{~}{f}f\epsilon `$, $`\beta \left(f_n\right)\beta _0`$ for all $`n`$, and $`\gamma \left(\left(f_n\right)\right)\gamma _0.`$
A couple more preparatory steps will allow us to improve Theorem 4.2 to an exact result (i.e., $`\epsilon =0`$) when $`\gamma _0`$ is of the right form.
###### Theorem 4.3 (\[3, Lemma 2.5\]).
If $`\left(f_n\right)`$ and $`\left(g_n\right)`$ are two sequences of real-valued functions on $`K`$ such that $`\gamma \left(\left(f_n\right)\right)\omega ^\xi `$ and $`\gamma \left(\left(g_n\right)\right)\omega ^\xi `$ for some $`\xi <\omega _1,`$ then $`\gamma \left(\left(f_n+g_n\right)\right)\omega ^\xi .`$
###### Proposition 4.4.
For $`1\xi <\omega _1,`$ $`𝔅_1^\xi \left(K\right)=\{f^K:\beta \left(f\right)\omega ^\xi \}`$ is a vector subspace of $`^K`$ that is closed under the topology uniform convergence.
We postpone the proof of Proposition 4.4 until the next section. We are now ready to prove the converse of Theorem 2.3 in certain cases.
###### Theorem 4.5.
If $`f𝔅_1\left(K\right)`$ and $`\beta \left(f\right)\beta _0\omega ^{\gamma _0}`$ for some $`1\beta _0<\omega _1`$ and $`\gamma _0<\omega _1,`$ then there exists $`\left(f_n\right)𝔅_1\left(K\right)`$ such that $`\left(f_n\right)`$ converges pointwise to $`f,`$ $`\beta \left(f_n\right)\beta _0`$ for all $`n`$ and $`\gamma \left(\left(f_n\right)\right)\omega ^{\gamma _0}.`$
###### Proof.
First we assume $`\beta _0`$ is of the form $`\omega ^{\alpha _0},`$ where $`\alpha _0<\omega _1.`$ By Theorem 4.2 there exist a sequence $`\left(f_n^1\right)𝔅_1\left(K\right)`$ and a function $`f^1𝔅_1\left(K\right)`$ such that, $`\beta \left(f_n^1\right)\omega ^{\alpha _0}`$ for all $`n,`$ $`\left(f_n^1\right)`$ converges pointwise to $`f^1,`$ $`f^1f{\displaystyle \frac{1}{2}},`$ and $`\gamma \left(\left(f_n^1\right)\right)\omega ^{\gamma _0}.`$ Then $`\beta \left(f^1\right)\omega ^{\alpha _0}\omega ^{\gamma _0}=\omega ^{\alpha _0+\gamma _0}`$ by Theorem 2.3. This implies that $`\beta \left(ff^1\right)\omega ^{\alpha _0+\gamma _0}`$ by Proposition 4.4. Hence there exist $`\left(f_n^2\right)𝔅_1\left(K\right)`$ and $`f^2`$ such that $`\beta \left(f_n^2\right)\omega ^{\alpha _0}`$ for all $`n`$, $`\left(f_n^2\right)`$ converges pointwise to $`f^2,ff^1f^2{\displaystyle \frac{1}{2^2}},`$ and $`\gamma \left(\left(f_n^2\right)\right)\omega ^{\gamma _0}.`$ We may assume that $`f_n^2{\displaystyle \frac{1}{2}}`$ for all $`n,`$ for otherwise, simply replace $`f_n^2`$ by $`\widehat{f}_n^2=\left(f_n^2\frac{1}{2}\right)\frac{1}{2}.`$ Continuing, we obtain $`f^m`$ and $`\left(f_n^m\right)_{n=1}^{\mathrm{}}`$ for each $`m`$ such that
* $`f_n^m{\displaystyle \frac{1}{2^{m1}}},`$
* $`\beta \left(f_n^m\right)\omega ^{\alpha _0}`$ for all $`m,n,`$
* $`\gamma \left(\left(f_n^m\right)_n\right)\omega ^{\gamma _0}`$ for all $`m,`$
* $`f^m=\underset{n}{lim}f_n^m`$ (pointwise) for all $`m,`$ and
* $`_{m=1}^{\mathrm{}}f^m`$ converges uniformly to $`f`$ on $`K.`$
Let $`g_n^m=f_n^1+f_n^2+\mathrm{}+f_n^m`$ and $`g_n=_{m=1}^{\mathrm{}}f_n^m.`$ By Theorem 4.3, $`\gamma \left(\left(g_n^m\right)_n\right)\omega ^{\gamma _0}`$ for all $`m.`$ Given $`\epsilon >0,`$ there exists $`m_0`$ such that for all $`n,`$ $`g_n^{m_0}g_n\epsilon .`$ Then $`K^{\omega ^{\gamma _0}}(\left(g_n\right),3\epsilon )K^{\omega ^{\gamma _0}}(\left(g_n^{m_0}\right),\epsilon )=\mathrm{}.`$ Therefore $`\gamma \left(\left(g_n\right)\right)\omega ^{\gamma _0}.`$ By Proposition 4.4, $`\beta \left(g_n^m\right)\omega ^{\alpha _0}`$ for all $`m,n.`$ Therefore, $`\beta \left(g_n\right)\omega ^{\alpha _0}`$ by Proposition 4.4. Moreover,
$`\underset{n}{lim}g_n`$ $`=\underset{n}{lim}\underset{m}{lim}g_n^m=\underset{m}{lim}\underset{n}{lim}g_n^m`$
$`=\underset{m}{lim}{\displaystyle \underset{k=1}{\overset{m}{}}}f^k=f\text{ pointwise.}`$
This proves the theorem in case $`\beta _0=\omega ^{\alpha _0}`$, with $`\left(g_n\right)`$ in place of $`\left(f_n\right).`$
For a general nonzero countable ordinal $`\beta _0,`$ write $`\beta _0`$ in Cantor normal form as
$$\beta _0=\omega ^{\beta _1}m_1+\omega ^{\beta _2}m_2+\mathrm{}+\omega ^{\beta _k}m_k,$$
where $`k,m_1,\mathrm{},m_k,`$ $`\omega _1>\beta _1>\beta _2>\mathrm{}>\beta _k.`$ If $`\gamma _00,`$ then $`\beta _0\omega ^{\gamma _0}=\omega ^{\beta _1}\omega ^{\gamma _0}.`$ By the previous case, there exists $`\left(f_n\right)𝔅_1\left(K\right)`$ such that $`\beta \left(f_n\right)\omega ^{\beta _1}\beta _0,`$ $`\gamma \left(\left(f_n\right)\right)\omega ^{\gamma _0}`$ and $`\left(f_n\right)`$ converges pointwise to $`f.`$ If $`\gamma _0=0,`$ take $`f_n=f`$ for all $`n.`$ Then $`\beta \left(f_n\right)\beta _0`$ for all $`n`$, $`\gamma \left(\left(f_n\right)\right)=1=\omega ^{\gamma _0}`$ and $`\left(f_n\right)`$ converges pointwise to $`f`$. ∎
The combination of Theorem 2.3 and Corollary 4.6 yields Theorem 1.2.
###### Corollary 4.6.
Let $`f`$ $`𝔅_1^\xi \left(K\right),`$ respectively, $`_1^\xi \left(K\right),`$ for some $`\xi <\omega _1.`$ For all countable ordinals $`\mu ,`$ $`\nu `$ such that $`\mu +\nu \xi ,`$ there exists a sequence $`\left(f_n\right)𝔅_1^\mu \left(K\right)`$, respectively, $`_1^\mu \left(K\right),`$ such that $`f_nf`$ pointwise, and $`\gamma \left(\left(f_n\right)\right)\omega ^\nu .`$
We do not know if Theorem 4.5 holds without the restriction on the form of the ordinal $`\gamma \left(\left(f_n\right)\right).`$
###### Problem 4.7.
Is it true that if $`f𝔅_1\left(K\right)`$ with $`\beta \left(f\right)\beta _0\gamma _0`$ for some countable ordinals $`\beta _0`$ and $`\gamma _0,`$ then there exists a sequence $`\left(f_n\right)`$ converging pointwise to $`f`$ so that $`\underset{n}{sup}\beta \left(f_n\right)\beta _0`$ and $`\gamma \left(\left(f_n\right)\right)\gamma _0\mathrm{?}`$
As another application of our results, we give the proof of another characterization of the classes $`_1^\xi \left(K\right)`$ due to Kechris and Louveau.
###### Definition 4.8 (\[3, Section 3\]).
A family $`\{\mathrm{\Phi }_\xi :0\xi <\omega _1\}`$ of real-valued functions on $`K`$ is defined as follows.
$$\mathrm{\Phi }_0=C\left(K\right),$$
$$\mathrm{\Phi }_{\xi +1}=\left\{\begin{array}{c}f:f\text{ is the pointwise limit of a bounded sequence}\\ \text{ }\left(f_n\right)\mathrm{\Phi }_\xi \text{ such that }\gamma \left(\left(f_n\right)\right)\omega .\end{array}\right\},$$
and for limit ordinals $`\lambda ,`$
$$\mathrm{\Phi }_\lambda =\left\{\begin{array}{c}f:f\text{ is the uniform limit of a bounded sequence}\\ \text{ }\left(f_n\right)_{\xi <\lambda }\mathrm{\Phi }_\xi \text{.}\end{array}\right\}.$$
###### Corollary 4.9 (\[3, Theorem 4.2\]).
For each $`\xi <\omega _1,`$ $`_1^\xi \left(K\right)=\mathrm{\Phi }_\xi .`$
###### Proof.
The case $`\xi =0`$ is trivial. Suppose the corollary holds for some $`\xi <\omega _1.`$ If $`f_1^{\xi +1}\left(K\right),`$ it follows from Corollary 4.6 that $`f`$ is the pointwise limit of a bounded sequence $`\left(f_n\right)`$ in $`_1^\xi \left(K\right)`$ such that $`\gamma \left(\left(f_n\right)\right)\omega .`$ Since $`_1^\xi \left(K\right)=\mathrm{\Phi }_\xi `$ by the inductive hypothesis, $`f\mathrm{\Phi }_{\xi +1}.`$ Conversely, if $`f\mathrm{\Phi }_{\xi +1},`$ then $`f`$ is the pointwise limit of a sequence $`\left(f_n\right)`$ in $`\mathrm{\Phi }_\xi `$ with $`\gamma \left(\left(f_n\right)\right)\omega .`$ Since $`\mathrm{\Phi }_\xi =_1^\xi \left(K\right),`$ $`\beta \left(f\right)\omega ^{\xi +1}`$ by Theorem 2.3. Thus $`f_1^{\xi +1}\left(K\right).`$
Now assume that the corollary holds for all $`\xi ^{}<\xi ,`$ where $`\xi `$ is a countable limit ordinal. Let $`f\mathrm{\Phi }_\xi .`$ By the inductive hypothesis, $`\mathrm{\Phi }_\xi ^{}=_1^\xi ^{}\left(K\right)_1^\xi \left(K\right)`$ for $`\xi ^{}<\xi .`$ Hence $`f`$ is the uniform limit of a sequence in $`_1^\xi \left(K\right)`$, and thus belongs to $`_1^\xi \left(K\right).`$ Conversely, assume that $`f_1^\xi \left(K\right).`$ For every $`n,`$ there exists $`\xi _n<\xi `$ such that $`\beta (f,\frac{1}{n})\omega ^{\xi _n}.`$ By Corollary 3.3, the exists $`f_n_1^{\xi _n}\left(K\right)=\mathrm{\Phi }_{\xi _n}`$ such that $`ff_n\frac{1}{n}.`$ Thus $`f\mathrm{\Phi }_\xi ,`$ as required. ∎
###### Remark 4.10.
*If a family* $`\{\mathrm{\Psi }_\xi :0\xi <\omega _1\}`$ *is defined in a similar way as the family* $`\{\mathrm{\Phi }_\xi :0\xi <\omega _1\}`$ *except for the removal of the boundedness condition on the sequence* $`\left(f_n\right),`$ *then* $`\mathrm{\Psi }_\xi =B_1^\xi \left(K\right)`$ *for all* $`\xi <\omega _1.`$
## 5. Optimal limit of continuous functions
In this section we prove the equivalence of the indices $`\beta `$ and $`\gamma `$ for functions in $`𝔅_1\left(K\right)`$ in the same sense that was established for $`_1\left(K\right)`$ in Theorem 2.3 of . Namely, it is shown that for all $`f𝔅_1\left(K\right),`$ $`\beta \left(f\right)`$ is the smallest ordinal $`\gamma _0`$ for which there exists a sequence $`\left(f_n\right)`$ in $`C\left(K\right)`$ converging pointwise to $`f`$ and satisfying $`\gamma \left(\left(f_n\right)\right)\gamma _0.`$ Let us note that this result is also the converse of Theorem 2.3 when $`\beta _0=1.`$
###### Definition 5.1.
Let $`\left(f_n\right)^K`$ and $`f^K.`$ We write
(a) $`\left(g_n\right)\left(f_n\right)`$ if $`\left(g_n\right)`$ is a convex block combination of $`\left(f_n\right),`$ i.e., there exists a sequence of non-negative real numbers $`\left(a_k\right)`$ and a strictly increasing sequence $`\left(p_n\right)`$ in $``$ such that $`_{k=p_{n1}+1}^{p_n}a_k=1`$ and $`g_n=_{k=p_{n1}+1}^{p_n}a_kf_k`$ for all $`n`$ $`\left(p_0=0\right).`$
(b) $`\left(g_n\right)\stackrel{𝑎}{}\left(f_n\right)`$ if there exists $`m`$ such that $`\left(g_n\right)_{n=m}^{\mathrm{}}\left(f_n\right),`$ and
(c) $`\left[f\right]_M^M=(fM)M,`$ where $`0M.`$
The easy proof of the next lemma is left to the reader.
###### Lemma 5.2.
If $`\left(g_n\right)\stackrel{𝑎}{}\left(f_n\right),`$ then $`\gamma (\left(g_n\right),\epsilon )\gamma (\left(f_n\right),\epsilon )`$ for all $`\epsilon >0.`$
###### Lemma 5.3.
Let $`f`$ be a Baire-1 function on $`K.`$ Suppose $``$ is a countable collection of compact subsets of $`K`$ such that $`f_H<\mathrm{}`$ for all $`H`$ and $`_HH=K.`$ Then there exists $`\left(f_n\right)C\left(K\right)`$ such that
(i) $`f_n`$ $``$ $`f`$ pointwise, and
(ii) $`\left(f_{n|H}\right)`$ is a bounded subset of $`C\left(H\right)`$ for all $`H.`$
###### Proof.
Write $``$ as a sequence $`\left(H_m\right)_{m=1}^{\mathrm{}}.`$ Without loss of generality, assume that $`H_mH_{m+1}`$ for all $`m.`$ Since $`f`$ is Baire-1, there exists $`\left(f_n^0\right)C\left(K\right)`$ such that $`\left(f_n^0\right)`$ converges pointwise to $`f.`$ Assume that $`\left(f_n^{m1}\right)_nC\left(K\right)`$ has been chosen so that $`\underset{n}{lim}f_n^{m1}=f`$ pointwise. If $`m,n`$, let $`U_n^m`$ be the $`{\displaystyle \frac{1}{n}}`$neighborhood of $`H_m`$ in $`K`$ and let $`M_m=f_{H_m}.`$ For all $`n,`$ the function $`\left[f_n^{m1}\right]_{M_m|H_m}^{M_m}f_{n|KU_n^m}^{m1}`$ is continuous on $`H_m\left(KU_n^m\right).`$ Let $`f_n^m`$ be a continuous extension of the function onto $`K.`$ Then $`\left(f_n^m\right)C\left(K\right).`$ If $`xH_m,`$ then $`\underset{n}{lim}f_n^m\left(x\right)=\underset{n}{lim}\left[f_n^{m1}\left(x\right)\right]_{M_m}^{M_m}=\left[f\left(x\right)\right]_{M_m}^{M_m}=f\left(x\right)`$ since $`f_{H_m}=M_m.`$ If $`xH_m,`$ then there exists $`n_0`$ such that $`xKU_{n_0}^m;`$ thus $`xKU_n^m`$ for all $`nn_0.`$ Therefore $`f_n^m\left(x\right)=f_n^{m1}\left(x\right)`$ for all $`nn_0.`$ Hence $`\underset{n}{lim}f_n^m\left(x\right)=f\left(x\right).`$ Thus $`\underset{n}{lim}f_n^m=f`$ pointwise. Now for each $`n,`$ let $`f_n=f_n^n.`$ Since $`H_mH_n`$ for all $`nm`$, on $`H_m`$ we have
$`f_n`$ $`=f_n^n=\left[f_n^{n1}\right]_{M_n}^{M_n}`$
$`=\left[\left[f_n^{n2}\right]_{M_{n1}}^{M_{n1}}\right]_{M_n}^{M_n}=\mathrm{}=\left[\mathrm{}\left[\left[f_n^{m1}\right]_{M_m}^{M_m}\right]_{M_{m+1}}^{M_{m+1}}\mathrm{}\right]_{M_n}^{M_n}`$
$`=\left[f_n^{m1}\right]_{M_m}^{M_m}\text{ as }M_mM_{m+1}\mathrm{}M_n.`$
Thus $`f_n=\left[f_n^{m1}\right]_{M_m}^{M_m}`$ on $`H_m`$ for all $`nm.`$ In particular, on the set $`H_m,`$
$$\underset{n}{lim}f_n=\left[\underset{n}{lim}f_n^{m1}\right]_{M_m}^{M_m}=\left[f\right]_{M_m}^{M_m}=f$$
since $`f_{H_m}=M_m.`$ As $`K=H_m,`$ we see that $`f_nf`$ pointwise. Also, for each $`m,`$ $`\left(f_{n|H_m}\right)_{n=m}^{\mathrm{}}`$ is bounded (by $`M_m`$) in $`C\left(H_m\right);`$ thus $`\left(f_{n|H_m}\right)_{n=1}^{\mathrm{}}`$ is bounded in $`C\left(H_m\right).`$
For the next lemma, recall that for a real-valued function $`f`$ defined on a set $`S,`$ $`osc(f,S)=sup\left\{\right|f\left(s_1\right)f\left(s_2\right)|:s_1,s_2S\}.`$
###### Lemma 5.4.
Let $`\left(f_n\right)`$ be bounded in $`C\left(H\right),`$ where $`H`$ is a compact metric space. Suppose $`\left(f_n\right)`$ converges pointwise to $`f`$ and $`H^1(f,\epsilon )=\mathrm{}`$ for some $`\epsilon >0,`$ then there exists $`\left(g_n\right)\left(f_n\right)`$ such that $`H^1(\left(g_n\right),7\epsilon )=\mathrm{}.`$
###### Proof.
By Corollary 3.3, there exists $`\stackrel{~}{f}C\left(H\right)`$ such that $`f\stackrel{~}{f}_H\epsilon .`$ Then $`\left(f_n\stackrel{~}{f}\right)`$ is bounded in $`C\left(H\right),`$ $`f_n\stackrel{~}{f}f\stackrel{~}{f}`$ pointwise and $`osc(f\stackrel{~}{f},H)2\epsilon `$. By the first statement in the proof of Theorem 2.3 in , there exists $`\left(h_n\right)\left(f_n\stackrel{~}{f}\right)`$ such that $`h_n(f\stackrel{~}{f})_H3\epsilon .`$ Let $`g_n=h_n+\stackrel{~}{f}`$ for all $`n`$. Then $`\left(g_n\right)\left(f_n\right)`$ and $`g_nf_H3\epsilon `$ for all $`n`$. It follows that $`H^1(\left(g_n\right),7\epsilon )=\mathrm{}.`$
###### Theorem 5.5.
Let $`f`$ be a Baire-1 function on $`K.`$ There exists a sequence $`\left(f_n\right)C\left(K\right)`$ such that $`\left(f_n\right)`$ converges pointwise to $`f`$ and $`\gamma \left(\left(f_n\right)\right)=\beta \left(f\right).`$
###### Proof.
Let $`\beta _0=\beta \left(f\right).`$ For each $`\alpha <\beta _0,`$ and all $`m,j`$, let $`U_{m,j}^\alpha `$ be the $`{\displaystyle \frac{1}{j}}`$neighborhood of $`K^\alpha (f,{\displaystyle \frac{1}{m}})`$ in $`K.`$ Define
$$=\{K^\alpha (f,\frac{1}{m})U_{m,j}^{\alpha +1}:\alpha <\beta _0,m,j\}.$$
Then $``$ is a countable collection of compact subsets of $`K`$ such that $`_HH=K.`$ If $`\alpha <\beta _0`$ and $`m,j`$, by Lemma 3.1, there is a continuous function $`g`$ on $`H=K^\alpha (f,{\displaystyle \frac{1}{m}})U_{m,j}^{\alpha +1}`$ such that $`gf_H{\displaystyle \frac{1}{m}}.`$ Hence $`f_H<\mathrm{}`$ for all $`H.`$ By Lemma 5.3, there exists $`\left(g_n\right)C\left(K\right)`$ such that $`\left(g_n\right)`$ converges pointwise to $`f`$ and $`\left(g_{n|H}\right)`$ is bounded in $`C\left(H\right)`$ for all $`H`$.
List the elements of $``$ in a sequence $`\left(H_k\right)_{k=1}^{\mathrm{}}.`$ Take $`\epsilon _k={\displaystyle \frac{1}{m}}`$ if $`H_k`$ is of the form $`K^\alpha (f,{\displaystyle \frac{1}{m}})U_{m,j}^{\alpha +1}`$ for some $`\alpha ,m,j.`$ Let $`\left(g_n^0\right)=\left(g_n\right).`$ Suppose $`\left(g_n^{k1}\right)_n\left(g_n\right)_n`$ has been chosen. Then $`\left(g_n^{k1}\right)_n`$ converges to $`f`$ pointwise, $`\left(g_{n|H_k}^{k1}\right)`$ is a bounded sequence in $`C\left(H_k\right),`$ and $`\left(H_k\right)^1(f,\epsilon _k)=\mathrm{}.`$ By Lemma 5.4, there exists $`\left(g_n^k\right)_n\left(g_n^{k1}\right)_n`$ such that $`\left(H_k\right)^1(\left(g_n^k\right)_n,7\epsilon _k)=\mathrm{}.`$ Let $`f_n=g_n^n`$ for all $`n`$. Then $`\left(f_n\right)\left(g_n\right).`$ Therefore $`\left(f_n\right)C\left(K\right)`$ and $`\left(f_n\right)`$ converges pointwise to $`f`$. We claim that for all $`m`$ and for all $`\alpha \beta _0,`$ $`K^\alpha (\left(f_n\right),{\displaystyle \frac{7}{m}})K^\alpha (f,{\displaystyle \frac{1}{m}}).`$ We prove the claim by induction on $`\alpha .`$ The claim is trivial if $`\alpha =0`$ or $`\alpha `$ is a limit ordinal. Assume that $`\alpha \beta _0`$ is a successor ordinal and that the claim holds for $`\alpha 1.`$ Let $`xK^\alpha (\left(f_n\right),{\displaystyle \frac{7}{m}}).`$ Then $`xK^{\alpha 1}(\left(f_n\right),{\displaystyle \frac{7}{m}})K^{\alpha 1}(f,{\displaystyle \frac{1}{m}}).`$ If $`xK^\alpha (f,{\displaystyle \frac{1}{m}}),`$ there exists $`j`$ such that $`d(x,K^\alpha (f,{\displaystyle \frac{1}{m}}))>{\displaystyle \frac{1}{j}}.`$ Choose $`k`$ such that $`H_k=K^{\alpha 1}(f,{\displaystyle \frac{1}{m}})U_{m,j}^\alpha .`$ Then $`\left(f_n\right)\stackrel{𝑎}{}\left(g_n^k\right)_n`$ and $`\gamma _{H_k}(\left(g_n^k\right)_n,7\epsilon _k)1`$ since $`\left(H_k\right)^1(\left(g_n^k\right)_n,7\epsilon _k)=\mathrm{}.`$ By Lemma 5.2, $`\left(H_k\right)^1(\left(f_n\right),7\epsilon _k)=\mathrm{}.`$ Thus $`\left(H_k\right)^1(\left(f_n\right),{\displaystyle \frac{7}{m}})=\mathrm{}.`$ But since $`d(x,K^\alpha (f,{\displaystyle \frac{1}{m}}))>{\displaystyle \frac{1}{j}},`$ there exists an open set $`U`$ in $`\stackrel{~}{K}=K^{\alpha 1}(f,{\displaystyle \frac{1}{m}})`$ such that $`xUH_k\stackrel{~}{K}.`$ By Lemma 2.1(d), $`\left(\stackrel{~}{K}\right)^1(\left(f_n\right),{\displaystyle \frac{7}{m}})U\left(H_k\right)^1(\left(f_n\right),{\displaystyle \frac{7}{m}})=\mathrm{}.`$ Therefore $`x\left(\stackrel{~}{K}\right)^1(\left(f_n\right),{\displaystyle \frac{7}{m}})=K^\alpha (\left(f_n\right),{\displaystyle \frac{7}{m}}),`$ a contradiction. This proves the claim. From the claim, $`K^{\beta _0}(\left(f_n\right),{\displaystyle \frac{7}{m}})K^{\beta _0}(f,{\displaystyle \frac{1}{m}})=\mathrm{}`$ for all $`m.`$ Therefore $`\gamma \left(\left(f_n\right)\right)\beta _0.`$ Since $`\gamma \left(\left(f_n\right)\right)\beta _0`$ by \[3, Proposition 2.1\], (or Theorem 2.3), $`\gamma \left(\left(f_n\right)\right)=\beta _0=\beta \left(f\right).`$
###### Remark 5.6.
*Unlike in Theorem 2.3 of , in general we cannot get a sequence* $`\left(g_n\right)\left(f_n\right)`$ *such that* $`\gamma \left(\left(g_n\right)\right)=\beta \left(f\right).`$ *Indeed, let* $`K=[0,1]`$ *and for each* $`nN`$ *let* $`f_n`$ *be a continuous function that vanishes outside* $`[{\displaystyle \frac{1}{n+1}},{\displaystyle \frac{1}{n}}]`$ *such that* $`_Kf_n=1.`$ *Then* $`\left(f_n\right)`$ *converges pointwise to* $`f=0`$*. Suppose* $`\left(g_n\right)\left(f_n\right),`$ *then* $`_Kg_n=1`$ *for all* $`nN.`$ *Thus* $`\left(g_n\right)`$ *does not converge uniformly to* $`f,`$ *i.e.,* $`\gamma \left(\left(g_n\right)\right)>1=\beta \left(f\right).`$
###### Proof of Proposition 4.4.
It is easy to see that for all $`f𝔅_1^\xi \left(K\right)`$ and $`a,`$ $`af𝔅_1^\xi \left(K\right).`$ If $`f`$, $`g𝔅_1^\xi \left(K\right),`$ then by Theorem 5.5 there exist two sequences of continuous functions $`\left(f_n\right)`$ and $`\left(g_n\right)`$ converging pointwise to $`f`$ and $`g`$ respectively such that $`\gamma \left(\left(f_n\right)\right)\omega ^\xi `$ and $`\gamma \left(\left(g_n\right)\right)\omega ^\xi .`$ According to Theorem 4.3, $`\gamma \left(\left(f_n+g_n\right)\right)\omega ^\xi .`$ Hence by Theorem 2.3, $`f+g𝔅_1^\xi \left(K\right).`$ Finally, given $`f\overline{𝔅_1^\xi \left(K\right)}`$ and $`\epsilon >0`$, choose $`g𝔅_1^\xi \left(K\right)`$ such that $`fg{\displaystyle \frac{\epsilon }{3}}.`$ Then $`K^{\omega ^\xi }(f,\epsilon )K^{\omega ^\xi }(g,{\displaystyle \frac{\epsilon }{3}})=\mathrm{}.`$ Thus $`f𝔅_1^\xi \left(K\right).`$
## 6. Product of Baire-1 functions
In , it is observed that the classes $`_1^\xi \left(K\right),`$ $`\xi <\omega _1`$ are closed under multiplication. However, it is relative easy to see that this fails for the classes $`𝔅_1^\xi \left(K\right).`$ In this section, we show that if $`f𝔅_1^{\xi _1}\left(K\right)`$ and $`g𝔅_1^{\xi _2}\left(K\right),`$ then $`fg𝔅_1^\xi \left(K\right),`$ where $`\xi =\mathrm{max}\{\xi _1+\xi _2,\xi _2+\xi _1\}.`$ It is also shown that the result is sharp. The proof of the next lemma is left to the reader.
###### Lemma 6.1.
If $`f`$ is bounded and $`\gamma \left(\left(g_n\right)\right)\xi ,`$ then $`\gamma \left(\left(fg_n\right)\right)\xi .`$
###### Lemma 6.2.
If $`f_1^{\xi _1}\left(K\right)`$ and $`g𝔅_1^{\xi _2}\left(K\right),`$ then $`fg𝔅_1^{\xi _1+\xi _2}\left(K\right).`$
###### Proof.
By Theorem 5.5, there exists a sequence $`\left(g_n\right)C\left(K\right)`$ converging to $`g`$ pointwise such that $`\gamma \left(\left(g_n\right)\right)=\omega ^{\xi _2}.`$ For each $`n,`$ $`g_nC\left(K\right)_1^{\xi _1}\left(K\right)`$ and $`f_1^{\xi _1}\left(K\right).`$ By (see the remark on \[3, p. 217\]), $`fg_n_1^{\xi _1}\left(K\right).`$ Lemma 6.1 implies that $`\gamma \left(\left(fg_n\right)\right)\omega ^{\xi _2}.`$ Since $`\left(fg_n\right)`$ converges to $`fg`$ pointwise, it follows from Theorem 2.3 that $`\beta \left(fg\right)\omega ^{\xi _1+\xi _2},`$ i.e., $`fg𝔅_1^{\xi _1+\xi _2}\left(K\right).`$
Now suppose $`f𝔅_1^{\xi _1}\left(K\right)`$ and $`g𝔅_1^{\xi _2}\left(K\right).`$ By Lemma 3.1, for all $`\alpha <\omega ^{\xi _2}`$, there is a continuous function $`g_\alpha :K^\alpha (g,1)K^{\alpha +1}(g,1)`$ such that
$$g_\alpha g_{K^\alpha (g,1)K^{\alpha +1}(g,1)}1.$$
Let $`h=_{\alpha <\omega ^{\xi _2}}g_\alpha .`$ It follows from the proof of Theorem 3.2 that $`\beta \left(h\right)\omega ^{\xi _2}.`$ Given a closed set $`HK,`$ we write
$$\underset{f}{d}\left(H\right)=\{xH:\underset{yH}{\underset{yx}{lim\; sup}}\left|f\left(y\right)\right|=\mathrm{}\}.$$
It is easy to see that $`\underset{f}{d}\left(H\right)`$ is a closed subset of $`H`$ such that $`\underset{f}{d}\left(H\right)H^1(f,\epsilon )`$ for any $`\epsilon >0.`$
###### Lemma 6.3.
Suppose that $`\alpha <\omega _1,`$ $`\delta >0`$ and $`s>2.`$ If $`x\left[KK^1(g,1)\right]K^\alpha (fh,\delta ),`$ then $`xK^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1).`$
###### Proof.
The proof is by induction on $`\alpha .`$ The result is clear if $`\alpha =0`$ or a limit ordinal. Assume that the lemma holds for some $`\alpha <\omega _1.`$ Suppose $`\delta >0`$ and $`s>2`$ are given. Let $`x\left[KK^1(g,1)\right]K^{\alpha +1}(fh,\delta ).`$ If $`x\underset{f}{d}\left(K^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1)\right),`$ then $`xK^{\alpha +1}(f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1)`$ and we are done. Otherwise, assume that $`x\underset{f}{d}\left(K^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1)\right).`$ Then there exist a neighborhood $`U_1`$ of $`x`$ in $`K`$ and $`M<\mathrm{}`$ such that $`\left|f\left(y\right)\right|M`$ for all $`yU_1K^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1).`$ Since $`h=g_0`$ on $`KK^1(g,1),`$ and $`g_0`$ is continuous on $`KK^1(g,1),`$ there exists a neighborhood $`U_2`$ of $`x`$ in $`K`$ such that $`\left|h\left(x_1\right)h\left(x_2\right)\right|{\displaystyle \frac{\delta }{2M}}`$ and $`2\left(\left|h\left(x_1\right)\right|+1\right)<s\left(\left|h\left(x\right)\right|+1\right)`$ for all $`x_1,x_2U_2.`$ Set $`U=\left(U_1U_2\right)K^1(g,1).`$ Then $`U`$ is a neighborhood of $`x.`$
Claim. $`K^\alpha (fh,\delta )UK^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1).`$
Note that if $`yU,`$ then $`yU_2.`$ Hence there exists $`t>2`$ such that $`t\left(\left|h\left(y\right)\right|+1\right)s\left(\left|h\left(x\right)\right|+1\right).`$ Also, $`yK^\alpha (fh,\delta )U`$ implies that $`y\left[KK^1(g,1)\right]K^\alpha (fh,\delta ).`$ Thus $`yK^\alpha (f,{\displaystyle \frac{\delta }{t\left(\left|h\left(y\right)\right|+1\right)}}1)`$ by the inductive hypothesis. Since
$$\frac{\delta }{t\left(\left|h\left(y\right)\right|+1\right)}\frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}1,$$
$`yK^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1),`$ as required.
Now if $`V`$ is a neighborhood of $`x`$ in $`K,`$ there exist $`x_1,x_2UV`$ $`K^\alpha (fh,\delta )`$ such that
$`\delta `$ $`\left|f\left(x_1\right)h\left(x_1\right)f\left(x_2\right)h\left(x_2\right)\right|`$
$`\left|f\left(x_1\right)f\left(x_2\right)\right|\left|h\left(x_1\right)\right|+\left|h\left(x_1\right)h\left(x_2\right)\right|\left|f\left(x_2\right)\right|`$
$`\left|f\left(x_1\right)f\left(x_2\right)\right|\left|h\left(x_1\right)\right|+{\displaystyle \frac{\delta }{2M}}M,`$
where, in the last inequality, $`\left|f\left(x_2\right)\right|M`$ since $`x_2UK^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1)`$ by the claim. Therefore,
$$\left|f\left(x_1\right)f\left(x_2\right)\right|\frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}1.$$
By the claim, $`x_1,x_2VK^\alpha (f,{\displaystyle \frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}}1).`$ Since $`V`$ is an arbitrary neighborhood of $`x,`$ this shows that
$$xK^{\alpha +1}(f,\frac{\delta }{s\left(\left|h\left(x\right)\right|+1\right)}1).$$
This completes the induction. ∎
It follows from Lemma 6.3 that
$$K^{\omega ^{\xi _1}}(fh,\delta )K^1(g,1).$$
Repeating the argument in Lemma 6.3 inductively yields
###### Lemma 6.4.
$`K^{\omega ^{\xi _1}\alpha }(fh,\delta )K^\alpha (g,1)`$ for all $`\alpha <\omega _1`$, and $`\delta >0.`$
In particular, $`K^{\omega ^{\xi _1}\omega ^{\xi _2}}(fh,\delta )=\mathrm{}`$ for all $`\delta >0,`$ i.e., $`fh𝔅_1^{\xi _1+\xi _2}\left(K\right).`$
###### Theorem 6.5.
If $`f𝔅_1^{\xi _1}\left(K\right)`$ and $`g𝔅_1^{\xi _2}\left(K\right),`$ then $`fg𝔅_1^\xi \left(K\right),`$ where $`\xi =\mathrm{max}\{\xi _1+\xi _2,\xi _2+\xi _1\}.`$
###### Proof.
From the above, we obtain a function $`h`$ in $`K`$ such that $`gh1,`$ $`\beta \left(h\right)\omega ^{\xi _2}`$ and $`fh𝔅_1^{\xi _1+\xi _2}\left(K\right).`$ Since $`g,h𝔅_1^{\xi _2}\left(K\right),`$ it follows from Proposition 4.4 that $`gh`$ $`𝔅_1^{\xi _2}\left(K\right).`$ As $`gh`$ is bounded, we see that $`gh_1^{\xi _2}\left(K\right).`$ By Lemma 6.2, $`\left(gh\right)f𝔅_1^{\xi _2+\xi _1}\left(K\right)𝔅_1^\xi \left(K\right).`$ Also, $`fh𝔅_1^{\xi _1+\xi _2}\left(K\right)𝔅_1^\xi \left(K\right).`$ Applying Proposition 4.4 again gives $`fg=f\left(gh\right)+fh𝔅_1^\xi \left(K\right).`$
Our final result shows that Theorem 6.5 is sharp. We omit the easy proof of the next lemma.
###### Lemma 6.6.
Suppose that $`h𝔅_1\left(K\right),`$ $`\alpha <\omega _1,`$ and $`\epsilon >0.`$ Let $`V=KK^\alpha (h,\epsilon ).`$ For any $`\eta <\omega _1,`$
$$K^\eta (h,\epsilon )K^\alpha (h,\epsilon )K^\eta (h\chi _V,\epsilon ).$$
###### Theorem 6.7.
Suppose that $`\xi _1`$, $`\xi _2`$ are countable ordinals, and let
$$\xi =\mathrm{max}\{\xi _1+\xi _2,\xi _2+\xi _1\}.$$
If $`K`$ is a compact metric space such that $`K^{(\omega ^\xi )}\mathrm{},`$ then
$$sup\{\beta \left(fg\right):f𝔅_1^{\xi _1}\left(K\right),g𝔅_1^{\xi _2}\left(K\right)\}=\omega ^\xi .$$
###### Proof.
We may of course assume that neither $`\xi _1`$ nor $`\xi _2`$ is $`0,`$ and that $`\xi =\xi _1+\xi _2.`$ The assumption on $`K`$ yields a $`\{0,1\}`$-valued function $`h`$ in $`𝔅_1\left(K\right)`$ such that $`K^{\omega ^\xi }(h,1)\mathrm{}.`$ Denote $`K^\alpha (h,1)`$ by $`K_\alpha ,`$ $`\alpha <\omega _1.`$ Choose a sequence of ordinals $`\left(\rho _k\right)_{k=0}^{\mathrm{}}`$ with $`\rho _0=0`$ that strictly increases to $`\omega ^{\xi _1}.`$ Let $`\lambda `$ be any ordinal that is less than $`\omega ^{\xi _2}.`$ Fix a function $`u:[0,\omega ^\lambda )`$ such that $`\{\alpha [0,\omega ^\lambda ):u\left(\alpha \right)k\}`$ is finite for all $`k.`$ Define real-valued functions $`f`$ and $`g`$ on $`K`$ as follows. If $`tK_{\omega ^{\xi _1}\lambda },`$ let $`f\left(t\right)=g\left(t\right)=0.`$ If $`t`$ $`K_{\omega ^{\xi _1}\alpha +\rho _{k1}}`$ $`K_{\omega ^{\xi _1}\alpha +\rho _k}`$ for some $`\alpha <\omega ^\lambda `$ and $`k,`$ let $`f\left(t\right)={\displaystyle \frac{h\left(t\right)}{ku\left(\alpha \right)}}`$ and $`g\left(t\right)=ku\left(\alpha \right).`$ Notice that $`fg=h\chi _V,`$ where $`V=KK^{\omega ^{\xi _1}\lambda }(h,1).`$ It follows from Lemma 6.6 that $`K^\eta (h,1)K^{\omega ^{\xi _1}\lambda }(h,1)K^\eta (fg,1)`$ for all $`\eta <\omega _1.`$ Since $`K^{\omega ^\xi }(h,1)\mathrm{},`$ and $`h𝔅_1\left(K\right),`$ $`K^\eta (h,1)K^{\omega ^{\xi _1}\lambda }(h,1)\mathrm{}`$ for all $`\eta <\omega ^{\xi _1}\lambda .`$ Thus $`K^\eta (fg,1)\mathrm{}`$ for all $`\eta <\omega ^{\xi _1}\lambda .`$ Hence $`\beta \left(fg\right)\omega ^{\xi _1}\lambda .`$
We now turn to the calculation of $`\beta \left(g\right)`$ and $`\beta \left(f\right).`$ First notice that the sets $`K_{\omega ^{\xi _1}\alpha +\rho _{k1}}K_{\omega ^{\xi _1}\alpha +\rho _k}`$, $`k`$, form a partition of $`K_{\omega ^{\xi _1}\alpha }`$ $`K_{\omega ^{\xi _1}\left(\alpha +1\right)}`$ into relatively open sets for any $`\alpha <\omega ^\lambda ,`$ and that $`g`$ is constant on each set $`K_{\omega ^{\xi _1}\alpha +\rho _{k1}}K_{\omega ^{\xi _1}\alpha +\rho _k}.`$ Hence the restriction of $`g`$ to $`K_{\omega ^{\xi _1}\alpha }`$ $`K_{\omega ^{\xi _1}\left(\alpha +1\right)}`$ is a continuous function for each $`\alpha <\omega ^\lambda .`$ It follows readily by induction that for any $`\epsilon >0,`$ $`K^\alpha (g,\epsilon )K_{\omega ^{\xi _1}\alpha }`$ for all $`\alpha \omega ^\lambda .`$ But $`g=0`$ on $`K_{\omega ^{\xi _1}\alpha }.`$ Thus $`K^{\omega ^\lambda +1}(g,\epsilon )=\mathrm{}.`$ Therefore $`\beta \left(g\right)\omega ^\lambda +1\omega ^{\xi _2}.`$
Finally, consider the function $`f`$. Let $`k_0`$ be given. The set
$$A=\{(\alpha ,k):k,\text{ }\alpha [0,\omega ^\lambda ),ku\left(\alpha \right)k_0\}$$
is finite. List the elements of $`A`$ in a finite sequence $`\left((\alpha _i,k_i)\right)_{i=1}^j`$ in lexicographical order. Then $`\left|f\left(t_1\right)f\left(t_2\right)\right|<{\displaystyle \frac{1}{k_0}}`$ for all $`t_1,`$$`t_2KK_{\omega ^{\xi _1}\alpha _1+\rho _{k_11}}.`$ Hence $`K^1(f,\frac{1}{k_0})K_{\omega ^{\xi _1}\alpha _1+\rho _{k_11}}.`$ Note that $`f=\frac{h}{k_1u\left(\alpha _1\right)}`$ on $`K_{\omega ^{\xi _1}\alpha _1+\rho _{k_11}}`$ $`K_{\omega ^{\xi _1}\alpha _1+\rho _{k_1}}.`$ Thus $`K^{1+\eta }(f,\frac{1}{k_0})K_{\omega ^{\xi _1}\alpha _1+\rho _{k_11}+\eta }`$ for all $`\eta `$ such that $`\omega ^{\xi _1}\alpha _1+\rho _{k_11}+\eta \omega ^{\xi _1}\alpha _1+\rho _{k_1}.`$ Let $`\eta _0`$ be such that $`\omega ^{\xi _1}\alpha _1+\rho _{k_11}+\eta _0=\omega ^{\xi _1}\alpha _1+\rho _{k_1}.`$ Then $`\eta _0\rho _{k_1}.`$ Therefore,
$$K^{1+\rho _{k_1}}(f,\frac{1}{k_0})K^{1+\eta _0}(f,\frac{1}{k_0})K_{\omega ^{\xi _1}\alpha _1+\rho _{k_1}}.$$
Repeating the argument, we see that
$$K^\rho (f,\frac{1}{k_j})K_{\omega ^{\xi _1}\alpha +\rho _{k_j}},$$
where $`\rho =1+\rho _{k_1}+1+\rho _{k_2}+\mathrm{}+1+\rho _{k_j}.`$ Since $`0f\left(t\right)<{\displaystyle \frac{1}{k_0}}`$ for all $`tK_{\omega ^{\xi _1}\alpha +\rho _{k_j}},`$
$$K^{\rho +1}(f,\frac{1}{k_j})=\mathrm{}.$$
As $`\left(\rho _k\right)`$ increases to $`\omega ^{\xi _1},`$ $`\rho +1<\omega ^{\xi _1}.`$ Hence $`K^{\omega ^{\xi _1}}(f,{\displaystyle \frac{1}{k_0}})=\mathrm{}`$ for any $`k_0.`$ It follows that $`\beta \left(f\right)\omega ^{\xi _1}.`$ Summarizing, we have functions $`f`$ and $`g`$ such that $`f𝔅_1^{\xi _1}\left(K\right)`$, $`g𝔅_1^{\xi _2}\left(K\right)`$ and $`\beta \left(fg\right)\omega ^{\xi _1}\lambda .`$ Since $`\lambda <\omega ^{\xi _2}`$ is arbitrary, the theorem is proved. ∎
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# 1 Introduction
## 1 Introduction
In the $`SU(2)_L\times U(1)`$ Standard Model (SM), the vertex factor for the weak neutral current interaction in the $`Z^0f\overline{f}`$ process is given by:
$$i\frac{g}{\mathrm{cos}\theta _W}\frac{\gamma ^\mu }{2}(v_fa_f\gamma ^5)$$
(1)
where $`v_f`$ and $`a_f`$ denote the vector and axial-vector couplings:
$$v_f=(c_L^f+c_R^f)=T_f^32Q_f\mathrm{sin}^2\theta _W,a_f=(c_L^fc_R^f)=T_f^3.$$
(2)
Here $`T_f^3`$ denotes the $`3^{rd}`$ component of the fermion weak isospin, $`Q_f`$ is the electric charge of the fermion, and $`\theta _W`$ represents the electroweak mixing angle.
From combinations of these couplings the $`Z^0`$ pole observables $`A_f`$ and $`R_f`$ can be formed. $`A_f`$ represents the extent of parity violation in the coupling of the $`Z^0`$ boson to the fermion of type $`f`$:
$$A_f=\frac{2v_fa_f}{v_f^2+a_f^2}=\frac{(c_L^f)^2(c_R^f)^2}{(c_L^f)^2+(c_R^f)^2},$$
(3)
and $`R_f`$ denotes the fraction of $`Z^0f\overline{f}`$ events in hadronic $`Z^0`$ decays:
$$R_f=\frac{\mathrm{\Gamma }(Z^0f\overline{f})}{\mathrm{\Gamma }(Z^0\text{hadrons})}(c_L^f)^2+(c_R^f)^2.$$
(4)
The precise experimental determination of $`A_f`$ and $`R_f`$ should ultimately verify whether or not the couplings for all generations and weak isospin states are described by the theory isospin assignments and a universal value of $`\mathrm{sin}^2\theta _W`$.
At Born level, the differential production cross section for $`e^+e^{}Z^0f\overline{f}`$ and longitudinally polarized electrons can be written as:
$$\sigma ^f(x)=\frac{d\sigma ^f}{dx}(1A_eP_e)(1+x^2)+2A_f(A_eP_e)x$$
(5)
where $`x`$ denotes the cosine of the polar angle of the outgoing fermion $`f`$ with respect to the incident electron beam direction, $`P_e`$ represents the longitudinal polarization of the electron beam, and the positron beam is assumed unpolarized. If one measures the polar angle distribution for a given final state $`f\overline{f}`$, one can derive the forward-backward production asymmetry:
$$A_{FB}^f(x)=\frac{\sigma ^f(x)\sigma ^f(x)}{\sigma ^f(x)+\sigma ^f(x)}=2A_f\frac{A_eP_e}{1A_eP_e}\frac{x}{1+x^2}$$
(6)
which depends on both the initial and final state coupling parameters as well as on the beam polarization. For zero polarization, one measures the product of couplings $`A_eA_f`$. If one measures the distributions in equal luminosity samples taken with negative $`(L)`$ and positive $`(R)`$ beam polarization of magnitude $`P_e`$, then one can derive the left-right-forward-backward asymmetry:
$$\stackrel{~}{A}_{FB}^f(x)=\frac{(\sigma _L^f(x)+\sigma _R^f(x))(\sigma _R^f(x)+\sigma _L^f(x))}{(\sigma _L^f(x)+\sigma _R^f(x))+(\sigma _R^f(x)+\sigma _L^f(x))}=2|P_e|A_f\frac{x}{1+x^2}$$
(7)
for which the dependence on the initial state coupling disappears, allowing a direct measurement of the final state coupling parameter $`A_f`$. Thus the presence of electron beam polarization permits unique $`A_f`$ measurements, not only independent of those inferred from the unpolarized forward-backward asymmetry which measures the combination $`A_eA_f`$, but also with a statistical advantage of $`(P_e/A_e)^225`$. The initial state coupling is determined most precisely via the left-right cross section asymmetry:
$$A_{LR}=\frac{1}{P_e}\frac{\sigma _L\sigma _R}{\sigma _L+\sigma _R}=A_e$$
(8)
which yields a very precise measurement of the electroweak mixing angle due to $`\delta A_e8\delta \mathrm{sin}^2\theta _W`$. The measurement and comparison of $`A_f`$ for the different charged lepton species also provides a direct test of lepton universality.
In addition, precise measurements of $`A_f`$ and $`R_f`$ can probe the effect of radiative corrections to the $`Z^0`$ propagator or the $`Z^0f\overline{f}`$ vertex. The radiative corrections depend on the masses of top and Higgs, and precise electroweak measurements can constrain these quantities. The coupling of the $`Z^0`$ boson to the $`b`$ quark is particularly interesting. Physics beyond the SM may couple more strongly to 3<sup>rd</sup> generation fermions, producing larger deviations in $`b`$ quark couplings than in other quark couplings. Since $`(c_L^b)^230(c_R^b)^2`$, $`R_b`$ has large sensitivity to possible deviations from the predicted left-handed coupling of the $`Z^0`$ boson to the $`b`$ quark, complementary to $`A_b`$ which has greater sensitivity to the right-handed coupling.
## 2 Unique Features of SLD/SLC
The performance of the SLC in the 1997-8 SLD data run has been excellent, with peak luminosities of $`3\times 10^{30}`$ cm<sup>-2</sup> $`s^1`$ (i.e. 20,000 $`Z^0`$ decays/week). Thus approximately 550,000 $`Z^0`$ decays have been collected during the 1993-8 data runs.
A general description of the SLD detector can be found in Ref. . Here we list several of the unique features which allow the SLD experiment to perform many competitive electroweak and heavy flavor measurements:
* A highly longitudinally polarized (average $``$73%) electron beam.
* A small and stable beam spot (1.5 $`\mu `$m $`\times `$ 0.8 $`\mu `$m $`\times `$ 700 $`\mu `$m) and a high precision 3D CCD-based pixel vertex detector allow the interaction point to be determined to 6 $`\mu `$m $`\times `$$`\mu `$m $`\times `$ 25 $`\mu `$m with an impact parameter resolution of 11 $`\mu `$m $`\times `$ 23 $`\mu `$m ($`r\varphi \times rz`$) for high momentum tracks.
* Good particle identification provided by the Cherenkov Ring Imaging Detector (CRID) .
## 3 Lepton Coupling Measurements
### A Left-right Cross Section Asymmetry ($`A_{LR}`$)
$`A_{LR}`$ provides a direct measurement of the initial state electron coupling, independent of the final state coupling. No efficiency or acceptance corrections are needed. The final state identification is relatively unsophisticated, and since practically all of the data can be used, $`A_{LR}`$ can be measured with high statistical precision. In fact, due to the high precision on the polarization measurement, the result is still statistically limited ($``$1.3% statistical error compared to $``$0.65% systematic error). The precision on $`A_{LR}`$ demands extensive cross checks to confirm the measurement, and three of the significant more recent checks have been the secondary, independent confirmations of the electron polarization measurement, the verification of the center-of-mass collision energy, and the measurement confirming the absence of positron polarization. The $`A_{LR}`$ measurement uses all hadronic events and is effectively a counting experiment; $`A_{LR}=\frac{1}{P_e}\frac{N_LN_R}{N_L+N_R}`$ where $`N_L`$($`N_R`$) denotes the number of $`Z^0`$ decays recorded with left (right) electron beam polarization. After applying additional corrections due to $`\gamma `$ exchange and $`\gamma Z^0`$ interference, we measure $`A_{LR}^0=0.15108\pm 0.00218`$ and $`\mathrm{sin}^2\theta _W^{eff}=0.23101\pm 0.00028`$ with the 1992-8 data sample. The $`A_{LR}`$ measurement provides the most precise single measurement of $`\mathrm{sin}^2\theta _W^{eff}`$ presently available.
### B $`A_{lepton}`$ from $`\stackrel{~}{A}_{FB}^f`$
The lepton couplings $`A_e`$, $`A_\mu `$ and $`A_\tau `$ are measured at SLD from leptonic decays of $`Z^0`$ bosons making use of the corresponding left-right-forward-backward asymmetry, $`\stackrel{~}{A}_{FB}^f`$, for each lepton type. $`A_e`$ and $`A_l`$ (with $`l=\mu ,\tau `$) are extracted simultaneously using a maximum likelihood fit. We obtain $`A_e=0.1558\pm 0.0064`$, $`A_\mu =0.137\pm 0.016`$, and $`A_\tau =0.142\pm 0.016`$. These results are consistent with lepton universality and can be combined to yield $`A_{e\mu \tau }=0.1523\pm 0.0057`$ corresponding to $`\mathrm{sin}^2\theta _W^{eff}=0.23085\pm 0.00073`$.
## 4 $`R_b`$ and $`R_c`$ Measurements
The $`R_c`$ and $`R_b`$ measurements heavily exploit the excellent vertexing capabilities of the SLD via a robust and efficient topological vertex algorithm . Following a standard hadronic event selection, each event is divided into two hemispheres where secondary (and tertiary) vertices are found. After calculating the $`p_t`$ corrected vertex invariant mass, $`M_{vtx}`$, hemispheres are tagged as containing a $`b`$ quark if $`M_{vtx}>2`$ GeV/c<sup>2</sup> and the secondary vertex is at least 5$`\sigma `$ away from the primary vertex. The hemisphere purity for $`b`$ events is $``$98% with $``$50% efficiency. Events are selected and tagged as $`b\overline{b}`$ if at least one of their hemispheres satisfies these conditions. Similarly, an event is tagged as $`c\overline{c}`$ if there is at least one track with 3D impact parameter more than 3$`\sigma `$ from the primary vertex, $`0.55`$ GeV/c$`{}_{}{}^{2}<M_{vtx}<2`$ GeV/c<sup>2</sup>, $`P_{vtx}>5`$ GeV/c and $`P_{vtx}>15M_{vtx}c10`$ GeV/c. The hemisphere purity for $`c`$ events is $``$70% with $``$16% efficiency. Figure 1 shows the $`M_{vtx}`$ and $`P_{vtx}`$ vs. $`M_{vtx}`$ distributions for $`uds`$, $`c`$ and $`b`$ hemispheres. The charm and bottom tagging efficiencies are self-calibrated directly from the experimental data using the single, double and mixed tag rates. The Monte Carlo is used as input for the $`c`$ and $`uds`$ efficiencies in the $`b`$ tag region. Therefore it is important to have high purity in the $`b`$ tagged event sample in order to reduce the systematic uncertainties due to the modeling of charm production and decay in the simulation. Hemisphere correlations are also derived from the simulation. We measure $`R_b=0.2159\pm 0.0014(stat.)\pm 0.0014(syst.)`$ and $`R_c=0.1685\pm 0.0047(stat.)\pm 0.0043(syst.)`$. Approximately 150,000 hadronic $`Z^0`$ decays from the last part of the 1998 run have not yet been included in the $`R_b`$ result.
## 5 $`A_q`$ Measurements
SLD provides several methods to measure the quark couplings $`A_q`$, and their statistical and systematic correlations are taken into account in the combinations of their results. All these measurements construct the left-right-forward-backward asymmetries defined previously and use, in different ways, the secondary vertex information. The four different techniques employed to measure $`A_b`$ make use of the jet charge, the kaon charge, the vertex charge, and the lepton charge to tag the $`b`$ hemisphere. To determine the $`c`$ hemisphere, the four $`A_c`$ measurements employ the kaon or vertex charge, the lepton charge, the soft pion charge or $`D^{()}`$ meson decays reconstructed exclusively. $`A_s`$ is measured at SLD using identified kaons.
### A $`A_b`$ with Kaon Tag
The $`bcs`$ decay chain is exploited in this measurement to tag the sign of the initial $`b`$ quark. This measurement heavily relies on the good $`K^\pm `$ identification provided by the CRID and the excellent separation between $`K^\pm `$ coming from the secondary vertex and those from the IP. The analyzing power for $`b`$ events is calibrated from the data. We obtain $`A_b=0.960\pm 0.040(stat.)\pm 0.069(syst.)`$.
### B $`A_b`$ with Vertex Charge
In this measurement , the sum of the charges of tracks attached to the reconstructed vertex is used to tag the initial $`b`$ quark sign. The analyzing power for $`b`$ events is calibrated from the data. We measure $`A_b=0.897\pm 0.027(stat.)\pm 0.034(syst.)`$ with the 1996-8 data sample.
### C $`A_b`$ with Momentum-Weighted Jet Charge
The momentum weighted jet charge is defined as:
$$Q_{diff}=Q_bQ_{\overline{b}}=\underset{tracks}{}q_i\text{sgn}(\stackrel{}{p}_i\widehat{T})|(\stackrel{}{p}_i\widehat{T})|^\kappa $$
(9)
where $`\stackrel{}{p}_i`$ and $`q_i`$ denote the $`i^{th}`$ track momentum and charge, respectively, and $`\widehat{T}`$ represents the direction of the thrust axis. The coefficient $`\kappa `$ was chosen to be 0.5 in order to maximize the analyzing power of the tag. $`Q_{diff}`$ is the difference between the momentum-weighted charges in the two hemispheres. The analyzing power for $`b`$ events is calibrated from the data. The hemisphere correlation is taken from the simulation. Figure 2 shows the polar angle distributions of the signed thrust axis for left-handed and right-handed electron beams. We measure $`A_b=0.882\pm 0.020(stat.)\pm 0.029(syst.)`$.
### D $`A_b`$ and $`A_c`$ with Lepton Tag
$`A_b`$ and $`A_c`$ can be measured by tagging bottom and charm hadrons using their semileptonic decays . The lepton total and transverse momenta with respect to the nearest jet are employed to calculate the probabilities that the lepton comes from each of the possible physics processes $`Z^0b\overline{b},bl`$; $`Z^0b\overline{b},\overline{b}\overline{c}l`$; $`Z^0b\overline{b},b\overline{c}l`$; $`Z^0c\overline{c},\overline{c}l`$; and background (leptons from light hadron decays, photon conversions, misidentified hadrons). The lepton charge provides quark-antiquark discrimination while the direction of the nearest jet to the lepton approximates the direction of the underlying quark.
Electrons are identified with both calorimeter and CRID information, and this information is incorporated in a neural network trained on Monte Carlo electrons. Muon identification uses information from tracking, the Warm Iron Calorimeter, and the CRID. Both electron and muon identification algorithms have been tested on control samples from data. If an electron is identified in an event, a secondary vertex is required in either hemisphere to reject $`uds`$ events. If a muon is identified, the vertex mass and the $`L/D`$ variable, illustrated in Figure 4, are included in the probability function together with the total and transverse lepton momenta. Figure 4 shows how $`L/D`$ distinguishes between muons from direct and cascade $`b`$ decays. The topological vertexing algorithm often finds only one vertex for both the $`B`$ and $`D`$ decay, thus $`L/D<1`$ ($`>1`$) indicates a muon from direct (cascade) $`b`$ decay. $`A_b`$ and $`A_c`$ are determined simultaneously from a maximum likelihood analysis. We measure $`A_b=0.924\pm 0.032(stat.)\pm 0.026(syst.)`$ and $`A_c=0.567\pm 0.051(stat.)\pm 0.064(syst.)`$.
### E $`A_c`$ with Exclusive $`D`$ mesons
This analysis exlusively reconstructs six modes to tag the charm quark: $`D^+K^{}\pi ^+\pi ^+`$, $`D^0K^{}\pi ^+`$, and $`D^+D^0\pi _{soft}^+`$ with $`D^0`$ decaying into $`K^{}\pi ^+`$, $`K^{}\pi ^+\pi ^0`$, $`K^{}\pi ^+\pi ^+\pi ^{}`$, $`K^{}l^+\nu _l`$ ($`l=e`$ or $`\mu `$). Both $`b`$ and $`uds`$ backgrounds are rejected with vertex information. The reconstruction efficiency is $``$4%, however, the high analyzing power and the good determination of the underlying charm quark direction lead to low systematic errors. Figure 6 shows the distribution of the mass difference between $`D^+`$ and $`D^0`$. The background under the signal is estimated from the sidebands. We obtain $`A_c=0.690\pm 0.042(stat.)\pm 0.022(syst.)`$.
### F $`A_c`$ with Inclusive Soft Pion
In this analysis the charm quark is tagged by the presence of a slow pion from the $`D^+D^0\pi _{soft}^+`$ decay. The soft pion in this decay is produced along the $`D^+`$ jet direction ($`P_T^20`$). Figure 6 illustrates the $`P_T^2`$ distribution for soft pion tracks. A signal to background ratio of $`1:2`$ is achieved for $`P_T^2<0.01`$ (GeV/c)<sup>2</sup>. This method yields $`A_c=0.683\pm 0.052(stat.)\pm 0.050(syst.)`$.
### G $`A_c`$ with Vertex Charge and Identified Kaons
This analysis uses the charm tag described for the $`R_c`$ analysis. Furthermore we require at least one hemisphere to pass the charm selection while neither hemisphere passes the bottom selection. The sign of the quark charge is determined by the charge of an identified $`K^\pm `$ (or by vertex charge), present in $``$25% ($``$50%) of the selected events with more than 90% correct sign fraction. Figure 7 shows the polar angle distributions of the signed thrust axis for left-handed and right-handed electron beams. The analyzing power for $`c`$ events is calibrated from the data. We obtain $`A_c=0.603\pm 0.028(stat.)\pm 0.023(syst.)`$.
### H $`A_s`$ Measurement
At SLD, $`s\overline{s}`$ events can be identified with relatively high purity due to the good separation of tracks from secondary vertices and the CRID particle identification. In the $`A_s`$ analysis , the charge of an identified $`K^\pm `$ is used to tag the sign of the initial $`s`$ quark. Information from the vertex detector is used to suppress the background from heavy flavor events. $`K^\pm `$ with $`p>9`$ GeV/c and $`K_s^0`$ with $`p>5`$ GeV/c are selected with 92% and 91% purity, respectively. Each thrust hemisphere of a light flavor tagged event is required to contain at least one identified strange particle. Strange hemispheres are tagged using the highest momentum strange particle present in the hemisphere. A tag is required in both hemispheres and at least one tag must be signed; if both are signed, signs must be opposite. The combined $`s\overline{s}`$ purity of the $`K^+K^{}`$ and $`K^\pm K_s^0`$ tagging modes is 66%. The initial $`s`$ quark direction is approximated by the thrust axis in the event, signed to point in the direction of negative strangeness. Figure 8 shows the polar angle distributions for left-handed and right-handed electron beams. The background from $`ud`$ events as well as the analyzing power of the method for $`s`$ events are constrained from the data. We measure $`A_s=0.895\pm 0.066(stat.)\pm 0.063(syst.)`$.
## 6 Conclusions
We have presented the results of several electroweak measurements performed by the SLD Collaboration. These results are summarized in Table 1. The combined $`A_{LR}`$ and $`A_{lepton}`$ SLD measurements yield $`\mathrm{sin}^2\theta _W^{eff}=0.23099\pm 0.00026`$. The LEP measurements of lepton forward-backward asymmetries and tau polarization have been combined into a LEP lepton based $`\mathrm{sin}^2\theta _W^{eff}=0.23151\pm 0.00033`$. These results are consistent. Both the $`R_c`$ and $`R_b`$ SLD results are in good agreement with the LEP results and the Standard Model prediction. The SLD $`R_c`$ measurement is the most precise single determination of this variable.
The SLD $`A_s`$ result is in agreement with LEP results and the Standard model prediction. It is also consistent with previous $`A_b`$ measurements performed by SLD and LEP, and therefore supports the predicted universality of the $`Z^0`$ to down-type quark couplings. The combined SLD $`A_c`$ measurements give $`A_c=0.634\pm 0.027`$, in agreement with both the corresponding LEP result and the Standard Model prediction. The combined SLD $`A_b`$ results yield $`A_b=0.905\pm 0.026`$ which is consistent with the Standard Model prediction. The corresponding LEP result is $`A_b=0.881\pm 0.020`$. A combined SLD and LEP average for $`A_b`$ is about 2.8 standard deviations below the Standard Model prediction.
## 7 Acknowledgements
I thank all SLD collaborators for their support and efforts and the SLAC accelerator department for its outstanding performance. I thank the organizers of Les Rencontres de la Vallee d’Aoste for inviting me and for the excellent hospitality. This work was supported in part by DOE Contract DE-AC03-76SF00515 (SLAC).
## <sup>∗∗</sup>List of Authors
Kenji Abe,<sup>(21)</sup> Koya Abe,<sup>(33)</sup> T. Abe,<sup>(29)</sup> I.Adam,<sup>(29)</sup> T. Akagi,<sup>(29)</sup> N. J. Allen,<sup>(5)</sup> W.W. Ash,<sup>(29)</sup> D. Aston,<sup>(29)</sup> K.G. Baird,<sup>(17)</sup> C. Baltay,<sup>(40)</sup> H.R. Band,<sup>(39)</sup> M.B. Barakat,<sup>(16)</sup> O. Bardon,<sup>(19)</sup> T.L. Barklow,<sup>(29)</sup> G. L. Bashindzhagyan,<sup>(20)</sup> J.M. Bauer,<sup>(18)</sup> G. Bellodi,<sup>(23)</sup> R. Ben-David,<sup>(40)</sup> A.C. Benvenuti,<sup>(3)</sup> G.M. Bilei,<sup>(25)</sup> D. Bisello,<sup>(24)</sup> G. Blaylock,<sup>(17)</sup> J.R. Bogart,<sup>(29)</sup> G.R. Bower,<sup>(29)</sup> J. E. Brau,<sup>(22)</sup> M. Breidenbach,<sup>(29)</sup> W.M. Bugg,<sup>(32)</sup> D. Burke,<sup>(29)</sup> T.H. Burnett,<sup>(38)</sup> P.N. Burrows,<sup>(23)</sup> A. Calcaterra,<sup>(12)</sup> D. Calloway,<sup>(29)</sup> B. Camanzi,<sup>(11)</sup> M. Carpinelli,<sup>(26)</sup> R. Cassell,<sup>(29)</sup> R. Castaldi,<sup>(26)</sup> A. Castro,<sup>(24)</sup> M. Cavalli-Sforza,<sup>(35)</sup> A. Chou,<sup>(29)</sup> E. Church,<sup>(38)</sup> H.O. Cohn,<sup>(32)</sup> J.A. Coller,<sup>(6)</sup> M.R. Convery,<sup>(29)</sup> V. Cook,<sup>(38)</sup> R. Cotton,<sup>(5)</sup> R.F. Cowan,<sup>(19)</sup> D.G. Coyne,<sup>(35)</sup> G. Crawford,<sup>(29)</sup> C.J.S. Damerell,<sup>(27)</sup> M. N. Danielson,<sup>(8)</sup> M. Daoudi,<sup>(29)</sup> N. de Groot,<sup>(4)</sup> R. Dell’Orso,<sup>(25)</sup> P.J. Dervan,<sup>(5)</sup> R. de Sangro,<sup>(12)</sup> M. Dima,<sup>(10)</sup> A. D’Oliveira,<sup>(7)</sup> D.N. Dong,<sup>(19)</sup> M. Doser,<sup>(29)</sup> R. Dubois,<sup>(29)</sup> B.I. Eisenstein,<sup>(13)</sup> V. Eschenburg,<sup>(18)</sup> E. Etzion,<sup>(39)</sup> S. Fahey,<sup>(8)</sup> D. Falciai,<sup>(12)</sup> C. Fan,<sup>(8)</sup> J.P. Fernandez,<sup>(35)</sup> M.J. Fero,<sup>(19)</sup> K.Flood,<sup>(17)</sup> R. Frey,<sup>(22)</sup> J. Gifford,<sup>(36)</sup> T. Gillman,<sup>(27)</sup> G. Gladding,<sup>(13)</sup> S. Gonzalez,<sup>(19)</sup> E. R. Goodman,<sup>(8)</sup> E.L. Hart,<sup>(32)</sup> J.L. Harton,<sup>(10)</sup> A. Hasan,<sup>(5)</sup> K. Hasuko,<sup>(33)</sup> S. J. Hedges,<sup>(6)</sup> S.S. Hertzbach,<sup>(17)</sup> M.D. Hildreth,<sup>(29)</sup> J. Huber,<sup>(22)</sup> M.E. Huffer,<sup>(29)</sup> E.W. Hughes,<sup>(29)</sup> X.Huynh,<sup>(29)</sup> H. Hwang,<sup>(22)</sup> M. Iwasaki,<sup>(22)</sup> D. J. Jackson,<sup>(27)</sup> P. Jacques,<sup>(28)</sup> J.A. Jaros,<sup>(29)</sup> Z.Y. Jiang,<sup>(29)</sup> A.S. Johnson,<sup>(29)</sup> J.R. Johnson,<sup>(39)</sup> R.A. Johnson,<sup>(7)</sup> T. Junk,<sup>(29)</sup> R. Kajikawa,<sup>(21)</sup> M. Kalelkar,<sup>(28)</sup> Y. Kamyshkov,<sup>(32)</sup> H.J. Kang,<sup>(28)</sup> I. Karliner,<sup>(13)</sup> H. Kawahara,<sup>(29)</sup> Y. D. Kim,<sup>(30)</sup> M.E. King,<sup>(29)</sup> R. King,<sup>(29)</sup> R.R. Kofler,<sup>(17)</sup> N.M. Krishna,<sup>(8)</sup> R.S. Kroeger,<sup>(18)</sup> M. Langston,<sup>(22)</sup> A. Lath,<sup>(19)</sup> D.W.G. Leith,<sup>(29)</sup> V. Lia,<sup>(19)</sup> C.Lin,<sup>(17)</sup> M.X. Liu,<sup>(40)</sup> X. Liu,<sup>(35)</sup> M. Loreti,<sup>(24)</sup> A. Lu,<sup>(34)</sup> H.L. Lynch,<sup>(29)</sup> J. Ma,<sup>(38)</sup> G. Mancinelli,<sup>(28)</sup> S. Manly,<sup>(40)</sup> G. Mantovani,<sup>(25)</sup> T.W. Markiewicz,<sup>(29)</sup> T. Maruyama,<sup>(29)</sup> H. Masuda,<sup>(29)</sup> E. Mazzucato,<sup>(11)</sup> A.K. McKemey,<sup>(5)</sup> B.T. Meadows,<sup>(7)</sup> G. Menegatti,<sup>(11)</sup> R. Messner,<sup>(29)</sup> P.M. Mockett,<sup>(38)</sup> K.C. Moffeit,<sup>(29)</sup> T.B. Moore,<sup>(40)</sup> M.Morii,<sup>(29)</sup> D. Muller,<sup>(29)</sup> V.Murzin,<sup>(20)</sup> T. Nagamine,<sup>(33)</sup> S. Narita,<sup>(33)</sup> U. Nauenberg,<sup>(8)</sup> H. Neal,<sup>(29)</sup> M. Nussbaum,<sup>(7)</sup> N.Oishi,<sup>(21)</sup> D. Onoprienko,<sup>(32)</sup> L.S. Osborne,<sup>(19)</sup> R.S. Panvini,<sup>(37)</sup> C. H. Park,<sup>(31)</sup> T.J. Pavel,<sup>(29)</sup> I. Peruzzi,<sup>(12)</sup> M. Piccolo,<sup>(12)</sup> L. Piemontese,<sup>(11)</sup> K.T. Pitts,<sup>(22)</sup> R.J. Plano,<sup>(28)</sup> R. Prepost,<sup>(39)</sup> C.Y. Prescott,<sup>(29)</sup> G.D. Punkar,<sup>(29)</sup> J. Quigley,<sup>(19)</sup> B.N. Ratcliff,<sup>(29)</sup> T.W. Reeves,<sup>(37)</sup> J. Reidy,<sup>(18)</sup> P.L. Reinertsen,<sup>(35)</sup> P.E. Rensing,<sup>(29)</sup> L.S. Rochester,<sup>(29)</sup> P.C. Rowson,<sup>(9)</sup> J.J. Russell,<sup>(29)</sup> O.H. Saxton,<sup>(29)</sup> T. Schalk,<sup>(35)</sup> R.H. Schindler,<sup>(29)</sup> B.A. Schumm,<sup>(35)</sup> J. Schwiening,<sup>(29)</sup> S. Sen,<sup>(40)</sup> V.V. Serbo,<sup>(29)</sup> M.H. Shaevitz,<sup>(9)</sup> J.T. Shank,<sup>(6)</sup> G. Shapiro,<sup>(15)</sup> D.J. Sherden,<sup>(29)</sup> K. D. Shmakov,<sup>(32)</sup> C. Simopoulos,<sup>(29)</sup> N.B. Sinev,<sup>(22)</sup> S.R. Smith,<sup>(29)</sup> M. B. Smy,<sup>(10)</sup> J.A. Snyder,<sup>(40)</sup> H. Staengle,<sup>(10)</sup> A. Stahl,<sup>(29)</sup> P. Stamer,<sup>(28)</sup> H. Steiner,<sup>(15)</sup> R. Steiner,<sup>(1)</sup> M.G. Strauss,<sup>(17)</sup> D. Su,<sup>(29)</sup> F. Suekane,<sup>(33)</sup> A. Sugiyama,<sup>(21)</sup> S. Suzuki,<sup>(21)</sup> M. Swartz,<sup>(14)</sup> A. Szumilo,<sup>(38)</sup> T. Takahashi,<sup>(29)</sup> F.E. Taylor,<sup>(19)</sup> J. Thom,<sup>(29)</sup> E. Torrence,<sup>(19)</sup> N. K. Toumbas,<sup>(29)</sup> T. Usher,<sup>(29)</sup> C. Vannini,<sup>(26)</sup> J. Va’vra,<sup>(29)</sup> E. Vella,<sup>(29)</sup> J.P. Venuti,<sup>(37)</sup> R. Verdier,<sup>(19)</sup> P.G. Verdini,<sup>(26)</sup> D. L. Wagner,<sup>(8)</sup> S.R. Wagner,<sup>(29)</sup> A.P. Waite,<sup>(29)</sup> S. Walston,<sup>(22)</sup> J.Wang,<sup>(29)</sup> S.J. Watts,<sup>(5)</sup> A.W. Weidemann,<sup>(32)</sup> E. R. Weiss,<sup>(38)</sup> J.S. Whitaker,<sup>(6)</sup> S.L. White,<sup>(32)</sup> F.J. Wickens,<sup>(27)</sup> B. Williams,<sup>(8)</sup> D.C. Williams,<sup>(19)</sup> S.H. Williams,<sup>(29)</sup> S. Willocq,<sup>(17)</sup> R.J. Wilson,<sup>(10)</sup> W.J. Wisniewski,<sup>(29)</sup> J. L. Wittlin,<sup>(17)</sup> M. Woods,<sup>(29)</sup> G.B. Word,<sup>(37)</sup> T.R. Wright,<sup>(39)</sup> J. Wyss,<sup>(24)</sup> R.K. Yamamoto,<sup>(19)</sup> J.M. Yamartino,<sup>(19)</sup> X. Yang,<sup>(22)</sup> J. Yashima,<sup>(33)</sup> S.J. Yellin,<sup>(34)</sup> C.C. Young,<sup>(29)</sup> H. Yuta,<sup>(2)</sup> G. Zapalac,<sup>(39)</sup> R.W. Zdarko,<sup>(29)</sup> J. Zhou.<sup>(22)</sup>
<sup>(1)</sup>Adelphi University, Garden City, New York 11530, <sup>(2)</sup>Aomori University, Aomori , 030 Japan, <sup>(3)</sup>INFN Sezione di Bologna, I-40126, Bologna Italy, <sup>(4)</sup>University of Bristol, Bristol, U.K., <sup>(5)</sup>Brunel University, Uxbridge, Middlesex, UB8 3PH United Kingdom, <sup>(6)</sup>Boston University, Boston, Massachusetts 02215, <sup>(7)</sup>University of Cincinnati, Cincinnati, Ohio 45221, <sup>(8)</sup>University of Colorado, Boulder, Colorado 80309, <sup>(9)</sup>Columbia University, New York, New York 10533, <sup>(10)</sup>Colorado State University, Ft. Collins, Colorado 80523, <sup>(11)</sup>INFN Sezione di Ferrara and Universita di Ferrara, I-44100 Ferrara, Italy, <sup>(12)</sup>INFN Lab. Nazionali di Frascati, I-00044 Frascati, Italy, <sup>(13)</sup>University of Illinois, Urbana, Illinois 61801, <sup>(14)</sup>Johns Hopkins University, Baltimore, MD 21218-2686, <sup>(15)</sup>Lawrence Berkeley Laboratory, University of California, Berkeley, California 94720, <sup>(16)</sup>Louisiana Technical University - Ruston,LA 71272, <sup>(17)</sup>University of Massachusetts, Amherst, Massachusetts 01003, <sup>(18)</sup>University of Mississippi, University, Mississippi 38677, <sup>(19)</sup>Massachusetts Institute of Technology, Cambridge, Massachusetts 02139, <sup>(20)</sup>Institute of Nuclear Physics, Moscow State University, 119899, Moscow Russia, <sup>(21)</sup>Nagoya University, Chikusa-ku, Nagoya 464 Japan, <sup>(22)</sup>University of Oregon, Eugene, Oregon 97403, <sup>(23)</sup>Oxford University, Oxford, OX1 3RH, United Kingdom, <sup>(24)</sup>INFN Sezione di Padova and Universita di Padova I-35100, Padova, Italy, <sup>(25)</sup>INFN Sezione di Perugia and Universita di Perugia, I-06100 Perugia, Italy, <sup>(26)</sup>INFN Sezione di Pisa and Universita di Pisa, I-56010 Pisa, Italy, <sup>(27)</sup>Rutherford Appleton Laboratory, Chilton, Didcot, Oxon OX11 0QX United Kingdom, <sup>(28)</sup>Rutgers University, Piscataway, New Jersey 08855, <sup>(29)</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309, <sup>(30)</sup>Sogang University, Seoul, Korea, <sup>(31)</sup>Soongsil University, Seoul, Korea 156-743, <sup>(32)</sup>University of Tennessee, Knoxville, Tennessee 37996, <sup>(33)</sup>Tohoku University, Sendai 980, Japan, <sup>(34)</sup>University of California at Santa Barbara, Santa Barbara, California 93106, <sup>(35)</sup>University of California at Santa Cruz, Santa Cruz, California 95064, <sup>(36)</sup>University of Victoria, Victoria, B.C., Canada, V8W 3P6, <sup>(37)</sup>Vanderbilt University, Nashville,Tennessee 37235, <sup>(38)</sup>University of Washington, Seattle, Washington 98105, <sup>(39)</sup>University of Wisconsin, Madison,Wisconsin 53706, <sup>(40)</sup>Yale University, New Haven, Connecticut 06511.
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# The 12C/13C-ratio in cool carbon starsPresented in this paper is observational data collected using the Swedish-ESO submillimetre telescope, La Silla, Chile, the 20 m telescope at Onsala Space Observatory, Chalmers Tekniska Högskola, Sweden, and the NRAO 12 m telescope located at Kitt Peak, USA.
## 1 Introduction
The <sup>12</sup>C/<sup>13</sup>C-ratio is an important measure of stellar evolution and nucleosynthesis. Current theoretical models of stellar evolution on the asymptotic giant branch (AGB) predict that stars with “normal” chemical composition transform into carbon stars, from spectral type M via the intermediate states MS, S and SC, through the dredge-up of freshly synthesized carbon in He-shell flashes (Busso et al. 1999). This scenario is also confirmed by observations (Smith & Lambert 1985, 1990; Dominy et al. 1986; Lambert et al. 1986). Since mainly the <sup>12</sup>C-isotope is synthesized, there should be a coeval evolution in the <sup>12</sup>C/<sup>13</sup>C-ratio from the value determined during the first red giant evolution. In addition, processes like hot bottom burning will affect this ratio (Boothroyd et al. 1993), as well as set an upper mass limit for carbon stars. Thus, an accurate estimate of the <sup>12</sup>C/<sup>13</sup>C-ratio should increase our understanding of the processes that lead to the formation of carbon stars (e.g., Forestini & Charbonnel 1997; Wallerstein & Knapp 1998). The <sup>12</sup>C/<sup>13</sup>C-ratio is also an important tracer of the past starformation rate and stellar mass function (Prantzos et al. 1996; Greaves & Holland 1997).
In a classical paper, Lambert et al. Lambert et al. (1986) determined the photospheric <sup>12</sup>C/<sup>13</sup>C-ratios of 30 optically bright carbon stars, by fitting stellar atmosphere models to near-IR data on the isotopomers of C<sub>2</sub>, CO, and CN. A decade later Ohnaka & Tsuji (1996, 1999), based on a different method and data on the CN red system around 8000 Å, presented <sup>12</sup>C/<sup>13</sup>C-ratios that, on the average, are about a factor of two lower than those of Lambert et al. Lambert et al. (1986) for the same stars. The activity in this field was further increased with the published results of Abia & Isern Abia and Isern (1997). They derived <sup>12</sup>C/<sup>13</sup>C-ratios of 44 carbon stars, using the CN red system, which fell in between the results obtained by Lambert et al. Lambert et al. (1986) and Ohnaka & Tsuji Ohnaka and Tsuji (1996). This is somewhat surprising since they used the same model atmospheres as Lambert et al. Lambert et al. (1986), suggesting that the derived <sup>12</sup>C/<sup>13</sup>C-ratios depend on the spectral features used. It is difficult to identify the main reason for the different results obtained by Lambert et al. Lambert et al. (1986) and Ohnaka & Tsuji (1996, 1999). The former used high-resolution near-infrared data, while the latter used data obtained closer to optical wavelengths. Also the atmospheric models differ. de Laverny & Gustafsson de Laverny and Gustafsson (1998) have investigated the method used by Ohnaka & Tsuji Ohnaka and Tsuji (1996) and found that it is very sensitive to model parameters and blends. Their conclusion is that the larger <sup>12</sup>C/<sup>13</sup>C-ratios determined by Lambert et al. Lambert et al. (1986) are more reliable, since this analysis is rather insensitive to the adopted model parameters, and also the effect of blends is less severe, but the discussion continues (Ohnaka & Tsuji 1998; de Laverny & Gustafsson 1999). Recently, Ohnaka et al. Ohnaka et al. (2000) have revised the values for three of the stars, which were originally published in Ohnaka & Tsuji Ohnaka and Tsuji (1996). Using new model atmospheres, they obtain <sup>12</sup>C/<sup>13</sup>C-ratios that are larger by about 40%, i.e., closer to those estimated by Lambert et al. Lambert et al. (1986).
To shed light on this disturbing controversy, we have performed independent estimates of the <sup>12</sup>C/<sup>13</sup>C-ratio, using CO radio line emission from the circumstellar envelopes (CSEs), for a sample of carbon stars showing large discrepancies between the sets of photospheric estimates. Due to the weakness of the circumstellar <sup>13</sup>CO lines, and the difficulties in the interpretation of the circumstellar emission, only a few attempts have been made to determine the <sup>12</sup>CO/<sup>13</sup>CO ratio in the CSEs of carbon stars (e.g., Knapp & Chang 1985; Sopka et al. 1989; Greaves & Holland 1997). In this paper we present new observational results, as well as a detailed modelling of circumstellar <sup>12</sup>CO and <sup>13</sup>CO radio line emission.
## 2 The observations
We have selected a sample of 11 optically bright carbon stars for which the photospheric <sup>12</sup>C/<sup>13</sup>C-ratios have been determined by Lambert et al. Lambert et al. (1986) and Ohnaka & Tsuji Ohnaka and Tsuji (1996). This sample is presented in Table 1, where we have also included S Sct, which was observed by Bergman et al. Bergman et al. (1993) and Olofsson el al. Olofsson et al. (1996), as well as the three J-type stars (i.e., stars with <sup>12</sup>C/<sup>13</sup>C-ratios $``$3) Y~CVn, RY~Dra, and T~Lyr (see Table 4 for references). These 15 stars were also included in the major survey of circumstellar molecular line emission by Olofsson et al. (1993a,b).
The <sup>13</sup>CO observations were carried out in October 1998 using the Swedish-ESO submillimetre telescope (SEST) located on La Silla, Chile, and in March 1999 using the Onsala 20 m telescope (OSO), Sweden. At SEST, a dual channel, heterodyne SIS receiver was used to simultaneously observe at 110 GHz (the $`J`$$`=`$$`1`$$``$$`0`$ line) and 220 GHz (the $`J`$$`=`$$`2`$$``$$`1`$ line). The single sideband temperatures of the receiver are about 110 K and 150 K at 110 GHz and 220 GHz, respectively. All the available acousto-optical spectrometers were used as backends. The two wideband (1 GHz), low resolution spectrometers (LRSs) were used to cover both lines. The third, narrow band (86 MHz), high resolution spectrometer (HRS) was used at 220 GHz, since the $`J`$$`=`$$`2`$$``$$`1`$ line is usually stronger than the lowest rotational line in CSEs. The LRS had 1440 channels separated by 0.7 MHz whereas the HRS used 2000 channels separated by 42 kHz.
At OSO, an SIS receiver, with a single sideband temperature of about 100 K at 110 GHz, was used for the observations. As backends, two filterbanks with bandwidths of 512 MHz (MUL A) and 64 MHz (MUL B) were used. The MUL A used 512 channels separated by 1 MHz and the MUL B filterbank used 256 channels with a separation of 250 kHz.
The observations were made in a dual beamswitch mode, where the source is alternately placed in the signal and the reference beam, using a beam throw of about 12$`\mathrm{}`$ at SEST and about 11$`\mathrm{}`$ at OSO. This method produces two spectra that are subtracted from each other, which results in very flat baselines, i.e., most of the frequency dependant response of the system is removed. The intensity scales are given in main beam brightness temperature, $`T_{\mathrm{mb}}`$$`=`$$`T_\mathrm{A}^{}/\eta _{\mathrm{mb}}`$, where $`T_\mathrm{A}^{}`$ is the antenna temperature corrected for atmospheric attenuation using the chopper wheel method, and $`\eta _{\mathrm{mb}}`$ is the main beam efficiency. For the SEST $`\eta _{\mathrm{mb}}`$ is 0.7 at 110 GHz and 0.5 at 220 GHz. At OSO $`\eta _{\mathrm{mb}}`$$`=`$0.5 at 110 GHz. The beamsizes of the SEST are 45$`\mathrm{}`$ and 24$`\mathrm{}`$ at 110 GHz and 220 GHz, respectively. At OSO the beam is 34$`\mathrm{}`$ at 110 GHz. Regular pointing checks were made on strong SiO masers and the pointing was usually better than $`\pm 3\mathrm{}`$ for both telescopes. The uncertainty in the intensity scale is estimated to be about $`\pm `$20%. The observational results are summarized in Table 2 and the detections of <sup>13</sup>CO line emission are shown in Fig. 1. These observations include first detections of circumstellar <sup>13</sup>CO emission towards R~Lep, UU~Aur, UX~Dra, and TX~Psc. The line parameters, i.e., the main beam brightness temperature at the line centre ($`T_{\mathrm{mb}}`$), the line centre velocity ($`v_{}`$), and half the full line width ($`v_\mathrm{e}`$), are obtained by fitting the data with an artificial line profile (Olofsson et al. 1993a)
$$T(v)=T_{\mathrm{mb}}\left[1\left(\frac{vv_{}}{v_\mathrm{e}}\right)^2\right]^{\gamma /2},$$
(1)
where $`\gamma `$ is a parameter describing the shape of the line. The integrated intensity ($`I_{\mathrm{mb}}`$) is obtained by integrating the emission between $`v_{}`$$`\pm `$$`v_\mathrm{e}`$. The upper limits are obtained as $`T_{\mathrm{pp}}2v_\mathrm{e}`$, where $`T_{\mathrm{pp}}`$ is the peak-to-peak noise obtained from spectra where the velocity resolution has been degraded to twice the expansion velocity. The corresponding <sup>12</sup>CO lines were also observed, simultaneously with the <sup>13</sup>CO lines, for all the sample stars. Spectra of these lines have been published in Olofsson et al. (1993a).
To complement these new observations we have collected additional <sup>13</sup>CO data and, in particular, <sup>12</sup>CO data on the sample stars, as well as for an additional seven stars for comparison, from various sources and included them in the modelling (see Table 4).
## 3 Radiative transfer
In order to model the circumstellar line emission, and to determine accurate <sup>12</sup>CO/<sup>13</sup>CO abundance ratios, we have used a non-LTE radiative transfer code based on the Monte Carlo method \[see Schöier Schöier (2000) for details\]. Assuming a spherically symmetric CSE expanding at a constant velocity, the code calculates the molecular excitation, i.e., the level populations, required to solve the radiative transfer exactly.
The excitation of the CO molecules were calculated taking into account 30 rotational levels in each of the ground and first vibrational states. The transition probabilities and energy levels are taken from Chandra et al. Chandra et al. (1996), and the rotational collisional rate coefficients (CO-H<sub>2</sub>) are based on the results in Flower & Launay Flower and Launay (1985) (they are extrapolated for $`J`$$`>`$11 and for temperatures higher than 250 K). Collisional transitions between vibrational levels are not important due to low densities and short radiative lifetimes.
The basic physical parameters of the CSE, e.g., the mass loss rate, expansion velocity, and temperature structure, are determined from the analysis of the observed <sup>12</sup>CO radio line emission. CO is well suited for this purpose since it is difficult to photodissociate and easy to excite through collisions, and thus is a very good tracer of the molecular gas and its temperature. The kinetic temperature of the gas is derived in a self-consistent manner, i.e., the calculations include the most important heating and cooling mechanisms for the gas, e.g., heating due to dust-gas collisions and molecular line cooling from CO. Once the characteristics of the CSE have been determined, the <sup>13</sup>CO excitation analysis is performed. The <sup>13</sup>CO abundance is varied until a satisfactory fit to the observations is obtained. In this way the circumstellar <sup>12</sup>CO/<sup>13</sup>CO-ratio is estimated.
The spatial extent of the molecular envelope is an important input parameter, and the derived mass loss rate and, to a lesser extent, the <sup>12</sup>CO/<sup>13</sup>CO-ratio will depend on this. The size of the circumstellar CO envelope, assumed to be the same for <sup>12</sup>CO and <sup>13</sup>CO, was estimated based on the modelling presented in Mamon et al. Mamon et al. (1988). It includes photodissociation, self-shielding and H<sub>2</sub>-shielding, and chemical exchange reactions. Here we assume the radial fractional abundance distribution to follow
$$f(r)=f_0\mathrm{exp}\left[\mathrm{ln}2\left(\frac{r}{r_\mathrm{p}}\right)^\alpha \right],$$
(2)
where $`f_0`$ is the initial (photospheric) abundance with respect to H<sub>2</sub>, $`r_\mathrm{p}`$ is the photodissociation radius (where the abundance has dropped to $`f_0/2`$), and $`\alpha `$ is a parameter describing the rate at which the abundance declines. Both $`r_\mathrm{p}`$ and $`\alpha `$ depend on the mass loss rate, the expansion velocity, and $`f_0`$. When modelling the <sup>12</sup>CO emission we assume $`f_0`$$`=`$1$`\times `$10<sup>-3</sup>. This is an average of the $`f_0`$:s estimated by Olofsson et al. Olofsson et al. (1993b) for a sample of optically bright carbon stars.
In our models we include both a central source of radiation and the cosmic background radiation at 2.7 K. The central radiation emanates from the star itself, which may be approximated by a blackbody. For heavily dust-enshrouded objects, like CW~Leo (a.k.a. IRC+10216), most of the stellar light is re-emitted by dust at longer wavelenghts. This emission source is also approximated by a blackbody. For low mass loss rate objects, the stellar blackbody temperature, $`T_{}`$, was estimated from a fit to the SED of the object. For stars of intermediate to high mass loss rates, two blackbodies were used, one representing the stellar contribution and one representing the dust. A fit to the SED gives the two blackbody temperatures, $`T_{}`$ and $`T_\mathrm{d}`$, and the relative luminosities of the two blackbodies, $`L_\mathrm{d}`$/$`L_{}`$. The method is described in Kerschbaum Kerschbaum (1999). The temperatures and luminosities used in the modelling are presented in Table 3. The inner boundary of the CSE was set to be outside the radius of the central blackbody(s), but never lower than 1$`\times `$10<sup>14</sup> cm.
The distances, presented in Table 3, were estimated using one of the following methods: the observed Hipparcos parallax, a period-luminosity relation (Groenewegen & Whitelock 1996), or an assumed bolometric luminosity. In the two former cases the luminosities were estimated using apparent bolometric magnitudes and the distances.
In the case of V~Hya and TX~Psc, stars with possible bipolar outflows (Heske et al. 1989; Kahane et al. 1996), no radiative transfer analysis was attempted due to the complexity of these outflows.
## 4 Model results
### 4.1 The <sup>12</sup>CO modelling
As explained above, the <sup>12</sup>CO line emission is used to derive the basic parameters of the CSE. The observational constraints for each source are presented in Table 4. The main part of the observational data used here are the $`J`$$`=`$$`1`$$``$$`0`$ and $`J`$$`=`$$`2`$$``$$`1`$ spectra obtained by Olofsson et al. Olofsson et al. (1993a) using the SEST, OSO, and the IRAM 30 m telescope at Pico Veleta, Spain (some of these observations have been remade and the intensities stated in Tab 4 may therefore be somewhat different from those originally given in the reference). Observations of the two lowest rotational transitions have also been performed using the NRAO 12 m telescope at Kitt Peak, USA, and of the $`J`$$`=`$$`3`$$``$$`2`$ line using the SEST (Schöier & Olofsson 2000, submitted). In addition, we have obtained publicly available data from the James Clerk Maxwell Telescope (JCMT) at Mauna Kea, Hawaii. The JCMT data are taken at face value. However, in the cases where there are more than one observation available, the derived line intensities are generally consistent within $`\pm `$20%. In addition, the good agreement with corresponding SEST observations lend further support for the reliability of the JCMT public data.
Thus, for all stars we use data from more than one <sup>12</sup>CO transition, in some cases four, to constrain the model. The intensities and overall line shapes of the circumstellar lines produced by the radiative transfer model generally agree well with those observed. This is illustrated here for two of our sample stars CW~Leo (Fig. 2) and U~Hya (Fig. 3). CW~Leo is a high mass loss rate Mira variable where the excitation of <sup>12</sup>CO is dominated by collisions. U~Hya, on the other hand, is a low mass loss rate object where radiation emitted by the central star plays a role in the excitation. The <sup>12</sup>CO model results presented in this paper constitute a sub-sample of the results of an analysis of a large survey of carbon stars presented in Schöier & Olofsson (2000, submitted). The reader is referred to this paper for a detailed description of the sensitivity of the model to the various input parameters. However, we point out here that tests made by Schöier & Olofsson show that the derived mass loss rate, for the majority of objects in this study, is mostly affected by the temperature structure and not by the assumed inner radius of the shell or the luminosity of the star, i.e., collisional excitation dominates over radiative excitation. In addition, Schöier & Olofsson have tested the molecular envelope size calculations by comparing radial brightness distributions obtained from the model with those observed. It was found that the envelope sizes estimated using the results from the model by Mamon et al. Mamon et al. (1988) are generally consistent with the observations. This is illustrated here in Fig. 2 for CW~Leo where the radial brightness distribution of the <sup>12</sup>CO($`J`$$`=`$$`1`$$``$$`0`$) line emission, observed at OSO, is compared to the distribution obtained from the radiative transfer model.
The derived envelope parameters, used as input for the <sup>13</sup>CO modelling, are presented in Table 3. We believe that, within the adopted circumstellar model, the estimated mass loss rates are generally accurate to within $`\pm `$50% (neglecting errors introduced by the uncertain CO abundance and the distance estimates). In the models presented here only <sup>12</sup>CO line cooling was included. For the J-stars, cooling from <sup>13</sup>CO will be important (scales with the isotope ratio for these low mass loss rate objects), but it will not affect the derived mass loss rate since the heating must be increased, i.e., the kinetic temperature structure will not change significantly, in order to maintain a good fit (Schöier & Olofsson 2000, submitted). See also the discussion in Ryde et al. Ryde et al. (1999) for the high mass loss rate object IRAS~15194-5115, where radio and far-IR <sup>12</sup>CO and <sup>13</sup>CO data are modelled.
When comparing our mass loss rate estimate for CW~Leo to those obtained from other detailed radiative transfer models (Kastner 1992; Crosas & Menten 1997; Groenewegen et al. 1998) we find a very good agreement, within 20%, when adjustments for differences in $`f_0`$ and distance have been made. Kastner Kastner (1992) also modelled the high mass loss rate object RW~LMi (a.k.a. CIT~6) obtaining (when corrected for the difference in distance) a mass loss rate in excellent agreement with our estimate. Sopka et al. Sopka et al. (1989) used a much simpler radiative transfer model to derive the mass loss rate for a number of AGB stars, of which four overlap with our survey. For RW~LMi they derived a mass loss rate in excellent agreement with that obtained from our model. For LP~And, V384~Per, and V~Cyg, however, there are discrepancies of about a factor of two to three. It should be noted though that Sopka et al. Sopka et al. (1989) based their mass loss rate estimates on observations of a single line, $`J`$$`=`$$`1`$$``$$`0`$, using a single telescope (OSO). For LP~And, the source with the largest discrepancy, we observe a <sup>12</sup>CO($`J`$$`=`$$`1`$$``$$`0`$) line that is almost a factor of two stronger using the same telescope.
### 4.2 The <sup>13</sup>CO modelling
Once the general characteristics of the CSE have been determined from the <sup>12</sup>CO analysis, the observed <sup>13</sup>CO line emission is modelled. Again, the observed intensities, as well as the line shapes, are generally well reproduced by the model. As for the <sup>12</sup>CO modelling, this is illustrated for two of our sample stars, CW~Leo (Fig. 2) and U~Hya (Fig. 3).
The abundance of <sup>13</sup>CO is usually (the J-stars provide an exception) much smaller than that of <sup>12</sup>CO, which leads to significantly different excitation conditions for the rarer isotopomer. For instance, the model <sup>13</sup>CO line intensities are more sensitive to the assumed properties of the central source of emission, since radiative pumping via the first excited vibrational state becomes important. Indeed, the <sup>13</sup>CO $`J`$$`=`$$`1`$$``$$`0`$ and $`J`$$`=`$$`2`$$``$$`1`$ transitions have inverted populations over parts of the envelope even for a high mass loss rate object as CW~Leo, Fig. 2. For the high mass loss rate objects, however, the part over which there is a weak maser acting is very small and the emission emanating from this region is not detected in our observations. For thinner envelopes the part of the envelope where the lowest rotational levels are inverted is larger due to the fact that the pumping emission can penetrate further out into the wind. In the case of U~Hya, changing the inner radius of the CSE or the luminosity of the star by $``$50% will change the estimated <sup>13</sup>CO abundance by $``$20%.
In our modelling of the <sup>13</sup>CO line emission we have assumed that the <sup>13</sup>CO envelope size is equal to that of <sup>12</sup>CO. This is based upon the model results by Mamon et al. Mamon et al. (1988). In this model both the effects of photodissociation and of chemical fractionation were included. Chemical fractionation of CO occurs through the exchange reaction (Watson et al. 1976) $`{}_{}{}^{13}\mathrm{C}_{}^{+}`$$`+`$$`{}_{}{}^{12}\mathrm{CO}^{12}\mathrm{C}^+`$$`+`$$`{}_{}{}^{13}\mathrm{CO}`$$`+`$$`\mathrm{\Delta }\mathrm{E}`$, where $`\mathrm{\Delta }\mathrm{E}/k`$$`=`$35 K. Below this temperature, the backward reaction is suppressed and production of <sup>13</sup>CO from <sup>12</sup>CO is favoured. This reaction can effectively produce <sup>13</sup>CO in the outer, cool parts of circumstellar envelopes. Mamon et al. Mamon et al. (1988) concluded that the difference between the <sup>12</sup>CO and <sup>13</sup>CO abundance distributions, when tested over a large mass loss rate interval, is always small, no more than 10 to 20%. Without the effect of chemical fractionation the <sup>13</sup>CO envelope would be significantly smaller. This is due to the fact that CO is photodissociated in lines and thereby exhibits considerable self-shielding. Thus, the shielding is higher the higher the optical depth, and hence it is less efficient for the less abundant <sup>13</sup>CO. Using OSO we have obtained a brightness distribution map of the <sup>13</sup>CO($`J`$$`=`$$`1`$$``$$`0`$) emission around CW~Leo. It is found that the model, with an assumed <sup>13</sup>CO envelope size equal to that of <sup>12</sup>CO, reproduces the observed radial brightness distribution within the observational errors, Fig. 2. A 20% smaller <sup>13</sup>CO envelope gives the same <sup>12</sup>CO/<sup>13</sup>CO-ratio but fails to reproduce the observed radial brightness distribution within the observational errors. This shows the importance of chemical fractionation, at least in the cool outer parts of dense CSEs. In the case of the thin molecular envelope around U~Hya a 20% smaller <sup>13</sup>CO envelope size would increase the derived <sup>13</sup>CO abundance by about 20%. This illustrates that in the case of a thin CSE the CO molecules are effectively excited to the photodissociation radius, i.e., photodissociation determines the size of the emitting region, whereas in a dense CSE excitation sets the size of the emitting region.
### 4.3 The circumstellar <sup>12</sup>CO/<sup>13</sup>CO-ratio
In Table 4 we list the observed integrated <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios. They are corrected for the differences in line strengths and beam-filling factors (assuming the sources to be unresolved the combined effect gives a $`\nu ^3`$ correction, i.e., the observed ratio is lowered by a factor of 0.87). Although we note that some of our CSEs are resolved, we have nevertheless treated all our sources in the same manner. Here we have chosen to integrate the emission over the entire line in order to achieve better signal-to-noise ratios. For optically thin lines this ratio should provide a first order estimate of the abundance ratio. A narrow velocity interval centered on the systemic velocity, where optical depth effects are smallest, should give somewhat higher line intensity ratios than those presented here. Nine of our sample stars have been detected in the <sup>13</sup>CO($`J`$$`=`$$`1`$$``$$`0`$) and/or the <sup>13</sup>CO($`J`$$`=`$$`2`$$``$$`1`$) line. A simple comparison of the <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios and the <sup>12</sup>C/<sup>13</sup>C-ratios derived by Lambert et al. Lambert et al. (1986) shows a tight correlation of the form (in the cases where more than one value is available in Table 4 we have used an average),
$$\frac{I(^{12}\mathrm{CO})}{I(^{13}\mathrm{CO})}=(0.6\pm 0.2)\times \frac{{}_{}{}^{12}\mathrm{C}}{{}_{}{}^{13}\mathrm{C}}.$$
(3)
Thus, a straightforward use of line intensity ratios would lead to <sup>12</sup>C/<sup>13</sup>C-ratios that, on the average, agree with those of Ohnaka & Tsuji Ohnaka and Tsuji (1996). However, any optical depth effects would lower the observed line intensity ratio.
In Table 3 we present the <sup>12</sup>CO/<sup>13</sup>CO-ratios derived using our radiative transfer code. They span a large range, from 2.5 to 90 (steps of 5 in the <sup>12</sup>CO/<sup>13</sup>CO-ratio was used to find the best fit <sup>13</sup>CO model, except for the J-stars where a smaller step-size of 0.5 was used). For most of our observed stars the isotope ratio obtained from the detailed radiative transfer is higher than those estimated from simple line intensity ratios. The discrepancy is fairly small for the objects with thin CSEs, but it increases with the thickness of the CSE ($``$$`\dot{M}`$/$`v_\mathrm{e}`$) and reaches a factor of six for the high mass loss rate objects, Fig. 4. This reflects that even for low to intermediate mass loss rate objects there are optical depth effects, and a detailed modelling is needed to derive reliable isotope ratios. In Table 3 we present, for each star, the maximum tangential optical depth in the <sup>12</sup>CO($`J`$$`=`$$`2`$$``$$`1`$) transition obtained in the modelling. The radial variation of the tangential optical depth is shown for the high mass loss rate object CW~Leo in Fig. 2 and for the low mass loss rate object U~Hya in Fig. 3. The optical depth in the <sup>12</sup>CO($`J`$$`=`$$`1`$$``$$`0`$) line is significantly lower.
The estimated <sup>13</sup>CO abundance depends on the adopted <sup>12</sup>CO abundance (assumed to be equal for all stars). We find that to a first approximation it scales with $`f_0`$, since a lower (higher) $`f_0`$ leads to a higher (lower) mass loss rate to fit the <sup>12</sup>CO line intensities and hence a lower (higher) <sup>13</sup>CO abundance to fit the <sup>13</sup>CO line intensities. This means that to a first approximation the estimated <sup>12</sup>CO/<sup>13</sup>CO-ratios are only very weakly dependant on the adopted <sup>12</sup>CO abundance. We have varied some of the other parameters in the model (see above), and conclude that in doing so the derived <sup>12</sup>CO/<sup>13</sup>CO-ratio changes by about 20%. With an additional uncertainty of $`\pm `$15% in the relative calibration of the <sup>12</sup>CO and <sup>13</sup>CO data (the relative calibration is usually better than the absolute calibration for any given telescope) , we estimate that the <sup>12</sup>CO/<sup>13</sup>CO-ratio is uncertain by about $`\pm `$30%. In the cases were the <sup>12</sup>CO/<sup>13</sup>CO-ratio estimate relies on observations of just one line the uncertainty may be as high as $`\pm `$50% depending on the quality of the data.
We compare first our derived <sup>12</sup>CO/<sup>13</sup>CO-ratio for CW~Leo with those found by others. Crosas & Menten Crosas and Menten (1997) derived, using a radiative transfer model similar to ours, an isotope ratio of 50, i.e., the same as we do. Greaves & Holland Greaves and Holland (1997) found a lower ratio of 32, and also a much lower ratio (24) for LP~And, using a simple radiative transfer model. Kahane et al. Kahane et al. (1992) found the <sup>12</sup>C/<sup>13</sup>C-ratio to be 44 based on observations of optically thin lines. Kahane et al. Kahane et al. (1992) also determined the <sup>12</sup>C/<sup>13</sup>C-ratio for RW~LMi obtaining a value of 31, in good agreement with our <sup>12</sup>CO/<sup>13</sup>CO estimate for this object. Sopka et al. Sopka et al. (1989) estimated <sup>12</sup>CO/<sup>13</sup>CO-ratios for four of our stars that are all lower than our derived values. This is most probably an effect of the difference in the treatment of the radiative transfer. We also note that Dufour et al. (2000, in prep.), derive somewhat lower <sup>12</sup>CO/<sup>13</sup>CO-ratios for the three J-stars.
For our sample stars we compare the circumstellar <sup>12</sup>CO/<sup>13</sup>CO-ratios obtained from the modelling with the photospheric <sup>12</sup>C/<sup>13</sup>C-ratios estimated by Lambert et al. Lambert et al. (1986) in Fig 5. A good correlation of the form
$$\frac{{}_{}{}^{12}\mathrm{CO}}{{}_{}{}^{13}\mathrm{CO}}=(1.0\pm 0.2)\times \frac{{}_{}{}^{12}\mathrm{C}}{{}_{}{}^{13}\mathrm{C}}$$
(4)
is obtained. Lambert et al. Lambert et al. (1986) give an uncertainty in their estimated <sup>12</sup>C/<sup>13</sup>C-ratios of about $`\pm `$40%. Included in our sample are stars that have had a drastic change in their mass loss rate. For R~Scl and S~Sct, stars with known detached CSEs (probably) produced during a period of intense mass loss, we have performed the analysis using the envelope parameters determined by Olofsson et al. Olofsson et al. (1996). These results are included in Fig 5. We note here that our results suggest that the <sup>12</sup>C/<sup>13</sup>C-ratios in the detached shells (with ages of about 10<sup>3</sup> and 10<sup>4</sup> yr for R~Scl and S~Sct, respectively) are the same as the present ones in the photospheres.
## 5 The <sup>12</sup>C/<sup>13</sup>C-ratio
Taken at face values our derived <sup>12</sup>CO/<sup>13</sup>CO-ratios support the <sup>12</sup>C/<sup>13</sup>C-ratios derived by Lambert et al. Lambert et al. (1986), Fig 5. However, there are a number of uncertainties in the extrapolation from the circumstellar isotopomer result to the stellar isotope ratio. The properties of the CO molecule suggests that the isotopomer ratio in the gas leaving the star is equal to the isotope ratio. We have based our circumstellar model upon the photodissociation/chemical fractionation results of Mamon et al. Mamon et al. (1988), and this has so far proven to give good results in our test cases. Our radiative transfer calculations are detailed and also provide good fits to multi-transition data. We also believe that estimates of isotopomer ratios are far less dependant on the adopted circumstellar model than are individual abundances. The two molecules have essentially the same energy level diagrams, the same transition strengths, and the same collisional cross sections, and hence their relative abundances are much less dependant on the adopted model than their absolute abundances (note that for <sup>12</sup>CO we actually adopt an abundance). However, one should note that the circumstellar estimates apply to time scales of 10<sup>2</sup> to 10<sup>3</sup> years, but there is no reason to expect that these stars have changed their surface composition over such a short time scale. In conclusion, we believe that our derived <sup>12</sup>CO/<sup>13</sup>CO-ratios are reliable estimates of the stellar <sup>12</sup>C/<sup>13</sup>C-ratios. Thus, we support the results obtained by Lambert et al. Lambert et al. (1986).
Also for the J-stars, whose origin is uncertain, we estimate <sup>12</sup>C/<sup>13</sup>C-ratios that are more consistent with those of Lambert et al. Lambert et al. (1986) than those of Ohnaka & Tsuji (1999), i.e., ratios that are, at least in principle, possible to obtain within the CNO-cycle. IRAS~15194-5115 also has a low <sup>12</sup>C/<sup>13</sup>C-ratio, but it differs from the J-stars in the sense that it has a much higher mass loss rate (see also Ryde et al. 1999). This star could be a borderline case with a mass of about 3.5 M, where the CNO-cycle produces a low <sup>12</sup>C/<sup>13</sup>C-ratio, while the temperature is not high enough to effectively convert <sup>12</sup>C to <sup>14</sup>N (Ventura et al. 1999).
## 6 Conclusions
We have determined the <sup>12</sup>CO/<sup>13</sup>CO-ratio in the molecular envelopes around 20 carbon stars, using a detailed non-LTE radiative transfer code. The <sup>12</sup>CO/<sup>13</sup>CO-ratios found range from 2.5 to 90, and we believe that these ratios accurately measure the stellar <sup>12</sup>C/<sup>13</sup>C-ratios.
Due to optical depth effects, present mainly in the <sup>12</sup>CO line, we find it necessary to treat in detail the radiative transfer in order to obtain reliable isotopomer ratios. For instance, for the two stars in common, we derive <sup>12</sup>CO/<sup>13</sup>CO-ratios that are almost a factor of two higher than those of Greaves & Holland Greaves and Holland (1997) in their study of the evolution of the local interstellar <sup>12</sup>C/<sup>13</sup>C-ratio.
An important ingredient in the model is the molecular abundance distribution in the CSE. We have here used the results of Mamon et al. Mamon et al. (1988) for <sup>12</sup>CO and <sup>13</sup>CO, and we find that for both molecules these are consistent with the observations, although for <sup>13</sup>CO we have only been able to test the results for a high mass loss rate object.
Our estimated <sup>12</sup>C/<sup>13</sup>C-ratios agree well with those estimated in the photosphere by Lambert et al. Lambert et al. (1986). Thus, we do not support the claims by other authors that the isotope ratios derived by Lambert et al. Lambert et al. (1986) were too high, by about a factor of two.
###### Acknowledgements.
We are grateful to Dr. F. Kerschbaum for generously providing estimates of some of the input parameters to the CO modelling. We would also like to thank Dr. C. Kahane for valuable comments. Financial support from the Swedish Natural Science Research Council (NFR) is gratefully acknowledged.
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# Contents
## 1 Introduction
Supersymmetry (for reviews, see ) has the ability to stabilize the hierarchy between the electroweak and Planck scales. However, the minimal supersymmetric standard model (MSSM) still requires an explanation for the magnitude of the supersymmetric Higgs mass parameter $`\mu `$. Assuming that there are no fine-tuned cancellations in the MSSM Higgs potential, $`\mu `$ should be of roughly the same magnitude as the soft supersymmetry-breaking masses. This suggests that $`\mu `$ arises as a vacuum expectation value (VEV) which is fixed by a potential with dimensionful parameters that are in turn determined by supersymmetry breaking.
Supersymmetry also requires some additional structure in order to solve the strong CP-problem. The $`\mu `$ parameter breaks the Peccei-Quinn (PQ) symmetry that is otherwise naturally present in the MSSM at the renormalizable level, so it is an attractive proposition that the dynamics which leads to the $`\mu `$ term simultaneously provide for an invisible axion. Astrophysical constraints on the axion leave open a window from roughly
$`10^9\mathrm{GeV}\stackrel{<}{}f\stackrel{<}{}10^{12}\mathrm{GeV}`$ (1.1)
for the VEV of the PQ-breaking field.
In this paper I will consider the phenomenology of a class of models with an invisible axion of the DFSZ type . In these models, the $`\mu `$ term arises from non-renormalizable terms in the superpotential, for example:
$`W={\displaystyle \frac{\lambda _\mu }{M_P}}X_1X_2H_uH_d.`$ (1.2)
(Here $`M_P=2.4\times 10^{18}`$ GeV is the reduced Planck mass, and $`\lambda _\mu `$ is a dimensionless coupling which I assume is not much smaller than unity.) The sum of the PQ charges of $`X_1`$ and $`X_2`$ must be equal and opposite to that of the MSSM Higgs superfields. When the scalar components of the neutral chiral superfields $`X_1`$ and $`X_2`$ acquire VEVs of order $`f`$, the approximate PQ symmetry is spontaneously broken, giving rise to a $`\mu `$ term of the right order of magnitude and an invisible axion. It is natural to assume that this occurs with $`X_1`$ and $`X_2`$ along a nearly flat direction in the potential. Then the low-energy degrees of freedom will typically include a pair of neutral chiral supermultiplets which are mixtures of the original $`X_1`$ and $`X_2`$. One of these contains as its imaginary scalar component the invisible axion of the model, and its fermion superpartner, the “axino”. In some models, $`X_1`$ and $`X_2`$ are the same field, so that the axino is the only new light singlet fermion. If $`X_1`$ and $`X_2`$ are distinct, then there will be another light Majorana fermion “singlino”. The axino and singlino both have odd $`R`$-parity and can mix. They obtain masses which are not much larger than the weak scale (but might be as small as of order a keV depending on the details of the model ). The upper limit can be understood from the facts that the (mass)<sup>2</sup> splittings between members of the same supermultiplet are bounded above by roughly $`m_{3/2}^2`$, the squared gravitino mass, and the axion is nearly massless.
The axino and the singlino are very weakly coupled to MSSM states, and cannot be directly produced in collider experiments at any significant rate. However, if either (or both) of these particles is lighter than all of the MSSM superpartners, then it will be the lightest supersymmetric particle (LSP) and can appear in decays from ordinary supersymmetric events.If $`R`$-parity is conserved and the singlino or the axino is the LSP, it will be absolutely stable and could dominate the energy density of the universe too soon. This potential problem can be solved (as in many similar cases) by invoking a low reheat temperature, at the cost of requiring in addition a low-scale baryogenesis mechanism. Furthermore, the NLSP decays will safely occur long before nucleosynthesis in the standard cosmology. In this paper I will argue there will be an opportunity in the era after supersymmetry is discovered to search for and measure the very long lifetime of the next-to-lightest supersymmetric particle (NLSP) into final states that include the singlino or axino, despite its very weak coupling.
In order to discuss the phenomenology in a general way, I will denote the relevant axino or singlino LSP by $`\stackrel{~}{S}`$, and refer to it generically as a singlino. It is part of a superfield $`S`$. In the low-energy theory it participates in a superpotential of the form
$`W=\mu \left(1+{\displaystyle \frac{ϵ}{v}}S\right)H_uH_d+{\displaystyle \frac{1}{2}}m_{\stackrel{~}{S}}S^2`$ (1.3)
in which
$`\mu ={\displaystyle \frac{\lambda _\mu }{M_P}}X_1X_2,`$ (1.4)
and I have introduced a dimensionless coupling parameter
$`ϵv/f`$ (1.5)
with $`v=175`$ GeV, the electroweak scale. (There are also soft mass terms for the scalar components of $`S`$, which will not concern us.) For example, if $`X_1`$ and $`X_2`$ are the same field, then one can read off from eq. (1.2) that $`ϵ=2v/f`$. For numerical purposes, this paper will use as a benchmark the value $`ϵ10^8`$ corresponding to $`f`$ (few)$`\times 10^{10}`$ GeV. There is also a dimensionless, holomorphic soft term in the lagrangian
$`_{\mathrm{SUSY}\mathrm{breaking}}={\displaystyle \frac{h_b}{M_P}}X_1X_2H_uH_d+\mathrm{c}.\mathrm{c}.`$ (1.6)
where $`h_b`$ is of order $`m_W`$. This gives rise to (among other terms) the necessary holomorphic soft (mass)<sup>2</sup> term for the Higgs bosons in the MSSM:
$`_{\mathrm{SUSY}\mathrm{breaking}}=bH_uH_d+\mathrm{c}.\mathrm{c}.`$ (1.7)
Note that $`b=h_bX_1X_2/M_P`$ is of order $`m_W^2`$, as required for proper electroweak symmetry breaking.
The coupling $`ϵ`$ and mass parameter $`m_{\stackrel{~}{S}}`$ parameterize our ignorance of the high-energy theory. The smallness of $`ϵ`$ means that $`S`$ nearly decouples. However, the conservation of $`R`$-parity implies that if the singlino is the LSP, then decays of the NLSP to $`\stackrel{~}{S}`$ will not suffer any competition and can be observed if they happen within a collider detector. These decays occur and are potentially observable because the “singlino” $`\stackrel{~}{S}`$ mixes slightly with the gauginos and Higgsinos, as well as couples directly to higgsino-Higgs pairs. In that sense, these models are similar to the well-studied - “next-to-minimal supersymmetric standard model” (NMSSM) . The differences include: the extremely small magnitude of $`ϵ`$; the fact that the field $`S`$ does not obtain a VEV; the absence (or at least weak-scale phenomenological irrelevance) of an $`S^3`$ term in the superpotential; and the presence of a tree-level supersymmetric mass term for $`\stackrel{~}{S}`$. Nevertheless, it is useful to compare the situation under study here to a very weakly coupled limit of the NMSSM. Indeed, the possibility of macroscopic decays involving a singlino have already been noted in ref. , but considering larger couplings (effectively $`ϵ\stackrel{>}{}10^6`$) and for smaller values of $`m_{\stackrel{~}{S}}`$ appropriate for LEP.
If the NLSP is the lightest of the ordinary MSSM neutralinos $`\stackrel{~}{N}_1`$, then it can decay according to
$`\stackrel{~}{N}_1f\overline{f}\stackrel{~}{S}`$ (1.8)
through virtual sleptons and squarks and virtual or on-shell $`Z`$ bosons and Higgs bosons. The decay width is estimated very roughly by
$`\mathrm{\Gamma }{\displaystyle \frac{m_{\stackrel{~}{N}_1}}{16\pi }}|ϵ\mu /v|^2\times (\mathrm{suppression}\mathrm{factors}).`$ (1.9)
The suppression factors include the effects of electroweak couplings, mixing angles, kinematic suppressions, and (if the mediating boson is not on-shell) three-body phase space. Without these effects, the rough estimate (for $`m_{\stackrel{~}{N}_1}100`$ GeV and $`|ϵ\mu /v|10^8`$) would be of order 1 meter<sup>-1</sup> for $`\mathrm{\Gamma }`$.
After including the suppression effects in realistic models, we will find that when $`\stackrel{~}{N}_1`$ is allowed to decay through an on-shell CP-even Higgs boson $`h^0`$, the inverse decay width is of order meters or tens of meters. Of course, the decay $`\stackrel{~}{N}_1h^0\stackrel{~}{S}`$ may not be kinematically allowed. In that case, there may still be allowed decays $`\stackrel{~}{N}_1Z^0\stackrel{~}{S}`$. These are typically further suppressed by a mixing angle, because the singlino must mix with the MSSM higgsinos in order to couple to the $`Z`$ boson. Nevertheless, we will find that these mixing angles are typically large enough so that the inverse decay widths can be of order hundreds of meters. Finally, it may be that $`m_{\stackrel{~}{N}_1}m_{\stackrel{~}{S}}<m_Z`$. In that case, there can still be three-body decays $`\stackrel{~}{N}_1f\overline{f}\stackrel{~}{S}`$ through virtual sleptons, squarks, and the $`Z^0`$. (Decays through off-shell Higgs bosons can also occur, but are typically very small because the MSSM Higgs boson widths are tiny unless they are heavy.) If the $`Z`$ boson is far off-shell, then with the usual model prejudices that sleptons are much lighter than squarks one finds that the smallest inverse decay widths are for $`\mathrm{}^+\mathrm{}^{}\stackrel{~}{S}`$ final states, and can be of order tens of kilometers. (All of these results assume $`ϵ10^8`$, and the decay widths must be scaled with $`ϵ^2`$.)
The majority of decays $`\stackrel{~}{N}_1f\overline{f}\stackrel{~}{S}`$ will evidently occur well outside of a typical collider detector. However, with a significant number of supersymmetric events available, a small but finite fraction will occur inside the detector where the displaced secondary vertex can be distinguished. Since the decaying $`\stackrel{~}{N}_1`$ and the resulting $`\stackrel{~}{S}`$ are invisible, the experimental signature will involve an energetic lepton-antilepton pair or dijet pair with a significant opening angle appearing “from nothing” (with no corresponding charged particle track pointing back to the interaction point) at a common point. This determination could be accomplished within the inner tracking volume of a detector, but might also be possible and perhaps even easier to distinguish if the decay occurs farther from the beam pipe. Thus, decays occurring within a meter to several meters from the interaction point could be a striking, if rare, signal.
At the Large Hadron Collider, the number of supersymmetric events expected per year with a low luminosity option of 10 fb<sup>-1</sup>/year can be roughly of order a few thousand to a few million or more for 200 GeV $`>m_{\stackrel{~}{N}_1}>`$ 50 GeV. (This assumes $`m_{\stackrel{~}{g}}m_{\stackrel{~}{q}}7m_{\stackrel{~}{N}_1}`$; of course this is quite model-dependent.) Every event gives two possible NLSP decays. Therefore one can aspire to detect rare decays with widths as long as hundreds of kilometers by searching within the supersymmetric event sample . In the limit of small decay widths, the probability that a particular $`\stackrel{~}{N}_1f\overline{f}\stackrel{~}{S}`$ decay will occur within a distance $`L`$ of the interaction point is given by
$`P(L)L\mathrm{\Gamma }/\beta \gamma `$ (1.10)
where $`\mathrm{\Gamma }`$ is the invariant partial width for that decay channel. Supersymmetric events will be “tagged” by the other particles from the sparticle decays and the presence of large $`E\text{/}_T`$ . Since the slowly decaying $`\stackrel{~}{N}_1`$ will be massive and not ultrarelativistic, one can use timing information together with the pointing information from the tracking detectors and drift chambers and perhaps vetoes from the muon system to eradicate backgrounds from cosmic rays and other sources. There have also been proposals motivated by gauge-mediated supersymmetry breaking models and by neutralino decays to axino and photon to build special detector components and instrumented tunnels to aid in the search for very slow decays . At future $`e^+e^{}`$ linear colliders, supersymmetric event rates are smaller, but the two-body decays mentioned above might occur often enough to be detected. The signal is somewhat more problematic at future runs of the Fermilab Tevatron collider, since the total supersymmetric event rates are likely to be considerably smaller. In this paper, I will simply remain optimistic and choose to present results for decay partial widths down to as small as (1000 km)<sup>-1</sup>.
In these models, there are also singlet scalars $`S`$ within the same supermultiplet as $`\stackrel{~}{S}`$. These very weakly coupled scalars have masses of order $`m_W`$ (or, in the case of the invisible axion, essentially 0). However, decays like $`\stackrel{~}{N}_1\stackrel{~}{S}S`$ depend on couplings that are effectively doubly suppressed by $`ϵ`$. Other decays involving only singlet scalars always have competition from ordinary unsuppressed MSSM decays, and so are not relevant for colliders.
The rest of this paper is organized as follows. In section 2, I examine some specific models which realize the idea outlined above. Section 3 discusses the couplings and mixings of the singlino/axino, and the relevant decays. Some representative numerical results for decays to the singlino are presented in section 4. Section 5 contains some concluding remarks. An Appendix contains complete formulas for decay widths, including the effects of arbitrary phases.
## 2 Singlino masses and couplings in models with an intermediate-scale solution to the $`\mu `$ problem
Let us consider several models which realize the basic idea outlined in the introduction. The magnitude of $`fX_{1,2}`$ is (up to dimensionless couplings) the geometric mean of the Planck scale and the weak scale in order to agree with eq. (1.1). One way that this could happen is if the soft supersymmetry-breaking (mass)<sup>2</sup> of $`X_1`$ is driven negative at an intermediate scale. More generally, $`X_1`$ and $`X_2`$ correspond to a nearly-flat direction in the potential, so that dimensionless supersymmetry breaking terms involving $`X_1`$ and $`X_2`$ will always favor a non-trivial minimum at an intermediate scale. The key things we want to show are that these models generically contain one or more singlet fermions which have electroweak scale (or smaller) masses, and couplings that are of order $`v/f`$ (with possible enhancements) to the MSSM Higgs fields.
For example, suppose that the superpotential contains, in addition to eq. (1.2), a term
$`W={\displaystyle \frac{\lambda _X}{6M_P}}X_1X_2^3,`$ (2.1)
as in . This fixes the PQ charges of the superfields, ensuring the presence of an invisible axion provided that no other terms break the PQ symmetry. Here $`\lambda _X`$ is a dimensionless coupling which is assumed to be of order unity. The supersymmetry breaking Lagrangian must then include
$`_{\mathrm{SUSY}\mathrm{breaking}}=m_1^2|X_1|^2+m_2^2|X_2|^2{\displaystyle \frac{h_X}{6M_P}}(X_1X_2^3+\mathrm{c}.\mathrm{c}.)`$ (2.2)
where $`h_X`$ is a mass parameter of order the electroweak scale and has been taken to be real and positive without loss of generality. A nontrivial global minimum will exist provided e.g. that $`m_1^2`$ is negative at the scale of the VEV. However, it is important to note that this is not necessary. The presence of a holomorphic coupling $`h_X`$ always favors spontaneous symmetry breaking at an intermediate scale. So VEVs for $`X_1`$ and $`X_2`$ can arise from a negative squared mass, and/or an $`h_X`$ which is sufficiently large . In any case, $`X_1`$ and $`X_2`$ are of order $`(m_WM_P)^{1/2}10^{10}`$ GeV, which is naturally within the invisible axion window.
Let us parameterize the VEVs of $`X_1`$ and $`X_2`$ by an overall magnitude $`f`$ and an angle $`\varphi `$, so that $`X_1=f\mathrm{cos}\varphi `$ and $`X_2=f\mathrm{sin}\varphi `$. Now expand the fields around their VEVs to obtain low-energy supermultiplet degrees of freedom $`S_1,S_2`$:
$`\left(\begin{array}{c}X_1\\ X_2\end{array}\right)=\left(\begin{array}{c}\mathrm{cos}\varphi \\ \mathrm{sin}\varphi \end{array}\right)f+\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{c}S_1\\ S_2\end{array}\right).`$ (2.3)
By requiring that the superpotential masses for the fermions $`\stackrel{~}{S}_1`$ and $`\stackrel{~}{S}_2`$ derived from eq. (2.1) are diagonal, one can solve for the mixing angle $`\theta `$ in terms of the VEV angle $`\varphi `$, with the result $`\theta =\varphi /2`$. \[This typically does not diagonalize the axion and other light scalar masses, and depends particularly on the choice of eq. (2.1).\] The resulting masses and couplings for $`\stackrel{~}{S}_1`$ and $`\stackrel{~}{S}_2`$ are then found to be
$`m_{\stackrel{~}{S}_1}={\displaystyle \frac{\lambda _Xf^2}{2M_P}}\mathrm{sin}\varphi (\mathrm{cos}\varphi 1);`$ $`ϵ_{\stackrel{~}{S}_1}={\displaystyle \frac{v}{f}}\left({\displaystyle \frac{\mathrm{cos}\varphi /2}{\mathrm{cos}\varphi }}{\displaystyle \frac{\mathrm{sin}\varphi /2}{\mathrm{sin}\varphi }}\right);`$ (2.4)
$`m_{\stackrel{~}{S}_2}={\displaystyle \frac{\lambda _Xf^2}{2M_P}}\mathrm{sin}\varphi (\mathrm{cos}\varphi +1);`$ $`ϵ_{\stackrel{~}{S}_2}={\displaystyle \frac{v}{f}}\left({\displaystyle \frac{\mathrm{sin}\varphi /2}{\mathrm{cos}\varphi }}+{\displaystyle \frac{\mathrm{cos}\varphi /2}{\mathrm{sin}\varphi }}\right).`$ (2.5)
\[Compare eq. (1.3).\] The scale $`f`$ and the angle $`\varphi `$ could also be computed, in principle, in terms of the parameters in the soft supersymmetry-breaking Lagrangian eq. (2.2) and the superpotential. The same parameters also determine the soft scalar mass of the saxion and other scalars with electroweak scale masses. However, these will not play any direct role in this paper, so I will not do this explicitly, and I will treat $`f`$ and $`\varphi `$ as free parameters.
One interesting limit is that of small $`\varphi `$ (i.e., $`X_2`$ small compared to $`X_1`$), in which the mass eigenstate $`\stackrel{~}{S}_1`$ is the axino. Then one finds that, up to phases,
$`m_{\stackrel{~}{S}_1}={\displaystyle \frac{\lambda _Xf^2}{4M_P}}\mathrm{sin}^3\varphi ;ϵ_{\stackrel{~}{S}_1}={\displaystyle \frac{v}{2f}};`$ (2.6)
$`m_{\stackrel{~}{S}_2}={\displaystyle \frac{\lambda _Xf^2}{M_P}}\mathrm{sin}\varphi ;ϵ_{\stackrel{~}{S}_2}={\displaystyle \frac{v}{f\mathrm{sin}\varphi }}.`$ (2.7)
Note that both masses become small in this limit. The coupling $`ϵ_{\stackrel{~}{S}_1}`$ must grow like $`1/\mathrm{sin}\varphi `$ in this parameterization in order for $`\mu `$ to not become much less than $`v`$. (LEP has not discovered a higgsino.) Another interesting limit is $`\varphi =\pi /4`$ (VEVs of equal magnitude), resulting in
$`m_{\stackrel{~}{S}_1}=0.10{\displaystyle \frac{\lambda _Xf^2}{M_P}};ϵ_{\stackrel{~}{S}_1}=0.77v/f;`$ (2.8)
$`m_{\stackrel{~}{S}_2}=0.60{\displaystyle \frac{\lambda _Xf^2}{M_P}};ϵ_{\stackrel{~}{S}_2}=1.85v/f,`$ (2.9)
again, up to phases. Both $`\stackrel{~}{S}_1`$ and $`\stackrel{~}{S}_2`$ have masses that are roughly of order the electroweak scale. In general they each contain an admixture of the axino. It is interesting to note that $`|m_{\stackrel{~}{S}_1}|/(\lambda _Xf^2/M_P)<0.1`$ over the whole range $`\theta <\pi /4`$. So it is not unlikely that one or both of $`\stackrel{~}{S}_1`$ and $`\stackrel{~}{S}_2`$ is lighter than all MSSM sparticles.
Other models can be obtained by choosing superpotentials
$`W={\displaystyle \frac{\lambda _X}{6M_P}}X_1X_2^3+{\displaystyle \frac{\lambda _\mu }{2M_P}}X_2^2H_uH_d,`$ (2.10)
as in ref. , or
$`W={\displaystyle \frac{\lambda _X}{6M_P}}X_1X_2^3+{\displaystyle \frac{\lambda _\mu }{2M_P}}X_1^2H_uH_d.`$ (2.11)
In both of these cases, the diagonalized singlino masses are still given as in eqs. (2.4) and (2.5), since they only depend on the $`\lambda _X`$ term in the superpotential. However, the couplings are modified to, respectively: and
$`ϵ_{\stackrel{~}{S}_1}={\displaystyle \frac{2v\mathrm{sin}\varphi /2}{f\mathrm{sin}\varphi }};ϵ_{\stackrel{~}{S}_2}={\displaystyle \frac{2v\mathrm{cos}\varphi /2}{f\mathrm{sin}\varphi }}.`$ (2.12)
for eq. (2.10), and
$`ϵ_{\stackrel{~}{S}_1}={\displaystyle \frac{2v\mathrm{cos}\varphi /2}{f\mathrm{cos}\varphi }};ϵ_{\stackrel{~}{S}_2}={\displaystyle \frac{2v\mathrm{sin}\varphi /2}{f\mathrm{cos}\varphi }}`$ (2.13)
for eq. (2.11).
Another similar model is obtained by assuming a different form of the $`\lambda _X`$ term used to stabilize the potential at large field strengths:
$`W={\displaystyle \frac{\lambda _X}{4M_P}}X_1^2X_2^2+{\displaystyle \frac{\lambda _\mu }{2M_P}}X_1^2H_uH_d.`$ (2.14)
Again using the mixing parameterization eq. (2.3), one finds in this class of models that now $`\theta =(1/2)\mathrm{tan}^1(2\mathrm{tan}2\varphi )`$ in order that $`\stackrel{~}{S}_1`$ and $`\stackrel{~}{S}_2`$ are mass eigenstates. In terms of the VEV angle $`\varphi `$, we find:
$`m_{\stackrel{~}{S}_1}`$ $`=`$ $`{\displaystyle \frac{\lambda _Xf^2}{4M_P}}\left[1\mathrm{cos}2\varphi \sqrt{1+4\mathrm{tan}^22\varphi }\right];`$ (2.15)
$`ϵ_{\stackrel{~}{S}_1}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}v}{f\mathrm{cos}\varphi }}\left[1+{\displaystyle \frac{1}{\sqrt{1+4\mathrm{tan}^22\varphi }}}\right]^{1/2};`$ (2.16)
$`m_{\stackrel{~}{S}_2}`$ $`=`$ $`{\displaystyle \frac{\lambda _Xf^2}{4M_P}}\left[1+\mathrm{cos}2\varphi \sqrt{1+4\mathrm{tan}^22\varphi }\right];`$ (2.17)
$`ϵ_{\stackrel{~}{S}_2}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}v}{f\mathrm{cos}\varphi }}\left[1{\displaystyle \frac{1}{\sqrt{1+4\mathrm{tan}^22\varphi }}}\right]^{1/2},`$ (2.18)
up to phases. In the limit $`\varphi 0`$, one finds
$`m_{\stackrel{~}{S}_1}={\displaystyle \frac{3\lambda _Xf^2}{2M_P}}\mathrm{sin}^2\varphi ;ϵ_{\stackrel{~}{S}_1}={\displaystyle \frac{2v}{f}};`$ (2.19)
$`m_{\stackrel{~}{S}_2}={\displaystyle \frac{\lambda _Xf^2}{2M_P}};ϵ_{\stackrel{~}{S}_2}={\displaystyle \frac{4v}{f}}\mathrm{sin}\varphi ,`$ (2.20)
which features a parametric suppression in the mass of $`\stackrel{~}{S}_1`$. In the case of equal VEVs $`\varphi =\pi /4`$, one obtains instead
$`m_{\stackrel{~}{S}_1}={\displaystyle \frac{\lambda _Xf^2}{4M_P}};ϵ_{\stackrel{~}{S}_1}={\displaystyle \frac{2v}{f}};`$ (2.21)
$`m_{\stackrel{~}{S}_2}={\displaystyle \frac{3\lambda _Xf^2}{4M_P}};ϵ_{\stackrel{~}{S}_2}={\displaystyle \frac{2v}{f}}.`$ (2.22)
Finally, the limit $`\varphi \pi /2`$ (i.e., $`X_1`$ small compared to $`X_2`$) yields a very light $`\stackrel{~}{S}_2`$ and a relative enhancement of the $`\stackrel{~}{S}_1`$ coupling:
$`m_{\stackrel{~}{S}_1}={\displaystyle \frac{\lambda _Xf^2}{2M_P}};ϵ_{\stackrel{~}{S}_1}={\displaystyle \frac{2v}{f\mathrm{cos}\varphi }};`$ (2.23)
$`m_{\stackrel{~}{S}_2}={\displaystyle \frac{3\lambda _Xf^2}{2M_P}}\mathrm{cos}^2\varphi ;ϵ_{\stackrel{~}{S}_2}={\displaystyle \frac{4v}{f}}.`$ (2.24)
There are clearly many possible more complicated variations on these models; for example, schemes with more than two fields $`X_i`$ participating in the PQ-breaking and $`\mu `$-generating dynamics (see, for example, ref. ). It is also possible to have a scheme in which there is only one field $`X`$, which obtains a VEV at an intermediate scale below where the soft mass term $`m_X^2`$ runs negative. This corresponds to the $`\varphi 0,\lambda _X0`$ limit of eq. (2.14) with $`\stackrel{~}{S}_2`$ removed, so there is just a light axino with $`ϵ=2v/f`$. The essential features of all these models are that they contain one or more singlino fields, with couplings naively of order $`v/f`$ but which can be significantly enhanced, and which can easily be lighter than all of the MSSM odd-$`R`$-parity sparticles.
## 3 Mixing of the singlino with MSSM neutralinos
As shown in the previous section, one or both singlino mass eigenstates $`\stackrel{~}{S}_1`$ or $`\stackrel{~}{S}_2`$ can be lighter than all MSSM sparticles. In this section, I will consider the relevant mixings and couplings of such a singlino to the MSSM states, and the ensuing decay partial widths. I will use $`\stackrel{~}{S}`$ to refer generically to either $`\stackrel{~}{S}_1`$ or $`\stackrel{~}{S}_2`$.
The properties of the singlino are determined by the superpotential eq. (1.3). At tree level, there are singlino-higgsino-Higgs boson couplings. Other couplings, including singlino-fermion-sfermion and singlino-higgsino-$`Z`$ boson, arise due to singlino-gaugino and singlino-higgsino mixing. In order to discover the couplings of the $`\stackrel{~}{S}`$ to the physical MSSM states, one must diagonalize the the $`5\times 5`$ neutralino mass matrix. In the ($`\stackrel{~}{S}`$, $`\stackrel{~}{B}`$, $`\stackrel{~}{W}^0`$, $`\stackrel{~}{H}_d^0`$, $`\stackrel{~}{H}_u^0`$) basis, it is given by:
$`M^{(5)}=\left(\begin{array}{ccccc}m_{\stackrel{~}{S}}& 0& 0& ϵ\mu s_\beta & ϵ\mu c_\beta \\ 0& M_1& 0& c_\beta s_Wm_Z& s_\beta s_Wm_Z\\ 0& 0& M_2& c_\beta c_Wm_Z& s_\beta c_Wm_Z\\ ϵ\mu s_\beta & c_\beta s_Wm_Z& c_\beta c_Wm_Z& 0& \mu \\ ϵ\mu c_\beta & s_\beta s_Wm_Z& s_\beta c_Wm_Z& \mu & 0\end{array}\right),`$ (3.1)
where $`s_\beta ,c_\beta `$ stand for $`\mathrm{sin}\beta `$, $`\mathrm{cos}\beta `$, and $`s_W,c_W`$ for $`\mathrm{sin}\theta _W`$, $`\mathrm{cos}\theta _W`$. In the following, I shall take $`m_{\stackrel{~}{S}}`$ to be real and positive without loss of generality. This allows $`ϵ`$ to have an arbitrary phase. Now, the off-diagonal terms proportional to $`ϵ`$ can be treated as a perturbation. Therefore, our procedure is to first diagonalize $`M^{(4)}`$, the lower right $`4\times 4`$ mass sub-matrix. This is accomplished with a unitary matrix $`Z_{ij}`$ $`(i,j=1,\mathrm{},4)`$ according to:
$`Z_{ik}^{}M_{kl}^{(4)}Z_{jl}^{}=\delta _{ij}m_{\stackrel{~}{N}_j}.`$ (3.2)
Here the masses $`m_{\stackrel{~}{N}_j}`$ are real and positive; this can always be done, regardless of the relative complex phases of $`\mu `$, $`M_1`$ and $`M_2`$. To the lowest order in a perturbative expansion in $`ϵ`$, the singlino $`\stackrel{~}{S}`$ is a mass eigenstate, and the ordinary MSSM neutralinos have the same masses that they would have had if $`\stackrel{~}{S}`$ were absent. I will choose an ordering scheme such that $`\stackrel{~}{S}=\stackrel{~}{N}_0`$, with $`m_{\stackrel{~}{S}}=m_{\stackrel{~}{N}_0}<m_{\stackrel{~}{N}_1}<m_{\stackrel{~}{N}_3}<m_{\stackrel{~}{N}_3}<m_{\stackrel{~}{N}_4}`$.
The full $`5\times 5`$ neutralino-singlino mass matrix can then be diagonalized according to
$`N_{ik}^{}M_{kl}^{(5)}N_{jl}^{}=\delta _{ij}m_{\stackrel{~}{N}_j},`$ (3.3)
where now $`i,j=0,1,\mathrm{},4`$. To lowest order in a perturbation in $`ϵ`$, one finds that $`N_{00}=1`$ and $`N_{ij}=Z_{ij}`$ for $`i,j=1,2,3,4`$. So one can write
$`N_{ij}=\left(\begin{array}{cc}1& N_{0j}\\ N_{i0}& Z_{ij}\end{array}\right).`$ (3.4)
The neutralino-singlino mixing elements can be determined in terms of the $`4\times 4`$ $`Z_{ij}`$, the mass eigenvalues, and the parameter $`ϵ`$:
$`N_{0j}`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{4}{}}}Z_{kj}\left[ϵ^{}\mu ^{}m_{\stackrel{~}{N}_k}(s_\beta Z_{k3}+c_\beta Z_{k4})+ϵ\mu m_{\stackrel{~}{S}}(s_\beta Z_{k3}^{}+c_\beta Z_{k4}^{})\right]/(m_{\stackrel{~}{N}_k}^2m_{\stackrel{~}{S}}^2);`$ (3.5)
$`N_{i0}`$ $`=`$ $`\left[ϵ\mu m_{\stackrel{~}{N}_i}(s_\beta Z_{i3}^{}+c_\beta Z_{i4}^{})+ϵ^{}\mu ^{}m_{\stackrel{~}{S}}(s_\beta Z_{i3}+c_\beta Z_{i4})\right]/(m_{\stackrel{~}{N}_i}^2m_{\stackrel{~}{S}}^2).`$ (3.6)
(Again, nothing has been assumed here about the complex phases of the parameters $`M_1`$, $`M_2`$, $`\mu `$ and $`ϵ`$.) This expansion is an excellent approximation because $`|ϵ|`$ is very small. In particular, one can check by numerical diagonalization of eq. (3.1) that the perturbative result is very accurate unless the denominator $`m_{\stackrel{~}{N}_1}^2m_{\stackrel{~}{S}}^2`$ in eqs. (3.5) and (3.6) is tuned to 0 to an extreme accuracy (comparable to $`|ϵ\mu m_{\stackrel{~}{S}}|`$), in which case the decay widths studied below are completely negligible anyway.
The relevant neutralino and singlino couplings for low energy phenomenology are all contained in the five eigenmasses $`m_{\stackrel{~}{N}_i}`$ and $`m_{\stackrel{~}{S}}`$ (for which corrections proportional to $`ϵ`$ are negligible) and the neutralino mixing quantities $`N_{ij}`$ and $`N_{0j}`$ with ($`i,j=1,\mathrm{},4`$). The decay rates with a singlino in the final state can now be worked out, as in the NMSSM . I present formulas for the relevant widths in an Appendix, taking care to allow for possible arbitrary phases in the neutralino mixing matrix and in $`ϵ`$. In general, the phase of $`ϵ`$ could be anything, since it is not constrained by low-energy experiments on CP violation and is not necessarily correlated with the phase of $`\mu `$. Because the results might be useful for other problems, I will provide the general results for a decay involving $`\stackrel{~}{N}_i`$ to $`\stackrel{~}{N}_j`$; the case needed in this paper is obtained by simply taking $`i=1`$ and $`j=0`$. Numerical illustrations of these results will be given in Section 4.
The necessary coupling for the decay
$`\stackrel{~}{N}_1h^0\stackrel{~}{S}`$ (3.7)
arises directly from the superpotential eq. (1.3) through the higgsino content of $`\stackrel{~}{N}_1`$. It also obtains contributions from the singlino-higgsino mixings $`N_{03}`$ and $`N_{04}`$ combined with the gaugino content of $`\stackrel{~}{N}_1`$, and from the singlino-gaugino mixing elements $`N_{01}`$ and $`N_{02}`$ combined with the higgsino content of $`\stackrel{~}{N}_1`$. The total coupling is given explicitly by eq. (A.2). The Higgs will then decay according to $`h^0b\overline{b}`$, $`WW^{}`$, $`\tau ^+\tau ^{}`$, $`c\overline{c}`$, or $`g\overline{g}`$, providing a visible product displaced from the original interaction point producing the event.
The decay
$`\stackrel{~}{N}_1Z\stackrel{~}{S}`$ (3.8)
relies both on singlino-higgsino mixings $`N_{03}`$ and $`N_{04}`$ and on the higgsino content of $`\stackrel{~}{N}_1`$, and is therefore somewhat more suppressed in models with a bino-like NLSP. The relevant coupling is given explicitly by eq. (A.6). The $`Z`$ boson then decays with Standard Model branching fractions to quark-antiquark and lepton-antilepton pairs.
When neither two-body decay is open, the neutralino NLSP will decay through off-shell sleptons, squarks, the $`Z`$, and Higgs bosons. The full expressions for these decays are given in the Appendix. The most important contribution for a bino-like NLSP is typically through sleptons, and therefore relies on the singlino-bino and singlino-wino mixing elements $`N_{01}`$ and $`N_{02}`$. Fortunately, these are not greatly suppressed in many models unless $`|\mu |`$ is very large.
There remains the possibility of a two-body decay $`\stackrel{~}{N}_1\gamma \stackrel{~}{S}`$, which arises at the one-loop level. However, using the results of (with appropriate changes for the couplings to correspond to the model under present consideration), I have verified that these decays are always suppressed compared to the two- and three-body decays considered here, with widths that cannot be much larger than a few times (1000 km)<sup>-1</sup> in the examples in the next section with $`ϵ=10^8`$. Since these decays are not competitive here, I will not present results for them.
A separate possibility is that the NLSP is a stau, or that all three lighter, mostly-right-handed, slepton mass eigenstates ($`\stackrel{~}{\tau }_1,\stackrel{~}{\mu }_R,\stackrel{~}{e}_R`$) have no open decays except to the singlino. In that case, one can hope to observe $`\stackrel{~}{\tau }_1\tau \stackrel{~}{S}`$ (and perhaps $`\stackrel{~}{\mu }_R\mu \stackrel{~}{S}`$ and $`\stackrel{~}{e}_Re\stackrel{~}{S}`$). The decaying slepton will appear in the detector as a muon-like charged particle track, or as a track with an anomalously high ionization rate. The rare decay to the singlino will yield a large-angle kink in the track leading either to a tau jet or an electron or muon. Since the decaying particle is heavy, there will be a significant angle at the kink. The decay widths are suppressed only by the singlino-gaugino mixing, so they can occur within the detector often enough to measure, even if $`ϵ`$ is significantly less than $`10^8`$.
## 4 Representative results for decays to the singlino
In this section, I will consider some illustrative numerical results for decays to the singlino, first for neutralino NLSP models and then for stau or slepton NLSP models. I will take $`ϵ=10^8`$, with the understanding that the results have to be scaled according to $`\mathrm{\Gamma }ϵ^2`$.
### 4.1 Neutralino decays
In order to study the decay partial widths of a neutralino, I will employ the concept of “model lines”, in which one supersymmetry-breaking parameter is allowed to vary, setting the overall scale for all sparticle masses. First, consider a typical model scenario with a bino-like NLSP and the LSP singlino mass fixed at $`m_{\stackrel{~}{S}}=50`$ GeV. The bino mass parameter $`M_1`$ is varied, with the wino mass parameter $`M_2`$ and the $`\mu `$ term then determined according to $`M_2=2.0M_1`$; $`\mu =3.0M_1`$. The $`5\times 5`$ neutralino mass matrix is then fully determined by also choosing fixed values of $`\mathrm{tan}\beta =3.0`$ and $`ϵ=10^8`$. The right-handed slepton masses $`m_{\stackrel{~}{e}_R}=m_{\stackrel{~}{\mu }_R}=m_{\stackrel{~}{\tau }_R}`$ are constrained to be the greater of $`1.2m_{\stackrel{~}{N}_1}`$ and 110 GeV. This assures that a slepton cannot be the NLSP and should not be found at LEP. Mixing in the stau sector is neglected. The left-handed slepton masses are determined by $`m_{\stackrel{~}{e}_L}^2=m_{\stackrel{~}{e}_R}^2+0.5M_2^2`$. I will assume that squarks are not light enough to give a significant contribution to the decay. Finally, the lightest Higgs boson mass is assumed to be $`m_{h^0}=120`$ GeV, safely out of the reach of LEP, and to obey the decoupling limit $`\alpha =\beta \pi /2`$.
The results for the partial decay widths to visible states (excluding neutrinos), as found from the equations in the Appendix, are shown in Figure 1 as a function of $`m_{\stackrel{~}{N}_1}`$.
For $`m_{\stackrel{~}{N}_1}\stackrel{<}{}120`$ GeV in this model line, the decays are dominated by the contributions of the virtual right-handed sleptons. The total inverse decay lengths are of order tens of kilometers, and are nearly democratic between $`e^+e^{}`$, $`\mu ^+\mu ^{}`$ and $`\tau ^+\tau ^{}`$ final states. The “knee” near $`m_{\stackrel{~}{N}_1}=92`$ GeV is merely an artifact of the constraint $`m_{\stackrel{~}{e}_R}>110`$ GeV; for smaller neutralino masses, the virtual slepton is necessarily more off-shell due to the constraint that it has not been discovered at LEP.
For $`m_{\stackrel{~}{N}_1}\stackrel{>}{}120`$ GeV, the contributions from the virtual $`Z`$ boson \[the terms $`W_Z`$, $`W_{Zt}`$, and $`W_{Zu}`$ in eq. (A.12)\] begin to be important. This increases the decay partial widths, with contributions to $`f\overline{f}`$ that are roughly proportional to the $`Z`$ branching fraction, so that dijet final states dominate. For $`m_{\stackrel{~}{N}_1}>142`$ GeV, the virtual Z boson is on-shell, and the decay becomes two-body $`\stackrel{~}{N}_1Z^0\stackrel{~}{S}`$. (The three-body formula is used in the vicinity of threshold, however, in order to correctly include interference effects with the virtual slepton diagrams in that regime.) This leads to a total visible decay width greater than (1000 meters)<sup>-1</sup>. For negative $`\mu `$, the decay widths tend to be somewhat smaller.
Finally, for $`m_{\stackrel{~}{N}_1}>m_{\stackrel{~}{S}}+m_{h^0}=170`$ GeV in this model line, the decay $`\stackrel{~}{N}_1h^0\stackrel{~}{S}`$ opens up and completely dominates. Since there is a direct higgsino-singlino-Higgs boson coupling, this is much larger than the two-body decay to $`Z\stackrel{~}{S}`$, even though the $`Z`$ boson is lighter. The results are shown assuming Standard Model branching fractions for $`h^0`$ into final states $`b\overline{b}`$, $`WW^{}`$, $`\tau ^+\tau ^{}`$, and (lumped together into the “$`jj`$” category) $`c\overline{c}`$ and $`gg`$. Here the partial decay width of $`\stackrel{~}{N}_1`$ to the $`b\overline{b}\stackrel{~}{S}`$ final state is found to be of order (10 meters)<sup>-1</sup>. Of course, if the $`h^0`$ mass is smaller, this mode will open up and dominate for smaller values of $`m_{\stackrel{~}{N}_1}`$.
In the era after supersymmetry is discovered, the situation will be rather different; we will presumably know the MSSM sparticle mass spectrum, but the singlino mass and coupling will be completely unknown. So, a more useful summary of the situation we could face might be something like that shown in Figure 2. This is a particular point along the same model line just discussed, with $`m_{\stackrel{~}{N}_1}=150`$ GeV, and with the horizontal axis representing the possible values of the singlino mass.
A somewhat different scenario ensues if the NLSP is a higgsino-like neutralino. To illustrate this, I choose a pair of model lines with $`\mu =\pm 0.8M_1`$, and all other parameter relationships as described above for Figure 1. The results are shown in Figure 3, but now only for the total visible decay width. Since $`\stackrel{~}{N}_1`$ has a smaller gaugino content, the decays through virtual sleptons are highly suppressed. Conversely, the large higgsino component of $`\stackrel{~}{N}_1`$ enhances the probability of decay through a virtual $`Z`$ boson. So, for $`m_{\stackrel{~}{N}_1}<m_{\stackrel{~}{S}}+m_{\stackrel{~}{h}^0}=170`$ GeV in this model line, the $`f\overline{f}`$ decays will obey $`Z`$ boson branching fractions. However, there turns out to be an accidental suppression of the matrix element for $`\mu <0`$ in the convention specified in eq. (3.1) (which is the same as in refs. ), particularly when the $`Z`$ boson is off-shell. For $`m_{\stackrel{~}{N}_1}>170`$ GeV, the decay $`\stackrel{~}{N}_1h^0\stackrel{~}{S}`$ length is of order several meters if $`\stackrel{~}{N}_1`$ is mostly higgsino.
In the above analyses, I have assumed that $`H^0`$ and $`A^0`$ are very heavy, and that three-body amplitudes involving them are negligible. Although this is appropriate throughout most of parameter space, it is possible that for large $`\mathrm{tan}\beta `$, the couplings of $`H^0`$ and $`A^0`$ could be large enough to make an appreciable contribution to $`\mathrm{\Gamma }(\stackrel{~}{N}_1b\overline{b}\stackrel{~}{S})`$, even if $`H^0`$ and $`A^0`$ are far off-shell.
Finally, I note that the decay widths discussed here depend on the phase of the parameter $`ϵ`$. This phase is not constrained by low-energy CP violating observables, and so might be considered completely arbitrary. I have checked, using the formulas in the Appendix, that varying Arg($`ϵ`$), while keeping all other parameters fixed, can change the $`\stackrel{~}{N}_1`$ decay widths by an order of magnitude or so. In models with a bino-like NLSP as in Figure 1, the largest decay widths tend to occur for real $`ϵ`$.
### 4.2 Slepton decays
It is also possible that the NLSP is a slepton $`\stackrel{~}{\tau }_1`$ or, effectively, all three mainly right-handed sleptons $`\stackrel{~}{e}_R`$, $`\stackrel{~}{\mu }_R`$ and $`\stackrel{~}{\tau }_1`$. The latter scenario is realized if $`\mathrm{tan}\beta `$ is not too large, so that the three sleptons are mass-degenerate to within less than $`m_\tau `$. These possibilities are familiar - in gauge-mediated supersymmetry breaking models but could also be realized in supergravity-mediated models if there is not a large universal contribution to scalar masses, or if $`D`$-term contributions (proportional to some exotic $`U(1)`$ quantum number) are large.
If $`\stackrel{~}{\tau }_1`$ is the NLSP, then the two-body decays $`\stackrel{~}{\tau }_1\tau \stackrel{~}{S}`$ are suppressed only by the bino-singlino mixing. In Figure 4, I show the results for this decay width for a typical model line with varying $`M_1`$ and fixed LSP mass $`m_{\stackrel{~}{S}}`$. In order to ensure that a stau is the NLSP, the constraint $`m_{\stackrel{~}{\tau }_1}=0.9m_{\stackrel{~}{N}_1}`$ (thick lines) and $`0.7m_{\stackrel{~}{N}_1}`$ (thin lines) are imposed. Other relevant model line parameters are $`M_2=2.0M_1`$, $`\mu =3.0M_1`$, $`\mathrm{tan}\beta =3.0`$. To a good approximation, the decay width depends only on the absolute value of $`s_{\stackrel{~}{\tau }}`$ in the stau mixing parameterization of eq. (A.24), so results are shown for $`s_{\stackrel{~}{\tau }}=0`$, $`0.25`$, and $`0.5`$. The solid line is also approximately true for $`\stackrel{~}{e}_Re\stackrel{~}{S}`$ and $`\stackrel{~}{\mu }_R\mu \stackrel{~}{S}`$ as a function of $`m_{\stackrel{~}{e}_R}`$ and $`m_{\stackrel{~}{\mu }_R}`$, by taking $`s_{\stackrel{~}{\tau }}=0`$, $`c_{\stackrel{~}{\tau }}=1`$.
Because there is relatively little suppression in this case, the inverse decay widths are of order tens of meters. This is discernible at a collider which can produce several hundred supersymmetric events. Each such event would contain a pair of quasi-stable stau or slepton highly ionizing tracks, which can have an anomalously high $`dE/dx`$ to distinguish them from muons. In a small fraction of events, one of the stau or slepton tracks will have a kink leading to a lepton or tau jet, corresponding to the decay. The resulting tau or lepton would have a significant angle with respect to the original highly ionizing track, yielding a potentially spectacular and nearly background-free signal.
## 5 Conclusions
The presence of the $`\mu `$ term in the MSSM and the solution to the strong CP problem may have a common explanation at an intermediate scale. Direct detection of the resulting axion is quite problematic. In this paper, I have argued that these models may nevertheless give rise to observable signals at colliders, through delayed decays to singlino fermions that include the axino as a mixture. These events will be a rare (perhaps very rare) occurrence within a large sample of supersymmetric events at the Large Hadron Collider or a future $`e^+e^{}`$ linear collider. There is also a possibility that the lightest MSSM sparticle could have slow decays into more than one singlino. Note that the LHC cross sections can be very large precisely when the NLSP decay widths are small.
The numerical estimate in section 4 of this paper have used $`ϵ=10^8`$ for the singlino-higgsino-Higgs coupling parameter. Of course, the actual value could be significantly smaller. On the other hand, I showed that in some models the couplings are parametrically enhanced, and the mass of one or more singlinos is reduced, if one of the VEVs giving rise to the $`\mu `$ term is relatively small. Furthermore, the magnitude of $`ϵ`$ can be significantly larger if the high scale $`M_P`$ is replaced by a somewhat lower scale that governs non-renormalizable operators, for example a string scale or a compactification scale that is not far above the apparent gauge coupling unification scale.
Future planning and analysis of collider physics experiments should take into account the possibility that the apparent LSP is actually unstable. Besides the models I have discussed here, there are at least two other plausible variations on the MSSM which can lead to delayed rare decays of what might appear, at first, to be the stable LSP.
First, gauge-mediated supersymmetry-breaking (GMSB) models with a supersymmetry-breaking scale $`\sqrt{F}`$ that is not too large will give rise to decays that could have macroscopic proper lengths -. It is interesting to compare the reason for this to that in the models discussed in the present paper. In GMSB models, an estimate for a decay width of the NLSP to the goldstino/gravitino $`\stackrel{~}{G}`$ is $`\mathrm{\Gamma }m_W^5/16\pi (\sqrt{F})^4`$, while in the decays to a singlino/axino LSP $`\stackrel{~}{S}`$, the estimate is $`\mathrm{\Gamma }m_W^3/16\pi f^2`$. So NLSP decays are suppressed by the 4th power of the supersymmetry-breaking scale in GMSB, but only by the square of the PQ scale in light axino/singlino models. GMSB models can, in fact, give rise to signals which might be very difficult to distinguish from those discussed here. For example, if the NLSP is a neutralino with a significant higgsino content, it can have decays $`\stackrel{~}{N}_1h^0\stackrel{~}{G}`$ and $`\stackrel{~}{N}_1Z^0\stackrel{~}{G}`$ that look like the decays discussed in this paper. Or, if a stau is the NLSP, it can appear quasi-stable with rare decays $`\stackrel{~}{\tau }_1\tau \stackrel{~}{G}`$ occurring within the detector. Second, one can have weak $`R`$-parity violating couplings in the MSSM which could also give decays like $`\stackrel{~}{N}_1\mathrm{}^+\mathrm{}^{}\nu `$ or $`\stackrel{~}{N}_1q\overline{q}^{}\nu `$. These signatures could mimic those discussed in the present paper.
If these signals appear, it will be interesting to try to establish the correct explanation from among the competing hypotheses. The prize for doing so will be that we will gain an understanding of physics at scales far above those probed by direct sparticle production at colliders.
Acknowledgements: I thank Howard Haber, Janusz Rosiek, and James Wells for helpful discussions. This work was supported in part by National Science Foundation grant number PHY-9970691.
## Appendix: Complete formulas for neutralino decay widths including the effects of arbitrary phases
In this appendix, I give formulas for the two- and three-body decays of a neutralino to another neutralino and a Higgs boson, $`Z`$ boson, or fermion-antifermion pair. The model is the extension of the MSSM with one singlet superfield as specified by eq. (1.3). This also includes the MSSM and NMSSM as special cases. The results needed in Section 4 of this paper are obtained by taking the neutralino mass eigenstate indices to be $`i=1`$ and $`j=0`$ in the following.
First, let us consider the two-body decay of a neutralino to another neutralino and the lightest CP-even neutral Higgs boson. The relevant decay width is equal to:
$`\mathrm{\Gamma }(\stackrel{~}{N}_ih^0\stackrel{~}{N}_j)`$ $`=`$ $`{\displaystyle \frac{m_{\stackrel{~}{N}_i}}{16\pi }}\sqrt{\lambda (1,r_j^2,r_{h^0}^2)}\left(|G_{ij}^{h^0}|^2(1+r_j^2r_{h^0}^2)+2\mathrm{R}\mathrm{e}[(G_{ij}^{h^0})^2]r_j\right)`$ (A.1)
where $`r_j=m_{\stackrel{~}{N}_j}/m_{\stackrel{~}{N}_i}`$; $`r_{h^0}=m_{h^0}/m_{\stackrel{~}{N}_i}`$; $`\lambda (a,b,c)=a^2+b^2+c^22ab2ac2bc`$; and the neutralino-neutralino-Higgs coupling is given by
$`G_{ij}^{h^0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(gN_{i2}^{}g^{}N_{i1}^{})(s_\alpha N_{j3}^{}+c_\alpha N_{j4}^{})+{\displaystyle \frac{ϵ\mu }{\sqrt{2}v}}(c_\alpha N_{i3}^{}s_\alpha N_{i4}^{})N_{j0}^{}+(ij).`$ (A.2)
Here $`c_\alpha `$ and $`s_\alpha `$ denote $`\mathrm{cos}\alpha `$ and $`\mathrm{sin}\alpha `$ with $`\alpha `$ the Higgs mixing angle in the notation of .
The results for the decays to the heavier neutral CP-even ($`H^0`$) and CP-odd ($`A^0`$) Higgs bosons can be obtained by instead using the couplings:
$`G_{ij}^{H^0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(gN_{i2}^{}g^{}N_{i1}^{})(c_\alpha N_{j3}^{}+s_\alpha N_{j4}^{})+{\displaystyle \frac{ϵ\mu }{\sqrt{2}v}}(s_\alpha N_{i3}^{}+c_\alpha N_{i4}^{})N_{j0}^{}+(ij)`$ (A.3)
for $`\stackrel{~}{N}_iH^0\stackrel{~}{N}_j`$, and
$`G_{ij}^{A^0}`$ $`=`$ $`{\displaystyle \frac{i}{2}}(gN_{i2}^{}g^{}N_{i1}^{})(s_\beta N_{j3}^{}c_\beta N_{j4}^{})+i{\displaystyle \frac{ϵ\mu }{\sqrt{2}v}}(c_\beta N_{i3}^{}+s_\beta N_{i4}^{})N_{j0}^{}+(ij)`$ (A.4)
for $`\stackrel{~}{N}_iA^0\stackrel{~}{N}_j`$, and substituting $`m_{h^0}`$ $`m_{H^0}`$ or $`m_{A^0}`$ in the obvious way. However, in the numerical analyses of Section 4, $`H^0`$ and $`A^0`$ are assumed to be heavy and decoupled, so these decays are neglected.
Similarly, the two-body decay of a neutralino to another neutralino and a $`Z`$ boson has a width given by
$`\mathrm{\Gamma }(\stackrel{~}{N}_iZ^0\stackrel{~}{N}_j)={\displaystyle \frac{m_{\stackrel{~}{N}_i}}{16\pi }}\sqrt{\lambda (1,r_j^2,r_Z^2)}\left(|G_{ij}^Z|^2\left[1+r_j^22r_Z^2+(1r_j^2)^2/r_Z^2\right]+6\mathrm{R}\mathrm{e}[(G_{ij}^Z)^2]r_j\right)`$ (A.5)
where $`r_Z=m_Z/m_{\stackrel{~}{N}_i}`$, and the neutralino-neutralino-$`Z`$ coupling is given by
$`G_{ij}^Z={\displaystyle \frac{g}{2c_W}}(N_{i3}N_{j3}^{}+N_{i4}N_{j4}^{}).`$ (A.6)
Finally we consider three-body decays $`\stackrel{~}{N}_if\overline{f}\stackrel{~}{N}_j`$, where $`f`$ is any Standard Model quark or lepton. Results for these decays have appeared in refs. , but here I include the effects of Higgs boson exchanges and arbitrary phases of the couplings . The differential partial widths can be expressed in terms of the dimensionless mass ratios $`r_f=m_f/m_{\stackrel{~}{N}_i}`$, $`r_j=m_{\stackrel{~}{N}_j}/m_{\stackrel{~}{N}_i}`$, $`r_Z=m_Z/m_{\stackrel{~}{N}_i}`$, $`r_{\mathrm{\Gamma }_Z}=\mathrm{\Gamma }_Z/m_{\stackrel{~}{N}_i}`$, and $`r_\varphi =m_\varphi /m_{\stackrel{~}{N}_i}`$ for each of $`\varphi =h^0,A^0`$, and $`H^0`$. I use dimensionless kinematic variables
$`\widehat{s}`$ $`=`$ $`(p_{\stackrel{~}{N}_i}p_{\stackrel{~}{N}_j})^2/p_{\stackrel{~}{N}_i}^2`$ (A.7)
$`\widehat{t}`$ $`=`$ $`(p_{\stackrel{~}{N}_i}p_f)^2/p_{\stackrel{~}{N}_i}^2`$ (A.8)
$`\widehat{u}`$ $`=`$ $`1+r_j^2+2r_f^2\widehat{s}\widehat{t}`$ (A.9)
with limits of integration
$`\widehat{t}_{\mathrm{min},\mathrm{max}}={\displaystyle \frac{1}{2}}[1+r_j^2\widehat{s}+2r_f^2\{(14r_f^2/\widehat{s})\lambda (1,\widehat{s},r_j^2)\}^{1/2}];`$ (A.10)
$`\widehat{s}_{\mathrm{min}}=4r_f^2;\widehat{s}_{\mathrm{max}}=(1r_j)^2.`$ (A.11)
The results for widths can be expressed asIn computing these results, I have neglected fermion masses arising from spinor algebra in matrix elements, but not in the kinematic limits of integration or the couplings.
$`d\mathrm{\Gamma }(\stackrel{~}{N}_if\overline{f}\stackrel{~}{N}_j)={\displaystyle \frac{n_cm_{\stackrel{~}{N}_i}}{512\pi ^3}}({\displaystyle W})d\widehat{t}d\widehat{s},`$ (A.12)
where $`n_c=1`$ (3) for leptons (quarks). The individual contributions to $`W`$ are:
$`W_Z`$ $`=`$ $`{\displaystyle \frac{4(a_f^2+b_f^2)}{(\widehat{s}r_Z^2)^2+r_Z^2r_{\mathrm{\Gamma }_Z}^2}}\{|G_{ij}^Z|^2[(1\widehat{u})(\widehat{u}r_j^2)+(1\widehat{t})(\widehat{t}r_j^2)]`$ (A.13)
$`+2\mathrm{R}\mathrm{e}[(G_{ij}^Z)^2]r_j\widehat{s}\}`$
$`W_t`$ $`=`$ $`{\displaystyle \underset{n,n^{}=1}{\overset{2}{}}}(a_j^na_j^n^{}+b_j^nb_j^n^{})(a_i^na_i^n^{}+b_i^nb_i^n^{}){\displaystyle \frac{(1\widehat{t})(\widehat{t}r_j^2)}{(r_{\stackrel{~}{f}_n}^2\widehat{t})(r_{\stackrel{~}{f}_n^{}}^2\widehat{t})}}`$ (A.14)
$`W_u`$ $`=`$ $`W_t(\widehat{t}\widehat{u})`$ (A.15)
$`W_{tu}`$ $`=`$ $`2\mathrm{R}\mathrm{e}{\displaystyle \underset{n,n^{}=1}{\overset{2}{}}}{\displaystyle \frac{1}{(r_{\stackrel{~}{f}_n}^2\widehat{t})(r_{\stackrel{~}{f}_n^{}}^2\widehat{u})}}[(a_j^nb_j^n^{}a_i^n^{}b_i^n+a_j^n^{}b_j^na_i^nb_i^n^{})(r_j^2\widehat{t}\widehat{u})`$ (A.16)
$`+(a_j^na_j^n^{}a_i^na_i^n^{}+b_j^nb_j^n^{}b_i^nb_i^n^{})\widehat{s}r_j]`$
$`W_{Zt}`$ $`=`$ $`{\displaystyle \frac{4(\widehat{s}r_Z^2)}{(\widehat{s}r_Z^2)^2+r_Z^2r_{\mathrm{\Gamma }_Z}^2}}\mathrm{Re}{\displaystyle \underset{n=1}{\overset{2}{}}}[(a_fa_i^na_j^n+b_fb_i^nb_j^n)\{G_{ij}^Z(1\widehat{t})(\widehat{t}r_j^2)`$ (A.17)
$`+G_{ij}^Z\widehat{s}r_j\}]/(r_{\stackrel{~}{f}_n}^2\widehat{t})`$
$`W_{Zu}`$ $`=`$ $`W_{Zt}(\widehat{t}\widehat{u})`$ (A.18)
$`W_{h^0,H^0}`$ $`=`$ $`{\displaystyle \underset{\varphi ,\varphi ^{}=h^0,H^0}{}}{\displaystyle \frac{4\widehat{s}\mathrm{Re}[G_f^\varphi G_f^\varphi ^{}]}{(r_\varphi ^2\widehat{s})(r_\varphi ^{}^2\widehat{s})}}\left\{(1+r_j^2\widehat{s})\mathrm{Re}[G_{ij}^\varphi G_{ij}^\varphi ^{}]+2\widehat{s}r_j\mathrm{Re}[G_{ij}^\varphi G_{ij}^\varphi ^{}]\right\}`$ (A.19)
$`W_{A^0}`$ $`=`$ $`{\displaystyle \frac{4\widehat{s}|G_f^{A^0}|^2}{(r_{A^0}^2\widehat{s})^2}}\left\{(1+r_j^2\widehat{s})|G_{ij}^{A_0}|^2+2\widehat{s}r_j\mathrm{Re}[(G_{ij}^{A_0})^2]\right\}`$ (A.20)
$`W_{\varphi t}`$ $`=`$ $`{\displaystyle \underset{\varphi =h^0,H^0,A^0}{}}{\displaystyle \underset{n=1,2}{}}{\displaystyle \frac{2}{(r_\varphi ^2\widehat{s})(r_{\stackrel{~}{f}_n}^2\widehat{t})}}\mathrm{Re}[\widehat{s}G_f^\varphi (\widehat{t}G_{ij}^\varphi +r_jG_{ij}^\varphi )(a_i^nb_j^n+a_j^nb_i^n)]`$ (A.21)
$`W_{\varphi u}`$ $`=`$ $`W_{\varphi t}(\widehat{t}\widehat{u}).`$ (A.22)
Quantities appearing in the above results are as follows. First,
$`a_f={\displaystyle \frac{g}{c_W}}(T_{3f}q_fs_W^2);b_f=q_fgs_W^2/c_W`$ (A.23)
are the $`Z`$ boson couplings to quarks and leptons with $`(T_{3f},q_f)=(1/2,2/3)`$ for up-type quarks, $`(1/2,1/3)`$ for down-type quarks, and $`(1/2,1)`$ for charged leptons. Left-right mixing and CP violation in the sfermion sector are parameterized by a unitary matrix, which I chooseThis parameterization has the feature that in the typical unmixed, CP-conserving case $`c_{\stackrel{~}{f}}=1`$, $`s_{\stackrel{~}{f}}=0`$, $`\stackrel{~}{f}_1=\stackrel{~}{f}_R`$ and $`\stackrel{~}{f}_2=\stackrel{~}{f}_L`$, with no minus signs. to write as
$`\left(\begin{array}{c}\stackrel{~}{f}_R\\ \stackrel{~}{f}_L\end{array}\right)=\left(\begin{array}{cc}c_{\stackrel{~}{f}}& s_{\stackrel{~}{f}}\\ s_{\stackrel{~}{f}}^{}& c_{\stackrel{~}{f}}^{}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{f}_1\\ \stackrel{~}{f}_2\end{array}\right),`$ (A.24)
where $`|c_{\stackrel{~}{f}}|^2+|s_{\stackrel{~}{f}}|^2=1`$, and $`m_{\stackrel{~}{f}_1}<m_{\stackrel{~}{f}_2}`$, The resulting couplings for down-type sfermions $`(\stackrel{~}{b},\stackrel{~}{\tau })`$ are:
$`a_i^{\stackrel{~}{f}_1}`$ $`=`$ $`\sqrt{2}s_{\stackrel{~}{f}}\left[gT_{3f}N_{i2}^{}+g^{}(q_fT_{3f})N_{i1}^{}\right]c_{\stackrel{~}{f}}^{}gN_{i3}^{}m_f/\sqrt{2}c_\beta m_W,`$ (A.25)
$`b_i^{\stackrel{~}{f}_1}`$ $`=`$ $`\sqrt{2}c_{\stackrel{~}{f}}^{}g^{}q_fN_{i1}+s_{\stackrel{~}{f}}gN_{i3}m_f/\sqrt{2}c_\beta m_W`$ (A.26)
$`a_i^{\stackrel{~}{f}_2}`$ $`=`$ $`\sqrt{2}c_{\stackrel{~}{f}}\left[gT_{3f}N_{i2}^{}+g^{}(q_fT_{3f})N_{i1}^{}\right]s_{\stackrel{~}{f}}^{}gN_{i3}^{}m_f/\sqrt{2}c_\beta m_W,`$ (A.27)
$`b_i^{\stackrel{~}{f}_2}`$ $`=`$ $`\sqrt{2}s_{\stackrel{~}{f}}^{}g^{}q_fN_{i1}c_{\stackrel{~}{f}}gN_{i3}m_f/\sqrt{2}c_\beta m_W`$ (A.28)
for $`i=0,1,\mathrm{}4`$. For up-type fermions one must replace $`N_{i3}^{()}/c_\beta `$ by $`N_{i4}^{()}/s_\beta `$ in eqs. (A.25)-(A.28). (However, for the cases of interest in this paper, $`\stackrel{~}{N}_1t\overline{t}\stackrel{~}{S}`$ is surely not kinematically allowed, and decays $`\stackrel{~}{N}_1\nu \overline{\nu }S`$ are not interesting.) In the expressions for the contributions to the widths, I have used the abbreviations $`a_i^n=a_i^{\stackrel{~}{f}_n}`$, etc. The various Higgs boson couplings to Standard Model fermions are given by, e.g.:
$`G_b^{h^0}={\displaystyle \frac{gm_bs_\alpha }{2m_Wc_\beta }};G_t^{h^0}={\displaystyle \frac{gm_tc_\alpha }{2m_Ws_\beta }};`$ (A.29)
$`G_b^{H^0}={\displaystyle \frac{gm_bc_\alpha }{2m_Wc_\beta }};G_t^{H^0}={\displaystyle \frac{gm_ts_\alpha }{2m_Ws_\beta }};`$ (A.30)
$`G_b^{A^0}=i{\displaystyle \frac{gm_b\mathrm{tan}\beta }{2m_W}};G_t^{A^0}=i{\displaystyle \frac{gm_t\mathrm{cot}\beta }{2m_W}},`$ (A.31)
for bottom and top quarks. The Higgs couplings for taus are obtained by $`m_bm_\tau `$, and for other quarks and leptons by the obvious substitutions.
The two-body decay width for a stau to a neutralino or singlino is given by:
$`\mathrm{\Gamma }(\stackrel{~}{\tau }_1\tau \stackrel{~}{N}_i)`$ $`=`$ $`{\displaystyle \frac{m_{\stackrel{~}{\tau }_1}}{16\pi }}\sqrt{\lambda (1,r_i^2,r_\tau ^2)}\left\{(|a_i^{\stackrel{~}{\tau }_1}|^2+|b_i^{\stackrel{~}{\tau }_1}|^2)(1r_i^2r_\tau ^2)4r_\tau r_i\mathrm{Re}[a_i^{\stackrel{~}{\tau }_1}b_i^{\stackrel{~}{\tau }_1}]\right\}`$ (A.32)
where now $`r_i=m_{\stackrel{~}{N}_i}/m_{\stackrel{~}{\tau }_1}`$ and $`r_\tau =m_\tau /m_{\stackrel{~}{\tau }_1}`$, and the couplings $`a_i^{\stackrel{~}{\tau }_1}`$, $`b_i^{\stackrel{~}{\tau }_1}`$ are given already by eqs. (A.25)-(A.28). In section 4, this formula is used with $`i=0`$, corresponding to $`\stackrel{~}{\tau }_1\tau \stackrel{~}{S}`$. The results for $`\stackrel{~}{\mu }_R\mu \stackrel{~}{S}`$ and $`\stackrel{~}{e}_Re\stackrel{~}{S}`$ are obtained by taking $`c_{\stackrel{~}{\tau }}1`$ and $`s_{\stackrel{~}{\tau }}0`$.
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# Quantum cosmology of 5D non-compactified Kaluza-Klein theory
## 1 Introduction
The subject of initial conditions is one of the most important questions in cosmological models. Unlike a classical system where the dynamical equations are solved subject to the initial conditions, for cosmological models there are no initial conditions external to the universe to be considered for solving the Einstein equations. This is because there is no time parameter external to the universe. We know that the issue of initial conditions in classical cosmology corresponds to a boundary condition problem in quantum cosmology. Therefore, it seems that the initial conditions should be introduced, from the outside, within some boundary conditions. Two well-known proposals commonly used in the literature are the Hartle-Hawking no boundary proposal -, and the Vilenkin tunneling proposal -. The first proposal is that the universe has no boundary in 4D Euclidean space and the second one states that only the outgoing modes of the wave function should be taken at the singular boundary of superspace. Some attempts to generalize these proposals to higher-dimensional Kaluza-Klein cosmologies, to find a reasonable explanation to the large separation of the scale of the observed three-dimensional universe and the scale of the extra dimensions have already been done , . In these works, the extra dimensions were supposed to be stable and compactified by a cyclic symmetry to a small size.
In this paper, we follow another approach and investigate the quantum cosmology of a non-compactified Kaluza-Klein theory developed by Wesson and co-workers -. Unlike the usual Kaluza-Klein theory in which a cyclic symmetry associated with the extra dimension is assumed, the new approach removes the cyclic condition on the extra-dimension, and derivatives of the metric with respect to the extra-coordinate are retained. This induces non-trivial matter on the hypersurfaces of $`l=`$ const.
Our goal is to investigate the effect of the $`l`$-dependence of the metric on the quantum cosmology of a simple model and obtain the relevant initial condition on the fifth coordinate. We find that while the Vilenkin wave function leads to probability distribution of quantum tunneling peaked around the $`l0`$ hypersurface, the Hartle-Hawking wave function leads to large $`l`$ values corresponding to highest probability for the birth of a Lorentzian universe. This leads us to choose the Vilenkin wave function as the more convenient one for the non-compactified Kaluza-Klein theory, since it seems unnatural that the 4-dimensional universe was born very far from the $`l0`$ constant hypersurface.
## 2 Wheeler-DeWitt equation
In this non-compactified Kaluza-Klein theory we choose the general 5-dimensional metric as:
$$ds^2=\widehat{N}^2(t,l)e^{b\sigma (t)}dt^2+\frac{\widehat{a}^2(t,l)e^{b\sigma (t)}}{[1+\frac{1}{4}kr^2]^2}dx^idx^i+ϵe^{2b\sigma (t)}dl^2.$$
(1)
Here $`l`$ is the fifth coordinate, and $`\widehat{N}(t,l)=N(t)f(l)`$, $`\widehat{a}(t,l)=a(t)\chi (l)`$ are the $`l`$-dependent separable lapse function and scale factor respectively. Also $`k=0,\pm 1`$ is the curvature, $`b`$ is a parameter, $`\sigma (t)`$ is a dilaton field and $`ϵ=\pm 1`$ which leaves the signature of the fifth dimension general. The $`l`$-dependence of the 4D geometry, namely $`\widehat{N}(t,l)`$ and $`\widehat{a}(t,l)`$ indicates that the cyclic condition on the fifth coordinate is removed. The Ricci curvature scalar is calculated to be
$$\widehat{}=\left[6\frac{\ddot{\widehat{a}}}{\widehat{N}^2\widehat{a}}+b\frac{\ddot{\sigma }}{\widehat{N}^2}6\frac{\dot{\widehat{a}}^2}{\widehat{N}^2\widehat{a}^2}\frac{3b^2}{2}\frac{\dot{\sigma }^2}{\widehat{N}^2}+3b\frac{\dot{\widehat{a}}\dot{\sigma }}{\widehat{N}^2\widehat{a}}6\frac{k}{\widehat{a}^2}\right]e^{b\sigma }+6ϵ\left[\frac{\widehat{N}^{}\widehat{a}^{}}{\widehat{N}\widehat{a}}+\frac{\widehat{N}^{\prime \prime }}{3\widehat{N}}+\frac{\widehat{a}^2}{\widehat{a}^2}+\frac{\widehat{a}^{\prime \prime }}{\widehat{a}}\right]e^{2b\sigma }.$$
(2)
Inserting the scalar curvature into the 5D vacuum Einstein-Hilbert action
$$S=𝑑x^5\sqrt{\widehat{g}}\widehat{},$$
leads to the following effective action
$$S=\widehat{L}𝑑t,$$
with
$$\widehat{L}=\widehat{N}\left[\frac{\widehat{a}\dot{\widehat{a}}^2}{2\widehat{N}^2}\frac{b^2}{8}\frac{\widehat{a}^3\dot{\sigma }^2}{\widehat{N}^2}\frac{1}{2}k\widehat{a}+\frac{ϵ}{2}e^{3b\sigma }\widehat{a}^3F(l)\right].$$
(3)
Here
$$F(l)=\frac{f^{}\chi ^{}}{f\chi }+\frac{f^{\prime \prime }}{3f}+\frac{\chi ^2}{\chi ^2}+\frac{\chi ^{\prime \prime }}{\chi },$$
(4)
and the dilaton field potential is
$$U(\sigma )=ϵe^{3b\sigma }F(l).$$
(5)
The Hamiltonian form of the action can be written
$$S=(\widehat{P}_a\dot{\widehat{a}}+P_\sigma \dot{\sigma }\widehat{N}\widehat{H})𝑑t,$$
where $`\widehat{N}`$ appears as the Lagrange multiplier. The variation of the action with respect to $`\widehat{N}`$ leads to the Hamiltonian constraint
$$\widehat{H}=\frac{\widehat{P}_a^2}{2\widehat{a}}\frac{2}{b^2}\frac{P_\sigma ^2}{\widehat{a}^3}+\frac{1}{2}k\widehat{a}\frac{1}{2}\widehat{a}^3U(\sigma )=0.$$
(6)
Since no lapse function $`N(t)`$ appears in the Hamiltonian, it is an indication of the invariance of the Hamiltonian under time reparametrization. This means $`N(t)`$ has no physical importance and we may choose the common gauge in cosmology, namely $`N(t)=1`$. The Wheeler-DeWitt equation in the minisuperspace of coordinates $`0<R<\mathrm{},\mathrm{}<\sigma <\mathrm{}`$ can then be written as
$$\left[\frac{1}{\widehat{a}^p}\frac{}{\widehat{a}}\widehat{a}^p\frac{}{\widehat{a}}+\frac{4}{b^2\widehat{a}^3}\frac{^2}{\sigma ^2}+k\widehat{a}\widehat{a}^3U(\sigma )\right]\mathrm{\Psi }(\widehat{a},\sigma )=0.$$
(7)
Here $`p`$ covers the ambiguity in factor-ordering, but since we will restrict ourselves to the semiclassical approximation omitting the first derivatives, this factor is not important . Then, the Wheeler-DeWitt equation can be rewritten as
$$\left[\frac{^2}{\widehat{a}^2}+\frac{4}{b^2\widehat{a}^2}\frac{^2}{\sigma ^2}+W(\widehat{a},\sigma )\right]\mathrm{\Psi }(\widehat{a},\sigma )=0,$$
(8)
where
$$W(\widehat{a},\sigma )=\widehat{a}^2[k\widehat{a}^2U(\sigma )]$$
(9)
is the superpotential. In the investigation of the Wheeler-DeWitt equation we need to know the properties of the superpotential $`W(\widehat{a},\sigma )`$. For fixed $`\sigma `$ the superpotential consists of two terms, a curvature term $`k\widehat{a}`$ and the term $`\widehat{a}^3U(\sigma )`$, where $`U(\sigma )`$ acts effectively as a cosmological term. By appropriate choices for $`k`$ and $`ϵ`$ the superpotential may have a maximum necessary for quantum tunneling. However, along the line of fixed $`\widehat{a}`$, the potential $`U(\sigma )`$ has no a maximum. This would suggest that we may concentrate on the $`\widehat{a}`$ coordinate as a viable dynamical variable in the investigation of quantum tunneling in the $`\widehat{a}`$ direction, and consider the coordinate $`\sigma `$ as a parameter. We will discuss this subject in the next section.
The relevant classical cosmology subject to quantization is a closed universe with $`k=+1`$. This is because the universes with $`k=1,0`$ lead to an infinite volume in the integration of the action and so the nucleation probability of the universe in the quantum creation procces would be zero. Thus, we take a closed universe, $`k=+1`$. The superpotential (9) exhibits a barrier if the dilaton potential (5) is greater than zero
$$U(\sigma )>0,$$
or (if $`ϵ=1`$)
$$F(l)>0.$$
In the semiclassical approximation we want to find the most probable initial conditions for the classical motion of the universe. This motion is controlled by the superpotential (9). To this end, we write
$$\widehat{a}=\widehat{a}_0(\sigma )=\frac{1}{\sqrt{U(\sigma )}},$$
(10)
and so define a surface of constant superpotential $`W=0`$ in the minisuperspace. In fact, equation (9) describes a superpotential barrier in the $`\widehat{a}`$ direction and equation (10) separates the under-barrier region $`0<\widehat{a}<\widehat{a}_0`$ from the outer region $`\widehat{a}>\widehat{a}_0`$. From (10) we can also obtain an equation for $`f(l)`$ and $`\chi (l)`$:
$$\chi \chi ^{}\frac{f^{}}{f}+\frac{\chi ^2}{3}\frac{f^{\prime \prime }}{f}+\chi _{}^{}{}_{}{}^{2}+\chi \chi ^{\prime \prime }=\text{const}.$$
(11)
The presence of a barrier region indicates that we can consider the nucleation of the universe through a quantum tunneling effect as discussed above. Therefore, we proceed to consider the well-known tunneling condition of Vilenkin . Our aim is to find the approximate analytic solutions of the Wheeler-DeWitt equation, under the barrier and beyond the barrier.
## 3 Wave function
First, we show a behaviour of the “nothing state” for
$$\widehat{a}^2\widehat{a}_0^2.$$
According to Halliwell , in order to have a regular solution for the wave function in the limit $`\widehat{a}0`$, the wave function should be $`\sigma `$-independent because the coefficient of $`_\sigma ^2`$ in the Wheeler-DeWitt equation (8) diverges. The Wheeler-DeWitt equation then reduces to
$$\left[\frac{d}{d\widehat{a}^2}+\widehat{a}^2\right]\mathrm{\Psi }(\widehat{a})=0.$$
(12)
Introducing the auxiliary variable $`\mathrm{\Gamma }(\widehat{a})=\mathrm{\Psi }(\widehat{a})/\widehat{a}^{1/2}`$ and the transformation $`\nu =\widehat{a}^2/2`$, equation (12) reduces to the modified Bessel equation
$$\nu ^2d_\nu ^2\mathrm{\Gamma }+\nu d_\nu \mathrm{\Gamma }\left(\nu ^2+\frac{1}{16}\right)\mathrm{\Gamma }=0,$$
(13)
whose independent solutions are the well-known modified Bessel functions of order $`1/4`$, $`I_{1/4}(\nu )`$, and $`K_{1/4}(\nu )`$. Transforming to the old variables, we find the growing solution $`\widehat{a}^{1/2}I_{1/4}(\widehat{a}^2/2)`$ and the decreasing solution $`\widehat{a}^{1/2}K_{1/4}(\widehat{a}^2/2)`$ in the $`\widehat{a}`$ direction. To select one of them we will impose a matching condition with the solution (18) for $`\widehat{a}\widehat{a}_0`$. This gives the decreasing solution
$$\mathrm{\Psi }(\widehat{a})=\widehat{a}^{1/2}K_{1/4}\left(\frac{\widehat{a}^2}{2}\right),$$
(14)
which is the well-known solution of nothing due to Vilenkin , and goes like $`e^{\frac{\widehat{a}^2}{2}}`$ for $`\widehat{a}0`$. It was obtained by Vilenkin in the limit of small $`\widehat{a}`$ in the 4-dimensional model with topology $`R\times S^3`$ and inflation, without a dilaton field. Usually, in 4-dimensional cosmology nothing is the nonsingular boundary of the superspace that includes three-geometries given through a slicing of a regular four-geometry , . In higher-dimensional cosmologies, the extra dimensions usually play the role of a scalar field $`\sigma `$ in the equivalent four-dimensional model, such that the non-singular boundary of the minisuperspace is the configuration $`\widehat{a}=0,|\sigma |<\mathrm{}`$. This configuration is called external nothing since the extra dimension is assumed to be nonzero .
Now, consider the common Wheeler-DeWitt equation in 4D
$$\left[\frac{^2}{a^2}+\frac{1}{a^2}\frac{^2}{\sigma ^2}+a^2a^4U(\sigma )\right]\mathrm{\Psi }(a,\sigma )=0.$$
(15)
The WKB solution of Eq.(15) with Vilenkin boundary condition is well known in the region of the minisuperspace where the potential $`U(\sigma )`$ is sufficiently flat, i.e. where
$$\left|\frac{U^{}}{U}\right|max\{U(\sigma ),a^2\}.$$
(16)
In fact, by assuming the condition (16) the wave function becomes a slowly-varying function of $`\sigma `$. Therefore, one can neglect the derivative with respect to the dilaton field $`\sigma `$. Thus, $`\sigma `$ plays the role of a parameter in the Wheeler-DeWitt equation (15), and the problem is reduced to the one-dimensional minisuperspace model.
In the present model, however, the potential $`U(\sigma )=ϵF(l)e^{3b\sigma }`$ has a strongly asymmetric form for $`\sigma <0`$ and $`\sigma >0`$, and so the condition (16) does not hold in the region of the minisuperspace where $`\widehat{a}^2U(\sigma )>1`$ and $`\sigma 0`$ with $`b>0`$. Nevertheless, we may use an approximation to cast the Wheeler-DeWitt equation (8) into a solvable equation. The main point in the non-compactified five-dimensional model is to find the most probable $`l=\text{const}`$ hypersurface for the 4D universe to tunnel from nothing. On the other hand, we are interested in the tunneling in the direction of $`\widehat{a}`$ which is the only $`l`$-dependent dynamical variable in the model. Therefore, any result about the most probable $`l=\text{const}`$ hypersurface will be obtained due to the $`l`$-dependence of $`\widehat{a}`$, and is expected not to be affected by the $`\sigma `$ dependence of $`\mathrm{\Psi }`$ since $`\sigma `$ is independent of $`l`$. So, we may solve the Wheeler-DeWitt equation in that region of the minisuperspace where the condition (16) holds and deduce the result about the most probable $`l=`$ const hypersurface for tunneling. This result is anticipated not to be changed if we solve the Wheeler-DeWitt equation in the whole minisuperspace. In other words, to the extent that we are concerned about finding the most probable $`l=\text{const}`$ hypersurface ( and the dependence of $`\mathrm{\Psi }`$ on the $`\sigma `$ field is not important to us ) we may suppose the wave function $`\mathrm{\Psi }`$ to be independent of $`\sigma `$. This assumption is valid at least in the positive sector $`\sigma 0`$ with $`b>0`$ where the potential $`U(\sigma )`$ is sufficiently flat. Therefore, we may drop the $`\sigma `$ derivative and the Wheeler-DeWitt equation (8) takes the form
$$\left[\frac{d}{d\widehat{a}^2}+\widehat{a}^2[1U(\sigma )\widehat{a}^2]\right]\mathrm{\Psi }(\widehat{a})=0,$$
(17)
where $`\sigma `$ is just a parameter and the problem is essentially identical to the one-dimensional minisuperspace model with the dynamical variable $`\widehat{a}`$. In this approximation, the superpotential barrier becomes wide for $`\sigma 0`$ which, at first glance, makes it inconvenient for tunneling. However, due to the function $`F(l)`$ we will show that the barrier becomes sufficiently narrow for quantum tunneling, even in the $`\sigma 0`$ sector of the minisuperspace. The solutions of equation (17) are the well-known Vilenkin tunneling wave functions :
$$\mathrm{\Psi }_T=\mathrm{exp}\left(\frac{1[1U(\sigma )\widehat{a}^2]^{\frac{3}{2}}}{3U(\sigma )}\right),\widehat{a}^2U(\sigma )<1,$$
(18)
and
$$\mathrm{\Psi }_T=\mathrm{exp}\left(\frac{1}{3U(\sigma )}\right)\mathrm{exp}\left(\frac{i[U(\sigma )\widehat{a}^21]^{\frac{3}{2}}}{3U(\sigma )}\right),\widehat{a}^2U(\sigma )>1,$$
(19)
where $`\mathrm{\Psi }`$ for the region $`\widehat{a}^2U(\sigma )<1`$ (underbarrier) has the regular behaviour $`\mathrm{\Psi }e^{\frac{\widehat{a}^2}{2}}`$ for $`\widehat{a}0`$ matching with the nothing solution (14). We notice that the potential term
$$U(\sigma )=ϵF(l)e^{3b\sigma }$$
plays the role of an effective 4D cosmological term whose properties merit attention. First, the presence of $`ϵ=\pm 1`$ corresponds to positive or negative cosmological term for a given $`F(l)`$. But as was discussed before, only the $`ϵ=+1`$ case is relevant to quantum cosmology with $`k=+1`$, since for $`ϵ=1`$ there is no maximum (barrier) for the superpotential. Second, the presence of the $`l`$-dependent term $`F(l)`$ indicates the contributions of the fifth dimension to the effective 4D cosmological term. Third, for $`b1`$, if the parameter $`\sigma `$ undergoes time evolution from negative to positive values, then a large cosmological term would become a small one very quickly. This is important, as it is in agreement with quantum tunneling and inflationary ideas where the very early universe might have experienced an extremely inflationary era for $`\sigma <0`$ which was abruptly switched off for $`\sigma >0`$. This would also be an alternative solution to the well-known cosmological constant problem in that an initially large cosmological term becomes small after the inflationary period.
Considering the function $`F(l)`$, and (4) and (11), we find that a non-zero $`F`$ may be achieved by taking linear behaviours for the functions $`\chi `$ or $`f`$ with respect to $`l`$
$$f(l)=\chi (l)=\frac{l}{L},$$
(20)
where a constant $`L`$ is introduced to preserve physical dimensions. The choice (20) corresponds to the so-called canonical metric , and gives
$$F(l)=\frac{2}{l^2}.$$
(21)
The cosmological term is then obtained as
$$\mathrm{\Lambda }U(\sigma )=\frac{2}{l^2}e^{3b\sigma },$$
(22)
with the right dimension of (length)<sup>-2</sup>.
One may obtain the probability distribution for the Vilenkin wave function as in :
$$\rho _T\mathrm{exp}\left[\frac{2}{3U(\sigma )}\right]\mathrm{exp}\left[\frac{l^2e^{3b\sigma }}{3}\right].$$
(23)
This is maximized for $`\sigma 0`$ when $`l0.`$ This condition shows that the 4+1 dimensional universe could have tunneled with a large probability if the fifth dimension was very small, and that a small 4D universe $`\widehat{a}_01`$ was born on a 4D hypersurface near to the $`l0`$ hypersurface<sup>1</sup><sup>1</sup>1We recall that although we have taken large values of $`\sigma `$ in this approximation to the Wheeler-DeWitt equation (8), the condition $`l0`$ makes it possible to narrow down the superpotential barrier ($`U(\sigma )1`$) and tunnel from nothing to a small size universe $`\widehat{a}_0=\frac{1}{\sqrt{U(\sigma )}}`$. .
Let us now consider the Hartle-Hawking wave function for this model. In the presence of matter fields $`\varphi `$ the Hartle-Hawking wave function of the universe is obtained through the functional integral over Euclidean 4-metrics $`g_{\alpha \beta }`$ ($`\alpha ,\beta =0,1,2,3`$):
$$\mathrm{\Psi }(\stackrel{~}{h}_{ij},\stackrel{~}{\varphi })=_𝒞𝑑g_{\alpha \beta }𝑑\varphi \mathrm{exp}[I(g_{\alpha \beta },\varphi )].$$
(24)
Here the domain $`𝒞`$ is defined by the “ no boundary” proposal as all regular compact Euclidean 4-geometries ( the boundary of which is $`S^3`$ with the induced 3-metric $`\stackrel{~}{h}_{i,j}`$ ($`i,j=1,2,3`$) ) and the regular matter field configurations ( the value of which is $`\stackrel{~}{\varphi }`$ on the 3-manifold). However, in our model there are no extra matter fields, other than $`\sigma `$ which appears as the dilaton field in the 5-metric. In this way, the corresponding Hartle-Hawking wave function is given as
$$\mathrm{\Psi }(\stackrel{~}{\widehat{h}}_{\alpha \beta })=_𝒞𝑑\widehat{g}_{AB}\mathrm{exp}[I(\widehat{g}_{AB})].$$
(25)
Here $`\widehat{g}_{AB}(A,B=0,1,2,3,4)`$ is the 5-metric, and the domain $`𝒞`$ is the class of all regular compact Euclidean 5-geometries whose boundary is $`S^3\times R`$ ($`R`$ denotes the fifth non-compact coordinate), with the induced 4-metric $`\stackrel{~}{\widehat{h}}_{\alpha \beta }(\alpha ,\beta =1,2,3,4)`$. On the other hand, the present 5D model is effectively equivalent to a 4D non-minimally coupled dilaton field on $`l=`$const hypersurfaces. Therefore, using the 5D metric (1), the above integral may be rewritten in its familiar 4D form in the gauge $`\dot{\widehat{N}}=0`$ as:
$$\mathrm{\Psi }_H(\stackrel{~}{\widehat{a}},\stackrel{~}{\sigma })|_{l=\text{const}}=𝑑\widehat{N}𝒟\widehat{a}𝒟\sigma \mathrm{exp}(I[\widehat{a}(\tau ),\sigma (\tau ),\widehat{N}]).$$
(26)
Here $`I`$ is the Euclidean action for the model, $`\tau `$ is the Euclidean time and $`\stackrel{~}{\widehat{a}}`$ and $`\stackrel{~}{\sigma }`$ are the final values of $`\widehat{a}`$ and $`\sigma `$ on the 3-geometry $`\stackrel{~}{\widehat{h}}_{i,j}`$ and $`l=`$const hypersurface. However, in practice we are usually interested in the semi-classical approximation to the above path integral
$$\mathrm{\Psi }_H(\stackrel{~}{\widehat{a}},\stackrel{~}{\sigma })|_{l=\text{const}}\mathrm{exp}(I_{cl}(\stackrel{~}{\widehat{a}},\stackrel{~}{\sigma })),$$
(27)
where $`I_{cl}(\stackrel{~}{\widehat{a}},\stackrel{~}{\sigma })`$ is the action for the instanton solutions to the Euclidean field equations. The Hartle-Hawking wave function for the model defined by the Lagrangian (3) is well-known and may be obtained by the following transformation :
$$\mathrm{\Psi }_H=\mathrm{\Psi }_T(Ue^{i\pi }U,\widehat{a}e^{i\pi /2}\widehat{a}).$$
This yields
$$\mathrm{\Psi }_H|_{l=\text{const}}\mathrm{exp}\left(\frac{1[1U(\sigma )\widehat{a}^2]^{\frac{3}{2}}}{3U(\sigma )}\right),\widehat{a}^2U(\sigma )<1,$$
(28)
$$\mathrm{\Psi }_H|_{l=\text{const}}\mathrm{exp}\left(\frac{1}{3U(\sigma )}\right)\mathrm{cos}\left(\frac{[U(\sigma )\widehat{a}^21]^{\frac{3}{2}}}{3U(\sigma )}\frac{\pi }{4}\right),\widehat{a}^2U(\sigma )>1.$$
(29)
Now, the probability distribution for the Hartle-Hawking wave function as given in is
$$\rho _H\mathrm{exp}\left[\frac{2}{3U(\sigma )}\right]\mathrm{exp}\left[\frac{l^2e^{3b\sigma }}{3}\right].$$
(30)
Contrary to Vilenkin’s case, this probability distribution is maximized for $`l0`$ .
This condition, on the other hand, indicates that a Lorentzian universe was born from a mother Euclidean universe with a large probability when the 4D Lorentzian hypersurface was very far from the $`l0`$ hypersurface.
## 4 Discussion
Usually, in quantum cosmology there is a debate on the choice between Hartle-Hawking and Vilenkin wave functions in concern to the issue of inflation. It is commonly believed that the Vilenkin wave function leads to inflation. However, for some particular models the Hartle-Hawking wave function claims to predict a period of inflation -. So, as far as inflation is concerned, there is no clear way to decide between the two proposals.
In this paper, we have introduced a new way to compare the two proposals in a higher-dimensional Kaluza-Klein model by using the extra dimension. We have found that if there is a non-compactified extra dimension, the Vilenkin proposal seems more reasonable than the Hartle-Hawking proposal. In the Hartle-Hawking proposal, a large value of $`l`$ gives rise to a rather large initial radius of the universe, which seems unnatural. That is, in the Hartle-Hawking proposal there is no good justification for a big 4D universe which was born on a constant hypersurface $`l0`$. However, in the Vilenkin proposal the universe starts naturally from a small radius $`\widehat{a}_0`$ with a large cosmological term $`\mathrm{\Lambda }l^2e^{3b\sigma }`$ ($`0\sigma <\mathrm{}`$) on the $`l0`$ hypersurface. Although we have solved the problem in the approximation of large positive values of $`\sigma `$, the result $`l0`$ hypersurface is expected to remain unchanged for the whole domain $`\mathrm{}<\sigma <\mathrm{}`$ since $`\sigma `$ is independent of $`l`$. Therefore, in the full theory the time evolution $`\sigma (t)<0\sigma (t)>0`$ will switch off the large cosmological term. This may be considered as a solution to the cosmological constant problem.
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# 1 Introduction
## 1 Introduction
The $`q`$-state Potts model has served as a valuable model for the study of phase transitions and critical phenomena . On a lattice, or, more generally, on a (connected) graph $`G`$, at temperature $`T`$, this model is defined by the partition function
$$Z(G,q,v)=\underset{\{\sigma _n\}}{}e^\beta $$
(1.1)
with the (zero-field) Hamiltonian
$$=J\underset{ij}{}\delta _{\sigma _i\sigma _j}$$
(1.2)
where $`\sigma _i=1,\mathrm{},q`$ are the spin variables on each vertex $`iG`$; $`\beta =(k_BT)^1`$; and $`ij`$ denotes pairs of adjacent vertices. The graph $`G=G(V,E)`$ is defined by its vertex set $`V`$ and its edge set $`E`$; we denote the number of vertices of $`G`$ as $`n=n(G)=|V|`$ and the number of edges of $`G`$ as $`e(G)=|E|`$. We use the notation
$$K=\beta J,a=u^1=e^K,v=a1$$
(1.3)
so that the physical ranges are (i) $`a1`$, i.e., $`v0`$ corresponding to $`\mathrm{}T0`$ for the Potts ferromagnet, and (ii) $`0a1`$, i.e., $`1v0`$, corresponding to $`0T\mathrm{}`$ for the Potts antiferromagnet. One defines the (reduced) free energy per site $`f=\beta F`$, where $`F`$ is the actual free energy, via
$$f(\{G\},q,v)=\underset{n\mathrm{}}{lim}\mathrm{ln}[Z(G,q,v)^{1/n}].$$
(1.4)
where we use the symbol $`\{G\}`$ to denote $`lim_n\mathrm{}G`$ for a given family of graphs.
Let $`G^{}=(V,E^{})`$ be a spanning subgraph of $`G`$, i.e. a subgraph having the same vertex set $`V`$ and an edge set $`E^{}E`$. Then $`Z(G,q,v)`$ can be written as the sum -
$`Z(G,q,v)`$ $`=`$ $`{\displaystyle \underset{G^{}G}{}}q^{k(G^{})}v^{e(G^{})}`$ (1.5)
$`=`$ $`{\displaystyle \underset{r=k(G)}{\overset{n(G)}{}}}{\displaystyle \underset{s=0}{\overset{e(G)}{}}}z_{rs}q^rv^s`$ (1.8)
where $`k(G^{})`$ denotes the number of connected components of $`G^{}`$ and $`z_{rs}0`$. Since we only consider connected graphs $`G`$, we have $`k(G)=1`$. The formula (1.5) shows that $`Z(G,q,v)`$ is a polynomial in $`q`$ and $`v`$ and enables one to generalize $`q`$ from $`_+`$ to $`_+`$ and, indeed, to $``$. The Potts model partition function $`Z(G,q,v)`$ on a graph $`G`$ is essentially equivalent to the Tutte polynomial - $`T(G,x,y)`$ and Whitney rank polynomial $`R(G,\xi ,\eta )`$ ,, for this graph, as discussed in the appendix.
One special case of the Potts model partition function that is of particular interest is the zero-temperature limit of the Potts antiferromagnet (AF). For sufficiently large $`q`$, on a given lattice or graph $`G`$, this exhibits nonzero ground state entropy (without frustration). This is equivalent to a ground state degeneracy per site (vertex), $`W>1`$, since $`S_0=k_B\mathrm{ln}W`$. The $`T=0`$ (i.e., $`v=1`$) partition function of the above-mentioned $`q`$-state Potts antiferromagnet (PAF) on $`G`$ satisfies
$$Z(G,q,1)=P(G,q)$$
(1.9)
where $`P(G,q)`$ is the chromatic polynomial (in $`q`$) expressing the number of ways of coloring the vertices of the graph $`G`$ with $`q`$ colors such that no two adjacent vertices have the same color . This is termed a proper vertex coloring of $`G`$. The minimum number of colors necessary for this coloring is the chromatic number of $`G`$, denoted $`\chi (G)`$. Thus<sup>1</sup><sup>1</sup>1At certain special points $`q_s`$ (typically $`q_s=0,1,..,\chi (G)`$), one has the noncommutativity of limits $`lim_{qq_s}lim_n\mathrm{}P(G,q)^{1/n}lim_n\mathrm{}lim_{qq_s}P(G,q)^{1/n}`$, and hence it is necessary to specify the order of the limits in the definition of $`W(\{G\},q_s)`$ .
$$W(\{G\},q)=\underset{n\mathrm{}}{lim}P(G,q)^{1/n}.$$
(1.10)
A second special case of the Potts model partition function is for infinite temperature, i.e., $`K=0`$ or equivalently $`v=0`$. In this case, as is clear from either (1.1) with (1.2) or from (1.5), $`Z`$ reduces to the single term
$$Z(G,q,0)=q^{n(G)}.$$
(1.11)
A third special case is the Potts ferromagnet in the limit of zero temperature, $`T0`$, i.e., $`v\mathrm{}`$.
Here we shall consider families of graphs that are cyclic strips of the square and triangular lattices, taken to be oriented horizontally, with free transverse boundary conditions, denoted $`FBC_y`$, and periodic longitudinal boundary conditions, denoted $`PBC_x`$. A given strip of the triangular lattice may be visualized by starting with the corresponding strip of the square lattice and adding diagonal bonds joining, say, the upper left to lower right vertices of each square. The length and width are taken to be $`L_x`$ and $`L_y`$ vertices, and the total number of vertices is $`n=L_xL_y`$. We shall also make some comments about the analogous strips with twisted periodic longitudinal, i.e., Möbius, boundary conditions, denoted $`TPBC_x`$ or $`Mb`$. These various families of lattice strip graphs are examples of recursive families of graphs, in the sense that a strip of length $`L_x+1`$ is constructed by the addition of a given subgraph (here, a transverse layer of the strip) to the strip of length $`L_x`$. As derived in , using a type of transfer matrix argument, a general form for the Potts model partition function for the strip graphs considered here, or more generally, for a recursively defined graph comprised of $`L_x=m`$ repeated subunits (transverse layers of the strip here), is
$$Z(G,q,v)=\underset{j=1}{\overset{N_{Z,G,\lambda }}{}}c_{Z,G,j}(\lambda _{Z,G,j})^m$$
(1.12)
where $`\lambda _{Z,G,j}`$ is a function of $`q`$ and $`v`$, and both $`\lambda _{Z,G,j}`$ and the coefficients $`c_{Z,G,j}`$ are independent of $`L_x`$. The coefficients $`c_{Z,G,j}`$ are functions only of $`q`$ and, indeed, are polynomials, for the cyclic strip graphs of interest here. (We shall comment on other types of strip graphs below.)
From (1.9), it also follows that the chromatic polynomial has the same structure,
$$P(G,q)=\underset{j=1}{\overset{N_{P,G,\lambda }}{}}c_{P,G,j}(\lambda _{P,G,j})^m$$
(1.13)
where $`N_{P,G,\lambda }`$ depends on $`G`$; $`c_{P,G,j}`$ and $`\lambda _{P,G,j}`$ are independent of $`L_x=m`$; and
$$N_{P,G,\lambda }N_{Z,G,\lambda }.$$
(1.14)
As will be shown below, for the strips under consideration here, this inequality is realized as an equality for the circuit graph $`C_n`$ and as a strict inequality for widths $`L_y2`$ (see Tables 1 and 3). For the same reason, it follows that when one sets $`v=1`$ in $`Z(G,q,v)`$, i.e. specializes to the $`T=0`$ Potts antiferromagnet, a certain number $`N_{Z,G,\lambda }N_{P,G,\lambda }`$ of the $`\lambda _{Z,G,j}`$’s vanish and the remaining, nonvanishing $`\lambda _{Z,G,j}(q,v)`$’s are precisely those occurring in $`P(G,q)`$:
$$\mathrm{If}\lambda _{Z,G,j}(q,v=1)0\mathrm{then}\lambda _{Z,G,j}(q,1)=\lambda _{P,G,j}(q).$$
(1.15)
Previous calculations of chromatic polynomials for recursive families of graphs of arbitrary length include -; we shall also use the results of calculations of Potts model partition functions (Tutte polynomials) for cyclic and Möbius strip graphs of arbitrary length in . Some relevant results on transfer matrices and the Temperley-Lieb algebra are in , .
Several basic questions about the structure of $`Z(G,q,v)`$ and $`P(G,q)`$ for these cyclic strip graphs are the following:
1. Can one obtain a general formula for the coefficients $`c_{Z,G,j}`$ and $`c_{P,G,j}`$ in (1.12) and (1.13)?
2. What are the respective sums of the coefficients $`c_{Z,G,j}`$ and $`c_{P,G,j}`$ in (1.12) and (1.13)?
3. For a given cyclic strip $`G`$, how many $`\lambda _{Z,G,j}`$’s in (1.12) have a particular coefficient, and how many $`\lambda _{P,G,j}`$’s in (1.13) have this coefficient?
4. What are the total numbers of different terms $`N_{Z,G,\lambda }`$ and $`N_{P,G,\lambda }`$ in (1.12) and (1.13)?
In this paper, we shall obtain a general formula for the coefficients $`c_{Z,G,j}`$ and $`c_{P,G,j}`$, an answer to the second pair of questions, and, based on these results, answers to the remaining two pairs of questions. We shall also comment on other types of lattice strips.
The chromatic polynomial, can be calculated by various methods, including iterative application of the deletion-contraction theorem (e.g., ), a certain matrix technique , a generating function method , or a coloring (compatibility) matrix method . Similarly, the Potts model partition function or equivalent Tutte polynomial can be calculated by iterative application of the generalized deletion-contraction theorem or by transfer matrix methods including a relation with the Temperley-Lieb algebra . Of course these different methods can be used to provide cross-checks; for example, iterated deletion-contraction and transfer matrix methods were used to check each other in and deletion-contraction and coloring matrix methods were used together in works such as . The last, the coloring matrix method, is also useful for deriving rigorous upper and lower bounds on $`W(\{G\},q)`$ ,,. In this method, one first selects a transverse slice of the strip, denoted $`_{L_y}`$. For our cyclic graphs, this is simply a line graph with $`L_y`$ vertices (oriented vertically, given that we orient the long direction of the strip horizontally). Denote an allowed $`q`$-coloring of this path as $`c(_{L_y})`$. The number of allowed colorings of the path $`_{L_y}`$ is $`𝒩=P(_{L_y},q)`$. For our chromatic polynomials of cyclic strips of the square and triangular lattices, $`P(_n,q)=P(T_{L_y},q)`$, where $`T_n`$ is the tree graph with $`n`$ vertices and
$$P(T_n,q)=q(q1)^{n1}.$$
(1.16)
Now focus on two adjacent paths $`_{L_y}`$ and $`_{L_y}^{}`$. Define compatible $`q`$-colorings of these paths as colorings such that no two adjacent vertices $`v_{L_y}`$ and $`v^{}_{L_y}^{}`$ (i.e. vertices connected by an edge = bond of the lattice strip graph) have the same color. One can then associate with this pair of paths an $`𝒩\times 𝒩`$ dimensional symmetric matrix $`𝒯`$ with entries $`𝒯_{c(_{L_y}),c(_{L_y}^{})}=1`$ or 0 if the $`q`$-colorings of $`_{L_y}`$ and $`_{L_y}^{}`$ are or are not compatible, respectively. Then the chromatic polynomial of the cyclic strip of the lattice $`\mathrm{\Lambda }`$, taken here to be square ($`sq`$) or triangular ($`tri`$), is (with $`L_x=m`$) given by
$$P(\mathrm{\Lambda },L_y\times L_x,FBC_y,PBC_x,q)=Tr(𝒯^m).$$
(1.17)
Since $`𝒯`$ is a symmetric real matrix (indeed, composed only of 0’s and 1’s, although we do not need this here), it can be diagonalized by an orthogonal transformation, so that the above trace is $`Tr(𝒯^m)=_j(\lambda _{P,G,j})^m`$, where the $`\lambda _{P,G,j}`$’s are the eigenvalues of $`𝒯`$. In the context of the $`q`$-coloring problem, if one denotes the multiplicity of the $`j`$’th distinct eigenvalue by $`c_{P,G,j}`$, one obtains the formula (1.13). That is, from the coloring matrix viewpoint, $`c_{P,G,j}`$ is the dimension of the invariant subspace in the full $`𝒩`$-dimensional space of coloring configurations of the transverse slice of the strip corresponding to the eigenvalue $`\lambda _{P,G,j}`$. For the full temperature-dependent Potts model partition function, one can define an analogous matrix, $`𝒯_Z`$ whose entries, rather than being 0 and 1, are appropriate Boltzmann weights for the given spin configurations on the successive transverse slices $`_{L_y}`$ and $`_{L_y}^{}`$. Then this partition function is (with $`L_x=m`$)
$$Z(\mathrm{\Lambda },L_y\times L_x,FBC_y,PBC_x,q,v)=Tr[(𝒯_Z)^m]$$
(1.18)
which was used in to derive the formula (1.12). Note that in this case all colorings, including those that yield the same color on adjacent vertices, are allowed. One can obtain these formulas (1.17) and (1.18) for sufficiently large integral $`q`$ that the multiplicities $`c_{P,G,j}`$ and $`c_{Z,G,j}`$ are positive integers; in both of these cases, the fact that the coefficients are multiplicities shows that for sufficiently large positive integer $`q`$ they are also positive integers. In the case of the chromatic polynomial, it is obvious that $`c_{P,G,j}`$ can depend only on $`q`$; in the case of the full Potts model partition function, for the cyclic strip graphs considered here, this also follows from (1.18) and (1.12), since the multiplicities of the eigenvalues cannot depend on the variable parameter $`v(1,\mathrm{})`$.
Note that, while the coefficients $`c_{P,G,j}`$ and $`c_{Z,G,j}`$ can be obtained as multiplicities of the distinct eigenvalues $`\lambda _{P,G,j}`$ and $`\lambda _{Z,G,j}`$ for sufficiently large integer $`q`$, when one considers positive $`q<4`$, they may be zero or negative (see eqs. (2.31)-(2.35) below). More generally, while they play the role of eigenvalue multiplicities for sufficiently large integer $`q`$, the domain of their definition may be generalized via (1.5) to $`q_+`$ or, indeed, to $`q`$. Some properties of determinants of coloring matrices for various strip graphs are given in the appendix.
The dimension of the space of coloring configurations, $`𝒩`$, is equal to the sum of the multiplicities of each distinct eigenvalue, i.e., the sum of the dimensions of the invariant subspaces corresponding to each of these distinct eigenvalues. For the chromatic polynomial, this is $`𝒩`$, which is equal to the sum
$$C_{P,G}=\underset{j=1}{\overset{N_{P,G,\lambda }}{}}c_{P,G,j}$$
(1.19)
while for the full Potts model partition function, we shall denote it as
$$C_{Z,G}=\underset{j=1}{\overset{N_{Z,G,\lambda }}{}}c_{Z,G,j}.$$
(1.20)
General results for these sums will be given below for the strips of interest.
Before proceeding, we note that for the chromatic polynomial, since two identical color configurations on $`_{L_y}`$ and $`_{}^{}{}_{L_y}{}^{}`$ are incompatible, the diagonal elements of the coloring matrix are zero and hence its trace is zero:
$$Tr(𝒯(G,q))=0\mathrm{for}P(G,q).$$
(1.21)
Since the coefficients $`c_{P,G,j}`$ are just the multiplicities of the given eigenvalues $`\lambda _{P,G,j}`$, we can write this as the sum over distinct $`\lambda _{P,G,j}`$’s:
$$\underset{j=1}{\overset{N_{P,G,\lambda }}{}}c_{P,G,j}\lambda _{P,G,j}=0.$$
(1.22)
In the following, since we shall present results that are applicable to cyclic strips of both the square and triangular lattices, we shall often leave the lattice type $`\mathrm{\Lambda }`$ and the boundary conditions $`(FBC_y,PBC_x)`$ implicit in the notation where they are obvious. The fact that these results will be applicable to cyclic strips of both the square and triangular lattices depends on the property that (having expressed the strip of the triangular lattice as a strip of the square lattice with diagonal bonds added to each square, as described above) the transverse slices, i.e., line graphs of $`L_y`$ vertices, of both the cyclic square-lattice and triangular-lattice strips are identical.<sup>2</sup><sup>2</sup>2Although the transverse slices are identical for the cyclic strips of the square and triangular lattice, they are different for other lattices, such as the honeycomb = brick lattice, and this difference was also significant for the earlier use of coloring matrix methods to obtain rigorous bounds on the Potts model ground state degeneracy per site, $`W`$ -. Similarly, for quantities that are independent of some part of $`G`$, such as $`c_{Z,G,j}`$ and $`\lambda _{Z,G,j}`$, which are independent of $`L_x`$, this will be incorporated in the notation.
## 2 Coefficients in Potts Model Partition Functions for Cyclic Lattice Strips
In this section we address the first pair of questions posed in the Introduction. From our exact solutions of the chromatic polynomials and, more generally, of the full temperature-dependent Potts model partition functions on cyclic strips of the square and triangular lattice, we have found that the coefficients $`c_{Z,G,j}`$ and $`c_{P,G,j}`$ are polynomials in $`q`$. Further, we have found that for a given cyclic strip of the square or triangular lattice with width $`L_y`$, the coefficients $`c_{P,G,j}`$ and $`c_{Z,G,j}`$ are of a limited set; only one type of polynomial of each degree in $`q`$ occurs<sup>3</sup><sup>3</sup>3From the studies of Möbius strips of the square lattice , we find that the coefficients are again polynomials, of the form $`\pm c^{(d)}`$; however, for Möbius strips of the triangular lattice, the coefficients are not, in general, polynomials, but algebraic functions of $`q`$, as was shown by the exact solutions for the chromatic polynomial in and for the full Potts model partition in for the $`L_y=2`$ Möbius strip. For strips with torus boundary conditions, exact solutions show that the coefficients are polynomials of degree $`d`$ in $`q`$, but it is not, in general, the case that there is a unique coefficient of degree $`d`$ and they are not, in general, of the form $`c^{(d)}`$. We shall comment further on this below., and, denoting this as $`c^{(d)}`$, there are coefficients of the form $`c^{(d)}`$ with $`0dL_y`$. We infer the following general formula for $`c^{(d)}`$, i.e. for the multiplicity of the corresponding eigenvalue of $`𝒯_Z`$:
$$c^{(d)}=U_{2d}\left(\frac{\sqrt{q}}{2}\right)$$
(2.1)
where $`U_n(x)`$ is the Chebyshev polynomial of the second kind, defined by (e.g. )
$`U_n(x)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1x^2}}}Im[(x+i\sqrt{1x^2})^{n+1}]`$ (2.2)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{[\frac{n}{2}]}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{j}}\right)(2x)^{n2j}`$ (2.4)
where in eq. (2.4) and similar equations below, the notation $`[\frac{n}{2}]`$ in the upper limit on the summand means the integral part of $`\frac{n}{2}`$. In the first line of (2.4), $`x`$ must be real, but the polynomial in the second line of (2.4) constitutes a definition of $`U_n(x)`$ for complex $`x`$. The first few of these coefficients are, with $`x=\frac{\sqrt{q}}{2}`$,
$$c^{(0)}=U_0(x)=1$$
(2.5)
$`c^{(1)}`$ $`=`$ $`U_2(x)=4x^21`$ (2.6)
$`=`$ $`q1`$ (2.8)
$`c^{(2)}`$ $`=`$ $`U_4(x)=16x^412x^2+1`$ (2.9)
$`=`$ $`q^23q+1`$ (2.11)
$`c^{(3)}`$ $`=`$ $`U_6(x)=64x^680x^4+24x^21`$ (2.12)
$`=`$ $`q^35q^2+6q1`$ (2.14)
$`c^{(4)}`$ $`=`$ $`U_8(x)=256x^8448x^6+240x^440x^2+1`$ (2.15)
$`=`$ $`q^47q^3+15q^210q+1`$ (2.17)
$`=`$ $`(q1)(q^36q^2+9q1)`$ (2.19)
and so forth for higher $`d`$. We have found, as an equivalent formula
$$c^{(d)}=\underset{k=1}{\overset{d}{}}(qq_{d,k})$$
(2.20)
where
$$q_{d,k}2+2\mathrm{cos}\left(\frac{2\pi k}{2d+1}\right)=4\mathrm{cos}^2\left(\frac{\pi k}{2d+1}\right),\mathrm{for}k=1,2,\mathrm{}d.$$
(2.21)
Note that the apparent square root in (2.1) is absent in the actual formula (2.1) for the $`c^{(d)}`$’s, which applies for arbitrary complex $`q`$. One way to derive (2.1) for cyclic strips is to use the fact that the partition function is given by the trace (1.18), which, with (1.12) makes clear the role of the coefficients $`c_{Z,G,j}`$ as multiplicities of eigenvalues of the transfer matrix $`𝒯_Z`$; then, one utilizes the connection with the Temperley-Lieb algebra to infer that these multiplicities are given by $`c^{(d)}`$ (for sufficiently large integer $`q`$ where the $`c^{(d)}`$ are positive integers), and finally, one continues this result to general complex $`q`$. Related discussions of dimensions of invariant subspaces of operators in the Temperley-Lieb algebra are given in ,. The form (2.20) is also related to the property that each coefficient $`c^{(d)}`$ has a simple zero at the Tutte-Beraha number $`B_{2d+1}`$; indeed, one could derive the form (2.1) by using this as a starting point. In this approach, one would start with the property that $`c^{(d)}`$ has the factor $`(qB_{2d+1})=(qq_{d,1})`$. Since $`B_{2d+1}`$ is, in general, irrational, it is necessary to symmetrize the product $`c^{(d)}=(qq_z)`$ over the zeros $`q_z`$ in order to obtain rational (indeed, integer) coefficients for the $`c^{(d)}`$. This symmetrization yields (2.20). To see this, we note that to incorporate this zero of $`c^{(d)}`$ at $`q=B_{2d+1}`$ one can start with the identity $`[\mathrm{exp}(2\pi i/(2d+1))]^{2d+1}1=0`$, express this as $`[\mathrm{cos}(\frac{2\pi }{2d+1})+i\mathrm{sin}(\frac{2\pi }{2d+1})]^{2d+1}`$, expand the resulting expression, and then impose the condition that the polynomial vanishes at $`q=q_{d,1}=B_{2d+1}`$ by setting $`\mathrm{cos}(2\pi /(2d+1))=q/21`$. This actually yields two equations, one for the real part and one for the imaginary part. Concentrating on the real part, we observe that the resultant equation has the form $`(1/2)(q4)(c^{(d)})^2`$. This yields the result
$$c^{(d)}=\left[\frac{2}{q4}\left(T_{2d+1}(\frac{q}{2}1)1\right)\right]^{1/2}$$
(2.22)
where $`T_n(x)`$ is the Chebyshev polynomial of the first kind, defined by (e.g., )
$`T_n(x)`$ $`=`$ $`Re\left[(x+i\sqrt{1x^2})^n\right]`$ (2.23)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=0}{\overset{[\frac{n}{2}]}{}}}(1)^j{\displaystyle \frac{n}{nj}}\left({\displaystyle \genfrac{}{}{0pt}{}{nj}{j}}\right)(2x)^{n2j}`$ (2.25)
As with $`U_n(x)`$, in the first line of (2.25), $`x`$ must be real, but the polynomial in the second line of (2.25) constitutes a definition of $`T_n(x)`$ for complex $`x`$; note also that the right-hand side of the second line is defined to be 1 for $`n=0`$. Simplifying eq. (2.22), one finally obtains eq. (2.1) in terms of the Chebyshev polynomial of the second kind. The starting point is equivalent to the identity $`[\mathrm{exp}(2\pi k/(2d+1))]^{2d+1}1=0`$, and hence it follows that $`c^{(d)}`$ vanishes at $`q_{d,k}`$ for $`2kd`$ as well as for $`k=1`$. This implies the formula (2.20).
The general formula (2.1) or the equivalent result (2.20) gives a deeper understanding of, the known exact solutions for chromatic polynomials and the full Potts model partition functions on cyclic strips of the square and triangular lattice. For reference, these include the degenerate case $`L_y=1`$ (circuit graph), and the cases $`L_y=2`$ in , $`L_y=3`$ in , and $`L_y=4`$ in for cyclic strips of the square lattice; and the cases $`L_y=2`$ in and $`L_y=3`$ and $`L_y=4`$ in for cyclic strips of the triangular lattice. For the full Potts model partition functions, the exact solutions include the elementary $`L_y=1`$ solution for the circuit graph, and the $`L_y=2`$ solutions in and for cyclic strips of the square and triangular lattices, respectively. Explicit examples will be given below. As noted, the coefficients with degree $`d`$ occur for $`0dL_y`$. In the following we shall work out implications of the formula (2.1) for the structure of the chromatic polynomial and full Potts model partition function for arbitrarily wide cyclic strips of the square and triangular lattices.
From the formula (2.1) or (2.20) and properties of the Chebyshev polynomials of the second kind, one can establish a number of useful properties of the coefficients $`c^{(d)}`$. First, we recall that a generating function for the $`U_n(x)`$ is (e.g. )
$$\frac{1}{12xz+z^2}=\underset{n=0}{\overset{\mathrm{}}{}}U_n(x)z^n$$
(2.26)
The formal sum in (2.26) converges if $`|x|<1`$ and $`|z|<1`$; however, for our purposes, we shall only use it as a means of extracting the coefficient of the $`z^n`$ term as $`U_n(x)`$, and this can be done for any $`x`$ and $`z`$. We infer the following properties of the $`c^{(d)}`$:
1. $`c^{(d)}`$ is a polynomial of degree $`d`$ in the variable $`q`$, with integer coefficients, highest-order term equal to $`q^d`$, and alternating signs for the subsequent terms of descending degree in $`q`$.
2. We recall the definition of a unimodal polynomial as one with the property that the magnitudes of its coefficients increase monotonically up to a point and then (with possible equality of the maximal two coefficients) decrease monotonically. From the property that the magnitudes of the coefficients of $`U_n(x)`$ are unimodal, it follows directly that the same is true of the magnitudes of the coefficients of $`c^{(d)}`$.
3. The zeros of $`c^{(d)}`$ are real, positive, simple, lie in the interval $`0<q<4`$ and are given as follows:
$$c^{(d)}=0\mathrm{at}q=q_{d,k},k=1,2,\mathrm{},d$$
(2.27)
In the limit $`d\mathrm{}`$, these zeros become dense in this interval.
4. If $`d1`$, then one of the zeros of $`c^{(d)}`$ is at a Tutte-Beraha number, namely,
$$q_{d,1}=B_{2d+1}$$
(2.28)
where the Tutte-Beraha number $`B_r`$ is defined as ,
$$B_r=2+2\mathrm{cos}\left(\frac{2\pi }{r}\right)=4\mathrm{cos}^2\left(\frac{\pi }{r}\right),r=1,2,\mathrm{}$$
(2.29)
The first Tutte-Beraha numbers are $`B_1=4`$, $`B_2=0`$, $`B_3=1`$, $`B_4=2`$, $`B_5=(3+\sqrt{5})/2=2.618..`$, $`B_6=3`$.
5. Depending on the value of $`d`$, $`c^{(d)}`$ may have other zero(s) at Tutte-Beraha number(s) distinct from $`B_{2d+1}`$ in (2.28). This cannot happen if $`2d+1`$ is prime. If $`2d+1`$ is not prime, consider the case where there exists $`d^{}`$ such that $`k(2d^{}+1)=2d+1`$. Since $`2d+1`$ and $`2d^{}+1`$ are both odd, it follows that $`k`$ is also odd, so we can write this equation in a symmetric manner as
$$(2d^{}+1)(2d^{\prime \prime }+1)=2d+1.$$
(2.30)
Then $`c^{(d)}`$ has additional zero(s) at $`B_{2d^{}+1}`$ and $`B_{2d^{\prime \prime }+1}`$. These degenerate into a single additional zero if $`2d+1=p^2`$ where $`p`$ is prime. Note that the factorization (2.30) is not unique. For example, consider $`c^{(d)}`$ for $`d=52`$. Then $`2d+1`$ has the full and unique factorization $`105=357`$. There are thus three ways of writing the twofold factorization in (2.30): $`335`$, $`521`$, and $`715`$, and hence, in addition to the zero at the Tutte-Beraha number $`B_{2d+1}=B_{105}`$, the coefficient $`c^{(105)}`$ also has zeros at the six other Tutte-Beraha numbers $`B_3`$, $`B_5`$, $`B_7`$, $`B_{15}`$, $`B_{21}`$, and $`B_{35}`$. As this example makes clear, in other cases there could be more than three twofold factorizations of the form (2.30).
6. As a particular case of the previous item, if $`d=1`$ mod 3, then, setting $`d=3j+1`$ and substituting in the above equation, we find $`(2d^{}+1)(2d^{\prime \prime }+1)=3(2j+1)`$. If $`j=0`$ so that $`d=1`$, i.e., $`c^{(d)}=q1`$, then there is only the single zero at $`B_3=1`$ given by (2.28). If $`j=1`$, i.e., $`d=4`$, then $`c^{(4)}`$ has a zero at $`B_3`$ and at $`B_9`$. If $`j2`$, then in addition to the zero at $`B_{2d+1}`$ given by (2.28), $`c^{(d)}`$ has zeros at the Tutte-Beraha numbers $`B_3`$ and $`B_{2j+1}`$ (as well as possible others). Note that since $`c^{(d)}`$ has only simple zeros, it follows that if $`d=1`$ mod 3, then has a simple factor $`(q1)`$. (The expression given above for $`c^{(4)}`$ illustrates this.)
7. Since $`U_{2n}(0)=(1)^n`$, it follows that
$$c^{(d)}=(1)^d\mathrm{for}q=0.$$
(2.31)
8. $$\mathrm{If}q=1\mathrm{then}c^{(d)}=\{\begin{array}{cc}1\hfill & \text{if }d=0\text{ mod 3}\hfill \\ 0\hfill & \text{if }d=1\text{ mod 3}\hfill \\ 1\hfill & \text{if }d=2\text{ mod 3}\hfill \end{array}$$
(2.32)
9. $$\mathrm{If}q=2\mathrm{then}c^{(d)}=\{\begin{array}{cc}1\hfill & \text{if }d=0,1\text{ mod 4}\hfill \\ 1\hfill & \text{if }d=2,3\text{ mod 4}\hfill \end{array}$$
(2.33)
10. $$\mathrm{If}q=3\mathrm{then}c^{(d)}=\{\begin{array}{cc}1\hfill & \text{if }d=0,2\text{ mod 6}\hfill \\ 2\hfill & \text{if }d=1\text{ mod 6}\hfill \\ 1\hfill & \text{if }d=3,5\text{ mod 6}\hfill \\ 2\hfill & \text{if }d=4\text{ mod 6}\hfill \end{array}$$
(2.34)
11. $$c^{(d)}=2d+1\mathrm{for}q=4.$$
(2.35)
12. As a consequence of the inequality $`|U_n(x)|n+1`$ for $`1x1`$, we have the inequality
$$|c^{(d)}|2d+1\mathrm{for}0q4.$$
(2.36)
This inequality is saturated at $`q=4`$, i.e. $`x=1`$.
13. Using the recursion relation for Chebyshev polynomials,
$$U_{n+1}(x)=2xU_n(x)U_{n1}(x)$$
(2.37)
iteratively, we find the resulting recursion for the $`c^{(d)}`$,
$$c^{(d+1)}=(q2)c^{(d)}c^{(d1)}.$$
(2.38)
If one lets
$$q=2+2\mathrm{cos}\theta =4\mathrm{cos}^2\left(\frac{\theta }{2}\right)$$
(2.39)
one sees that the argument of the Chebyshev polynomial of the second kind in (2.1) is given by $`x=\frac{\sqrt{q}}{2}=\mathrm{cos}(\theta /2)`$. A relevant identity is (with $`\omega =\theta /2`$ here)
$$U_n(\mathrm{cos}\omega )=\frac{\mathrm{sin}((n+1)\omega )}{\mathrm{sin}\omega }.$$
(2.40)
The fact that the transformation (2.39) applies, with $`\theta `$ real, for $`q[0,4]`$ shows the special role of this interval for Potts model partition functions on the strips under consideration here. This makes an intriguing connection with the 2D Potts ferromagnet, which has a second-order phase transition for the values $`q=2`$ (Ising), $`q=3`$, and $`q=4`$; and a first-order phase transition for $`q5`$. Of course, the physical thermodynamic behavior of the Potts ferromagnet on the infinite-length, finite width strips is quite different from that of the model on 2D lattices; in the former cases, since these are quasi-one-dimensional, the model has only a zero-temperature critical point rather than a finite-temperature phase transition. For the antiferromagnetic Potts model, the transformation means that the value $`q=4`$ is special (as the positive value of $`q`$ where the angle $`\theta `$ changes from being real to imaginary) and is in accord with the special role of $`q=4`$ in this model on the infinite two-dimensional lattice (see also ). This is somewhat similar to earlier situations in which studies of lower-dimensional realizations of spin models gave insight into properties of higher-dimensional realizations (e.g. - for $`O(N)`$ models and for Potts/Ising models). However, in assessing the connection of the special role of the value $`q=4`$ in the coefficients $`c^{(d)}`$ with the $`T=0`$ critical point at $`q=4`$ for the Potts antiferromagnet on the triangular lattice, one must also taken into account the fact that the formula (2.1) also applies for the cyclic strips of the square lattice that have been studied; however, again taking the limit $`L_y\mathrm{}`$, the Potts antiferromagnet on the square lattice has a zero-temperature critical point at $`q=3`$ rather than at $`q=4`$ .
One of the interesting aspects of the zeros $`q_{d,k}`$ of $`c^{(d)}`$ is that at these values of $`q`$, the critical exponents of the 2D Potts model are rational. We recall that the thermal, magnetic, and tricritical exponents for the paramagnetic to ferromagnetic transition in the 2D $`q`$-state Potts ferromagnet are, for $`q4`$, -,
$$y_t=\frac{3(1u)}{2u}$$
(2.41)
$$y_h,y_{h,tric.}=\frac{(3u)(5u)}{4(2u)}$$
(2.42)
where
$$u=\frac{2}{\pi }\mathrm{arccos}\left(\frac{\sqrt{q}}{2}\right)$$
(2.43)
so that $`0u1`$ for $`4q0`$. The angle $`\theta `$ in (2.39) is equal to $`\pi u`$ in (2.43), and hence
$$q=q_{d,k}u=\frac{2k}{2d+1}.$$
(2.44)
Note that the converse does not hold; that is, there are rational values of $`u`$ that are not of the form (2.44) and hence do not correspond to any of the $`q_{d,k}`$. An example is $`u=1/2`$.
## 3 Determination of $`n_P(L_y,d)`$ for Cyclic Strips of the Square and Triangular Lattices
In this section and the next we use (2.1) together with two theorems (eqs. (3.10 and (4.1) below) to determine structural properties of the chromatic polynomial and the full Potts model partition function for cyclic strip graphs $`G`$ of the square and triangular lattices. Let us define $`n_P(L_y,d)`$ as the number of terms $`\lambda _{P,G,j}`$ in $`P(G,q)`$ that have as their coefficients $`c_{P,G,j}=c^{(d)}`$ and $`n_Z(L_y,d)`$ as the number of terms $`\lambda _{Z,G,j}`$ in $`Z(G,q,v)`$ that have as their coefficients $`c_{Z,G,j}=c^{(d)}`$. For the cyclic strip graphs under consideration here, these coefficients are independent of $`L_x`$ and depend on $`L_y`$ and $`d`$; furthermore, they are the same for both square and triangular strips, as discussed further below. While the individual $`\lambda _{P,G,j}`$’s in (1.13) are, in general, different for the cyclic strips of the square and triangular lattices, the total number of $`\lambda _{P,G,j}`$’s is the same, so we use the short notation $`N_{P,L_y,\lambda }`$. The same is true for the full partition function, so we use the notation $`N_{Z,L_y,\lambda }`$.
Since the chromatic polynomial $`P(G,q)`$ is the special case $`v=1`$ (i.e. zero-temperature antiferromagnet) of the general (finite-temperature, $`J`$ positive or negative) Potts model partition function, (1.9), it follows that
$$n_P(L_y,d)n_Z(L_y,d).$$
(3.1)
The total number, $`N_{P,G,\lambda }`$, of different terms $`\lambda _{P,G,j}`$ in (1.13) is given by
$$N_{P,L_y,\lambda }=\underset{d=0}{\overset{L_y}{}}n_P(L_y,d).$$
(3.2)
Since each term $`\lambda _{P,G,j}`$ is a distinct eigenvalue of the coloring matrix $`𝒯`$, the number $`N_{P,L_y,\lambda }`$ is the number of different invariant subspaces in the full $`𝒩`$–dimensional space of coloring configurations of the transverse slices of the strips. Similarly, for the full Potts model partition function,
$$N_{Z,L_y,\lambda }=\underset{d=0}{\overset{L_y}{}}n_Z(L_y,d).$$
(3.3)
and analogously this represents the number of different invariant subspaces for the matrix $`𝒯_Z`$. It was shown in that the $`\lambda _{P,G(L_y),j}`$’s are the same for the cyclic and Möbius strips (although the corresponding $`c_{P,G,j}`$’s are different)
$$\lambda _{P,G(L_y),FBC_y,PBC_x,j}=\lambda _{P,G(L_y),FBC_y,TPBC_x,j}j$$
(3.4)
and hence the total number of terms is also the same:
$$N_{P,G(L_y),FBC_y,PBC_x,\lambda }=N_{P,G(L_y),FBC_y,TPBC_x,\lambda }.$$
(3.5)
The argument in relied upon the local nature of the deletion-contraction operations used in calculating $`P(G,q)`$, and the same property is true of the deletion-contraction theorem used for calculating the Tutte polynomial, or equivalently, the Potts model partition function, so that the following generalizations of (3.5) holds :
$$\lambda _{Z,G(L_y),FBC_y,PBC_x,j}=\lambda _{Z,G(L_y),FBC_y,TPBC_x,j}j$$
(3.6)
and hence
$$N_{Z,G(L_y),FBC_y,PBC_x,\lambda }=N_{Z,G(L_y),FBC_y,TPBC_x,\lambda }.$$
(3.7)
For the sum of the coefficients in (1.13), i.e., (1.19), we have
$$C_{P,L_y}=\underset{j=1}{\overset{N_{P,L_y,\lambda }}{}}c_{P,L_y,j}=\underset{d=0}{\overset{L_y}{}}n_P(L_y,d)c^{(d)}$$
(3.8)
and for the corresponding sum of coefficients in (1.12),
$$C_{Z,L_y}=\underset{j=1}{\overset{N_{Z,L_y,\lambda }}{}}c_{Z,L_y,j}=\underset{d=0}{\overset{L_y}{}}n_Z(L_y,d)c^{(d)}.$$
(3.9)
We first recall a theorem specifying $`C_{P,L_y}`$ for cyclic strips of the square and triangular lattice :
Theorem 1.
$$C_{P,L_y}=P(T_{L_y},q)=q(q1)^{L_y1}.$$
(3.10)
Proof. Using coloring matrix methods , one has that $`C_{P,L_y}`$ is equal to the chromatic polynomial for the coloring of the transverse slice of the strip, which is a line with $`L_y`$ vertices. This is a special case of a tree graph, for which the elementary general result in eq. (1.16) holds. $`\mathrm{}`$
Next, we have
Theorem 2. The $`n_P(L_y,d)`$, $`d=0,1,..L_y`$ are determined as follows. One has
$$n_P(L_y,d)=0\mathrm{for}d>L_y,$$
(3.11)
$$n_P(L_y,L_y)=1$$
(3.12)
and
$$n_P(1,0)=1$$
(3.13)
with all other numbers $`n_P(L_y,d)`$ being determined by the two recursion relations
$$n_P(L_y+1,0)=n_P(L_y,1)$$
(3.14)
and
$$n_P(L_y+1,d)=n_P(L_y,d1)+n_P(L_y,d)+n_P(L_y,d+1)\mathrm{for}L_y1\mathrm{and}1dL_y+1.$$
(3.15)
Proof. We substitute for $`c^{(d)}`$ from eq. (2.1) in eq. (3.10). We obtain another equation by differentiating this with respect to $`q`$ once; another by differentiating twice, and so forth up to $`L_y`$-fold differentiations. This yields $`L_y+1`$ linear equations in the $`L_y+1`$ unknowns, $`n_P(L_y,d)`$, $`d=0,1,\mathrm{},L_y`$. We solve this set of equations to get the $`n_P(L_y,d)`$. $`\mathrm{}`$
A corollary is that
$$n_P(L_y,L_y1)=L_y.$$
(3.16)
The numbers $`n_P(L_y,d)`$ can be viewed as integer sequences in $`L_y`$ for a given value of $`d`$. We find that (for $`L_y1`$ where our strips are defined) $`n_P(L_y,0)`$ is a Motzkin number -
$$n_P(L_y,0)=M_{L_y1}$$
(3.17)
where the Motzkin number $`M_n`$ is given by
$$M_n=\underset{j=0}{\overset{n}{}}(1)^jC_{n+1j}\left(\genfrac{}{}{0pt}{}{n}{j}\right)$$
(3.18)
where
$$C_n=\frac{1}{n+1}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)$$
(3.19)
is the Catalan number. (We also use the symbol $`C_n`$ to denote the circuit graph with $`n`$ vertices; the meaning will be clear from context.) The Catalan and Motzkin numbers occur in many combinatoric applications -. Among these is the construction of non-intersecting chords on a circle; the number of ways of connecting a subset of $`n`$ points on a circle by non-intersecting chords is $`M_n`$, while the number of ways of completely connecting $`2n`$ points on the circle by such chords is $`C_n`$. Summing over subsets of points connected by chords, this yields the well-known relation
$$M_n=\underset{k=0}{\overset{[\frac{n}{2}]}{}}\left(\genfrac{}{}{0pt}{}{n}{2k}\right)C_k.$$
(3.20)
As a consequence of our relation (3.14), the equation (3.17) also implies that
$$n_P(L_y,1)=M_{L_y}.$$
(3.21)
A generating function is
$$G_{n_P(L_y,0)}(x)=\frac{1x\sqrt{12x3x^2}}{2x^2}=\underset{L_y=1}{\overset{\mathrm{}}{}}n_P(L_y,0)x^{L_y1}.$$
(3.22)
In eq. (5.1.3) below we shall give an exact determination of the total number, $`N_{P,L_y,\lambda }`$, of types of terms $`\lambda _{P,L_y,j}`$. In Table 1 we list the numbers $`n_P(L_y,d)`$ and the sums $`N_{P,L_y,\lambda }`$ for the first several widths of cyclic strips of the square and triangular lattice, $`1L_y10`$. Others can easily be calculated using our general formulas.
Now let us consider a random walk on the nonnegative integers such that in each step the walker moves by $`+1`$, $`1`$, or 0 units. Denoting $`m(n,k)`$ as the number of walks of length $`n`$ steps starting at $`0`$ and ending at $`k`$, we obtain the Motzkin triangle in Table 2. The first column, corresponding to $`k=0`$ is the number of walks defined above that return to the origin after $`n`$ steps. This is given by the Motzkin number defined in (3.18); $`m(n,0)=M_n`$. The second column, $`m(n,1)`$, is given by the first differences of Motzkin numbers:
$$m(n,1)=\underset{j=0}{\overset{[\frac{n+1}{2}]}{}}\frac{1}{j+1}\left(\genfrac{}{}{0pt}{}{n}{2j1}\right)\left(\genfrac{}{}{0pt}{}{2j}{j}\right).$$
(3.23)
The row sums in Table 2 will be important below; we denote them as
$$𝒮_n=\underset{k=0}{\overset{n}{}}m(n,k).$$
(3.24)
Notice that Table 1 can be viewed as the combination of two Motzkin triangles, i.e. Table 2, as follows,
$$n_P(L_y,d)=m(L_y1,d)+m(L_y1,d1).$$
(3.25)
In light of the property that the coefficients (= multiplicities of eigenvalues of the transfer matrix) $`c^{(d)}`$ are Chebyshev polynomials of the second kind for the cyclic strips of the square and triangular lattices and for Möbius strips of the square lattice (discussed below), and the finding in eq. (3.17) for the numbers of $`\lambda _{P,G,j}`$’s with coefficients $`c^{(0)}`$ and $`c^{(1)}`$, it is interesting to note that in a connection has been shown between what is called a Motzkin polynomial and this Chebyshev polynomial. To see this, one defines $`s_{h,0}=M_h`$, $`s_{h,1}=M_hM_{h1}`$ and shows that, with $`s_{h,n}=s_{h,n1}s_{h1,n1}s_{h2,n2}`$ for $`1nh`$, it is possible to express $`s_{h,n}`$ as $`s_{h,n}=a_nM_h+a_{n1}M_{h1}+\mathrm{}+a_0M_{hn}`$, where the coefficients $`a_j`$ are independent of $`h`$. The Motzkin polynomial is defined as $`S_n(x)=a_nx^n+\mathrm{}+a_1x+a_0`$, and the connection is that $`S_n(x)=U_n(\frac{x1}{2})`$ .
Certain ($`q`$-independent) relations between the $`n_P(L_y,d)`$ can be derived by evaluating the sum (3.8) and its result (3.10) for $`q=0,1`$, and 2. Setting $`q=0`$ in (3.8) and (3.10), and using (2.31), we have
$$\underset{d=0}{\overset{L_y}{}}n_P(L_y,d)c^{(d)}(q=0)=\underset{d=0}{\overset{L_y}{}}(1)^dn_P(L_y,d)=0.$$
(3.26)
Next we evaluate (3.8) and (3.10) for $`q=1`$. If $`L_y=1`$, this only involves one term, $`n_P(1,0)=1`$. So consider $`L_y2`$ and set $`L_y=3k+r`$ with integer $`k0`$ and $`r=0,1`$, or 2. To express the equation compactly, it is convenient to use a Heaviside step function $`\theta (z)`$, where $`\theta (z)=1`$ if $`z>0`$ and $`\theta (z)=0`$ if $`z0`$. Using (2.32), we have
$`{\displaystyle \underset{d=0}{\overset{L_y}{}}}n_P(L_y,d)c^{(d)}(q=1)`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{k}{}}}n_P(L_y,3j){\displaystyle \underset{j=0}{\overset{k1+\theta (r\frac{3}{2})}{}}}n_P(L_y,3j+2)`$ (3.27)
$`=`$ $`0`$ (3.29)
where here and in other summations, it is understood that if the upper limit on the summation is negative, the sum is zero. A similar relation can be obtained in a straightforward manner by substituting $`q=2`$ in (3.8) and (3.10). Other similar relations among the $`n_P(L_y,d)`$ can be obtained by evaluating (3.8) and (3.10) for $`q=3`$ and $`q=4`$; for these relations, the right-hand sides grow with $`L_y`$, in contrast to (3.26)-(3.29), where the right-hand sides are constant.
## 4 Determination of $`n_Z(L_y,d)`$ for Cyclic Strips of the Square and Triangular Lattices
We next generalize these results to the full Potts model partition function. We determine the numbers $`n_Z(L_y,d)`$ by observing that in the cyclic strips under consideration here, the only difference between the coloring for the $`T=0`$ Potts antiferromagnet and for the general Potts model at finite temperature is the relaxation of the constraint that no two adjacent vertices can have the same color (indeed, for ferromagnetic coupling, there is a preference for these vertices to have the same color). Hence, one simply counts all possible colorings, and obtains
Theorem 3.
$$\underset{d=0}{\overset{L_y}{}}n_Z(L_y,d)c^{(d)}=q^{L_y}.$$
(4.1)
Next, we have
Theorem 4. The $`n_Z(L_y,d)`$ are determined as follows. One has
$$n_Z(L_y,d)=0\mathrm{for}d>L_y,$$
(4.2)
$$n_Z(L_y,L_y)=1$$
(4.3)
and
$$n_Z(1,0)=1.$$
(4.4)
All other numbers $`n_Z(L_y,d)`$ are then determined by the two recursion relations
$$n_Z(L_y+1,0)=n_Z(L_y,0)+n_Z(L_y,1)$$
(4.5)
and
$`n_Z(L_y+1,d)`$ $`=`$ $`n_Z(L_y,d1)+2n_Z(L_y,d)+n_Z(L_y,d+1)`$ (4.8)
$`\mathrm{for}1dL_y+1.`$
Proof. The proof is similar to that for Theorem 2: we substitute for $`c^{(d)}`$ from eq. (2.1) in eq. (4.1). We obtain another equation by differentiating this with respect to $`q`$ once; another by differentiating twice, and so forth up to $`L_y`$-fold differentiations. This yields $`L_y+1`$ linear equations in the $`L_y+1`$ unknowns, $`n_Z(L_y,d)`$, $`d=0,1,\mathrm{},L_y`$. We solve this set of equations to get the $`n_Z(L_y,d)`$. $`\mathrm{}`$
A corollary is that
$$n_Z(L_y,L_y1)=2L_y1.$$
(4.9)
From Theorems 3 and 4 we find a general formula for the numbers $`n_Z(L_y,d)`$:
$$n_Z(L_y,d)=\frac{(2d+1)}{(L_y+d+1)}\left(\genfrac{}{}{0pt}{}{2L_y}{L_yd}\right)$$
(4.10)
for $`0dL_y`$, with $`n_Z(L_y,d)=0`$ for $`d>L_y`$. For fixed $`d`$ in the range $`0dL_y`$ where $`n_Z(L_y,d)`$ is nonvanishing, it has the leading asymptotic behavior
$$n_Z(L_y,d)(2d+1)\pi ^{1/2}L_y^{3/2}4^{L_y}\left[1+O(L_y^1)\right]\mathrm{as}L_y\mathrm{}\mathrm{for}\mathrm{fixed}d.$$
(4.11)
(The formal notation $`1+O(L_y^1)`$ in (4.11) and in other asymptotic expansions is not intended to indicate the sign of the coefficient of the $`O(L_y^1)`$ term; here it is actually negative.) As a measure of the asymptotic behavior of $`n_Z(L_y,d)`$ when $`d`$, rather than being fixed, is a finite fraction of $`L_y`$, we take the central value $`d=L_y/2`$ and calculate that
$$n_Z(L_y,\frac{L_y}{2})\pi ^{1/2}L_y^{1/2}4^{2L_y+1}3^{\frac{3}{2}(L_y+1)}\left[1+O(L_y^1)\right]\mathrm{as}L_y\mathrm{}.$$
(4.12)
Note that for the special case $`d=0`$, this reduces to
$$n_Z(L_y,0)=C_{L_y}$$
(4.13)
where the Catalan number $`C_n`$ was defined in (3.19). In Table 3 we list the first few numbers $`n_Z(L_y,d)`$ and the total sums $`N_{Z,L_y,\lambda }`$ that will be calculated below in (5.2.4).
Similar to Table 2, we can also make up a Catalan triangle, and shall find that Table 3 is just the combination of two triangles.
Certain ($`q`$-independent) relations between the $`n_Z(L_y,d)`$ can be derived by evaluating the sum (3.9) and its result (4.1) for $`q=0`$ and $`q=1`$. Setting $`q=0`$ in this sums and using (2.31), we have
$$\underset{d=0}{\overset{L_y}{}}(1)^dn_Z(L_y,d)=0.$$
(4.14)
Next we evaluate (3.9) and (4.1) for $`q=1`$. Let $`L_y=3k+r`$ with $`k0`$ and $`r=0,1`$, or 2. Using (2.32), we have
$$\underset{j=0}{\overset{k}{}}n_Z(L_y,3j)\underset{j=0}{\overset{k1+\theta (r\frac{3}{2})}{}}n_Z(L_y,3j+2)=1.$$
(4.15)
We have obtained the following relations involving both $`n_Z(L_y,d)`$ and $`n_P(L_y,d)`$, as well as the total numbers $`N_{Z,L_y,\lambda }`$ and $`N_{P,L_y,\lambda }`$
$$n_P(L_y,d)=\underset{j=0}{\overset{L_y1}{}}(1)^j\left(\genfrac{}{}{0pt}{}{L_y1}{j}\right)n_Z(L_yj,d)$$
(4.16)
and
$$n_Z(L_y,d)=\underset{j=0}{\overset{L_y1}{}}\left(\genfrac{}{}{0pt}{}{L_y1}{j}\right)n_P(L_yj,d).$$
(4.17)
Since the total numbers of terms $`N_{P,L_y,\lambda }`$ and $`N_{Z,L_y,\lambda }`$ are sums of the $`n_P(L_y,d)`$ and $`n_Z(L_y,d)`$ given, respectively, by (3.2) and (3.3), it follows that
$$N_{P,L_y,\lambda }=\underset{j=0}{\overset{L_y1}{}}(1)^j\left(\genfrac{}{}{0pt}{}{L_y1}{j}\right)N_{Z,L_yj,\lambda }$$
(4.18)
and
$$N_{Z,L_y,\lambda }=\underset{j=0}{\overset{L_y1}{}}\left(\genfrac{}{}{0pt}{}{L_y1}{j}\right)N_{P,L_yj,\lambda }.$$
(4.19)
Note that for $`d=0`$, our eq. (4.16) reduces to the relation (3.18) expressing the Motzkin number as a certain weighted sum of Catalan numbers, while eq. (4.17) reduces to the relation expressing the Catalan number as a weighted sum of Motzkin numbers (e.g., )
$$C_n=\underset{j=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n1}{j}\right)M_{nj1}=\underset{k=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n1}{k}\right)M_k$$
(4.20)
for $`n1`$ (with $`C_0=1`$). Further, for the even and odd widths $`L_y=2\mathrm{}`$ and $`L_y=2\mathrm{}+1`$,
$$N_{Z,\mathrm{},\lambda }=\underset{j=0}{\overset{2\mathrm{}1}{}}(1)^j\left(\genfrac{}{}{0pt}{}{2\mathrm{}1}{j}\right)N_{P,2\mathrm{}j,\lambda }$$
(4.21)
and
$$2N_{Z,\mathrm{},\lambda }=\underset{j=0}{\overset{2\mathrm{}}{}}(1)^j\left(\genfrac{}{}{0pt}{}{2\mathrm{}}{j}\right)N_{P,2\mathrm{}+1j,\lambda }.$$
(4.22)
## 5 Determination of $`N_{P,L_y,\lambda }`$ and $`N_{Z,L_y,\lambda }`$ for Cyclic and Möbius Strips of the Square and Triangular Lattices and Relation with Directed Lattice Animals
In this section we shall use our results above to calculate the total number, $`N_{P,G,\lambda }`$, of $`\lambda _{P,G,j}`$’s in the chromatic polynomial (1.13) and the total number, $`N_{Z,G,\lambda }`$, of $`\lambda _{Z,G,j}`$’s in the full Potts model partition function for cyclic strips $`G`$ of the square and triangular lattice. Since the individual numbers $`n_P(L_y,d)`$ were shown to be the same for the cyclic strips of the square and triangular lattices, and similarly for $`n_Z(L_y,d)`$, clearly it is also true that the respective total numbers $`N_{P,G,\lambda }`$ and $`N_{Z,G,\lambda }`$ are the same for cyclic strips of the square and triangular lattices. Furthermore, as a consequence of eqs. (3.4)-(3.7), our use of cyclic strips to calculate the total numbers $`N_{P,G,\lambda }`$ and $`N_{Z,G,\lambda }`$ also yields these respective numbers for the Möbius strips of the square and triangular lattices. This is useful since, as will be seen in the next sections, the individual coefficients involve a larger set for Möbius strips of the square lattice, namely not just $`c^{(d)}`$ but also $`c^{(d)}`$; moreover, the coefficients for Möbius strips of the triangular lattice are not, in general, polynomials in $`q`$ .
### 5.1 $`N_{P,L_y,\lambda }`$
For the total number of terms $`N_{P,L_y,\lambda }`$, from our theorems above we obtain the recursion relation
$$N_{P,L_y+1,\lambda }=3N_{P,L_y,\lambda }2n_P(L_y,0).$$
(5.1.1)
We find that
$$N_{P,L_y,\lambda }=2𝒮_{L_y}=2(L_y1)!\underset{j=0}{\overset{[\frac{L_y}{2}]}{}}\frac{(L_yj)}{(j!)^2(L_y2j)!}$$
(5.1.2)
where $`𝒮_n`$ was given in (3.24). A generating function for these numbers is
$$G_{N_{P,L_y,\lambda }}(x)=\left(\frac{1+x}{13x}\right)^{1/2}1=\underset{L_y=1}{\overset{\mathrm{}}{}}N_{P,L_y,\lambda }x^{L_y}.$$
(5.1.3)
From this, it follows that the number $`N_{P,L_y,\lambda }`$ grows exponentially fast with the width $`L_y`$ of the cyclic or Möbius strip of the square or triangular lattice, with the leading asymptotic behavior
$$N_{P,L_y,\lambda }L_y^{1/2}3^{L_y}\mathrm{as}L_y\mathrm{}.$$
(5.1.4)
### 5.2 $`N_{Z,G,\lambda }`$
A corollary of Theorem 4 is that
$$N_{Z,L_y+1,\lambda }=4N_{Z,L_y,\lambda }2n_Z(L_y,0).$$
(5.2.1)
Summing the individual numbers $`n_Z(L_y,d)`$ to evaluate the total, (3.3), we have, for cyclic and Möbius strips of width $`L_y`$ of the square and triangular lattices,
$`N_{Z,L_y,\lambda }`$ $`=`$ $`{\displaystyle \underset{d=0}{\overset{L_y}{}}}{\displaystyle \frac{(2d+1)}{(L_y+d+1)}}\left({\displaystyle \genfrac{}{}{0pt}{}{2L_y}{L_yd}}\right)`$ (5.2.2)
$`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{2L_y}{L_y}}\right).`$ (5.2.4)
As $`L_y\mathrm{}`$, $`N_{Z,L_y,\lambda }`$ has the leading asymptotic behavior
$$N_{Z,L_y,\lambda }\pi ^{1/2}L_y^{1/2}4^{L_y}\left[1+O(L_y^1)\right]\mathrm{as}L_y\mathrm{}.$$
(5.2.5)
We give two comparisons of the asymptotic behavior of individual numbers $`n_Z(L_y,d)`$ with the total, $`N_{Z,L_y,\lambda }`$. First, we compare the growth of $`n_Z(L_y,d)`$ for fixed $`d`$ (in the range $`0dL_y`$ where $`n_Z(L_y,d)`$ is nonzero) with the growth of $`N_{Z,L_y,\lambda }`$:
$$\frac{n_Z(L_y,d)}{N_{Z,L_y,\lambda }}\frac{(2d+1)}{L_y}\left[1+O(L_y^1)\right]\mathrm{as}L_y\mathrm{}\mathrm{for}\mathrm{fixed}d.$$
(5.2.6)
Second, we compare the relative growths of the central number $`n_Z(L_y,\frac{L_y}{2})`$ and the total, $`N_{Z,L_y,\lambda }`$:
$`{\displaystyle \frac{n_Z(L_y,\frac{L_y}{2})}{N_{Z,L_y,\lambda }}}`$ $``$ $`\left({\displaystyle \frac{4}{3^{3/2}}}\right)^{L_y+1}\left[1+O(L_y^1)\right]`$ (5.2.7)
$``$ $`(0.769800\mathrm{})^{L_y+1}\left[1+O(L_y^1)\right]`$ (5.2.11)
$`\mathrm{as}L_y\mathrm{}.`$
### 5.3 Connection with Directed Lattice Animals
For an arbitrary $`G`$, the Potts model partition function $`Z(G,q,v)`$ gives information about certain graph-theoretic quantities describing $`G`$, as is clear from the Kasteleyn-Fortuin representation (1.5) and the equivalence in eqs. (12.1.6) and (12.1.9) to the Tutte polynomial and Whitney rank polynomial. An example of this information is illustrated by the zero-temperature antiferromagnetic special case $`P(G,q)=Z(G,q,1)`$ in eq. (1.9), i.e., the chromatic polynomial, counting the number of proper vertex colorings of $`G`$. Other examples are that the Tutte polynomials $`T(G,x,y)`$ for various special values of its arguments, counts the number of spanning trees, spanning forests, connected spanning subgraphs, and spanning subgraphs of $`G`$ (see appendix). It is also well known that the $`q1`$ limit of the Potts model is related to bond percolation and the $`q1`$ limit of a certain multisite Potts model is related to site percolation, as reviewed in .
Here we would like to report two new and very interesting relations with graph-theoretic quantities that we have found from our calculations. These involve directed lattice animals. To explain this, we recall that an animal $`A`$ on a lattice or more generally a graph $`G`$ is defined as a (finite) set of vertices in $`G`$ with the property that any two vertices in $`A`$ are connected by means of a path in $`G`$ having all its vertices in $`A`$. Thus, any vertex in $`A`$ is adjacent to another vertex in $`A`$. If the graph is a lattice, the animal is termed a lattice animal, and on an infinite lattice, the lattice animals are defined up to an overall lattice translation. A directed lattice animal $`A_{dir}`$ is an animal on a lattice $`\mathrm{\Lambda }`$ together with a special origin or root point $`O`$ such that any vertex of $`A`$ can be reached starting from $`O`$ by an oriented path on $`\mathrm{\Lambda }`$ having all of its vertices in $`A`$. For example, a directed lattice animal on the square lattice could be defined by restricting the orientations of each step of the oriented path to be “eastward” or “northward”, so that, starting at the origin, this animal would extend in the northeast quadrant and would not contain any backtracking steps going in the south or west directions. (The animals that we consider here are also called site animals to distinguish them from bond animals; we shall leave this implicit.) Animals and directed animals are related, respectively, to percolation and directed percolation. Let us denote $`N_{A,\mathrm{\Lambda },n}`$ and $`N_{DA,\mathrm{\Lambda },n}`$ as the total number of lattice animals and directed lattice animals with $`n`$ vertices, respectively, on the lattice $`\mathrm{\Lambda }`$. For $`n\mathrm{}`$, the numbers of lattice animals and directed lattice animals grow asymptotically like $`n^\theta a^n`$, where $`a`$ depends on the lattice and, e.g., for 2D lattices, $`\theta =1/2`$ (e.g. - and references therein). Having given this background, we now state the relations that we have obtained.
First, we find that the total number of distinct eigenvalues of the coloring matrix, i.e. the total number, $`N_{P,L_y,\lambda }`$, of different $`\lambda _{P,G,j}`$’s in the chromatic polynomial (1.13) for cyclic strips with width $`L_y`$ of the square and triangular lattice (equal by (3.5) to the same number for the corresponding Möbius strips of the square and triangular lattices) is twice the number of directed lattice animals with $`n=L_y`$ vertices on the square (sq) lattice:
$$N_{P,L_y,\lambda }=2N_{DA,sq,L_y}.$$
(5.3.1)
We established this by observing that the generating function for directed lattice animals on the square lattice, which is known exactly , is precisely $`(1/2)`$ times the generating function (5.1.3).
Second, we find that the analogous total number, $`N_{Z,L_y,\lambda }`$, of different eigenvalues $`\lambda _{Z,G,j}`$’s appearing in the full Potts model partition function (1.12) for these cyclic and Möbius strips with width $`L_y`$ of the square and triangular lattices is twice the number of directed lattice animals with $`n=L_y`$ vertices on the triangular lattice:
$$N_{Z,L_y,\lambda }=2N_{DA,tri,L_y}.$$
(5.3.2)
We established this by recalling the known result
$$N_{DA,tri,n}=\frac{1}{2}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)$$
(5.3.3)
and using our calculation (5.2.4). It should be emphasized that, as we have shown above, the numbers $`N_{P,L_y,\lambda }`$ and $`N_{Z,L_y,\lambda }`$ on the respective left-hand sides of (5.3.1) and (5.3.2) apply for cyclic and Möbius strips of both the square and triangular lattice, whereas the right-hand side of (5.3.1) is specific to the square lattice and the right-hand side of (5.3.2) is specific to the triangular lattice.
Parenthetically, we mention some other relations. The enumeration of directed lattice animals on the square and triangular lattices was shown in to be connected with the hard-square lattice gas model . As discussed in , no exact result has been obtained, analogous to those in , for the number of directed lattice animals on the honeycomb lattice. Directed lattice animals are connected with directed percolation, and the latter has been related to a kind of chiral Potts model . Compact lattice animals are the subset of lattice animals with no unoccupied interior sites; directed compact lattice animals have been related to the $`q\mathrm{}`$ limit of the Potts model and have been shown to have numbers that grow less rapidly than the leading $`a^n`$ growth in the numbers of regular directed lattice animals .
## 6 Determination of $`n_{P,Mb}(L_y,d,\pm )`$ for Möbius Strips of the Square Lattice
Möbius strips differ in their global topology from cyclic strips. For the Möbius strips of the square lattice that have been studied, the coefficients are polynomials in $`q`$, but the polynomials that occur for the strip of width $`L_y`$ arise from the set $`\pm c^{(d)}`$, where $`0dL_y/2`$ for even $`L_y`$ and $`0d(L_y+1)/2`$ for odd $`L_y`$, rather than the set $`c^{(d)}`$, $`0dL_y`$ as in the case of the cyclic strips. In passing, we note that for Möbius strips of the triangular lattice, we have found, originally for width $`L_y=2`$ in and later for width $`L_y=3`$ in , that in the case of chromatic polynomials, some of the coefficients are not polynomials, but rather, algebraic functions of $`q`$ ; similarly, some of the coefficients in the full Potts model partition function for the Möbius strip of the triangular lattice are not polynomials in $`q`$ . This is related to the fact, discussed before , that the Möbius strip of the triangular lattice involves a seam, i.e., it is not translationally homogeneous in the longitudinal direction, whereas the Möbius strip of the square lattice does not have a seam and is translationally homogeneous in the longitudinal direction. Hence, for Möbius strips, we only consider the square lattice here.
Since both signs of the $`c^{(d)}`$ occur, it is necessary to define a larger set of numbers for Möbius strips of the square lattice. Let $`n_{P,Mb}(L_y,d,\pm )`$ denote the number of terms $`\lambda _j`$’s in the expression (1.13) for $`P(sq,L_y\times L_x,FBC_y,TPBC_x,q)`$ with coefficients $`c_{P,L_y,Mb,j}=\pm c^{(d)}`$, respectively, where $`Mb`$ denotes Möbius. The notation $`c_{P,L_y,Mb,j}`$ reflects the fact that these coefficients depend on $`L_y`$ and the Möbius boundary conditions, but are independent of $`L_x`$. The sums and differences for each $`d`$ are defined as
$$n_{P,Mb,tot}(L_y,d)=n_{P,Mb}(L_y,d,+)+n_{P,Mb}(L_y,d,)$$
(6.1)
and
$$\mathrm{\Delta }n_{P,Mb}(L_y,d)=n_{P,Mb}(L_y,d,+)n_{P,Mb}(L_y,d,).$$
(6.2)
We first obtain a general theorem for the sum of the coefficients:
Theorem 5.
$$C_{P,L_y,Mb}=\underset{j=1}{\overset{N_{P,L_y,\lambda }}{}}c_{P,L_y,Mb,j}=\underset{d=0}{\overset{d_{max}}{}}\mathrm{\Delta }n_{P,Mb}(L_y,d)c^{(d)}=\{\begin{array}{cc}0\hfill & \text{for even }L_y\hfill \\ P(T_{(\frac{L_y+1}{2})},q)\hfill & \text{for odd }L_y\hfill \end{array}$$
(6.3)
where $`P(T_n,q)`$ was given in (1.16) and
$$d_{max}=\{\begin{array}{cc}\frac{L_y}{2}\hfill & \text{for even }L_y\hfill \\ \frac{(L_y+1)}{2}\hfill & \text{for odd }L_y\hfill \end{array}.$$
(6.4)
Proof. Our method for proving this theorem is inspired by coloring matrix methods . The Möbius strip involves a reversed-orientation periodic longitudinal boundary condition. We can think of constructing such a strip by cutting a cyclic strip, reversing the orientation of one of the ends, and gluing these ends together again. For a strip with $`L_y`$ even, let us label the vertices on the two ends as $`1,2,..,L_y`$; then the Möbius boundary condition means identifying vertex 1 with vertex $`L_y`$, vertex 2 with vertex $`L_y1`$, and so forth. As regards the coloring constraint, in the case $`L_y=2`$, this effectively produces a subgraph consisting of a single vertex with an edge forming a loop. In the case $`L_y=4`$, this produces a vertex $`1=4`$ with two edges connecting to a vertex $`2=3`$, which in turn has a loop connected to it, and so forth for higher even values of $`L_y`$. Because of the loop that appears in each even-$`L_y`$ case, the chromatic polynomial for coloring this subgraph vanishes identically. (The value of $`d_{max}`$ given in (6.4) follows from this graphical construction.) This proves the theorem for the case of even $`L_y`$. For odd $`L_y`$, consider first the case $`L_y=3`$; here the Möbius boundary condition identifies vertex 1 with vertex 3 and leaves vertex 2 invariant. Hence, as regards the coloring, it effectively produces a subgraph consisting of the vertex $`1=3`$ connected by two edges with the vertex 2. The chromatic polynomial for the coloring of this subgraph is not sensitive to the multiple edges and hence is $`P(T_2,q)=q(q1)`$. Again, this subgraph construction yields the value of $`d_{max}`$ for odd $`L_y`$ in (6.4). In a similar manner, for higher odd values of $`L_y`$, the Möbius boundary condition leads to a subgraph with $`(L_y+1)/2`$ vertices forming a chain, with each interior pair connected to the next by two edges. The chromatic polynomial for the coloring of this subgraph is $`P(T_{(L_y+1)/2},q)=q(q1)^{(L_y1)/2}`$. This completes the proof of the theorem. $`\mathrm{}`$
Two corollaries of this theorem are as follows.
Corollary 1
$$n_{P,Mb}(L_y,d,+)=n_{P,Mb}(L_y,d,)\mathrm{for}\mathrm{even}L_y$$
(6.5)
Proof. If $`L_y`$ is even, then in order for the terms of the highest power in $`q`$ to cancel so as to yield a sum of 0 as in (6.3), it is necessary that $`n_{P,Mb}(L_y,d_{max},+)=n_{P,Mb}(L_y,d_{max},)`$, where $`d_{max}`$ was given above in (6.4). But then in turn, in order for the terms of degree $`d_{max}1`$ to sum to zero, it is necessary that $`n_{P,Mb}(L_y,d_{max}1,+)=n_{P,Mb}(L_y,d_{max}1,)`$, and so forth for all powers. $`\mathrm{}`$
Corollary 2
$$\mathrm{\Delta }n_{P,Mb}(L_y,d)=n_P(\frac{L_y+1}{2},d)\mathrm{for}\mathrm{odd}L_y3$$
(6.6)
Proof. This follows from eq. (6.3) by the same kind of argument that was used to obtain the $`n_P(L_y,d)`$. Given that the sum $`C_{P,L_y,Mb}`$ satisfies (6.3), this uniquely determines the $`\mathrm{\Delta }n_{P,Mb}(L_y,d)`$’s just as (3.10) determined the $`n_P(L_y,d)`$. $`\mathrm{}`$
In order to determine the numbers of coefficients $`c^{(d)}`$, i.e., $`n_{P,Mb}(L_y,d,\pm )`$, for a Möbius strip of width $`L_y`$, we start with the chromatic polynomial for the cyclic $`L_y`$ strip and determine how the coefficients $`c^{(d)}`$ change when one changes the longitudinal boundary condition from cyclic to Möbius. The next two theorems determines this (the proofs are given after the second of these two theorems).
Theorem 6 Consider an $`L_y\times L_x`$ strip of the square lattice with $`(FBC_y,PBC_x)`$. As before, the coefficients $`c_{P,G,j}`$ are made up from the set $`c^{(d)}`$ with $`0dL_y`$. When one changes the longitudinal boundary condition from cyclic to Möbius, the following respective changes of coefficients of even degree $`d=2k`$ and of odd degree $`2k+1`$ occur:
$$c^{(0)}\pm c^{(0)}$$
(6.7)
$$c^{(2k)}\pm c^{(k1)},1k\left[\frac{L_y}{2}\right]$$
(6.8)
$$c^{(2k+1)}\pm c^{(k+1)},0k\left[\frac{L_y1}{2}\right]$$
(6.9)
where in eqs. (6.8) and (6.9), we again use the notation $`[\nu ]`$ to denote the integral part of $`\nu `$. Thus, if there are $`n_P(L_y,0)`$ coefficients $`c_{P,G,j}`$ of the terms $`(\lambda _j)^m`$ of the form $`c^{(0)}`$ in eq. (1.13) for $`P(sq(L_y\times L_x;FBC_y,PBC_x),q)`$, then the respective coefficients multiplying the terms $`(\lambda _j)^m`$ in the chromatic polynomial for the corresponding Möbius strip, $`P(sq(L_y\times L_x;FBC_y,TPBC_x),q)`$, are either of the form $`+c^{(0)}`$ or $`c^{(0)}`$, and so forth for the $`c^{(d)}`$ with $`d>0`$, as specified by eqs. (6.8) and (6.9).
Theorem 7. With the same premise as in Theorem 6, we have
$$n_{P,Mb}(L_y,0,\pm )=\frac{1}{2}n_{P,Mb,tot}(L_y,0)=\frac{1}{2}\left[n_P(L_y,0)+n_P(L_y,2)\right]\mathrm{for}L_y\mathrm{even}$$
(6.10)
$`n_{P,Mb}(L_y,d,\pm )`$ $`=`$ $`{\displaystyle \frac{1}{2}}n_{P,Mb,tot}(L_y,d)={\displaystyle \frac{1}{2}}\left[n_P(L_y,2d1)+n_P(L_y,2d+2)\right]`$ (6.13)
$`\mathrm{for}L_y\mathrm{even}\mathrm{and}1d{\displaystyle \frac{L_y}{2}}`$
$$n_{P,Mb}(L_y,d,\pm )=0\mathrm{for}L_y\mathrm{even}\mathrm{and}d\frac{L_y+2}{2}$$
(6.14)
$$n_{P,Mb}(L_y,0,\pm )=\frac{1}{2}\left[n_P(L_y,0)+n_P(L_y,2)\pm n_P(\frac{L_y+1}{2},0)\right]\mathrm{for}L_y\mathrm{odd}$$
(6.15)
$`n_{P,Mb}(L_y,d,\pm )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[n_P(L_y,2d1)+n_P(L_y,2d+2)\pm n_P({\displaystyle \frac{L_y+1}{2}},d)\right]`$ (6.18)
$`\mathrm{for}L_y\mathrm{odd}\mathrm{and}1d{\displaystyle \frac{L_y+1}{2}}`$
$$n_{P,Mb}(L_y,d,\pm )=0\mathrm{for}L_y\mathrm{odd}\mathrm{and}d\frac{L_y+3}{2}$$
(6.19)
Proofs. We prove these theorems using an inductive argument. From the chromatic polynomials for the cyclic and Möbius strips of the square lattice with width $`L_y=2`$ , one knows the transformation rule for the coefficients $`c^{(d)}`$ with $`0d2`$ as one goes from cyclic to Möbius longitudinal boundary conditions, namely, $`c^{(0)}\pm c^{(0)}`$, $`c^{(1)}\pm c^{(1)}`$, and $`c^{(2)}\pm c^{(0)}`$, where the sign information will not be needed. For $`L_y=3`$, we know the numbers $`n_P(L_y,d)`$ for the cyclic strip from the previous theorem. We also know the differences of the numbers of coefficients of each degree for the $`L_y=3`$ Möbius strip, $`\mathrm{\Delta }n_{P,Mb}(3,d)`$; these are determined by the relation (6.3). In particular, we have $`\mathrm{\Delta }n_{P,Mb}(3,0)=n_P(2,0)=1`$, $`\mathrm{\Delta }n_{P,Mb}(3,1)=n_P(2,1)=2`$, and $`\mathrm{\Delta }n_{P,Mb}(3,2)=n_P(2,2)=1`$. From the transformation rules obtained so far, it follows that the sources for coefficients of degree 0 for the $`L_y=3`$ Möbius strip are the coefficients of degree 0 and 2 for the corresponding cyclic strip. Adding these, we have $`n_{P,Mb,tot}(3,0)=n_P(3,0)+n_P(3,2)=5`$. This relation for the sum of $`n_{P,Mb}(3,0,+)`$ and $`n_{P,Mb}(3,0,)`$, together with the relation $`\mathrm{\Delta }n_{P,Mb}(3,0)=1`$ for their difference, enables us to solve for each of these numbers, and we get $`n_{P,Mb}(3,0,+)=3`$, $`n_{P,Mb}(3,0,)=2`$. From the transformation rule $`c^{(1)}\pm c^{(1)}`$ found previously, we infer that for the Möbius strip, $`n_{P,Mb,tot}(3,1)=n_P(3,1)=4`$. Combining this with the relation $`\mathrm{\Delta }n_{P,Mb}(3,1)=2`$, we solve to get $`n_{P,Mb}(3,1,+)=3`$ and $`n_{P,Mb}(3,1,)=1`$. We next consider $`n_{P,Mb}(3,2,\pm )`$. We know that $`\mathrm{\Delta }n_{P,Mb}(3,2)=1`$, and, given the transformation rules obtained so far, there is only one source for terms with degree coefficients of degree 2 in the chromatic polynomial for the $`L_y=2`$ Möbius strip, namely the single term with a $`c^{(3)}`$ coefficient. Hence, $`n_{P,Mb,tot}(3,2)=1`$, and so, knowing the sum and difference, we get $`n_{P,Mb}(3,2,+)=1`$ and $`n_{P,Mb}(3,2,)=0`$. Since there are no sources for any higher-degree coefficients, we have $`n_{P,Mb}(3,d,\pm )=0`$ for $`d3`$. Thus, we have both determined all of the numbers $`n_{P,Mb}(3,d,\pm )`$ and the transformation rule for the coefficients of next higher degree, namely $`c^{(3)}\pm c^{(2)}`$.
Proceeding to $`L_y=4`$, we start with the knowledge of the $`n_P(4,d)`$ for the cyclic strip and of the differences $`\mathrm{\Delta }n_{P,Mb}(4,d)=0`$. For coefficients of degree 0 in the chromatic polynomial for the Möbius strip, the transformation rules indicate two sources, namely the coefficients with degree 0 or 2 in the corresponding cyclic strip. We shall show, it a posteriori, that these are the only sources, so we infer that $`n_{P,Mb,tot}(4,0)=n_P(4,0)+n_P(4,2)=12`$. Combining this with the relation $`\mathrm{\Delta }n_{P,Mb}(4,0)=0`$, we obtain $`n_{P,Mb}(4,0,\pm )=(1/2)n_{P,Mb,tot}(4,0)=6`$. Next, for the coefficients of degree 1, we start with those from the cyclic strip, $`n_P(4,1)=9`$. But we know that there must be another source because this is an odd number and, by itself, would give the unacceptable non-integral result for $`n_{P,Mb}(4,0,\pm )=(1/2)n_{P,Mb,tot}(4,0)`$. The only possible source is the coefficients of degree 4 in the cyclic strip, so we infer the next transformation rule: $`c^{(4)}\pm c^{(1)}`$. This yields $`n_{P,Mb,tot}(4,1)=n_P(4,1)+n_P(4,4)=10`$ and hence $`n_{P,Mb}(4,0,\pm )=5`$. Finally, for the coefficients of degree 2 in the chromatic polynomial for the $`L_y=4`$ Möbius strip, the source is the coefficients of degree 3 in the cyclic $`L_y=4`$ strip, so $`n_{P,Mb,tot}(4,2)=n_P(4,3)=4`$, whence $`n_{P,Mb}(4,2,\pm )=2`$. Since there are no sources for any higher-degree coefficients, we have $`n_{P,Mb}(4,d,\pm )=0`$ for $`d3`$. Thus again we have determined all of the numbers $`n_{P,Mb}(4,d,\pm )`$ and also the next higher transformation rule for $`c^{(4)}`$.
It should now be clear how one proceeds iteratively: at each higher width $`L_y`$, one uses the previously established transformation rules, and obtains the transformation rule for the next higher degree coefficient, to determine all of the numbers $`n_{P,Mb}(L_y,d,\pm )`$. This completes the proof of Theorems 5 and 6. $`\mathrm{}`$.
Two corollaries of these theorems are
$$n_{P,Mb}(L_y,\frac{L_y}{2},\pm )=\frac{L_y}{2}\mathrm{for}\mathrm{even}L_y$$
(6.20)
$$n_{P,Mb}(L_y,\frac{L_y+1}{2},+)=1,n_{P,Mb}(L_y,\frac{L_y+1}{2},)=0\mathrm{for}\mathrm{odd}L_y$$
(6.21)
Values of the numbers $`n_{P,Mb}(L_y,d,\pm )`$ for the first several widths, $`2L_y10`$, are given in Table 4.
Further corollaries involve ($`q`$-independent) relations between the $`n_P(L_y,d)`$. As before for the cyclic strips, we derive these by evaluating (6.3) for $`q=0,1`$, and 2. For example, setting $`q=0`$ and using (2.31), we have (for $`L_y2`$ where Möbius strips are defined)
$$\underset{d=0}{\overset{d_{max}}{}}(1)^d\mathrm{\Delta }n_{P,Mb}(L_y,d)=0\mathrm{where}d_{max}=\{\begin{array}{cc}\frac{L_y}{2}\hfill & \text{for even }L_y\hfill \\ \frac{(L_y+1)}{2}\hfill & \text{for odd }L_y\hfill \end{array}$$
(6.22)
It is straightforward to derive similar relations among the $`\mathrm{\Delta }n_{P,Mb}(L_y,d)`$ by evaluating (6.3) for $`q=1`$ and $`q=2`$.
Another result pertains to the detailed structure of the chromatic polynomials for the cyclic, as compared with Möbius strips of a given width $`L_y`$, as the length $`L_x`$ gets large. As discussed in , because the cyclic and Möbius strips of a given width have the same number of vertices and edges, the coefficients of the leading powers of $`q`$ are the same. We recall this result. For $`m`$ greater than a minimal value<sup>4</sup><sup>4</sup>4 For example, for $`m4`$ ($`m3`$) the $`L_y=2`$ cyclic (Möbius) square strips exhibit special behavior regarding $`g`$ and $`k_g`$; for $`m`$ larger than these respective values, they both have $`g=4`$ and $`k_g=m`$., the cyclic and Möbius strips of a given $`(G_s)_m`$ have the same number of vertices $`n`$, edges (bonds) $`e`$, girth $`g`$ (length of minimum closed circuit on $`G_s`$) and number $`k_g`$ of circuits of length $`g`$. One has $`n=t_sm`$ where $`t_s`$ depends on $`G_s`$. Writing
$$P((G_s),q)=\underset{j=0}{\overset{n1}{}}(1)^jh_{nj}q^{nj}$$
(6.23)
and using the results that $`h_{nj}=\left(\genfrac{}{}{0pt}{}{e}{j}\right)`$ for $`0j<g1`$ (whence $`h_n=1`$ and $`h_{n1}=e`$) and $`h_{n(g1)}=\left(\genfrac{}{}{0pt}{}{e}{g1}\right)k_g`$, it follows that for $`m`$ greater than the above-mentioned minimal value, these $`h_j`$’s are the same for the cyclic and Möbius strips of each type $`G_s`$. For a given $`G_s`$, as $`m`$ increases, the $`h_{nj}`$’s for the cyclic and Möbius strips become equal for larger $`j`$. Our relations (6.7)-(6.9) make clear which coefficients $`c_{G_s(L_y),j}`$ in (1.13), and hence which coefficients $`h_{G_s(L_y),j}`$ in (6.23), remain equal for the cyclic versus Möbius strips.
## 7 Determination of the numbers $`n_{Z,Mb}(L_y,d,\pm )`$ for Möbius Strips of the Square Lattice
In this section we generalize the calculations of the previous section to the full Potts model partition function on Möbius strips of the square lattice. (The restriction to the square lattice will be implicit in the notation.) Let $`n_{Z,Mb}(L_y,d,\pm )`$ denote the number of terms $`\lambda _{Z,L_y,j}`$ in the expression (1.12) for $`Z(L_y\times L_x,FBC_y,TPBC_x,q,v)`$ with coefficients $`c_{Z,L_y,Mb,j}=\pm c^{(d)}`$, respectively, where $`Mb`$ denotes Möbius. The notation $`c_{Z,L_y,Mb,j}`$ reflects the fact that these coefficients are independent of $`L_x`$. The sums and differences for each $`d`$ are defined as
$$n_{Z,Mb,tot}(L_y,d)=n_{Z,Mb}(L_y,d,+)+n_{Z,Mb}(L_y,d,)$$
(7.1)
and
$$\mathrm{\Delta }n_{Z,Mb}(L_y,d)=n_{Z,Mb}(L_y,d,+)n_{Z,Mb}(L_y,d,).$$
(7.2)
We first have the two theorems
Theorem 8.
$$C_{Z,L_y,Mb}\underset{j=1}{\overset{N_{Z,L_y,\lambda }}{}}c_{Z,L_y,Mb,j}=\underset{d=0}{\overset{d_max}{}}\mathrm{\Delta }n_{Z,Mb}(L_y,d)c^{(d)}=\{\begin{array}{cc}q^{L_y/2}\hfill & \text{for even }L_y\hfill \\ q^{(L_y+1)/2}\hfill & \text{for odd }L_y\hfill \end{array}$$
(7.3)
Theorem 9. For the Möbius strips of the square lattice,
$$n_{Z,Mb}(L_y,0,\pm )=\frac{1}{2}\left[n_Z(L_y,0)+n_Z(L_y,2)\pm n_Z(\frac{L_y}{2},0)\right]\mathrm{for}L_y\mathrm{even}$$
(7.4)
$`n_{Z,Mb}(L_y,d,\pm )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[n_Z(L_y,2d1)+n_Z(L_y,2d+2)\pm n_Z({\displaystyle \frac{L_y}{2}},d)\right]`$ (7.7)
$`\mathrm{for}L_y\mathrm{even}\mathrm{and}1d{\displaystyle \frac{L_y}{2}}`$
$$n_{Z,Mb}(L_y,0,\pm )=\frac{1}{2}\left[n_Z(L_y,0)+n_Z(L_y,2)\pm n_Z(\frac{L_y+1}{2},0)\right]\mathrm{for}L_y\mathrm{odd}$$
(7.8)
$`n_{Z,Mb}(L_y,d,\pm )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[n_Z(L_y,2d1)+n_Z(L_y,2d+2)\pm n_Z({\displaystyle \frac{L_y+1}{2}},d)\right]`$ (7.11)
$`\mathrm{for}L_y\mathrm{odd}\mathrm{and}1d{\displaystyle \frac{L_y+1}{2}}.`$
$$n_{Z,Mb}(L_y,d,\pm )=0\mathrm{for}d>\left[\frac{L_y+1}{2}\right]$$
(7.12)
The proofs of these theorems are analogous to those for the analogous Theorems 5 and 7 for chromatic polynomials for Möbius strips of the square lattice.
Two corollaries of these theorems are
$$n_{Z,Mb}(L_y,\frac{L_y}{2},+)=L_y,n_{Z,Mb}(L_y,\frac{L_y}{2},)=L_y1\mathrm{for}\mathrm{even}L_y$$
(7.13)
$$n_{Z,Mb}(L_y,\frac{L_y+1}{2},+)=1,n_{Z,Mb}(L_y,\frac{L_y+1}{2},)=0\mathrm{for}\mathrm{odd}L_y$$
(7.14)
Values of the numbers $`n_{Z,Mb}(L_y,d,\pm )`$ for the first several widths, $`2L_y10`$, are given in Table 5.
Two relations among the $`n_{Z,Mb}(L_y,d)`$ can be derived as before, by evaluating the sum (7.3) for $`q=0`$ and $`q=1`$. For example, from the evaluation for $`q=0`$ we have (for $`L_y2`$)
$$\underset{d=0}{\overset{d_{max}}{}}(1)^d(\mathrm{\Delta }n_{Z,Mb}(L_y,d))=0\mathrm{where}d_{max}=\{\begin{array}{cc}\frac{L_y}{2}\hfill & \text{for even }L_y\hfill \\ \frac{(L_y+1)}{2}\hfill & \text{for odd }L_y\hfill \end{array}$$
(7.15)
It is straightforward to obtain a corresponding relation via the evaluation of (7.3) for $`q=1`$.
## 8 Some Examples
Since some of the notation above is complicated, it is worthwhile to illustrate our general calculations with specific exact solutions. These will also be useful in a discussion of the special cases of the infinite-temperature limit of the Potts model and the zero-temperature limit of the Potts ferromagnet below.
### 8.1 Chromatic Polynomials
For the $`L_y=3`$ cyclic strip of the square lattice, $`N_{P,L_y=3,\lambda }=10`$ and, with the abbreviation $`sq3FP`$ for $`sq,L_y=3,FBC_y,PBC_x`$, the chromatic polynomial is
$$P(sq,3\times L_x,FBC_y,PBC_x,q)=\underset{j=1}{\overset{10}{}}c_{sq3FP,j}(\lambda _{sq3FP,j})^{L_x}$$
(8.1.1)
where
$$\lambda _{sq3FP,1}=1$$
(8.1.2)
$$\lambda _{sq3FP,2}=q1$$
(8.1.3)
$$\lambda _{sq3FP,3}=q2$$
(8.1.4)
$$\lambda _{sq3FP,4}=q4$$
(8.1.5)
$$\lambda _{sq3FP,5}=(q2)^2$$
(8.1.6)
$$\lambda _{sqFP,(6,7)}=\frac{1}{2}\left[(q2)(q^23q+5)\pm \left\{(q^25q+7)(q^45q^3+11q^212q+8)\right\}^{1/2}\right]$$
(8.1.7)
and $`\lambda _{sqFP,j}`$, $`j=8,9,10`$, are the roots of the cubic equation
$`\xi ^3+(2q^29q+12)\xi ^2+(q^410q^3+36q^256q+31)\xi `$ (8.1.8)
(8.1.9)
$`(q1)(q^49q^3+29q^240q+22)=0.`$ (8.1.10)
The coefficients are
$$c_{sq3FP,1}=c^{(3)}$$
(8.1.11)
$$c_{sq3FP,j}=c^{(2)}\mathrm{for}j=2,3,4$$
(8.1.12)
$$c_{sq3FP,j}=c^{(1)}\mathrm{for}j=5,8,9,10$$
(8.1.13)
$$c_{sq3FP,j}=c^{(0)}=1\mathrm{for}j=6,7.$$
(8.1.14)
For the numbers $`n_P(L_y,d)`$, one sees that $`n_P(3,0)=2`$, $`n_P(3,1)=4`$, $`n_P(3,2)=3`$, $`n_P(3,3)=1`$, and the total number $`N_{P,L_y,\lambda }=10`$, in agreement with Table 1. The exact solution for the chromatic polynomial for the $`L_y=3`$ cyclic strip of the triangular lattice, given in , has the same $`N_{P,L_y,\lambda }`$ and $`n_P(L_y,d)`$’s, although eight of the ten $`\lambda _{P,G,j}`$’s are different.
As an illustration of our results for chromatic polynomials for Möbius strips of the square lattice, we again take the $`L_y=3`$ case. The exact solution for the chromatic polynomial (with the abbreviation $`sq3Mb`$ for $`sq,L_y=3,FBC_y,TPBC_x`$) is
$$P(sq,3\times L_x,FBC_y,TPBC_x,q)=\underset{j=1}{\overset{10}{}}c_{P,sq3Mb,j}(\lambda _{P,sq3Mb,j})^{L_x}$$
(8.1.15)
where $`\lambda _{P,sq3Mb,j}=\lambda _{P,sq3FP,j}j`$, in accordance with the general result (3.4) and
$$c_{P,sq3Mb,1}=c^{(2)}$$
(8.1.16)
$$c_{P,sq3Mb,j}=c^{(1)}\mathrm{for}j=8,9,10$$
(8.1.17)
$$c_{P,sq3Mb,5}=c^{(1)}$$
(8.1.18)
$$c_{P,sq3Mb,j}=c^{(0)}\mathrm{for}j=3,6,7$$
(8.1.19)
$$c_{P,sq3Mb,j}=c^{(0)}\mathrm{for}j=2,4.$$
(8.1.20)
Thus, for the numbers $`n_P(L_y,d,\pm )`$, we have $`n_P(3,0,+)=3`$, $`n_P(3,0,)=2`$, $`n_P(3,1,+)=3`$, $`n_P(3,1,)=1`$, $`n_P(3,2,+)=1`$, and $`n_P(3,2,)=0`$, with $`n_P(3,d,\pm )=0`$ for $`d3`$, in agreement with Table 4.
The fact that the coefficients in the chromatic polynomial for the Möbius strip of the triangular lattice are not, in general, polynomials in $`q`$, was shown via the exact solution for the lowest width, $`L_y=2`$ in (with $`L_x=m`$):
$`P(tri,2\times L_x,FBC_y,TPBC_x,q)=1+[(q2)^2]^m{\displaystyle \frac{(q1)(q3)}{\sqrt{94q}}}\left[(\lambda _{t2Mb,3})^m(\lambda _{t2Mb,4})^m\right]`$ (8.1.21)
$`=`$ $`1+(q2)^{2m}2^{1m}(q1)(q3){\displaystyle \underset{s=0}{\overset{[(m1)/2]}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{m}{2s+1}}\right)(52q)^{m2s1}(94q)^s`$ (8.1.23)
where
$$\lambda _{t2Mb,(3,4)}=\frac{1}{2}\left[52q\pm \sqrt{94q}\right].$$
(8.1.24)
As we have discussed in , this can be attributed to the feature that the Möbius strip of the triangular lattice has a seam, and hence is not translationally homogeneous in the longitudinal direction, in contrast to the Möbius strip of the square lattice, which is translationally homogeneous in the longitudinal direction.
### 8.2 Potts Model Partition Function
As an illustration for the full Potts model partition function, we recall the calculation of this function for the $`L_y=2`$ cyclic strip of the square lattice. Here, for $`N_{Z,L_y,\lambda }`$, we have $`N_{Z,2,\lambda }=6`$, and (with the shorthand $`sq2FP`$ for $`sq,L_y=2,FBC_y,PBC_x`$)
$$Z(sq,2\times L_x,FBC_y,PBC_x,q,v)=\underset{j=1}{\overset{6}{}}c_{Z,sq2FP,j}(\lambda _{Z,sq2FP,j})^{L_x}$$
(8.2.1)
where
$$\lambda _{Z,sq2FP,1}=v^2$$
(8.2.2)
$$\lambda _{Z,sq2FP,2}=v(v+q)$$
(8.2.3)
$$\lambda _{Z,sq2FP,(3,4)}=\frac{v}{2}\left[q+v(v+4)\pm (v^4+4v^3+12v^22qv^2+4qv+q^2)^{1/2}\right]$$
(8.2.4)
$$\lambda _{Z,sq2FP,(5,6)}=\frac{1}{2}(T_{S12}\pm \sqrt{R_{S12}})$$
(8.2.5)
with
$$T_{S12}=v^3+4v^2+3qv+q^2$$
(8.2.6)
$$R_{S12}=v^6+4v^52qv^42q^2v^3+12v^4+16qv^3+13q^2v^2+6q^3v+q^4.$$
(8.2.7)
and
$$c_{Z,sq2FP,1}=c^{(2)}$$
(8.2.8)
$$c_{Z,sq2FP,j}=c^{(1)}\mathrm{for}j=2,3,4$$
(8.2.9)
$$c_{Z,sq2FP,j}=c^{(0)}\mathrm{for}j=5,6.$$
(8.2.10)
Thus, for the individual numbers, $`n_Z(2,0)=2`$, $`n_Z(2,1)=3`$, $`n_Z(2,2)=1`$, and $`n_Z(2,d)=0`$ for $`d3`$, in agreement with Table 3.
As an example of our general results for the structure of the Potts model partition function on Möbius strips of the square lattice, we display the exact solution for width $`L_y=2`$ (with the abbreviation $`sq2Mb`$ for $`sq,L_y=2,FBC_y,TPBC_x`$):
$$Z(sq,2\times L_x,FBC_y,TPBC_x,q,v)=\underset{j=1}{\overset{6}{}}c_{Z,sq2Mb,j}(\lambda _{Z,sq2Mb,j})^{L_x}$$
(8.2.11)
where $`\lambda _{Z,sq2Mb,j}=\lambda _{Z,sq2FP,j}j`$, in accord with the result (3.6) and
$$c_{Z,sq2Mb,j}=c^{(1)}\mathrm{for}j=3,4$$
(8.2.12)
$$c_{Z,sq2Mb,2}=c^{(1)}$$
(8.2.13)
$$c_{Z,sq2Mb,j}=c^{(0)}\mathrm{for}j=5,6$$
(8.2.14)
$$c_{Z,sq2Mb,1}=c^{(0)}$$
(8.2.15)
Thus, $`n_{Z,Mb}(2,0,+)=2`$, $`n_{Z,Mb}(2,0,)=1`$, $`n_{Z,Mb}(2,1,+)=2`$, and $`n_{Z,Mb}(2,1,+)=1`$, in agreement with Table 5.
### 8.3 Other Limits
In the text above we have given results for structural properties of the general Potts model partition function $`Z(G,q,v)`$ and for the special case of the zero-temperature Potts antiferromagnet, $`Z(G,q,v=1)=P(G,q)`$. Another special case is infinite-temperature, i.e., $`v=0`$. In this case, for any graph $`G`$, the Potts model partition function reduces to $`Z(G,q,0)=q^{n(G)}`$, as indicated above in (1.11). For the strip graphs under consideration here, the equality in (1.11) corresponds to the property that at $`v=0`$ all but one of the $`\lambda _{Z,G,j}`$’s in eq. (1.12) vanish, and the nonvanishing $`\lambda _{Z,G,j}`$ necessarily has the value $`q^{L_y}`$ and the coefficient $`c^{(0)}=1`$, so that the right-hand side of (1.12) becomes $`q^{L_yL_x}=q^n`$. One can see this explicitly in the exact solution given above for the Potts model partition function for the $`L_y=2`$ strip of the square lattice.
Secondly, one may consider the zero-temperature limit of the Potts ferromagnet, $`v\mathrm{}`$. This is quite different from the $`T=0`$ limit of the Potts antiferromagnet; among other differences, the ground state entropy per site is zero for the ferromagnet rather than being nonzero as it is for the antiferromagnet with sufficiently large $`q`$. In this limit $`T0`$, the partition function of the Potts ferromagnet grows like $`a^{e(G)}`$ where $`a`$ was defined in (1.3) and $`e(G)`$ was also defined in the introduction as the number of edges (bonds) in $`G`$. Hence, it is convenient to use the reduced partition function $`Z_r`$, defined by
$$Z_r(G,q,v)=a^{e(G)}Z(G,q,v)=u^{e(G)}Z(G,q,v)$$
(8.3.1)
which has the finite limit $`Z_r1`$ as $`T0`$. For the cyclic or Möbius strip graphs of the lattice $`\mathrm{\Lambda }`$ of interest here we thus write
$`Z_r(\mathrm{\Lambda },FBC_y,(T)PBC_x,q,v)`$ $`=`$ $`u^{e(G)}{\displaystyle \underset{j=1}{\overset{N_{Z,G,\lambda }}{}}}c_{Z,G,j}(\lambda _{Z,G,j})^{L_x}`$ (8.3.2)
$``$ $`{\displaystyle \underset{j=1}{\overset{N_{Z,G,\lambda }}{}}}c_{Z,G,j}(\lambda _{Z,G,j,u})^{L_x}`$ (8.3.4)
where as before we use the shorthand notation $`G=\mathrm{\Lambda },FBC_y,(T)PBC_x`$, and
$$\lambda _{Z,G,j,u}=u^{e(G)/L_x}\lambda _{Z,G,j}.$$
(8.3.5)
For example, for cyclic or Möbius strips of the square lattice, $`e=(2L_y1)L_x`$, so the prefactor in (8.3.5) is $`u^{2L_y1}`$, while for the corresponding strips of the triangular lattice, $`e=(3L_y2)L_x`$, so the prefactor is $`u^{3L_y2}`$. For the $`L_y=1`$ circuit graph,
$$\lambda _{Z,C,1,u}=1u,\lambda _{Z,C,2,u}=1+(q1)u.$$
(8.3.6)
so that both of these two $`\lambda _{Z,G,j,u}`$’s remain important in the $`T=0`$ limit. In contrast, for the $`L_y=2`$ cyclic or Möbius strips of the square and triangular lattices, our exact solutions in show that, of the six $`\lambda _{Z,G,j,u}`$’s, two are finite and nonzero for $`u0`$ while the other four vanish. For example, for the cyclic/Möbius strip of the square lattice, in the vicinity of the zero-temperature point $`u=0`$, one has
$$\lambda _{Z,sq2FP,1,u}=u2u^2+u^3$$
(8.3.7)
$$\lambda _{Z,sq2FP,2,u}=u+(q2)u^2+(1q)u^3$$
(8.3.8)
and the Taylor series expansions
$$\lambda _{Z,sq2FP,3,u}=1u^2+2(q2)u^3+O(u^4)$$
(8.3.9)
$$\lambda _{Z,sq2FP,4,u}=u+(q4)u^2+(73q)u^3+O(u^4)$$
(8.3.10)
$$\lambda _{Z,sq2FP,5,u}=1+(q1)u^2\left[1+4u+O(u^2)\right]$$
(8.3.11)
$$\lambda _{Z,sq2FP,6,u}=u+2(q2)u^2+(q^27q+7)u^3+O(u^4).$$
(8.3.12)
Hence,
$$lim_{u0}\frac{\lambda _{Z,sq2FP,j,u}}{\lambda _{Z,sq2FP,3,u}}=lim_v\mathrm{}\frac{\lambda _{Z,sq2FP,j}}{\lambda _{Z,sq2FP,3}}=\{\begin{array}{cc}1\hfill & \text{for }j=5\hfill \\ 0\hfill & \text{for }j=1,2,4,6\hfill \end{array}.$$
(8.3.13)
In this sense, one can define an effective, reduced, $`N_{Z,G,\lambda }`$, and this has the value 2, somewhat analogous to the value $`N_{P,G,\lambda }=4`$ for the $`T=0`$ limit of the Potts antiferromagnet (1.13); however, a difference is that when one reinserts the prefactor in (8.3.4) to get back the actual partition function, all of the six $`\lambda _{Z,G,j}`$’s do contribute.
## 9 Coefficients for Other Strip Graphs
### 9.1 Strips with Torus or Klein Bottle Boundary Conditions
Like the cyclic and Möbius strips, the strips with torus or Klein bottle boundary conditions have periodic or twisted periodic longitudinal boundary conditions, respectively. We first give two theorems for the sum of the coefficients for these types of strip graphs (which are the analogues of the theorems yielding eqs. (3.10) and (6.3) for the cyclic and Möbius strips). We have
Theorem 10. The sum of the coefficients in (1.19) for a strip of the square or triangular lattice of width $`L_y`$ and $`(PBC_y,PBC_x)=`$ torus boundary conditions is (independent of $`L_x`$)
$$C_{P,L_y,torus}=\underset{j=1}{\overset{N_{P,L_y,torus,\lambda }}{}}c_{P,L_y,torus,j}=P(C_{L_y},q)$$
(9.1.1)
where the chromatic polynomial of the circuit graph $`C_n`$ is
$$P(C_n,q)=(q1)^n+(q1)(1)^n.$$
(9.1.2)
Proof. This follows by noting that the transverse slice of the strip graphs of the square and triangular lattices of width $`L_y`$ is the circuit graph $`C_{L_y}`$, and the number of ways of proper colorings of this graph is given by the chromatic polynomial in (9.1.2). $`\mathrm{}`$
Theorem 11. The sum of the coefficients in (1.19) for a strip of the square or triangular lattice of width $`L_y`$ is (independent of $`L_x`$) and $`(PBC_y,TPBC_x)=`$ Klein bottle (KB) boundary conditions is
$$C_{P,L_y,KB}=\underset{j=1}{\overset{N_{P,L_y,KB,\lambda }}{}}c_{P,L_y,KB,j}=0$$
(9.1.3)
Proof. This follows by observing that this sum corresponds to the proper coloring of an effective graph (defined as in our study of Möbius strips) formed from the transverse slice of the strip, and in this case the effective graph always involves at least one loop, i.e. an edge that connects a vertex to itself (specifically, the effective graph involves one loop if $`L_y`$ is odd and two loops if $`L_y`$ is even). But the chromatic polynomial of a graph with a loop vanishes identically. $`\mathrm{}`$
Theorems 10 and 11 generalize the observations made for these sums of coefficients for the exact solutions of the chromatic polynomials of the $`L_y=3`$ torus and Klein bottle graphs of the square lattice in and to arbitrary $`L_y`$ values where these graphs are defined, i.e. $`L_y3`$.
On the basis of the exact solutions for the chromatic polynomials for the width $`L_y=3`$ strips of the square and triangular lattices with $`(PBC_y,PBC_x)=`$ torus and $`(PBC_y,TPBC_x)=`$ Klein bottle boundary conditions, we can make several further remarks. First, since it was found that $`N_{P,G,\lambda }=8`$ for $`G=sq(3\times L_x,PBC_y,PBC_x)`$ but the different value, $`N_{P,G,\lambda }=11`$ for $`tri(3\times L_x,PBC_y,PBC_x)`$, these strips behave fundamentally differently than the cyclic and Möbius strips of these two lattices, for which the value of $`N_{P,G,\lambda }`$ was the same for a given $`L_y`$. Secondly, the exact solutions in show that the coefficients that enter in the chromatic polynomial for the $`L_y=3`$ strips of the square and triangular lattice with torus or Klein bottle boundary conditions are polynomials in $`q`$ but are not of the form $`c^{(d)}=U_{2d}(\frac{\sqrt{q}}{2})`$. Indeed, in contrast to the situation for cyclic and Möbius square-lattice strips and cyclic triangular-lattice strips, the coefficient of degree $`d`$ is not unique, i.e., there can be more than one type of coefficient of a given degree $`d`$ in $`q`$. For example, for the $`L_y=3`$ square-lattice strip with torus boundary conditions, the coefficients that are quadratic functions of $`q`$ are $`q(q3)`$, $`(1/2)q(q3)`$, $`(q1)(q2)`$, and $`(1/2)(q1)(q2)`$. Another difference, again illustrated by the chromatic polynomial for the $`L_y=3`$ strips of the square and triangular lattice is that the coefficients can have zeros outside the range $`0<q<4`$, in contrast to the $`c^{(d)}`$. For example, one term in the chromatic polynomial for the $`L_y=3`$ strip with toroidal boundary conditions is $`(q^36q^2+8q1)(1)^{L_x}`$ for the square lattice and $`(1/3)(q1)(q^25q+3)(2)^{L_x}`$ for the triangular lattice ; the first coefficient has a zero at $`q4.11`$ and the other coefficient has a zero at $`q4.30`$. Given that the exact solutions for the chromatic polynomials on torus graphs in show that the coefficients, which are multiplicities of eigenvalues of the coloring matrix, are not of the form $`c^{(d)}`$ for strip graphs of the square and triangular lattices, it is not clear to us how to relate this with the remark in that certain related eigenvalue multiplicities $`d_r`$ for torus boundary conditions are of the form $`\mathrm{sin}(r\theta /2)/\mathrm{sin}(\theta /p)`$ where $`p`$ is 1 (2) if $`r`$ is even (odd) and $`\theta `$ satisfies (2.39).
The exact solutions for the chromatic polynomial in for the square lattice and in for the triangular lattice showed that if one starts with torus boundary conditions, cuts the tube graph and reglues the ends with reversed orientation to form Klein bottle boundary conditions, the number of $`\lambda _{P,G,j}`$’s does change (is reduced):
$$N_{P,\mathrm{\Lambda },L_y\times L_x,PBC_y,PBC_x}N_{P,\mathrm{\Lambda },L_y\times L_x,PBC_y,TPBC_x}.$$
(9.1.4)
Specifically, it was found that while
$$N_{P,sq,3\times L_x,PBC_y,PBC_x,\lambda }=8,N_{P,tri,3\times L_x,PBC_y,PBC_x,\lambda }=11$$
(9.1.5)
in contrast,
$$N_{P,sq,3\times L_x,PBC_y,TPBC_x,\lambda }=N_{P,tri,3\times L_x,PBC_y,PBC_x,\lambda }=5.$$
(9.1.6)
Thus, at least for this width, if one uses Klein bottle boundary conditions, then the value of $`N_{P,G,\lambda }`$ is the same, namely five, for both square and triangular strips. As is evident from the exact solutions , the five coefficients are the same, although the $`\lambda _{P,G,j}`$’s are different. These coefficients are
$$\{1,q1,(q1),\frac{(q1)(q2)}{2},\frac{q(q3)}{2}\}$$
(9.1.7)
Since the coloring matrix applies directly for the lattice strip graphs with torus boundary conditions, eq. (1.21) applies for the chromatic polynomials of these graphs.
### 9.2 Strips with Free Longitudinal Boundary Conditions
Previously, exact solutions for chromatic polynomials on strip graphs of regular lattices with $`(FBC_y,FBC_x)=`$ open boundary conditions were given in ; in these cases, as discussed before, the coefficients $`c_{P,G,j}`$ are not, in general, polynomials in $`q`$. This can be seen immediately from eqs. (2.14) or (2.19) of . An analogous statement holds for the full Potts model partition functions on these open strips . This property can be understood as a consequence of the fact that the chromatic polynomial is not the trace of an $`m`$’th power of the coloring matrix for these strips, but rather is given by
$$P(G,q)=c|𝒯^m|c^{}$$
(9.2.1)
where $`c`$ and $`c^{}`$ denote colorings of the ends of the strip. A similar statement holds for the full Potts partition function $`Z(G,q,v)`$ in terms of the corresponding matrix $`𝒯_Z`$.
As discussed before , for a strip graph of some lattice with given transverse boundary conditions, the term $`\lambda _{P,G,j}`$ that is dominant in the physical region of $`q`$ (including sufficiently positive integers and denoted as region $`R_1`$ in our previous work ) is independent of the longitudinal boundary conditions. (The same applies for the $`\lambda _{Z,G,j}`$ in $`Z(G,q,v)`$ that is dominant in the physical paramagnetic phase .) In particular, this means that the $`\lambda _{P,G,j}`$ that is dominant in region $`R_1`$ is the same for the cyclic (or Möbius) and open strip graphs of a given lattice type $`\mathrm{\Lambda }`$, thereby establishing a certain connection between these cyclic (or Möbius) and open strip graphs. Similarly, it means that the $`\lambda _{P,G,j}`$ that is dominant in region $`R_1`$ is the same for the torus and Klein bottle strips of this lattice $`\mathrm{\Lambda }`$ on the one hand, and the cylindrical strip graphs on the other. Now it was also shown that for graphs with periodic or twisted periodic boundary conditions, the coefficient of this dominant term is unity. If this term is a root $`r_{\mathrm{},s}`$ of an (irreducible) algebraic equation of degree $`d_{\mathrm{}}`$, then the theorem on symmetric functions of roots of an algebraic equation implies that for the cyclic strips, these roots $`r_{\mathrm{},s}`$ appear in $`P(G,q)`$ in the form of the sum $`_{s=1}^d_{\mathrm{}}r_{\mathrm{},s}^{L_x}`$, i.e., they all have the coefficient unity . This sets an upper bound on $`d_{\mathrm{}}(\mathrm{\Lambda },L_y)`$, namely
$$n_P(L_y,0)d_{\mathrm{}}(\mathrm{\Lambda },L_y)$$
(9.2.2)
where we have included the dependent of $`d_{\mathrm{}}`$ on the lattice type $`\mathrm{\Lambda }`$ and the strip width, and we recall from eq. (3.4) that $`d_{\mathrm{}}(\mathrm{\Lambda },L_y)`$ is the same for the cyclic and Möbius strips of a given lattice $`\mathrm{\Lambda }`$. For cyclic and Möbius strips of the square lattice, $`d_{\mathrm{}}(sq,1)=1`$, $`d_{\mathrm{}}(sq,2)=1`$ , $`d_{\mathrm{}}(sq,3)=2`$ , and $`d_{\mathrm{}}(sq,4)=3`$ , while for the cyclic and Möbius strips of the triangular lattice, $`d_{\mathrm{}}(tri,2)=1`$ and $`d_{\mathrm{}}(tri,3)=2`$, $`d_{\mathrm{}}(tri,4)=4`$ . This information is summarized in tables given in along with related results. Our calculations above yield, for the relevant values of $`n_P(L_y,0)`$ (see Table 1): $`n_P(1,0)=n_P(2,0)=1`$, $`n_P(3,0)=2`$, and $`n_P(4,0)=4`$. Hence, evidently, the inequality (9.2.2) is realized as an equality for the triangular strips and the square strips for $`L_y=1,2,3`$, but as a strict inequality for the $`L_y=4`$ square strips. Further, for the cases that we studied, we found that
$$d_{\mathrm{}}(\mathrm{\Lambda },L_y)_{cyc.,Mb}=d_{\mathrm{}}(\mathrm{\Lambda },L_y)_{open}$$
(9.2.3)
where $`d_{\mathrm{}}(\mathrm{\Lambda },L_y)_{open}d_{\mathrm{}}(\mathrm{\Lambda },FBC_y,FBC_x,L_y)`$. When we obtained the results in the present paper for $`n_P(L_y,d)`$, we noticed the interesting connection that for the cyclic and Möbius strips of the triangular lattice, for which (9.2.2) holds as an equality, it is also true that one has the fourfold equality
$$n_P(L_y,0)=d_{\mathrm{}}(tri,L_y)_{cyc.,Mb.}=d_{\mathrm{}}(tri,L_y)_{open}=N_{P,tri(L_y),open,\lambda }$$
(9.2.4)
i.e., the number of $`\lambda _{P,G,j}`$’s with coefficient $`c^{(0)}=1`$ in the chromatic polynomial for the cyclic or Möbius strips of the triangular lattice, which was equal to the degree of the algebraic equation one of whose roots was the dominant $`\lambda _{P,G,j}`$ in $`R_1`$, was also equal to the total number of terms in the chromatic polynomial for the open strip of the triangular lattice with the same width. We confirmed eq. (9.2.4) for widths $`L_y=2,3,4`$, based on the exact solutions in . In terms of the generating function formalism of , the denominator $`𝒟`$ of the generating function, as a polynomial in the auxiliary variable, has the maximal degree, $`n_P(L_y,0)`$, allowed by (9.2.2) for these open strips of the triangular lattice but does not, in general, for the open strips of the square lattice. As we noted in , this difference in behavior is related to the fact that the Möbius strips of the square and triangular lattices are different; in particular, the Möbius strip of the triangular lattice involves algebraic nonpolynomial coefficients. Further, even for the strips of the square lattice, where the first two terms of (9.2.4) are not, in general, equal (as our solutions for chromatic polynomials for $`L_y=4`$ showed), we still found
$$d_{\mathrm{}}(sq,L_y)_{cyc.,Mb.}=N_{P,sq(L_y),open,\lambda }.$$
(9.2.5)
That is, we find that the equality (9.2.5) holds for cyclic or Möbius and open strips of the square lattice for $`1L_y4`$, as well as for the corresponding strips of the triangular lattice with $`2L_y4`$; in particular, it holds even when $`n_P(L_y,0)`$ is not equal to (is greater than) $`d_{\mathrm{}}(\mathrm{\Lambda },L_y)`$. For reference, for open strips of the triangular lattice with $`L_y2`$, one can infer from calculations in that $`N_{Z,tri(L_y),open,\lambda }=C_{L_y}`$; furthermore, $`N_{P,tri(L_y),open,\lambda }=M_{L_y1}`$ . Another generalization would be that for triangular strips, the value of the degree in eq. (9.2.3) satisfies $`d_{\mathrm{}}(tri,L_y)_{cyc.,Mb}=d_{\mathrm{}}(tri,L_y)_{open}=M_{L_y1}`$. This would require that, in terms of the generating function , the denominator $`𝒟`$ does not factorize.
Previously, coloring matrix methods have been applied to the honeycomb ($`hc`$) lattice to derive rigorous bounds on $`W`$ . In the present context of strips, we observe that, in addition to the elementary result $`N_{P,hc,L_y=2,open,\lambda }=d_{\mathrm{}}(hc,L_y=2)_{open}=1`$, our exact solutions have yielded the results
$$N_{P,hc,L_y=3,open\lambda }=d_{\mathrm{}}(hc,L_y=3)_{open}=3$$
(9.2.6)
from and
$$N_{P,hc,L_y=2,cyc.,Mb.,\lambda }=d_{\mathrm{}}(hc,L_y=2)_{cyc.,Mb}=1$$
(9.2.7)
from . In view of eq. (9.2.6), for this strip of the honeycomb lattice, $`n_P(L_y=3,0)d_{\mathrm{}}(hc,L_y=3)_{open}`$, in contrast to the case with the square and triangular lattice strips.
### 9.3 Strips with Cylindrical Boundary Conditions
Exact solutions for chromatic polynomials for cylindrical strips, i.e. with boundary conditions of the form $`(PBC_y,FBC_x)`$ of the square and triangular lattices were presented in , and, again, the coefficients are, in general, not polynomials in $`q`$, for the same reason as in the case of the open strips. And again, the same statement holds for the full Potts model partition function. (Recent calculations of chromatic polynomials for wider strips with open and cylindrical boundary conditions include .)
## 10 Some General Geometric Identities
In this section we present some useful identities between Potts model partition functions on different types of lattice strips that follow from basic geometrical considerations. We have applied these as checks in our previous calculations. Consider a family of strip graphs of the square or triangular lattice with fixed width $`L_y`$ and arbitrary length $`L_x`$ with some set of longitudinal ($`x`$) and transverse ($`y`$) boundary conditions, denoted as $`BC=(BC_y,BC_x)`$. By the transfer matrix argument given in , it follows that the Potts model partition function for this strip is of the form
$$Z(\mathrm{\Lambda },L_y\times L_x,BC,q,v)=\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_y,BC}}{}}c_{Z,\mathrm{\Lambda },L_y,BC,j}(\lambda _{Z,\mathrm{\Lambda },L_y,BC,j})^{L_x}.$$
(10.1)
We recall that the total number of terms, $`N_{Z,\mathrm{\Lambda },L_y,BC}`$, the coefficients $`c_{Z,\mathrm{\Lambda },L_y,BC,j}`$, and the $`\lambda _{Z,\mathrm{\Lambda },L_y,BC,j}`$’s are independent of $`L_x`$. In our previous works, we have studied the $`L_x\mathrm{}`$ behavior of these Potts model partition functions. However, clearly, for fixed $`L_y`$ and $`L_x`$, by switching the $`x`$ and $`y`$ axes, one can equivalently view the strip as the length–$`L_y`$ member of a family of strip graphs with fixed width $`L_x`$ and length $`L_y`$.<sup>5</sup><sup>5</sup>5 Note that for other lattices such simple relations do not, in general, hold. For example, if one considers strips of the honeycomb lattice, constructed as a brick lattice, then rotating a open horizontal $`L_y\times L_x`$ strip of bricks by $`90^{}`$, one does not get an equivalent $`L_x\times L_y`$ strip of horizontally oriented bricks. This yields a number of useful identities. We proceed to describe these
### 10.1 Strips with $`(FBC_y,FBC_x)`$
For strips of the square and triangular lattices with $`(FBC_y,FBC_x)`$ (open) boundary conditions, the interchange of the $`x`$ and $`y`$ axes leaves the boundary conditions the same, and one gets the following identity (where we write $`F_yF_x(FBC_y,FBC_x)`$ to save space)
$$\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_y,F_y,F_x}}{}}c_{Z,\mathrm{\Lambda },L_y,F_y,F_x,j}(\lambda _{Z,\mathrm{\Lambda },L_y,F_y,F_x,j})^{L_x}=\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_x,F_x,F_y}}{}}c_{Z,\mathrm{\Lambda },L_x,F_x,F_y,j}(\lambda _{Z,\mathrm{\Lambda },L_x,F_x,F_y,j})^{L_y}.$$
(10.1.1)
Taking $`T=0`$ for the Potts antiferromagnet, one obtains the corresponding identity for chromatic polynomials
$$\underset{j=1}{\overset{N_{P,\mathrm{\Lambda },L_y,F_y,F_x}}{}}c_{P,\mathrm{\Lambda },L_y,F_y,F_x,j}(\lambda _{P,\mathrm{\Lambda },L_y,F_y,F_x,j})^{L_x}=\underset{j=1}{\overset{N_{P,\mathrm{\Lambda },L_x,F_y,F_x}}{}}c_{P,\mathrm{\Lambda },L_x,F_x,F_y,j}(\lambda _{P,\mathrm{\Lambda },L_x,F_x,F_y,j})^{L_y}.$$
(10.1.2)
Using the results in , one can construct various illustrations of this. One of the simplest is to take the open strip of the square lattice with width $`L_y=2`$, for which $`N_{P,sq,2,FBC_y,FBC_x}=1`$ and $`P(sq,2\times L_x,FBC_y,FBC_x,q)=q(q1)(q^23q+3)^{L_x1}`$. For $`L_x=3`$, this chromatic polynomial must be equal to that for the width $`L_y=3`$ open strip (which has $`N_{P,sq,3,FBC_y,FBC_x}=2`$) with length $`L_x=2`$, for which (using and eq. (2.15) of ) for general $`L_x`$,
$`P(sq,3\times L_x,FBC_y,FBC_y,q)={\displaystyle \frac{1}{(\lambda _{sq3FF,1}\lambda _{sq3FF,2})}}\times `$ (10.1.3)
(10.1.4)
$`\left[(A_{sq3FF,0}\lambda _{sq3FF,1}+A_{sq3FF,1})(\lambda _{sq3FF,1})^{L_x2}(A_{sq3FF,0}\lambda _{sq3FF,2}+A_{sq3FF,1})(\lambda _{sq3FF,2})^{L_x2}\right]`$ (10.1.5)
(10.1.6)
(10.1.7)
where
$$A_{sq3FF,0}=q(q1)(q^23q+3)^2$$
(10.1.8)
$$A_{sq3FF,1}=q(q1)^3(q^36q^2+13q11)$$
(10.1.9)
and
$$\lambda _{sq3FF,(1,2)}=\frac{1}{2}\left[(q2)(q^23q+5)\pm \left[(q^25q+7)(q^45q^3+11q^212q+8)\right]^{1/2}\right].$$
(10.1.10)
### 10.2 Cyclic and Cylindrical Strips
We have denoted strips of fixed width $`L_y`$ and arbitrary length $`L_x`$ vertices with $`(FBC_y,PBC_x)`$ and $`(PBC_y,FBC_x)`$ boundary conditions as cyclic and cylindrical, respectively. Topologically, these are the same; the distinction was made because we were, in particular, interested in the large $`L_x\mathrm{}`$ limit, which, for cyclic strips means that the length of a circuit goes to infinity, while for the cylindrical strips, the length of the circumference of the cylinder is fixed while its length goes to infinity. Clearly, a cyclic strip of the square or triangular lattice with a fixed width $`L_y`$ and length $`L_x`$ is the same as a cylindrical strip of the given lattice with width $`L_x`$ and length $`L_y`$. Letting $`\mathrm{\Lambda }`$ denote the type of lattice, square or triangular as before, and writing
$$Z(\mathrm{\Lambda },L_y\times L_x,FBC_y,PBC_x)=\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_y,F_y,P_x}}{}}c_{Z,\mathrm{\Lambda },L_y,F_y,P_x,j}(\lambda _{Z,\mathrm{\Lambda },L_y,F_y,P_x,j})^{L_x}$$
(10.2.1)
we have the identity
$$\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_y,F_y,P_x}}{}}c_{Z,\mathrm{\Lambda },L_y,F_y,P_x,j}(\lambda _{Z,\mathrm{\Lambda },L_y,F_y,P_x,j})^{L_x}=\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_x,P_y,F_x}}{}}c_{Z,\mathrm{\Lambda },L_x,P_y,F_x,j}(\lambda _{Z,\mathrm{\Lambda },L_x,P_y,F_x,j})^{L_y}.$$
(10.2.2)
Again, for the special case of the $`T=0`$ Potts antiferromagnet, we have the resultant identity for the chromatic polynomial
$$\underset{j=1}{\overset{N_{P,\mathrm{\Lambda },L_y,F_y,P_x}}{}}c_{P,\mathrm{\Lambda },L_y,F_y,P_x,j}(\lambda _{P,\mathrm{\Lambda },L_y,F_y,P_x,j})^{L_x}=\underset{j=1}{\overset{N_{P,\mathrm{\Lambda },L_x,P_y,F_x}}{}}c_{P,\mathrm{\Lambda },L_x,P_y,F_x,j}(\lambda _{P,\mathrm{\Lambda },L_x,P_y,F_x,j})^{L_y}.$$
(10.2.3)
Perhaps the simplest illustration of this general identity is the relation between the chromatic polynomial for the width $`L_y=2`$ cyclic strip of the square lattice and the width $`L_y=3`$ cylindrical strip of this lattice. For the $`L_y=2`$ cyclic strip, one has $`N_{P,sq,2\times L_x,F_y,P_x}=4`$ and
$$P(sq,2\times L_x,FBC_y,PBC_x,q)=(q^23q+1)+(q1)\left[(3q)^{L_x}+(1q)^{L_x}\right]+(q^23q+3)^{L_x}.$$
(10.2.4)
For the $`L_y=3`$ cylindrical strip, one has
$$P(sq,3\times L_x,PBC_y,FBC_x,q)=q(q1)(q2)(q^36q^2+14q13)^{L_x1}.$$
(10.2.5)
Thus, the identity (10.2.3) yields
$`(q^23q+1)+(q1)\left[(3q)^3+(1q)^3\right]+(q^23q+3)^3=`$ (10.2.6)
(10.2.7)
$`q(q1)(q2)(q^36q^2+14q13).`$ (10.2.8)
While this illustration is quite simple, more complicated cases involve identities between powers of algebraic roots of different types.
### 10.3 Strips with Torus Boundary Conditions
For a strip graph of the square or triangular lattice with $`(PBC_y,PBC_x)(P_y,P_x)`$ (torus) boundary conditions, the identities read
$$\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_y,P_y,P_x}}{}}c_{Z,\mathrm{\Lambda },L_y,P_y,P_x,j}(\lambda _{Z,\mathrm{\Lambda },L_y,P_y,P_x,j})^{L_x}=\underset{j=1}{\overset{N_{Z,\mathrm{\Lambda },L_x,P_x,P_y}}{}}c_{Z,\mathrm{\Lambda },L_x,P_x,P_y,j}(\lambda _{Z,\mathrm{\Lambda },L_x,P_x,P_y,j})^{L_y}$$
(10.3.1)
$$\underset{j=1}{\overset{N_{P,\mathrm{\Lambda },L_y,P_y,P_x}}{}}c_{P,\mathrm{\Lambda },L_y,P_y,P_x,j}(\lambda _{P,\mathrm{\Lambda },L_y,P_y,P_x,j})^{L_x}=\underset{j=1}{\overset{N_{P,\mathrm{\Lambda },L_x,P_x,P_y}}{}}c_{P,\mathrm{\Lambda },L_x,P_x,P_y,j}(\lambda _{P,\mathrm{\Lambda },L_x,P_x,P_y,j})^{L_y}$$
(10.3.2)
The simplest illustration of this is for the family of strips of the square lattice with width $`L_y=2`$, for which the torus degenerates to a cyclic strip, and one has $`P(sq,2\times L_x,PBC_y,PBC_x,q)=P(sq,2\times L_x,FBC_y,PBC_x,q)`$, given in (10.2.4). Let us next consider the strip of the square lattice with width $`L_y=3`$ and length $`L_x=2`$. This graph is related by the identity to two other strip graphs: (i) the cyclic strip of the square lattice with $`L_y=2`$ and $`L_x=3`$ (see eq. (10.2.4)) and (ii) the equivalent cylindrical strip of the square lattice with $`L_y=3`$ and $`L_x=2`$ (see eq. (10.2.5)). Indeed, when one evaluates the exact solution given in for $`L_x=2`$, one finds that it is equal to the right-hand side of eq. (10.2.8).
## 11 Conclusions
In this paper, using a general formula, in terms of Chebyshev polynomials of the second kind, for the coefficients that occur in the $`q`$-state Potts model partition function on cyclic strips of the square and triangular lattices and on Möbius strips of the square lattice, together with general formulas for sums of these coefficients, we have determined the number, $`n_Z(L_y,d)`$, of $`\lambda _{Z,G,j}`$’s with coefficient $`c^{(d)}`$ in $`Z(G,q,v)`$ and the total number, $`N_{Z,G,\lambda }`$, for these strips. Results are also given for the analogous numbers $`n_P(L_y,d)`$ and $`N_{P,L_y,\lambda }`$ for the zero-temperature Potts antiferromagnet partition functions, i.e., chromatic polynomials for these strips. The results for the total numbers, $`N_{Z,L_y,\lambda }`$ and $`N_{P,L_y,\lambda }`$, apply for both cyclic and Möbius strips of both the square and triangular lattices. Among other connections, we found that $`n_P(L_y,0)=n_P(L_y1,1)=M_{L_y1}`$, the Motzkin number; $`n_Z(L_y,0)=C_{L_y}`$, the Catalan number; the exact expression for $`N_{P,L_y,\lambda }`$ in eq. (5.1.2); and the relations $`N_{P,L_y,\lambda }=2N_{DA,sq,L_y}`$; and $`N_{Z,L_y,\lambda }=2N_{DA,tri,L_y}`$, where $`N_{DA,\mathrm{\Lambda },n}`$ denotes the number of directed lattice animals on the lattice $`\mathrm{\Lambda }`$. We also found the asymptotic growths $`N_{Z,L_y,\lambda }L_y^{1/2}4^{L_y}`$ and $`N_{P,L_y,\lambda }L_y^{1/2}3^{L_y}`$ as $`L_y\mathrm{}`$. Some commens about other lattice strips were made. In addition, we presented some useful general geometric identities for Potts model partition functions.
Acknowledgment: We thank Lee-Peng Teo for helpful discussions. The research of R. S. was supported in part by the NSF grant PHY-9722101 and at Brookhaven by the DOE contract DE-AC02-98CH10886.<sup>6</sup><sup>6</sup>6Accordingly, the U.S. government retains a non-exclusive royalty-free license to publish or reproduce the published form of this contribution or to allow others to do so for U.S. government purposes.
## 12 Appendix
### 12.1 Connection Between Potts Model Partition Function and Tutte Polynomial
The Potts model partition function $`Z(G,q,v)`$ is related to the Tutte polynomial $`T(G,x,y)`$ as follows. The graph $`G`$ has vertex set $`V`$ and edge set $`E`$, denoted $`G=(V,E)`$. A spanning subgraph $`G^{}`$ is defined as a subgraph that has the same vertex set and a subset of the edge set: $`G^{}=(V,E^{})`$ with $`E^{}E`$. The Tutte polynomial of $`G`$, $`T(G,x,y)`$, is then given by -
$$T(G,x,y)=\underset{G^{}G}{}(x1)^{k(G^{})k(G)}(y1)^{c(G^{})}$$
(12.1.1)
where $`k(G^{})`$, $`e(G^{})`$, and $`n(G^{})=n(G)`$ denote the number of components, edges, and vertices of $`G^{}`$, and
$$c(G^{})=e(G^{})+k(G^{})n(G^{})$$
(12.1.2)
is the number of independent circuits in $`G^{}`$ (sometimes called the co-rank of $`G^{}`$). Note that the first factor can also be written as $`(x1)^{r(G)r(G^{})}`$, where
$$r(G)=n(G)k(G)$$
(12.1.3)
is called the rank of $`G`$. The graphs $`G`$ that we consider here are connected, so that $`k(G)=1`$. Now let
$$x=1+\frac{q}{v}$$
(12.1.4)
and
$$y=a=v+1$$
(12.1.5)
so that $`q=(x1)(y1)=(x1)v`$. Then
$$Z(G,q,v)=(x1)^{k(G)}(y1)^{n(G)}T(G,x,y).$$
(12.1.6)
There is also a connection with the Whitney rank polynomial, $`R(G,\xi ,\eta )`$, defined as
$$R(G,\xi ,\eta )=\underset{G^{}G}{}\xi ^{r(G^{})}\eta ^{c(G^{})}$$
(12.1.7)
where the sum is again over spanning subgraphs $`G^{}`$ of $`G`$. Then
$$T(G,x,y)=(x1)^{r(G)}R(G,\xi =(x1)^1,\eta =y1)$$
(12.1.8)
and
$$Z(G,q,v)=q^{n(G)}R(G,\xi =\frac{v}{q},\eta =v).$$
(12.1.9)
Note that the chromatic polynomial is a special case of the Tutte polynomial:
$$P(G,q)=q^{k(G)}(1)^{k(G)+n(G)}T(G,x=1q,y=0)$$
(12.1.10)
(recall eq. (1.9)).
For a recursive family of graphs, such as the strip graphs considered in this paper, comprised of $`m`$ repetitions of a basic subgraph, the Tutte polynomial has the form
$$T(G,x,y)=\underset{j=1}{\overset{N_{T,G,\lambda }}{}}c_{T,G,j}(\lambda _{T,G,j})^m$$
(12.1.11)
where, from (12.1.6), one has the relation
$$\lambda _{T,G,j}=v^{L_y}\lambda _{Z,G,j}$$
(12.1.12)
so that
$$N_{T,G,\lambda }=N_{Z,G,\lambda }.$$
(12.1.13)
It is convenient to extract a common factor from the coefficients:
$$c_{T,G,j}\frac{\overline{c}_{T,G,j}}{x1}.$$
(12.1.14)
Of course, although the individual terms contributing to the Tutte polynomial are thus rational functions of $`x`$ rather than polynomials in $`x`$, the full Tutte polynomial is a polynomial in both $`x`$ and $`y`$. Given the relation (12.1.6), if one defines $`n_T(L_y,d)`$ as the number of terms $`\lambda _{T,L_y,j}`$ in $`T(L_y,q,v)`$ that have as their reduced coefficient $`\overline{c}_{T,L_y,j}=c^{(d)}`$, then these are the same numbers:
$$n_T(L_y,d)=n_Z(L_y,d).$$
(12.1.15)
Thus, Tables 3 and 5 apply equally to the structure of the Tutte polynomials for the cyclic and Möbius strips of the square and triangular lattices.
For a given graph $`G=(V,E)`$, at certain special values of the arguments $`x`$ and $`y`$, the Tutte polynomial $`T(G,x,y)`$ yields quantities of basic graph-theoretic interest -. We recall some definitions: a spanning subgraph was defined at the beginning of the paper; a tree is a connected graph with no cycles; a forest is a graph containing one or more trees; and a spanning tree is a spanning subgraph that is a tree. We recall that the graphs $`G`$ that we consider are connected. Then the number of spanning trees of $`G`$, $`N_{ST}(G)`$, is
$$N_{ST}(G)=T(G,1,1),$$
(12.1.16)
the number of spanning forests of $`G`$, $`N_{SF}(G)`$, is
$$N_{SF}(G)=T(G,2,1),$$
(12.1.17)
the number of connected spanning subgraphs of $`G`$, $`N_{CSSG}(G)`$, is
$$N_{CSSG}(G)=T(G,1,2),$$
(12.1.18)
and the number of spanning subgraphs of $`G`$, $`N_{SSG}(G)`$, is
$$N_{SSG}(G)=T(G,2,2).$$
(12.1.19)
These connections have been used in ,.
### 12.2 Determinants of Coloring Matrices
In this section we list the determinants of coloring matrices for some lattice strips. For the chromatic polynomial this is
$$detT_P(G)=\underset{j=1}{\overset{𝒩}{}}\lambda _{P,G,j}=\underset{j=1}{\overset{N_{P,G,\lambda }}{}}(\lambda _{P,G,j})^{c_{P,G,j}}.$$
(12.2.1)
For the Potts model partition function the analogous determinant is
$$detT_Z(G)=\underset{j=1}{\overset{𝒩_𝒵}{}}\lambda _{Z,G,j}=\underset{j=1}{\overset{N_{Z,G,\lambda }}{}}(\lambda _{Z,G,j})^{c_{Z,G,j}}.$$
(12.2.2)
Let us define the shorthand notation
$$D_P(G)=detT_P(G),D_Z(G)=detT_Z(G).$$
(12.2.3)
For the products of eigenvalues contributing to (1.13) for the chromatic polynomials and (1.12) for the Potts model partition function for cyclic and Möbius strips of the square lattice we find
$$D_P(sq,L_y=1,FBC_y,PBC_x)=D_P(\{C\},q)=(1)^{q1}(q1)$$
(12.2.4)
(where $`\{C\}`$ refers to the circuit graph) For the $`L_y=2`$ cyclic and Möbius strips of the square lattice, from ,
$$D_P(sq,L_y=2,FBC_y,PBC_x)=\left[(3q)(1q)\right]^{q1}(q^23q+3)$$
(12.2.5)
and
$$D_P(sq,L_y=2,FBC_y,TPBC_x)=(3q)^{q1}(1q)^{(q1)}(q^23q+3).$$
(12.2.6)
For the $`L_y=3`$ cyclic and Möbius strips of the square lattice, from ,
$`D_P(sq,L_y=3,FBC_y,PBC_x)=(1)^{(q2)(q^23q+1)}(q1)^{(q1)^2}(q2)^{q^2q1}\times `$ (12.2.7)
(12.2.8)
$`(q4)^{q^23q+1}(q^36q^2+13q11)(q^49q^3+29q^240q+22)^{q1}`$ (12.2.9)
and
$`D_P(sq,L_y=3,FBC_y,TPBC_x)=(1)^{q^24q+2}(q1)^{q1}(q2)^{(2q3)}\times `$ (12.2.10)
(12.2.11)
$`(q4)^1(q^36q^2+13q11)(q^49q^3+29q^240q+22)^{q1}.`$ (12.2.12)
The corresponding determinants for $`L_y=4`$ can be obtained from but are quite lengthy, so we do not list them here.
For the $`D_Z`$ determinants we have
$$D_Z(sq,L_y=1,FBC_y,PBC_x)=D_Z(\{C\},q)=v^{q1}(q+v)$$
(12.2.13)
and, from ,
$$D_Z(sq,L_y=2,FBC_y,PBC_x)=v^{2q(q1)}(v+1)^q(v+q)^{2q}$$
(12.2.14)
and
$$D_Z(sq,L_y=2,FBC_y,TPBC_x)=v^{2(q1)}(v+1)^q(v+q)^2.$$
(12.2.15)
Note that, except for the case of the circuit graph, $`sq,L_y=1,PBC_x`$, for the value $`v=1`$, where the Potts model partition function reduces to the chromatic polynomial, $`Z(G,q,v=1)=P(G,q)`$, some eigenvalues in the product contributing to $`Z`$ vanish, and hence this product vanishes. If one extracts these vanishing eigenvalues, then, of course, the rest yield the same product as for $`P(G,q)`$.
For cyclic strips of the triangular lattice we have, from
$$D_P(tri,L_y=2,FBC_y,PBC_x)=(q2)^{2q}$$
(12.2.16)
and from ,
$$D_P(tri,L_y=3,FBC_y,PBC_x)=(1)^{(q2)(q^23q+1)}(q2)^{q(2q1)}(q3)^{q(q1)}$$
(12.2.17)
$$D_Z(tri,L_y=2,FBC_y,PBC_x)=v^{2q(q1)}(v+1)^{2q}(v+q)^{2q}.$$
(12.2.18)
The determinant $`D_P(tri,L_y=4,FBC_y,PBC_x)`$ can be calculated from our exact solution in ; however, it is sufficiently lengthy that we do not include it here. The analogous determinants for the Möbius strips of the triangular lattice can also be calculated from the exact solutions that we have given ; however, they are more complicated, since some coefficients are algebraic, rather than polynomial, functions of $`q`$.
For the strip graphs of the square and triangular lattices with $`(PBC_y,PBC_x)=`$ torus and $`(PBC_y,TPBC_x)=`$ Klein bottle boundary conditions, we have, from the exact solutions in and ,
$`D_P(sq,L_y=3,PBC_y,PBC_x)=(1)^{q^36q^2+11q4}(q1)^{\frac{(q1)(q2)}{2}}\times `$ (12.2.19)
(12.2.20)
$`(q2)^{q^2+q4}(q4)^{(q1)(q2)}(q5)^{\frac{q(q3)}{2}}(q^27q+13)^{q1}\times `$ (12.2.21)
(12.2.22)
$`(q^36q^2+14q13)`$ (12.2.23)
and
$`D_P(sq,L_y=3,PBC_y,TPBC_x)`$ $`=`$ $`(q1)^{\frac{(q1)(q2)}{2}}(q5)^{\frac{q(q3)}{2}}(q^27q+13)^{q1}\times `$ (12.2.26)
$`(q^36q^2+14q13).`$
### 12.3 Some Other Coloring Matrix Results for Lattice Strip Graphs
Since the trace of a coloring matrix vanishes (recall (1.21)), the sum of eigenvalues, each multiplied by its multiplicity also vanishes, as indicated in eq. (1.22), for the cyclic and torus strips, for which (1.17) applies directly. Using coloring methods, we obtain the following general formulas for sums for other types of strips. These agree with our previous exact solutions in .
$$\underset{j=1}{\overset{N_{P,sq(L_y),Mb,\lambda }}{}}c_{P,sq(L_y),Mb,j}\lambda _{P,sq(L_y),Mb,j}=\{\begin{array}{cc}q(q1)(q^23q+3)^{\frac{L_y}{2}1}\hfill & \text{for even }L_y\hfill \\ 0\hfill & \text{for odd }L_y\hfill \end{array}$$
(12.3.1)
$$\underset{j=1}{\overset{N_{P,tri(L_y),Mb,\lambda }}{}}c_{P,tri(L_y),Mb,j}\lambda _{P,tri(L_y),Mb,j}=0L_y2.$$
(12.3.2)
The equations also apply to the corresponding strips with Klein bottle (KB) boundary conditions, with $`MbKB`$, and $`L_y3`$. The polynomial $`q^23q+3`$ in (12.3.1) is $`D_4`$ in our previous notation, where
$$D_k=\underset{s=0}{\overset{k2}{}}(1)^s\left(\genfrac{}{}{0pt}{}{k1}{s}\right)q^{k2s}.$$
(12.3.3)
From the exact solutions for the chromatic polynomials in , we have, for the cyclic strip of the kagomé lattice consisting of a succession of hexagons with two interleaved triangles per hexagon of arbitrary length (denoted as having width $`L_y=2`$ in ) we have
$$\underset{j=1}{\overset{6}{}}c_{P,kag2,j}\lambda _{P,kag2,j}=q(q1)^2(q2)^2.$$
(12.3.4)
Additional families of graphs are provided by the homeomorphic expansions of the cyclic and Möbius $`L_y=2`$ strips of the square lattice, in which one adds $`k2`$ degree-2 vertices to the upper and lower horizontal edges of each square, for $`k3`$. These may be viewed as strips of $`p`$-sided polygons, with $`p=2k`$, with each successive pair of polygons sharing one edge. For $`k=2`$ and $`k=3`$, these families are the $`L_y=2`$ cyclic and Möbius strips of the square and honeycomb lattice, respectively. (In the latter case, the honeycomb lattice is represented as a brick lattice with the bricks oriented horizontally.) Following , we denote the cyclic (cyc) and Möbius (Mb) strips of this type, of length $`L_x=m`$ as $`(Ch)_{k,m,cyc}`$ and $`(Ch)_{k,m,Mb}`$, where $`Ch`$ denotes “chain”. Exact solutions for the chromatic polynomials of these families were given in . From the construction of these graphs, it is clear that $`_jc_{P,(Ch)_{k,cyc,j}}=q(q1)`$ and $`_jc_{P,(Ch)_{k,Mb,j}}=0`$ hold, as in eqs. (3.10) and (6.3). For the other sums (with $`\lambda _{(Ch),k,cyc,j}=\lambda _{(Ch),k,Mb,j}`$), we have
$$\underset{j=1}{\overset{4}{}}c_{(Ch)_{k,cyc,j}}\lambda _{(Ch)_{k,cyc,j}}=q^23q+1+2(1)^{k+1}(q1)D_{k+1}+D_{2k}$$
(12.3.5)
where $`D_k`$ was defined in (12.3.3), and
$$\underset{j=1}{\overset{4}{}}c_{(Ch)_{k,Mb,j}}\lambda _{(Ch)_{k,Mb,j}}=1+2(1)^k(q1)D_k+D_{2k}.$$
(12.3.6)
For $`k=2`$, the sum (12.3.5) is zero, as in (1.21), and the sum (12.3.6) has the value $`q(q1)`$ as in the $`L_y=2`$ special case of (12.3.1). As an example of the values for higher–$`p`$ polygonal strips, for $`k=3`$, i.e., the strip of the honeycomb lattice, (12.3.5) has the value $`q(q1)^3`$ and (12.3.6) has the value $`q(q1)(q2)^2`$.
For the full Potts model partition functions, from our exact solutions in we find
$$\underset{j=1}{\overset{6}{}}c_{Z,sq2,cyc,j}\lambda _{Z,sq2,cyc,j}=q(v+1)^2(v+q)$$
(12.3.7)
$$\underset{j=1}{\overset{6}{}}c_{Z,sq2,Mb,j}\lambda _{Z,sq2,Mb,j}=q(v^3+3v^2+3v+q)$$
(12.3.8)
$$\underset{j=1}{\overset{6}{}}c_{Z,tri2,cyc,j}\lambda _{Z,tri2,cyc,j}=q(v+1)^2(v^2+2v+q)$$
(12.3.9)
$$\underset{j=1}{\overset{6}{}}c_{Z,tri2,Mb,j}\lambda _{Z,tri2,Mb,j}=q(v+1)(v^3+3v^2+3v+q).$$
(12.3.10)
As is evident, for the special case $`v=1`$ where the Potts model partition function reduces to the chromatic polynomial, these equations reduce to their analogues for the respective chromatic polynomials.
In terms of the Tutte polynomials, these formulas read
$$\underset{j=1}{\overset{6}{}}c_{T,sq2,cyc,j}\lambda _{T,sq2,cyc,j}=xy^2$$
(12.3.11)
$$\underset{j=1}{\overset{6}{}}c_{T,sq2,Mb,j}\lambda _{T,sq2,Mb,j}=x+y+y^2$$
(12.3.12)
(which is the Tutte polynomial for the graph known as the “thick link” with three edges, the planar dual to the circuit graph $`C_3`$)
$$\underset{j=1}{\overset{6}{}}c_{T,tri2,cyc,j}\lambda _{T,tri2,cyc,j}=(x+y)y^2$$
(12.3.13)
and
$$\underset{j=1}{\overset{6}{}}c_{T,tri2,Mb,j}\lambda _{T,tri2,Mb,j}=y(x+y+y^2).$$
(12.3.14)
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# Géométrie de l’espace tangent sur l’hyperboloïde quantique
## 1 introduction
Il est bien connu depuis les années 60 que toutes les constructions principales de la géométrie classique se généralisent dans le cadre de la superthéorie. Dans les années 80, plusieurs notions et constructions de la géométrie classique et de la superthéorie s’étendent au cas de la géométrie tressée (“braided geometry”). En particulier il existe une généralisation naturelle de la notion d’algèbre commutative (ou super-commutative) liée aux opérateurs de tresses involutifs (appelés aussi symétries). En effet soit $`A`$ une algèbre associative de produit noté $``$ et munie d’un opérateur de tresse<sup>1</sup><sup>1</sup>1Nous appelons opérateur de tresse (ou un tressage) une solution de l’équation de Yang-Baxter quantique (EYBQ) :
$$S^{12}S^{23}S^{12}=S^{23}S^{12}S^{23}$$
$`S^{12}=Sid,S^{23}=idS`$. $`S`$ involutif (i.e. $`S^2=id`$)
$$S:A^2A^2,$$
$`A`$ est appelée une algèbre S-commutative si :
1. $`=S,`$
2. $`S^{12}=^{23}S^{12}S^{23}`$.
Dans ce cas, certains aspects du calcul différentiel tressé ont été développés dans \[GRR\]. Plus particulièrement la notion de champ de vecteurs se généralise (dans l’esprit de la superthéorie) de la façon suivante : une application linéaire $`X:AA`$ est appelée un champ de vecteurs tressés si elle vérifie la $`S`$-analogue de la règle de Leibniz
$$X(fg)=X(f)g+ev^{23}S^{12}(Xfg)$$
(1)
ev est l’application d’évaluation $`XfX(f)`$. Nous supposons ici que $`S`$ peut être étendu à un tressage (noté encore $`S`$) “transposant” les fonctions et les opérateurs et qu’en outre l’application d’évaluation commute avec $`S`$ dans le sens suivant
$$ev^{23}S^{12}S^{23}(Xfg)=Sev^{12}(Xfg)$$
(il en est de même en remplaçant la fonction $`g`$ par un opérateur $`Y`$).
Ainsi, lorsque $`A`$ est une algèbre $`S`$-commutative, les champs de vecteurs sur une telle algèbre peuvent être définis par la $`S`$-analogue de la règle de Leibniz (1).
Si l’algèbre $`A`$ est plutôt munie d’une application de tresse non involutive (c’est juste le cas de l’algèbre que nous considérons ici) il n’existe aucune définition générale de l’analogue tressée d’une algèbre commutative. (Comme le montre plusieurs exemples y compris celui de l’hyperboloïde quantique, l’application naïve des relations de la $`S`$-commutativité (çi-dessus) n’est plus raisonnable). Autrement dit, il n’est pas évident de dire ce qu’est une algèbre $`S`$-commutative ni de décrire (dans l’esprit de la superthéorie) les éléments de la géométrie tressée. Par exemple la S-analogue de la règle de Leibniz définie par la relation (1) n’est plus valable et donc, il n’est pas trivial de décrire l’espace tangent d’une variété quantique. (Soulignons que pour les variétés quantiques liées aux applications de tresses involutives, l’espace tangent peut être décrit en termes de champs de vecteurs tressés.)
Le but de cet article est de définir l’espace tangent sur l’hyperboloïde quantique, d’étudier ses propriétés et d’introduire les analogues tressés de certaines structures classiques (métrique, connexion). Soulignons que nous traitons (dans l’esprit de \[S1\]) l’espace tangent comme un module projectif de rang fini sur l’algèbre des “fonctions polynômes sur l’hyperboloïde quantique”. Rappelons que dans \[S1\], Serre a établit une correspondance biunivoque entre les fibrés vectoriels et les modules projectifs de rang fini sur l’anneau des coordonnées.(Une analogue de cette correspondance pour les variétés lisses compactes a été établie par R.G.Swan.)
L’hyperboloïde quantique (ou tressé) est le plus simple exemple d’orbite quantique. Par orbite quantique, nous entendons les algèbres qui vérifient les propriétés suivantes :
1. elles sont $`U_q(𝗀)`$-covariantes (i.e. leur produit $``$ est $`U_q(𝗀)`$-covariant
$$X(ab)=\mathrm{\Delta }X(ab),XU_q(𝗀)$$
(2)
où a et b sont des éléments de l’algèbre considérée),
2. elles représentent une déformation plate <sup>2</sup><sup>2</sup>2Rappelons que l’algèbre $`𝒜_{\mathrm{}}`$ (où $`\mathrm{}`$ est un paramètre formel) est une déformation plate de l’algèbre $`𝒜_0,`$ si on a
$$𝒜_0=𝒜_{\mathrm{}}/\mathrm{}𝒜_{\mathrm{}}$$
et $`𝒜_0𝐊[[\mathrm{}]]`$ est isomorphe à $`𝒜_{\mathrm{}}`$ comme des $`𝐊[[\mathrm{}]]`$-modules (ici le produit tensoriel est complété en topologie $`\mathrm{}`$-adique). de leurs analogues classiques : i.e. des orbites habituelles dans $`𝗀^{}`$ (plus précisément, les algèbres de fonctions sur de telles orbites),
3. elles sont en un certain sens commutatives.
La dernière propriété est la plus difficile à traiter, car les orbites quantiques font partie d’une famille d’algèbres $`U_q(𝗀)`$-covariantes et ce n’est pas toujours facile de distinguer dans cette famille une algèbre “commutative” au sens tressé (voir \[DGK\] où ce problème est discuté). Pour expliquer le problème décrivons d’abord les objets infinitésimaux sur de telles algèbres.
Soit $``$ une variété lisse (en particulier une orbite dans $`𝗀^{}`$) et
$$\rho :𝗀\mathrm{Vect}()$$
la représentation de $`𝗀`$ sur l’espace des champs de vecteurs sur $``$. Associons à la $`R`$-matrice
$$R=\underset{\alpha \mathrm{\Delta }^+}{}X_\alpha X_\alpha ^2(𝗀)$$
$`\{H_\alpha ,X_\alpha ,X_\alpha \}`$ est une base de Cartan-Weyl en normalisation de Chevalley de l’algèbre de Lie $`𝗀`$ et $`\mathrm{\Delta }^+`$ désigne son système de racines positives (en supposant une décomposition triangulaire fixée de $`𝗀`$), le crochet de $`R`$-matrice
$$\{f,g\}_R=\mu <\rho ^2(R),dfdg>,f,g\mathrm{Fun}()$$
$`\mu `$ est le produit dans $`\mathrm{Fun}()`$.
Dans le cas général, ce crochet ne satisfait pas à la relation de Jacobi. Par conséquent, il n’est pas un crochet de Poisson. Mais il l’est sur certaines variétés dites de type $`R`$-matrices dans \[GP\]. C’est juste le cas des orbites dans $`𝗀^{}`$ qui admettent comme déformations plates, des algèbres $`U_q(𝗀)`$-covariantes. (Mais il existe quand même des orbites dans $`𝗀^{}`$ qui sont des déformations plates $`U_q(𝗀)`$-covariantes et sur lesquelles le crochet de Poisson correspondant est un peu différent de $`\{,\}_R`$, voir \[DGS\].) ll est facile de voir que sur de telles orbites, le crochet de Kirillov-Kostant-Souriau (KKS) noté $`\{,\}_{KKS}`$ (qui est la restriction sur $`𝗀^{}`$ tout entier du crochet linéaire dit de Lie-Poisson) et le crochet de $`R`$-matrice sont compatibles et donc sur elles, ces deux crochets engendrent ce qu’on appelle le pinceau de Poisson. Ainsi la famille de crochets définie par
$$\{,\}_{a,b}=a\{,\}_{KKS}+b\{,\}_R,$$
(3)
est une famille de crochets de Poisson.
Si l’on peut dire qu’un crochet de Poisson définit “une direction de déformation”, un pinceau de Poisson nous donne un plan de ces directions. La quantification double (i.e. simultanée) de ce pinceau de Poisson conduit aux algèbres non commutatives tressées. Ces algèbres dépendent de deux paramètres. Leur produit est toujours $`U_q(𝗀)`$-covariant.
Dans cet article, un cas particulier de telles algèbres est étudié. Notamment celle qui découle de l’hyperboloïde plongé comme une orbite (par rapport à l’action coadjointe du groupe $`SL(2)`$) d’un élément semi-simple de $`sl(2)^{}`$ (nous regardons plus précisément une famille d’orbites). Etant l’orbite d’un élément semi-simple la variété initiale est une variété algébrique affine fermée. Notons que dans ce cas particulier il n’est pas difficile de distinguer dans la famille des algèbres $`U_q(sl(2))`$-covariantes, une algèbre “commutative” tressée.
Soulignons que l’hyperboloïde quantique a été introduit par Podlès \[P\] sous le nom de la sphère quantique. En fait, la sphère quantique est l’hyperboloïde quantique munie d’une involution (nous ne regardons pas ici le problème d’une définition raisonnable d’une involution sur l’hyperboloïde quantique). Nous nous contentons de considérer simplement l’hyperboloïde quantique sans involution. Toutefois la forme compacte ou non de la variété initiale n’a aucune importance, puisque nous ne regardons que les fonctions polynômes restreintes sur cette variété. (Par un changement de variables convenable nous pouvons passer de l’algèbre des fonctions sur la sphère à celle des fonctions sur l’hyperboloïde correspondant et réciproquement.)
Certains aspects de l’algèbre provenant de la quantification du pinceau de Poisson (3) sur l’hyperboloïde ont été examinés dans \[DG\], \[DGR\], \[GV\]. En particulier il a été montré que cette algèbre est sans multiplicité (i.e. chaque module irréductible de dimension finie entre dans la décomposition en $`U_q(sl(2))`$-modules irréductibles de cette algèbre au maximum une fois).
Désignons par $`𝒜_{\mathrm{},q}^c`$ cette algèbre à deux paramètres qui est le résultat de la quantification double du pinceau de Poisson sur l’hyperboloïde plongé comme une orbite dans $`sl(2)^{}`$. En gros, nous disons que le paramètre $`\mathrm{}`$ est celui de la quantification du crochet de KKS, $`q`$ celui de tressage et $`c`$ numérote les orbites quantiques. Dans le cas où $`\mathrm{}=0`$, nous obtenons l’algèbre $`𝒜_{0,q}^c`$ qui est considérée comme la $`q`$-analogue de l’algèbre commutative correspondante. (Elle est parfois appelée sphère quantique standard.) C’est notre algèbre principale : “l’algèbre des fonctions polynômes restreintes sur l’hyperboloïde quantique” (où “l’anneau des coordonnées quantiques”). Ainsi dans la suite, tous les modules considérés (dans le cas quantique) seront des modules sur cette algèbre. Nous introduisons alors les notions de champs de vecteurs, d’espace tangent sur l’hyperboloïde quantique. Nous traitons cet espace tangent comme un $`𝒜_{0,q}^c`$-module et nous l’appellons “module tangent” sur l’hyperboloïde quantique.
Commençons par le module tangent. Dans le cas classique, il est formé par des champs de vecteurs tangents à la variété sous considération (ici la sphère ou l’hyperboloïde). Quels sont les analogues quantiques de ces champs ?
Les champs de vecteurs auxquels l’on peut spontanément penser, sont les générateurs $`X,Y,H`$ du groupe quantique (GQ) $`U_q(sl(2))`$ (voir Section 2). Mais si nous introduisons le module tangent comme toutes les combinaisons linéaires (ayant pour coefficients les éléments de l’algèbre $`𝒜_{0,q}^c`$) de ces opérateurs, nous n’obtenons pas une déformation plate du module tangent initial. En effet, notons $`\mathrm{Fun}(\mathrm{S}^2)`$ l’algèbre des fonctions polynômes restreintes sur la sphère ($`S^2`$). On peut décrire le “module tangent” sur la sphère comme toutes les combinaisons linéaires aux coefficients-fonctions de trois rotations infinitésimales
$$X=z_yy_z,Y=x_zz_x,Z=y_xx_y.$$
(4)
Ces rotations correspondent aux générateurs standards $`x,y,z`$ de $`𝗀=so(3)=su(2)`$ qui opèrent sur l’algèbre de Lie $`𝗀`$ elle-même par l’action adjointe, et leur extension sur les éléments de
$$\mathrm{Fun}(𝗀^{})=\mathrm{Sym}(𝗀)$$
s’effectue par la règle de Leibniz. Il est facile de voir que sur l’algèbre $`\mathrm{Fun}(\mathrm{S}^2)`$, les opérateurs $`X,Y,Z`$ satisfont à la relation
$$xX+yY+zZ=0.$$
En passant à la forme non compacte (i.e. à l’hyperboloïde $`\mathrm{H}`$), nous avons la relation
$$xY+yX+\frac{h}{2}H=\mathrm{\hspace{0.17em}0}$$
(5)
avec les générateurs standards $`x,y,h`$ de l’algèbre $`sl(2)`$ et $`X,Y,H`$ sont ici les rotations hyperboliques infinitésimales correspondantes. Notons $`\mathrm{Fun}(\mathrm{H})`$ l’algèbre des fonctions polynômes restreintes sur l’hyperboloïde. Finalement nous pouvons définir l’espace tangent sur l’hyperboloïde comme un module facteur d’un module libre $`A^3`$ de rang 3 (ici $`A=𝒜_{0,1}^c=\mathrm{Fun}(\mathrm{H})`$) de la façon suivante :
$$aX+bY+cH\mathrm{modulo}f.(xY+yX+\frac{h}{2}H)$$
$`a,b,c,f\mathrm{Fun}(\mathrm{H})`$. Notons $`T(\mathrm{H})`$ le $`\mathrm{Fun}(\mathrm{H})`$-module tangent sur l’hyperboloïde. Nous avons en outre l’action
$$T(\mathrm{H})AA$$
qui signifie que les éléments du $`A`$-module (disons gauche) $`T(\mathrm{H})`$ sont présentés comme des opérateurs sur $`A`$ et on appelle le plongement
$$sl(2)T(\mathrm{H}).$$
(6)
une ancre.
Notre but est de définir le module tangent sur l’hyperboloïde quantique muni d’une ancre quantique.
Contrairement au cas classique, les générateurs $`X,Y,H`$ du GQ $`U_q(sl(2))`$ ne satisfont à aucune relation du type (5) et donc la platitude de la déformation du module tangent (engendré par les opérateurs $`X,Y,H`$) n’a pas lieu.
Nous suggérons d’autres candidats pour le rôle des $`q`$-analogues des opérateurs $`X,Y,H(U(sl(2)))`$, de sorte qu’ils satisfassent à la $`q`$-analogue de la relation (5). Donc l’analogue quantique du module tangent peut être introduit de façon analogue au cas classique. La construction se fait en deux étapes :
$``$ Primo, nous définissons l’analogue tressé du crochet de Lie de $`sl(2)`$ (il a été introduit pour la première fois dans \[DG\] et généralisé après dans \[LS\]). Cela nous permet d’introduire l’ad-action de l’algèbre “$`sl(2)`$ tressée” sur elle-même. Il est facile de vérifier que les opérateurs correspondants aux générateurs de l’algèbre $`sl(2)`$ tressée, satisfont à une relation qui est la $`q`$-analogue de la relation (5). Cela incite alors à définir l’analogue quantique du module tangent comme un module quotient d’un module libre en utilisant cette relation.
$``$ Secundo, si nous voulons avoir une action de ce module sur l’algèbre $`𝒜_{0,q}^c`$, nous devons étendre l’action des opérateurs adjoints qui est bien définie pour le moment sur les éléments de degré un de l’algèbre $`𝒜_{0,q}^c`$ à toute fonction polynômes de degré supérieur à un, sur l’hyperboloïde quantique.
Malheureusement pour le cas que nous traitons, il n’existe aucune forme tressée de la règle de Leibniz. Nous présentons une autre méthode pour étendre l’action des opérateurs adjoints provenant de l’ad-action de l’algèbre $`sl(2)`$ tressée, sur toute l’algèbre $`𝒜_{0,q}^c`$ de sorte que la $`q`$-analogue de la relation (5) soit toujours satisfaite.
En fait notre construction nous amène à deux modules tangents : celui traité comme un $`𝒜_{0,q}^c`$-module gauche noté $`T(\mathrm{H}_q)_l`$ et celui traité comme un $`𝒜_{0,q}^c`$-module droit noté $`T(\mathrm{H}_q)_r`$. Le problème est de les identifier, c’est-à-dire construire un isomorphisme (dans la catégorie considérée<sup>3</sup><sup>3</sup>3Par définition, un morphisme dans cette catégorie est introduit comme une application qui commute à l’action de $`U_q(sl(2))`$.) entre ces deux modules. Notons qu’une telle identification est une étape indispensable pour définir par exemple une (pseudo)métrique tressée sur l’hyperboloïde quantique. Le problème est que notre méthode pour introduire une telle métrique consiste à définir d’abord un “couplage” $`<,>`$ sur $`T(\mathrm{H}_q)_l_𝐊T(\mathrm{H}_q)_r`$
$$<,>:T(\mathrm{H}_q)_l_𝐊T(\mathrm{H}_q)_r𝒜_{0,q}^c.$$
(Ici, $`T(\mathrm{H}_q)_l`$ est un $`𝒜_{0,q}^c`$-module gauche et $`T(\mathrm{H}_q)_r`$ un $`𝒜_{0,q}^c`$-module droit.)
Le passage à une (pseudo)métrique sur
$$T(\mathrm{H}_q)_r_𝐊T(\mathrm{H}_q)_r𝒜_{0,q}^c\text{ou}T(\mathrm{H}_q)_l_𝐊T(\mathrm{H}_q)_l𝒜_{0,q}^c$$
ne peut être effectué qu’après l’identification des modules $`T(\mathrm{H}_q)_l`$ et $`T(\mathrm{H}_q)_r`$ :
$$T(\mathrm{H}_q)_lT(\mathrm{H}_q)_r.$$
Soulignons que sans une telle identification, il n’est pas évident de sortir le facteur $`f`$ dans le terme
$$<X,fY>X,YT(\mathrm{H}_q)_l,f𝒜_{0,q}^c.$$
Plus précisément sans cette identification on ne peut pas définir le produit tensoriel $`T(\mathrm{H}_q)_ϵ_{𝒜_{0,q}^c}T(\mathrm{H}_q)_ϵ`$$`ϵ=l,r`$.
Nous construisons également une connexion (partiellement définie, i.e. sur un sous espace de $`T(\mathrm{H}_q)_ϵT(\mathrm{H}_q)_ϵ`$) tressée. Le problème dans ce cas est que, pour étendre cette connexion sur $`T(\mathrm{H}_q)_ϵT(\mathrm{H}_q)_ϵ`$, il faut pouvoir étendre le crochet de Lie tressé sur cet espace. Malheureusement nous ne savons pas prolonger (par les méthodes existantes) sur $`T(\mathrm{H}_q)_ϵT(\mathrm{H}_q)_ϵ`$ le crochet de Lie tressé initialement défini sur $`sl(2)`$.
Précison que les critères de raison d’être pour les objets et les opérateurs que nous introduisons sont de deux sortes : la platitude de la déformation et la $`U_q(sl(2))`$-covariance.
Le contenu de cet article est le suivant. Dans la Section 2, nous présentons notre algèbre principale : celle des fonctions sur l’hyperboloïde quantique. Nous rappelons dans la Section 3, l’algèbre enveloppante tressée de $`sl(2)`$. Dans la Section 4, nous définissons les champs de vecteurs tressés et l’espace tangent (considéré comme un $`𝒜_{0,q}^c`$-module) sur l’hyperboloïde quantique muni d’une ancre quantique. Dans la Section 5 nous montrons que ce $`𝒜_{0,q}^c`$-module tangent est projectif. Dans la Section 6 nous définissons et construisons une (pseudo)<sup>4</sup><sup>4</sup>4Pseudo signifie que son analogue classique n’est pas définie positive.métrique et une connexion (partiellement définie) tressées sur l’hyperboloïde quantique. Dans la dernière Section, nous proposons une façon canonique pour identifier les modules tangents gauche et droit sur l’hyperboloïde quantique.
Dans toute la suite, le corps de base $`𝐊`$ est $`\mathrm{I}\mathrm{R}`$ ou $`\mathrm{I}\mathrm{C}`$ et le paramètre $`q`$ ($`𝐊`$) est générique.
Remerciemments : D. Gurevich m’a soumis ce problème et m’a constamment guidé dans sa résolution et dans sa rédaction. Je l’en remercie.
## 2 Hyperboloïde quantique
Rappelons d’abord la présentation de l’hyperboloïde classique.
Soient :
$`sl(2)`$ l’algèbre de Lie du groupe de Lie $`SL(2)`$,
$`[,]`$ le crochet de Lie sur $`sl(2)`$,
$`sl(2)^{}`$ l’espace (vectoriel) dual de $`sl(2)`$ et
$`(X,Y,H)`$ la base de Cartan-Weyl de $`sl(2)`$.
La représentation “adjointe gauche” (Ad) de $`SL(2)`$ :
$$Ad_gX=gXg^1,gSL(2),Xsl(2)$$
fait de $`sl(2)`$ un $`SL(2)`$-module gauche. (Notons que Ad désigne souvent la représentation adjointe droite. Mais nous préférons réaliser $`sl(2)`$ comme un $`SL(2)`$-module gauche.)
Par la représentation coadjointe associée et notée $`Ad^{}`$, l’espace $`sl(2)^{}`$ est un $`SL(2)`$-module droit.
Soit $`\omega `$ l’élément de $`sl(2)^{}`$ défini par :
$$\omega (H)=a,\omega (X)=\omega (Y)=0,a𝐊,a0.$$
Désignons par $`𝒪_\omega `$ l’orbite de l’élément $`\omega `$ par la représentation coadjointe $`Ad^{}`$ :
$$𝒪_\omega =\{Ad_{g}^{}{}_{}{}^{}(\omega )/gSL(2)\}.$$
Le stabilisateur de l’élément $`\omega `$ est juste le sous-groupe de Cartan $`𝐇`$ de $`SL(2)`$. Par conséquent comme un espace homogène :
$`𝒪_\omega =SL(2)/𝐇.`$
Présentons maintenant l’orbite $`𝒪_\omega `$ comme une variété algébrique affine.
Considérons $`\mathrm{Fun}(\mathrm{sl}(2)^{})`$, l’algèbre des fonctions polynômes sur l’espace $`sl(2)^{}`$. Soit $`\mathrm{Sym}(\mathrm{sl}(2))`$ l’algèbre symétrique de l’espace $`sl(2)`$. On a alors de manière naturelle :
$$\mathrm{Fun}(\mathrm{sl}(2)^{})=\mathrm{Sym}(\mathrm{sl}(2)).$$
Soit $`U(sl(2))`$ l’algèbre enveloppante de $`sl(2)`$. Elle est une algèbre filtrée. Soit $`\mathrm{Gr}U(sl(2))`$ l’algèbre graduée associée à l’algèbre filtrée $`U(sl(2))`$. On a : $`\mathrm{Gr}U(sl(2))\mathrm{Sym}(\mathrm{sl}(2))`$ par le théorème de Poincaré-Birkhoff-Witt (PBW).
Soit $`\mathrm{C}`$ l’élément de Casimir de $`U(sl(2))`$
$$\mathrm{C}=XY+YX+\frac{H^2}{2},$$
$`\mathrm{C}`$, est un générateur du centre de l’algèbre $`U(sl(2))`$. Associons à chaque élément $`Z`$ de $`U(sl(2))`$ son image (notée $`z`$) dans $`\mathrm{Gr}U(sl(2))`$ ($`\mathrm{Sym}(\mathrm{sl}(2))`$). Alors par cette correspondance, l’image (que nous notons encore $`\mathrm{C}`$) de l’élément de Casimir dans $`\mathrm{Sym}(\mathrm{sl}(2))`$ est
$$\mathrm{C}=\mathrm{\hspace{0.17em}2}xy+\frac{h^2}{2}.$$
Il est bien connu que toute orbite de la représentation coadjointe $`Ad^{}`$ est contenue dans la variété algébrique affine définie par :
$$\mathrm{C}=\mathrm{\hspace{0.17em}2}xy+\frac{h^2}{2}=c\text{}c\text{est une constante dans}𝐊.$$
En particulier si $`𝐊=\mathrm{I}\mathrm{C}`$ et $`c0`$, l’orbite $`𝒪_\omega `$ coïncide avec l’hyperboloïde (classique) $`\mathrm{H}`$ d’équation
$$2xy+\frac{h^2}{2}=c=\mathrm{C}(\omega )=\frac{a^2}{2}.$$
(7)
Si $`𝐊=\mathrm{I}\mathrm{R}`$, l’hyperboloïde contient parfois une orbite, parfois deux.
Pour $`c=0`$, (7) définit le cône qui est composé de deux orbites $`\{0\}`$ et tout le reste si $`𝐊=\mathrm{I}\mathrm{C}`$; et de trois orbites si $`𝐊=\mathrm{I}\mathrm{R}`$. Fixons $`c0`$ et considérons $`\mathrm{Fun}(\mathrm{H})`$, l’algèbre des fonctions polynômes sur l’hyperboloïde $`\mathrm{H}`$. Par définition $`\mathrm{Fun}(\mathrm{H})`$ est la restriction des fonctions polynômes de $`\mathrm{Fun}(\mathrm{sl}(2)^{})`$ sur $`\mathrm{H}`$ i.e.
$$\mathrm{Fun}(\mathrm{H})=\mathrm{Fun}(\mathrm{sl}(2)^{})/\{\mathrm{C}\mathrm{c}\},$$
(8)
$`\{\mathrm{C}c\}`$ désigne l’idéal bilatère engendré par l’élément $`\mathrm{C}c`$.
Il est en outre facile de voir que la multiplication dans l’algèbre $`\mathrm{Fun}(\mathrm{H})`$ est covariante par l’action de $`U(sl(2))`$.
Par analogie avec le cas classique précédemment décrit, nous présentons l’analogue quantique noté $`\mathrm{H}_q`$ de l’hyperboloïde classique $`\mathrm{H}`$, sous forme de son algèbre des “fonctions quantiques”. La multiplication dans cette algèbre doit être en outre $`U_q(sl(2))`$-covariante. Pour ce faire, nous donnons dans la suite un analogue “tressé” de l’élément de Casimir $`\mathrm{C}`$ qui participera à nos constructions.
Soit $`U_q(sl(2))`$ le GQ associé au groupe $`SL(2)`$. Le groupe $`U_q(sl(2))`$ est une algèbre de Hopf. Dans le modèle de Drinfel’d-Jimbo, elle est engendrée par les éléments $`X,Y,H`$ satisfaisant aux relations de commutation (pour $`q0,q^21`$) :
$$[H,X]=2X,[H,Y]=2Y,[X,Y]=\frac{q^Hq^H}{qq^1}.$$
(9)
On peut choisir le coproduit ($`\mathrm{\Delta }`$) défini par exemple par :
$$\mathrm{\Delta }(X)=X1+q^HX,\mathrm{\Delta }(Y)=1Y+Yq^H,\mathrm{\Delta }(H)=H1+1H.$$
(10)
Alors l’antipode $`\gamma `$ est donnée par :
$$\gamma (X)=q^HX,\gamma (H)=H,\gamma (Y)=Yq^H.$$
(11)
(Pour $`q=1`$, les relations (9) et (10) correspondent à celles de l’algèbre de Hopf $`U(sl(2))`$).
Désignons par $`U_q(sl(2))Mod`$ la catégorie des $`U_q(sl(2))`$-modules de dimension finie qui sont analogues quantiques (c’est-à-dire des déformations) de $`U(sl(2))`$-modules irréductibles de dimension finie. Tout objet de $`U_q(sl(2))Mod`$ est appelé $`q`$-analogue (ou analogue tressé) de l’objet classique correspondant.
Le centre de l’algèbre $`U_q(sl(2))`$ est engendré (voir \[M\]) par l’opérateur de Casimir quantique
$$\mathrm{C}_q=(\frac{q^{\frac{H+1}{2}}q^{\frac{H+1}{2}}}{qq^1})^2+YX.$$
(12)
(On peut remarquer que pour $`q=1`$ on a $`\mathrm{C}_1=\frac{\mathrm{C}}{2}+\frac{id}{4}\mathrm{C}`$.)
Donnons à présent un autre analogue de l’élément de Casimir $`\mathrm{C}`$ qui nous servira dans la suite. Pour cela, considérons une seconde copie de l’espace $`sl(2)`$ que nous notons $`𝖵`$ pour la différencier de l’espace initiale $`sl(2)`$ et désignons par ($`u,v,w`$) une base de l’espace $`𝖵`$ :
$$𝖵=Span(u,v,w),\text{notons}sl(2):=(𝖵,[,]).$$
Munissons $`𝖵`$ de l’action de $`U_q(sl(2))`$ qui coïncide pour $`q=1`$ avec celle de la représentation adjointe (ad) de $`sl(2)`$. Elle est notée . et définie par :
$`X.u`$ $`=`$ $`0,X.v=(q+q^1)u,X.w=v,`$
$`Y.u`$ $`=`$ $`v,Y.v=(q+q^1)w,Y.w=0,`$
$`H.u`$ $`=`$ $`2u,H.v=0,H.w=2w.`$
(Pour $`q=1`$, on vérifie qu’on a bien l’“ad-action gauche” de $`sl(2)`$.)
La structure de coalgèbre de $`U_q(sl(2))`$ permet d’étendre cette action sur $`𝖵^2`$. Par le théorème de Clebsch-Gordan quantique (voir \[K\]), $`𝖵^2`$ se décompose en trois $`U_q(sl(2))`$-modules irréductibles de dimension finie $`𝖵_0^q,𝖵_1^q,𝖵_2^q`$ respectivement de spins $`\mathrm{\hspace{0.17em}0},1,2`$. Fixons respectivement dans les espaces $`𝖵,𝖵_0^q,𝖵_1^q,𝖵_2^q`$ leurs éléments de plus haut poids, notés :
$$x_0,𝒞_q,x_1,x_2$$
$`𝒞_q=(q^3+q)uw+vv+(q+q^1)wu`$.
La compatibilité de la $`U_q(sl(2))`$-action sur l’espace $`𝖵^2𝖵𝐊`$ impose les relations
$$𝒞_q=c,x_1=\mathrm{}x_0$$
(13)
$`c`$ et $`\mathrm{}`$ sont des constantes dans $`𝐊`$. $`𝒞_q`$ est appelé Casimir tressé. Il n’est pas à confondre avec le Casimir quantique $`\mathrm{C}_q`$ (défini par (12)) qui appartient à l’algèbre $`U_q(sl(2))`$.
Faisons remarquer que pour $`q=1`$, $`𝒞_1=4uw+v^2=2\mathrm{C}`$. $`𝒞_q`$ (plus précisément $`\frac{𝒞_q}{2}`$) est la $`q`$-analogue de $`\mathrm{C}`$, dont nous nous servons dans toute la suite pour nos constructions.
En opérant avec l’opérateur $`YU_q(sl(2))`$ sur la deuxième équation de (13), on en déduit en plus deux autres équations. (Voir appendice A pour leurs formes explicites). En posant :
$$x_0=u,x_1=q^2uvvu,x_2=uu$$
on en déduit l’expression des deux autres éléments de base de $`𝖵_1^q`$ en opérant avec $`Y(U_q(sl(2)))`$ sur l’élément $`x_1`$.
###### Définition 2.1
L’algèbre $`𝒜_{\mathrm{},q}^c`$ est le quotient de l’algèbre tensorielle libre $`T(𝖵)`$ par l’idéal bilatère $`I_{\mathrm{}}`$ engendré par les éléments :
$`q^2uvvu+2u\mathrm{},(q^3+q)(uwwu)+(1q^2)vv2v\mathrm{},`$
$`q^2vw+wv2w\mathrm{},𝒞_qc.`$
Pour $`\mathrm{}=0\text{et}c0`$, l’algèbre $`𝒜_{0,q}^c`$ est celle des fonctions sur l’hyperboloïde quantique (qui est aussi notée $`\mathrm{H}_q`$).
Remarque 2.1 Dans cette définition $`\mathrm{}`$ et $`q`$ sont des constantes fixées. Nous pouvons les considérer comme des paramètres. Il suffit pour cela de remplacer dans la définition précédente $`T(𝖵)`$ par $`T(𝖵)𝐊[[\mathrm{},q,q^1]]`$. Le paramètre orbital $`c`$ est une constante.
Bien que l’hyperboloïde classique soit une variété non compacte, en un certain sens, l’hyperboloïde quantique est plutôt un analogue tressé de la sphère. Ceci vient du fait que nous ne considérons que les fonctions polynomiales sur l’objet classique et leurs analogues tressés. En outre, toutes les représentations que nous considérons sont de dimension finie.
###### Proposition 2.1
$`𝒜_{\mathrm{},q}^c`$ est une algèbre associative $`U_q(sl(2))`$-covariante.
$`Preuve:`$ C’est par construction de l’algèbre $`𝒜_{\mathrm{},q}^c`$.
$``$ $`𝒜_{0,1}^c`$ (pour $`c0`$) est l’algèbre des fonctions polynômes sur $`sl(2)`$ restreintes à la variété algébrique affine définie par : $`\mathrm{\hspace{0.17em}\hspace{0.17em}4}uw+v^2=c`$, c’est l’algèbre des fonctions sur l’hyperboloïde (classique) qui est commutative.
$``$ L’algèbre $`𝒜_{\mathrm{},1}^c`$ (pour $`c0`$) est l’analogue non commutative de l’algèbre $`𝒜_{0,1}^c`$, mais elle est toujours $`sl(2)`$-invariante.
Remarque 2.2 L’algèbre $`𝒜_{\mathrm{},1}^c`$ est une déformation plate de l’algèbre $`𝒜_{0,1}^c`$. Cela découle du théorème de PBW. Ainsi pour $`c0`$, comme dans la décomposition en $`sl(2)`$-modules irréductibles de dimension finie ($`𝖵_k`$ de spin $`k𝐍`$) de l’algèbre $`𝒜_{0,1}^c`$, toute composante apparaît sans multiplicité, il en est de même pour les algèbres $`𝒜_{\mathrm{},1}^c`$ et $`𝒜_{\mathrm{},q}^c`$ ($`q1`$) (voir \[GV\],\[A\]).
$``$ Le cas $`c=0`$ correspond au cône (dit “quantique” si $`q1`$).
$``$ $`𝒜_{\mathrm{},q}^c`$ est une famille (dépendante de $`\mathrm{}`$) d’algèbres $`U_q(sl(2))`$-covariantes. La sphère quantique de Podlès \[P\] est en effet une autre présentation de telles algèbres, et munies d’une involution. Nous n’avons pas besoin dans la suite d’une quelconque involution sur l’algèbre $`𝒜_{0,q}^c`$.
$``$ Dans la famille $`𝒜_{\mathrm{},q}^c`$, nous traitons l’algèbre $`𝒜_{0,q}^c`$ comme la $`q`$-analogue de l’algèbre commutative $`𝒜_{0,1}^c`$.
## 3 $`sl(2)`$ tressé
### 3.1 Crochet de Lie tressé de $`sl(2)`$
Le crochet de Lie $`[,]`$ sur $`sl(2)`$ :
$`[,]:𝖵^2𝖵,`$
$``$ est une application $`𝐊`$-linéaire,
$``$ et $`sl(2)`$-invariante.
Nous allons définir de façon analogue le crochet de Lie tressé (noté $`[,]_q`$) sur $`sl(2)`$.
Nous introduisons les $`q`$-analogues notées $`I_\pm ^q`$ des sous-espaces symétriques et antisymétriques de l’espace $`sl(2)^2`$ de façon similaire au cas classique en posant :
$$I_+^q=𝖵_0^q𝖵_2^q\text{et}I_{}^q=𝖵_1^q.$$
(Notons que les algèbres correspondantes $`T(𝖵)/\{I_\pm ^q\}`$ sont des déformations plates de leurs analogues classiques. En outre comme l’algèbre $`𝒜_{0,q}^c`$ est une algèbre quotient de l’algèbre “$`q`$-symétrique” $`T(𝖵)/\{I_{}^q\}`$, $`𝒜_{0,q}^c`$ est également une algèbre $`q`$-symétrique.)
Notons que $`I_{}^q`$ est engendré par les trois tenseurs
$`q^2uvvu,(q^3+q)(uwwu)+(1q^2)vv,q^2vw+wv.`$
###### Définition 3.1
Le crochet de Lie tressé de $`sl(2)`$, est l’opérateur
$`[,]_q:𝖵^2𝖵\text{vérifiant}`$
1. $`[,]_qI_+^q=\mathrm{\hspace{0.17em}0}`$,
2. $``$ $`[,]_q(q^2uvvu)=\tau u`$,
$``$ $`[,]_q((q^3+q)(uwwu)+(1q^2)vv)=\tau v`$,
$``$ $`[,]_q(q^2vw+wv)=\tau w`$.
$`\tau `$ est une constante non nulle.
$`𝖵`$ muni du crochet $`[,]_q`$ est appelé algèbre de Lie tressée et noté $`sl(2)_q`$. Plus précisément
$$sl(2)_q:=(𝖵,[,]_q).$$
De fait ce crochet de Lie tressé dépend du facteur $`\tau `$. Mais nous négligeons cette dépendance en supposant que $`\tau `$ est fixé.
###### Proposition 3.1
1. $`[,]_q`$ est un $`U_q(sl(2))`$-morphisme (i.e. une application $`U_q(sl(2))`$-covariante).
2. La table de commutation de $`[,]_q`$ est :
$`[u,u]_q=0,[u,v]_q=q^2Mu,[u,w]_q=(q+q^1)^1Mv,`$
$`[v,u]_q=Mu,[v,v]_q=(1q^2)Mv,[v,w]_q=q^2Mw,`$
$`[w,u]_q=(q+q^1)^1Mv,[w,v]_q=Mw,[w,w]_q=0,`$
$`M=(1+q^4)^1\tau .`$
$`Preuve:`$ 1. C’est la propriété 2 de la définition 3.1
2. C’est un calcul direct.
(Si $`q=1`$ et $`\tau =4`$ (donc $`M=2`$), nous obtenons le crochet de Lie sur $`sl(2)`$.)
Remarque 3.1 C’est le fait que l’espace $`sl(2)`$ apparaisse une seule fois dans la décomposition en $`sl(2)`$-modules irréductibles de $`sl(2)^2`$ qui a permit de définir de façon unique le crochet de Lie tressé $`[,]_q`$. Pour les algèbres de Lie $`𝗀=sl(n)(n>\mathrm{\hspace{0.17em}2})`$, la multiplicité de l’espace $`sl(n)`$ dans $`sl(n)^2`$ est deux : l’une appartenant à la partie symétrique de $`sl(n)^2`$ et l’autre à sa partie antisymétrique. Il n’est donc pas évident de décrire les $`q`$-analogue des algèbres symétriques et antisymétriques de $`𝗀`$. Cependant, il existe un sous-espace $`I_{}^q𝗀_q^2`$$`𝗀_q`$ est l’espace $`sl(n)`$ muni de la $`U_q(sl(n))`$-action, telle que l’algèbre quadratique $`T(𝗀_q)/\{I_{}^q\}`$ soit une déformation plate de l’algèbre symétrique de $`𝗀`$ (voir \[D\]). Une description explicite du sous-espace $`I_{}^q`$ peut être donnée par l’équation appelée “reflection equation”
$$SL_1SL_1=L_1SL_1S$$
(14)
$`S`$ est une solution de l’équation de Yang-Baxter quantique de type Hecke (voir \[G1\]), $`L_1=Lid`$ et $`L`$ est une matrice dont les coefficients matriciels sont les éléments $`l_i^j,1i,jn`$. L’algèbre quadratique définie par l’équation (14) est habituellement appelée “reflection equation algebra” (REA). En considérant l’algèbre $`𝗀_q`$ introduite dans \[LS\] qui est la $`q`$-analogue de l’algèbre $`𝗀`$, on peut décrire à partir de la REA, l’algèbre enveloppante de l’algèbre $`𝗀_q`$ qui est aussi une déformation plate à deux paramètres de l’algèbre symétrique de $`𝗀`$. (voir \[AG\])
Si $`𝗀`$ est une algèbre de Lie simple différente de $`sl(n)`$, dans $`𝗀^2`$ toute composante qui apparaît est sans multiplicité. Donc on peut définir la $`q`$-analogue du crochet de Lie, en imposant qu’il soit un morphisme non trivial dans la catégorie des $`U_q(𝗀)`$-modules irréductibles de dimension finie (il est donc défini de façon unique à un facteur constant près) puis, on introduit son algèbre enveloppante, son algèbre symétrique et antisymétrique comme dans le cas classique, mais dans la catégorie sous considération. (Ici le fait que $`𝗀^2`$ soit sans multiplicité joue le rôle principale.) Cependant, ces algèbres ne sont pas des déformations plates de leurs analogues classiques (voir \[G2\]).
### 3.2 Algèbre enveloppante tressée de $`sl(2)`$
$``$ Rappelons que dans le cas classique ($`q=1`$), l’algèbre enveloppante $`U(sl(2))`$ est définie de la façon suivante
$$\begin{array}{ccc}U(sl(2))& =& T(sl(2))/\{ABBA[A,B]\}\\ & =& T(sl(2))/\{Im(id\frac{1}{2}[,])I_{}\},\end{array}$$
(15)
$`A,B`$ sont des éléments de $`sl(2)`$.
$``$ Par analogie avec l’algèbre enveloppante de $`sl(2)`$, nous définissons l’algèbre enveloppante tressée de $`sl(2)`$ notée $`U(sl(2)_q)`$ comme suit :
$$U(sl(2)_q)=T(sl(2)_q)/\{Im(id\kappa [,]_q)I_{}^q\}.$$
L’idéal $`\{Im(id\kappa [,]_q)I_{}^q\}`$ est celui engendré par les éléments :
$$\begin{array}{cc}& q^2uvvu\kappa (q^2[u,v]_q[v,u]_q),\\ & (q^3+q)(uwwu)+(1q^2)vv\kappa ((q^3+q)([u,w]_q\\ & [w,u]_q)(1q^2)[v,v]_q),\\ & q^2vw+wv\kappa (q^2[v,w]_q+[w,v]_q).\end{array}$$
(16)
Le choix de $`\kappa `$ sera précisé par la suite (voir la relation (23)).
Remarque 3.2 En fait on a
$$𝒜_{\mathrm{},q}^c=U(sl(2)_q)/\{𝒞_qc\}\text{avec}\mathrm{}=\frac{\kappa \tau }{2}.$$
###### Lemme 3.1
(\[DG\]) Le Casimir tressé $`𝒞_q`$ est un élément centrale de l’algèbre enveloppante tressée $`U(sl(2)_q)`$ c’est-à-dire :
$$X𝒞_q=𝒞_qXU(sl(2)_q).$$
Rappelons que nous voulons en fait d’une part, introduire le module tangent sur l’hyperboloïde quantique et d’autre part décrire les analogues tressées de certaines notions de la géométrie différentielle classique sur ce module. Pour ce faire, dans toute la suite nous travaillons principalement avec l’algèbre $`𝒜_{0,q}^c`$ et tous les modules considérés sont des modules sur cette algèbre.
## 4 Espace tangent quantique
Dans cette section, nous définissons l’espace tangent (considéré comme un $`𝒜_{0,q}^c`$-module) sur l’hyperboloïde quantique noté $`T(\mathrm{H}_q)`$. Puis nous construisons des champs de vecteurs tressés sur l’hyperboloïde quantique tels qu’ils nous permettent de munir $`T(\mathrm{H}_q)`$ d’une structure d’ancre quantique.
Pour mieux présenter le formalisme de la construction de ce module et de cette ancre dans le cas quantique, commençons par l’exemple du cas classique.
Considérons pour cela la sphère $`S^2`$ de dimension deux. Elle a pour équation
$$x^2+y^2+z^2=R^2$$
$`R`$ est une constante strictement positive. Posons
$$\mathrm{Fun}(\mathrm{S}^2)=𝐊[\mathrm{x},\mathrm{y},\mathrm{z}]/\{\mathrm{x}^2+\mathrm{y}^2+\mathrm{z}^2\mathrm{R}^2\}.$$
Nous donnons ici, trois descriptions (globales) de l’espace tangent sur la sphère noté $`T(S^2)`$.
a) Comme un champ de vecteurs c’est : $`\mathrm{Vect}(\mathrm{S}^2)`$ i.e. l’espace des champs de vecteurs sur $`S^2`$ .
b) Comme un $`\mathrm{Fun}(\mathrm{S}^2)`$-module : d’abord $`\mathrm{Vect}(\mathrm{S}^2)`$ est engendré par les trois rotations infinitésimales $`X,Y,Z`$ définies par (4) et qui vérifient dans l’algèbre $`\mathrm{Fun}(\mathrm{S}^2)`$ la relation
$$xX+yY+zZ=0.$$
(17)
(Dans l’expression (17) $`x,y,z`$ désignent en fait des opérateurs de multiplication). Donc comme un $`\mathrm{Fun}(\mathrm{S}^2)`$-module, $`T(S^2)`$ peut être réalisé comme le module quotient $`M/N`$
$$M=\{aX+bY+cZ,a,b,c\mathrm{Fun}(\mathrm{S}^2)\},$$
$$N=\{f(xX+yY+zZ),f\mathrm{Fun}(\mathrm{S}^2)\}.$$
c) Comme une variété algébrique affine : elle est plongée dans l’espace de dimension 6
$$(span(x,y,z,X,Y,Z))^{}$$
et définie par l’équation de la sphère et la relation (17).
Notons que l’espace tangent $`T(S^2)`$ est un cas particulier de fibré vectoriel (sur la sphère). Habituellement $`T(S^2)`$ est défini en termes de cartes locales. Nous n’utilisons pas cette description locale ici, car dans le cas quantique nous n’avons pas de localisation.
En passant à l’hyperboloïde $`\mathrm{H}`$ (i.e. à l’analogue non compacte de la sphère), nous avons de façon analogue trois descriptions (globales) de l’espace tangent sur l’hyperboloïde noté $`T(\mathrm{H})`$.
a’) Comme un champ de vecteurs c’est juste $`\mathrm{Vect}(\mathrm{H})`$.
b’) Comme un $`𝒜_{0,1}^c`$-module, il est réalisé comme le module quotient $`M/N`$
$$M=\{aU+bV+cW,a,b,c𝒜_{0,1}^c\},$$
$$N=\{f(2uW+vV+2wU),f𝒜_{0,1}^c=\mathrm{Fun}(\mathrm{H})\}.$$
Notons qu’ici $`U,V,W`$ sont les rotations hyperboliques infinitésimales associées respectivement aux générateurs $`u,v,w`$ de l’algèbre de Lie $`sl(2)`$. De même l’analogue non compacte de l’équation (17) s’écrit :
$$2uW+vV+2wU=\mathrm{\hspace{0.17em}0}.$$
(18)
Nous pouvons mettre l’équation (18) sous la forme symbolique
$$(𝖵𝖵^{})_0=0,$$
(19)
$`𝖵`$ désigne (encore) l’espace vectoriel engendré par $`u,v,w`$ et la marque $`^{}`$ désigne l’espace vectoriel engendré par les rotations hyperboliques infinitésimales. En outre $`𝖵`$ et $`𝖵^{}`$ sont des $`U(sl(2))`$-modules. La composante $`(𝖵𝖵^{})_i`$ dénote celle de spin $`i`$ dans la décomposition en $`U(sl(2))`$-modules irréductibles de dimension finie de $`𝖵𝖵^{}`$.
Ici, nous regardons le module tangent $`T(\mathrm{H})`$ comme un $`𝒜_{0,1}^c`$-module gauche. Comme un $`𝒜_{0,1}^c`$-module droit il est donné par l’équation $`(𝖵^{}𝖵)_0=0`$. Il est bien connu que ces deux $`𝒜_{0,1}^c`$-modules s’identifient naturellement. Nous discutons dans la section 7, du problème d’identification des modules tangents gauche et droit sur l’hyperboloïde quantique (réalisés tous deux comme des $`𝒜_{0,q}^c`$-modules).
c’) Enfin comme une variété algébrique affine, elle est plongée dans l’espace de dimension 6
$$(span(u,v,w,U,V,W))^{}$$
et définie par l’équation de l’hyperboloïde et la relation (18).
Laquelle des trois descriptions précédentes admet une “bonne” (i.e. plate) $`q`$-analogue ?
Il est évident que si nous voulons définir sur l’hyperboloïde quantique le $`𝒜_{0,q}^c`$-module tangent gauche noté $`T(\mathrm{H}_q)_l`$ comme une déformation plate de son analogue classique, nous devons utiliser la même formule (19) mais dans la catégorie $`U_q(sl(2))Mod`$. Cela nous amène à l’équation symbolique :
$$(𝖵𝖵^q)_0=0.$$
(20)
Dans l’identité (20), $`𝖵`$ est muni de la $`U_q(sl(2))`$-action et $`𝖵^q`$ est l’analogue tressé de $`𝖵^{}`$.
Interessons nous d’abord à la $`q`$-analogue des descriptions a’) et b’).
### 4.1 $`T(\mathrm{H}_q)`$ comme $`𝒜_{0,q}^c`$-module tressé
Soit $`𝒩`$ un élément de la catégorie $`U_q(sl(2))Mod`$ et qui soit en outre un $`𝒜_{0,q}^c`$-module. Soit l’application
$$\mu :𝒜_{0,q}^c𝒩𝒩$$
qui désigne l’action de $`𝒜_{0,q}^c`$ sur $`𝒩`$.
###### Définition 4.1
$`𝒩`$ est appelé un $`𝒜_{0,q}^c`$-module tressé si $`\mu `$ est un morphisme dans la catégorie $`U_q(sl(2))Mod`$ i.e. :
$$z.\mu (an)=\mu (z_{(1)}.az_{(2)}.n)$$
$`zU_q(sl(2)),\mathrm{\Delta }(z)=z_{(1)}z_{(2)},`$ $`a𝒜_{0,q}^c`$ et $`n𝒩`$.
Notons $`U^q,V^q,W^q`$ les générateurs de $`𝖵^q`$. Le GQ $`U_q(sl(2))`$ agit sur ces générateurs comme il opérait sur les éléments de l’espace $`𝖵`$ i.e.
$`X.U^q=\mathrm{\hspace{0.17em}0},X.V^q=(q+q^1)U^q,X.W^q=V^q,`$
$`Y.U^q=V^q,Y.V^q=(q+q^1)W^q,Y.W^q=0,`$
$`H.U^q=\mathrm{\hspace{0.17em}2}U^q,H.V^q=0,H.W^q=2W^q.`$
Alors l’identité (20) devient en forme explicite :
$$(q^3+q)uW^q+vV^q+(q+q^1)wU^q=0$$
(21)
Ainsi le $`𝒜_{0,q}^c`$-module tangent gauche $`T(\mathrm{H}_q)_l`$ est réalisé comme un $`𝒜_{0,q}^c`$-module facteur du $`𝒜_{0,q}^c`$-module
$$(𝒜_{0,q}^c)^3=M_l^q=\{aU^q+bV^q+cW^q,a,b,c,𝒜_{0,q}^c\}$$
par le $`𝒜_{0,q}^c`$ sous-module
$$N_l^q=\{f((q^3+q)uW^q+vV^q+(q+q^1)wU^q),f𝒜_{0,q}^c\}.$$
Précisons que les $`𝒜_{0,q}^c`$-modules
$$M_l^q,\mathrm{N}_l^q,T(\mathrm{H}_q)_l=M_l^q/\mathrm{N}_l^q$$
sont des modules tressés au sens de la définition 4.1.
De façon similaire le $`𝒜_{0,q}^c`$-module tangent droit sur l’hyperboloïde quantique noté $`T(\mathrm{H}_q)_r`$ est réalisé comme le module quotient $`M_r^q/N_r^q`$
$$M_r^q=\{\overline{U}^qa+\overline{V}^qb+\overline{W}^qc,a,b,c,𝒜_{0,q}^c\},$$
$$N_r^q=\{f((q^3+q)\overline{U}^qw+\overline{V}^qv+(q+q^1)\overline{W}^qu),f𝒜_{0,q}^c\},$$
$`\overline{U}^q,\overline{V}^q,\overline{W}^q`$ sont les générateurs du $`𝒜_{0,q}^c`$-module $`T(\mathrm{H}_q)_r`$. (L’action du groupe quantique $`U_q(sl(2))`$ sur ces générateurs est la même que sur $`𝖵^q`$.) De même les $`𝒜_{0,q}^c`$-modules
$$M_r^q,\mathrm{N}_r^q,T(\mathrm{H}_q)_r=M_r^q/\mathrm{N}_r^q$$
sont des modules tressés au sens de la définition (4.1).
###### Proposition 4.1
(\[A\]) Le $`𝒜_{0,q}^c`$-module tangent $`T(\mathrm{H}_q)`$ est une déformation plate de son analogue classique.
### 4.2 $`T(\mathrm{H}_q)`$ comme champs de vecteurs tressés
Nous considérons les générateurs $`u,v,w`$ de l’algèbre $`𝒜_{0,q}^c`$ comme des opérateurs de multiplication (à gauche) dans cette même algèbre, puis nous définissons les opérateurs $`U^q,V^q,W^q`$ qui sont les $`q`$-analogues des rotations hyperboliques infinitésimales $`U,V,W`$. Finalement nous obtenons la $`q`$-analogue des champs de vecteurs qui est engendré par les opérateurs $`U^q,V^q,W^q`$ et les éléments de l’algèbre $`𝒜_{0,q}^c`$ sont traités comme des opérateurs.
Soulignons une fois de plus que la tentation est de faire jouer aux générateurs $`X,Y,H`$ du GQ $`U_q(sl(2))`$ le rôle de $`q`$-analogue des rotations hyperboliques infinitésimales $`U,V,W`$. Mais alors dans ce cas, on a aucune relation de la forme (20). Par conséquent ils ne peuvent être considérer comme les $`q`$-analogues des champs de vecteurs $`U,V,W`$. Il faut donc trouver un autre moyen pour réaliser les générateurs $`U^q,V^q,`$
$`W^q`$ de $`𝖵^q`$ comme des opérateurs sur l’algèbre $`𝒜_{0,q}^c`$.
Par analogie avec le cas classique, à partir du crochet de Lie tressé de $`sl(2)`$ précédemment défini, on peut définir $`U^q,V^q,W^q`$ comme des opérateurs sur l’espace $`𝖵`$ (identifié à l’espace des fonctions linéaires sur $`𝖵^{}`$). En effet, associons au vecteur de base $`u`$ de l’espace $`𝖵`$, l’opérateur
$$\begin{array}{cccc}U^q:& 𝖵& & 𝖵\hfill \\ & z& & U^qz=ad^qu(z)=[u,z]_q,U^q1=0,\hfill \end{array}$$
défini sur les vecteurs de base de $`𝖵`$. Les opérateurs $`V^q,W^q`$ associés respectivement aux vecteurs de base $`v,w`$ sont définis de façon analogue à l’opérateur $`U^q`$.
Imposons maintenant que les opérateurs $`U^q,V^q,W^q`$ ainsi définis sur $`𝖵`$ vérifient les relations de définition de l’algèbre enveloppante tressée $`U(sl(2)_q)`$. C’est-à-dire que l’équation (16) soit encore satisfaite si nous remplaçons les générateurs $`u,v,w`$ (de l’espace $`𝖵`$) respectivement par leurs images par $`ad^q`$. Autrement dit, pour l’opérateur $`ad^q`$ nous avons pour tout vecteur de base $`z`$ de l’espace $`𝖵`$ :
$$\begin{array}{cc}& q^2[u,[v,z]_q]_q[v,[u,z]_q]_q=\kappa [q^2[u,v]_q[v,u]_q,z]_q\\ & (q^3+q)([u,[w,z]_q]_q[w,[u,z]_q]_q)+(1q^2)[v,[v,z]_q]_q=\\ & =\kappa (q^3+q)[[u,w]_q[w,u]_q,z]_q+(1q^2)[[v,v]_q,z]_q,\\ & q^2[v,[w,z]_q]_q+[w,[v,z]_q]_q=\kappa [q^2[v,w]_q+[w,v]_q,z]_q.\end{array}$$
(22)
Les relations du (22) fixent le choix de la constante $`\kappa `$ de (16). Par exemple en se servant de la première relation du (22), nous obtenons
$$\kappa =1(q^2+q^2)^1.$$
(23)
Ainsi pour ce choix de $`\kappa `$ fixé par l’identité (23), $`ad^q`$ est une représentation de $`sl(2)_q`$ (sur $`𝖵`$). Nous regardons les relations du (22) comme l’analogue tressé de l’identité de Jacobi.
Remarque 4.2.1 Il a été montré dans \[LS\] qu’il existe également un analogue tressé de l’identité de Jacobi pour les algèbres $`sl(n),n>2`$. Pour d’autres algèbres de Lie simples, il n’en existe (apparemment) pas.
Puisque $`ad^q`$ est une représentation, les opérateurs $`U^q,V^q,W^q`$ sont des opérateurs adjoints (gauches) tressés. Ils sont donc définis par la table suivante :
$`U^qu=[u,u]_q=0,U^qv=q^2Mu,U^qw=(q+q^1)^1Mv,`$
$`V^qu=Mu,V^qv=(1q^2)Mv,V^qw=q^2Mw,`$
$`W^qu=(q+q^1)^1Mv,W^qv=Mw,W^qw=0.`$
###### Lemme 4.1
$$(q^3+q)uW^q+vV^q+(q+q^1)wU^q=\mathrm{\hspace{0.17em}0},$$
(24)
sur les éléments de degré un de l’algèbre $`𝒜_{0,q}^c`$.
$`Preuve:`$ Il suffit de la vérifier sur les éléments $`u,v,w`$ de l’algèbre $`𝒜_{0,q}^c`$. Ce qui est immédiat. Par exemple pour $`u`$ on a :
$$(q^3+q)uW^q(u)+vV^q(u)+(q+q^1)wU^q(u)=$$
$$=M(q^2uv+vu)=\mathrm{\hspace{0.17em}0}\text{dans l’algèbre}𝒜_{0,q}^c.$$
Il reste maintenant à résoudre le problème qui consiste à étendre les opérateurs $`U^q,V^q,W^q`$ (bien définis pour le moment sur l’espace $`𝖵`$) sur tout élément de l’algèbre $`𝒜_{0,q}^c`$ de telle manière que l’équation (24) soit vérifiée sur $`𝒜_{0,q}^c`$ et les opérateurs ainsi prolongés satisfassent aux relations de définition de l’algèbre $`U(sl(2)_q)`$.
Dans le cas classique cette extension est effectuée par la règle de Leibniz. Il existe aussi une forme de cette règle pour les solutions involutives de l’équation de Yang-Baxter quantique (voir (1)). Mais dans notre cas, où l’opérateur de Yang-Baxter quantique provenant du GQ $`U_q(sl(2))`$ n’est pas involutif, ce n’est pas évident d’étendre de façon naturelle ces opérateurs sur les éléments de degré supérieur (à un) de l’algèbre $`𝒜_{0,q}^c`$.
Le fait que l’algèbre $`𝒜_{0,q}^c`$ se décompose en $`U_q(sl(2))`$-modules irréductibles de spin $`k`$, $`𝖵_k^q(𝖵^k)`$ nous permet de définir cette extension sur les composantes $`𝖵_k^q`$.
###### Théorème 4.1
Il existe une application
$$\beta :T(\mathrm{H}_q)𝒜_{0,q}^c𝒜_{0,q}^c$$
telle que :
1) Le diagramme suivant soit commutatif
$$\begin{array}{ccc}𝒜_{0,q}^cT(\mathrm{H}_q)𝒜_{0,q}^c& & T(\mathrm{H}_q)𝒜_{0,q}^c\\ & & \\ 𝒜_{0,q}^c𝒜_{0,q}^c& & 𝒜_{0,q}^c\end{array}$$
(où dans les lignes horizontales de ce diagramme $`𝒜_{0,q}^c`$ agit sur $`𝒜_{0,q}^c`$ par son produit habituel et dans les lignes verticales on opère par l’application $`\beta `$).
2) L’application $`\beta `$ restreinte à l’espace $`𝖵^q`$ est une représentation de l’algèbre $`U(sl(2)_q)`$. En outre les générateurs de $`\beta (𝖵^q)`$ vérifie l’identité (24).
$`Preuve:`$
2) Elle se fait en deux étapes :
Etape 1. Nous allons réaliser les opérateurs $`U^q,V^q,W^q`$ comme une série de représentations de l’algèbre $`𝒜_{0,q}^c`$ sur les composantes $`𝖵_k^q`$.
Soit
$$\begin{array}{cccc}P_k^q:& 𝖵^k& & 𝖵_k^q\hfill \\ & z& & P_k^qz\hfill \end{array}$$
le projecteur tressé correspondant à son analogue classique noté $`P_k`$. $`P_k^q`$ est un morphisme dans la catégorie $`U_q(sl(2))Mod`$. Considérons également l’application
$$\rho _k^q(z)v=\alpha _kP_k^q(\rho ^q(z)id_{k1})v,zsl(2)_q,v𝖵_k^q,\alpha _k𝐊$$
(25)
$`id_{k1}`$ est l’opérateur identité sur l’espace $`𝖵^{(k1)}`$. Notons que dans la formule (25), $`\alpha _k`$ est une constante quelconque. En supposant que l’application $`\rho _k^q`$ soit une représentation de l’algèbre $`𝒜_{0,q}^c`$, cela impose le choix de la constante $`\alpha _k`$ donnée par la propositon suivante
###### Proposition 4.2
(\[G2\])
$$\alpha _k=(q^1+q^3)(1+q^2+\mathrm{}+q^{2(k1)})(q^1+q^{2k+1})^1.$$
Ainsi à l’élément de base $`u`$ de l’espace $`𝖵`$, on associe la famille d’opérateurs $`U_k^q`$ définie sur le $`U_q(sl(2))`$-module irréductible de spin $`k`$, $`𝖵_k^q𝖵^k`$ :
$$U_k^q:𝖵_k^q𝖵_k^q$$
$$U_k^q(v)=\alpha _kP_k^q(U^qid_{k1})v,v𝖵_k^q$$
$`\alpha _k`$ est donnée par la proposition 4.2. La famille d’opérateurs $`U_k^q`$ ($`k𝐍`$) définie l’extension (encore notée $`U^q`$) de l’opérateur $`U^q`$ (initialement définit sur $`𝖵`$) sur l’algèbre $`𝒜_{0,q}^c`$
$$U^q:𝒜_{0,q}^c𝒜_{0,q}^c.$$
On prolonge de la même manière, les opérateurs $`V^q,W^q`$. Nous montrons dans la deuxième étape que (pour toute valeur de la constante $`\alpha _k`$) l’équation (24) est encore satisfaite par les opérateurs prolongés $`U^q,V^q,W^q`$.
Pour le faire montrons d’abord que l’élément $`(U^qid_{k1})v_k^q`$ (où $`v_k^q𝖵_k^q`$), réduit en forme de base dans l’algèbre $`𝒜_{0,q}^c`$ est identique (à un facteur près) à l’élément $`P_k^q(U^qid_{k1})v_k^q`$. Pour cela, il suffit en fait de montrer que dans la forme réduite de l’élément $`(U^qid_{k1})v_k^q`$, il n’y a pas d’éléments appartenant aux composantes $`𝖵_i^q,i<k`$. L’élément $`(U^qid_{k1})v_k^q`$ résulte de l’application de l’opérateur $`[,]_q`$ à l’élément $`uv_k^qsl(2)_q𝖵_k^q`$. Cette opération commute avec l’action du GQ $`U_q(sl(2))`$.
Dans la décomposition en $`U_q(sl(2))`$-modules du produit $`sl(2)_q𝖵_k^q`$ il y a trois composantes irréductibles : $`𝖵_{k+1}^q,𝖵_k^q,𝖵_{k1}^q`$. En tenant compte du fait que l’élément $`(U^qid_{k1})v_k^q𝖵^k`$ et que dans la réduction de base d’un élément de $`𝖵^k`$ seulement les composantes $`𝖵_k^q,𝖵_{k2}^q,𝖵_{k4}^q`$ peuvent apparaître, nous en concluons que l’élément $`(U^qid_{k1})v_k^q𝖵_k^q`$. Dans le raisonnement précédent en remplaçant $`U^q`$ par $`V^q`$ ou par $`W^q`$ on obtient la même conclusion.
Etape 2. Ayant finalement réalisé $`U^q,V^q,W^q`$ comme des opérateurs sur l’algèbre $`𝒜_{0,q}^c`$, nous montrons à présent que l’équation (24) reste encore vraie sur toute l’algèbre $`𝒜_{0,q}^c`$. En effet, soit $`u_i\{u,v,w\}`$ et $`g𝒜_{0,q}^c`$. L’élément $`g`$ s’écrit :
$$g=\underset{k}{}P_k^q(g)=\underset{k}{}g_k,g_k=P_k^q(g)\text{avec}g_k𝖵_k^q𝖵^k.$$
Ecrivons également $`g_k`$ sous la forme $`u_ig_{k1}`$. Pour $`k`$ fixé, nous avons :
$$[(q^3+q)uW^q+vV^q+(q+q^1)wU^q](g_k)=()$$
$$=\alpha _k[(q^3+q)uP_k^q(W^qid)+vP_k^q(V^qid)+(q+q^1)wP_k^q(U^qid)](g_k)$$
$`()`$ (écrit dans la base de $`𝒜_{0,q}^c`$) donne (à un facteur constant près) :
$$()=\alpha _k[(q^3+q)P_k^q(uW^qid)+P_k^q(vV^qid)+(q+q^1)P_k^q(wU^qid)](g_k)$$
$$=\alpha _kP_k^q([(q^3+q)uW^qu_i+vV^qu_i+(q+q^1)wU^qu_i]id_{k1}(g_{k1})=\mathrm{\hspace{0.17em}0}.$$
La dernière égalité étant due à l’équation (24).
1) Elle est une conséquence immédiate de la façon dont les opérateurs $`U^q,V^q,W^q`$ ont été construits dans le 2).
###### Définition 4.2
Les opérateurs $`U^q,V^q,W^q`$ et toutes leurs combinaisons linéaires à coefficients dans l’algèbre $`𝒜_{0,q}^c`$ sont appelés les champs de vecteurs tressés gauches. Ils sont de la forme
$$aU^q+bV^q+cW^q,a,b,c𝒜_{0,q}^c.$$
Remarque 4.2.2 Pour $`q=1`$, les opérateurs $`U^1,V^1,W^1`$ coïncident respectivement avec les champs de vecteurs (ou les rotations hyperboliques infinitésimales) $`U,V,W`$ sur l’hyperboloïde.
De la même manière, on définit $`\overline{U}^q,\overline{V}^q,\overline{W}^q`$ les champs de vecteurs tressés droits associés respectivement aux éléments de base $`u,v,w`$ de $`𝖵`$. Il est facile de vérifier que dans l’algèbre $`𝒜_{0,q}^c`$, ces opérateurs satisfont à une relation analogue au (24) :
$$(q^3+q)\overline{U}^qw+\overline{V}^qv+(q+q^1)\overline{W}^qu=\mathrm{\hspace{0.17em}0}.$$
Par le théorème 4.1, $`T(\mathrm{H}_q)`$ est le $`𝒜_{0,q}^c`$-module tangent des champs de vecteurs tressés sur l’hyperboloïde quantique.
Remarque 4.2.3 Par la méthode développée dans \[LS\] qui consiste à décrire l’algèbre enveloppante tressée à partir de la REA, on aurait pû traiter facilement les opérateurs adjoints tressés correspondants aux générateurs $`u,v,w`$ comme des opérateurs sur l’algèbre $`𝒜_{0,q}^c`$. Mais la différence fondamentale avec notre méthode est que les opérateurs provenant de la méthode suggérée dans \[LS\] ne permettent pas de contrôler l’identité (24) sur l’algèbre $`𝒜_{0,q}^c`$. À notre connaissance à part notre façon de définir les opérateurs $`U^q,V^q,W^q`$ il n’en existe apparement pas d’autre qui puisse contrôler (24).
###### Définition 4.3
Le plongement
$$sl(2)_qT(\mathrm{H}_q)$$
est appelé une ancre quantique.
Notons que le plongement de la définition 4.3 est la $`q`$-analogue de celui définit par (6), qui est l’exemple le plus simple d’une ancre. Rappellons qu’une ancre est constituée d’une variété $``$ d’une algèbre de Lie $`𝗀`$ et d’un plongement de $`𝗀`$ dans l’espace des champs de vecteurs sur $``$. C’est la raison principale pour laquelle nous appellons le plongement de la définition (4.3) ancre quantique et cela en dépit du fait que le $`𝒜_{0,q}^c`$-module $`T(\mathrm{H}_q)`$ n’est muni d’aucun crochet de Lie tressé. Nous considérons également le couple ($`T(\mathrm{H}_q),𝒜_{0,q}^c`$) comme une $`q`$-analogue d’une algèbre de Lie-Rinehart \[R\] partielle (partielle est associé au fait que $`T(\mathrm{H}_q)`$ n’est muni d’aucune structure de Lie tressé).
Remarque 4.2.4 Quant à la description c) qui décrivait l’espace tangent sur l’hyperboloïde comme une variété algébrique affine, pour la $`q`$-déformé il est nécéssaire de trouver la $`q`$-analogue de l’algèbre symétrique de l’espace tangent $`T(\mathrm{H}_q)`$ tel qu’il soit une déformation plate de son analogue classique. Le problème fondamentale qui se pose pour l’existence de cette algèbre symétrique déformée est de trouver une façon raisonnable de transposer les éléments de l’algèbre $`𝒜_{0,q}^c`$ et de l’espace $`𝖵^q`$. Le seul bon candidat susceptible de pouvoir réaliser une telle transposition est l’opérateur de Yang-Baxter (YB) quantique provenant de la R-matrice universelle du GQ $`U_q(sl(2))`$. Malheureusement cette méthode conduit à une déformation non plate de l’algèbre symétrique (classique).
Remarque 4.2.5 Une façon d’introduire dans le cas classique les champs de vecteurs sur une variété algébrique affine, consiste à définir d’abord les champs de vecteurs dans l’espace ambiant comme toutes les combinaisons linéaires à coefficients-fonctions (ici les fonctions dans l’espace ambiant) de dérivées partielles. Ensuite, on définit les champs de vecteurs sur la variété donnée comme de tels champs de vecteurs qui respectent les équations définissant la variété en question. L’autre façon consiste au passage aux cartes, i.e. au considération locale.
Malheureusement nous ne connaissons pas de $`q`$-analogues des dérivées partielles. Par conséquent nous ne savons pas définir les champs de vecteurs tressés sur l’espace $`sl(2)_{q}^{}{}_{}{}^{}`$ tout entier, puisque les dérivées partielles tressées ne sont pas définies. C’est la raison principale pour laquelle nous avons introduit les champs de vecteurs tressés sur l’hyperboloïde quantique à partir du crochet de Lie tressé.
Nous montrons maintenant que sur la sphère (ou l’hyperboloïde), la notion de champs de vecteurs définie à partir des dérivées partielles et celle définie à partir des champs de vecteurs adjoints sont équivalentes.
En effet, il est clair que les champs de vecteurs $`X,Y,Z`$ sur la sphère définis par (4) tels que
$$X(x^2+y^2+z^2R^2)=Y(x^2+y^2+z^2R^2)=Z(x^2+y^2+z^2R^2)=\mathrm{\hspace{0.17em}0}$$
s’expriment en fonction des dérivées partielles $`_x,_y,_z`$. Montrons la réciproque. Cela revient à montrer que pour tout champ de vecteurs $`𝒳`$ sur la sphère de la forme
$$𝒳=\alpha _x+\beta _y+\gamma _z,$$
$`\alpha ,\beta ,\gamma \mathrm{Fun}(\mathrm{S}^2)`$ sont tels que
$$𝒳(x^2+y^2+z^2R^2)=0\text{i.e.}\alpha x+\beta y+\gamma z=\mathrm{\hspace{0.17em}0},$$
(26)
$`𝒳`$ est une combinaison linéaire à coefficients-fonctions (ici dans
$`\mathrm{Fun}(\mathrm{S}^2)`$) des champs de vecteurs $`X,Y,Z`$. Notons que la condition (26) signifie que le champ de vecteurs $`𝒳`$ est tangent à la sphère.
###### Proposition 4.3
Il existe $`k,l,m\mathrm{Fun}(\mathrm{S}^2)`$ tels que :
$$𝒳=\alpha _x+\beta _y+\gamma _z=kX+lY+mZ.$$
(27)
$`Preuve:`$ En appliquant (27) respectivement à $`x,y,z`$ les fonctions $`\alpha ,\beta ,\gamma `$ se mettent sous la forme
$$\alpha =zl+ym,\beta =zkxm,\gamma =yk+xl.$$
(28)
Considérons les fonctions $`k,l,m`$ définies par
$$k=\frac{1}{R^2}(\beta z\gamma y),l=\frac{1}{R^2}(\alpha z+\gamma x),m=\frac{1}{R^2}(\alpha y\beta x).$$
(29)
Par l’équation (26), il est facile de vérifier que $`(k,l,m)`$ ainsi donné vérifie (27).
Par le changement de base :
$$u=i(x+iy),v=\sqrt{2}z,w=i(xiy),$$
on se ramène au cas de l’hyperboloïde. On peut remarquer qu’un champ de vecteurs
$$𝒳=\alpha _u+\beta _v+\gamma _w$$
“respecte” l’équation
$$2uw+\frac{v^2}{2}c=0$$
(i.e. $`\mathrm{\hspace{0.17em}\hspace{0.17em}2}\alpha w+\beta v+2\gamma u=0`$) si et seulement si $`𝒳`$ peut être présenté sous la forme :
$$𝒳=kU+lV+mW$$
$`U,V,W`$ sont les rotations hyperboliques infinitésimales.
## 5 Projectivité du $`𝒜_{0,q}^c`$-module $`T(\mathrm{H}_q)`$
Il existe plusieurs définitions équivalentes d’un module projectif (voir par exemple \[L\]). En remplaçant homomorphisme par homomorphisme dans la catégorie $`U_q(sl(2))Mod`$ et module par module tressé, nous adoptons ces définitions et propriétés pour notre cas tressé.
Pour montrer que le module tangent sur l’hyperboloïde quantique est un module projectif, nous traitons juste le cas classique, i.e. le cas de la sphère puis nous en déduisons les résultats dans notre cas tressé.
Comme nous l’avons déja fait remarquer, bien que le module tangent $`T(S^2)`$ ne soit pas un $`\mathrm{Fun}(\mathrm{S}^2)`$-module libre, il est par contre un $`\mathrm{Fun}(\mathrm{S}^2)`$-module projectif. En effet posons ici $`A=\mathrm{Fun}(\mathrm{S}^2)`$. Soit ($`X,Y,Z`$) la base du $`A`$-module $`A^3`$ et $`N`$ le $`A`$-module de type fini engendré par l’élément
$$xX+yY+zZ:=(x,y,z)$$
$`N`$ est un sous-module de $`A^3`$. Notons $`\overline{N}`$ le sous-module de $`A^3`$ engendré par les éléments
$$yXxY:=(y,x,0),zYyZ:=(0,z,y),$$
$$xZzX:=(z,0,x).$$
###### Proposition 5.1
En sens de $`A`$-module on a :
$$A^3=N\overline{N}$$
$`Preuve:`$ Notons $`\mathrm{Q}`$ le projecteur de $`A^3`$ tel que $`\mathrm{Im}\mathrm{Q}=\mathrm{N}`$ et défini par : $`(f,g,h)A^3`$
$`\mathrm{Q}(f,g,h)`$ $`=`$ $`R^2(fx+gy+hz)(x,y,z)`$
$`=`$ $`R^2(fx+gy+hz)(xX+yY+zZ).`$
Montrer que l’intersection des deux $`A`$-modules $`N`$ et $`\overline{N}`$ est réduite à zéro, revient à montrer que $`\mathrm{Ker}\mathrm{Q}=\overline{\mathrm{N}}`$. Nous avons :
$$\mathrm{Q}(y,x,0)=\mathrm{Q}(0,z,y)=\mathrm{Q}(z,0,x)=0$$
par conséquent $`\overline{N}\mathrm{Ker}\mathrm{Q}`$. Il reste à prouver que $`\mathrm{Ker}\mathrm{Q}\overline{\mathrm{N}}`$.
Soit $`𝒳\mathrm{Ker}\mathrm{Q}`$, alors il existe $`(\alpha ,\beta ,\gamma )A^3`$ tel que
$$𝒳=\alpha X+\beta Y+\gamma Z\text{et}\mathrm{Q}(𝒳)=0.$$
La deuxième condition ($`\mathrm{Q}(𝒳)=0`$) entraîne
$$\alpha x+\beta y+\gamma z=\mathrm{\hspace{0.17em}0}.$$
Ainsi si $`𝒳\mathrm{Ker}\mathrm{Q}`$, cela revient à dire que $`𝒳`$ est de la forme :
$$𝒳=\alpha X+\beta Y+\gamma Z\text{avec}\alpha x+\beta y+\gamma z=\mathrm{\hspace{0.17em}0}.$$
Donc montrer qu’un tel champ $`𝒳`$ appartient à $`\overline{N}`$ revient à montrer la proposition 4.3 (il suffit de remplacer dans cette proposition, les dérivées partielles $`_x,_y,_z`$ respectivement par les champs $`X,Y,Z`$). Par conséquent on a bien $`\mathrm{Ker}\mathrm{Q}\overline{\mathrm{N}}`$.
En outre comme $`(f,g,h)=f(1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0})+g(0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}0})+h(0,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1}),`$ il est nécéssaire et suffisant de connaître cette décomposition pour les éléments $`X:=(1,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0});Y:=(0,\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}0});Z:=(0,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}1})`$ en vue de pouvoir décomposer tout élément du $`A`$-module $`A^3`$ comme somme d’un élément du $`A`$-module $`N`$ et du $`A`$-module $`\overline{N}`$. Par un calcul directe nous avons :
$`\mathrm{Q}(X)`$ $`=`$ $`R^2x(xX+yY+zZ)`$
$`\mathrm{Q}(Y)`$ $`=`$ $`R^2y(xX+yY+zZ)`$
$`\mathrm{Q}(Z)`$ $`=`$ $`R^2z(xX+yY+zZ).`$
Considérons en outre le projecteur noté $`\mathrm{P}`$, de $`A^3`$ sur le $`A`$-module $`\overline{N}`$ :
$$\mathrm{P}=id\mathrm{Q},$$
nous obtenons également par un calcul directe :
$`\mathrm{P}(X)`$ $`=`$ $`R^2(z(xZzX)y(yXxY))`$
$`\mathrm{P}(Y)`$ $`=`$ $`R^2(x(yXxY)z(zYyZ))`$
$`\mathrm{P}(Z)`$ $`=`$ $`R^2(y(zYyZ)x(xZzX)).`$
D’où la preuve de la proposition.
Nous avons alors les identifications suivantes :
$$A^3/\{N\}\overline{N},A^3/\{\overline{N}\}N.$$
Donc le $`A`$-module $`T(S^2)=A^3/\{N\}`$ est réalisé comme un sous-module ($`\overline{N}`$) de $`A^3`$ ayant un module supplémentaire ($`N`$).
Faisons enfin remarquer qu’en prenant $`A=\mathrm{Fun}(\mathrm{H}),`$ $`N`$ le $`A`$-module engendré par $`\mathrm{\hspace{0.17em}2}wU+vV+2uW`$ et $`\overline{N}`$ celui engendré par les générateurs de $`(𝖵𝖵^{})_1,`$ la proposition 5.1 reste également valable sur l’hyperboloïde.
Passons à présent au cas tressé. L’analogue tressé des $`\mathrm{Fun}(\mathrm{H})`$-modules $`N`$ et $`\overline{N}`$ sont respectivement le $`𝒜_{0,q}^c`$-module (disons gauche) $`N_l^q`$ et $`\overline{N}_l^q`$. Rappelons que $`N_l^q`$ est engendré par
$$(q^3+q)uW^q+vV^q+(q+q^1)wU^q$$
et précisons que $`\overline{N}_l^q`$ est engendré par les éléments :
$$q^2uV^qvU^q,(q^3+q)(uW^qwU^q)+(1q^2)vV^q,q^2vW^q+wV^q.$$
###### Proposition 5.2
En sens de $`𝒜_{0,q}^c`$-module on a :
$$(𝒜_{0,q}^c)^3=N_l^q\overline{N}_l^q$$
$`Preuve:`$ Pour montrer que l’intersection de $`N_l^q`$ et $`\overline{N}_l^q`$ est réduite à zéro, notons $`\mathrm{Q}_q`$ l’analogue tressé du projecteur $`\mathrm{Q}`$. $`\mathrm{Q}_q`$ est la projection de $`(𝒜_{0,q}^c)^3`$ telle que $`\mathrm{Im}\mathrm{Q}_\mathrm{q}=\mathrm{N}_\mathrm{l}^\mathrm{q}`$ et définie par $`g,h,k𝒜_{0,q}^c`$
$`\mathrm{Q}_q(g,h,k)=c^1[gu+hv+kw][(q^3+q)uW^q+vV^q+(q+q^1)wU^q]`$
Considérons “l’analogue tressé de la proposition 4.3” i.e. en remplaçant dans l’équation (27) les dérivées partielles par les générateurs de $`𝖵^q`$ et les champs $`X,Y,Z`$ par les générateurs de $`(𝖵𝖵^q)_1`$. En outre l’analogue tressé de la condition (26) s’écrit
$$(\alpha U^q+\beta V^q+\gamma W^q)(𝒞_qc)=\mathrm{\hspace{0.17em}0}$$
avec $`\alpha ,\beta ,\gamma 𝒜_{0,q}^c`$. Alors les coefficients $`k,l,m𝒜_{0,q}^c`$ de l’analogue tressé de la proposition 4.3 sont donnés (au facteur constant $`c^1`$ près) par les générateurs de l’espace $`(𝐕𝖵)_1`$ où l’espace $`𝐕`$ est celui engendré par les éléments $`\alpha ,\beta ,\gamma `$.
Le fait que $`\mathrm{Ker}(\mathrm{Q}_\mathrm{q})=\overline{\mathrm{N}}_\mathrm{l}^\mathrm{q}`$ découle (comme dans le cas classique) de “l’analogue tressé de la proposition 4.3”.
Comme pour la sphère (ou l’hyperboloïde classique) tout élément du $`𝒜_{0,q}^c`$-module $`(𝒜_{0,q}^c)^3`$ est la somme d’un élément de $`N_l^q`$ et de $`\overline{N}_l^q`$. En effet nous avons :
$`\mathrm{Q}_q(U^q)`$ $`=`$ $`c^1u((q^3+q)uW^q+vV^q+(q+q^1)wU^q)`$
$`\mathrm{Q}_q(V^q)`$ $`=`$ $`c^1v((q^3+q)uW^q+vV^q+(q+q^1)wU^q)`$
$`\mathrm{Q}_q(W^q)`$ $`=`$ $`c^1w((q^3+q)uW^q+vV^q+(q+q^1)wU^q).`$
En considérant en outre l’analogue tressé noté $`\mathrm{P}_q`$ du projecteur $`\mathrm{P}`$ qui est tel que $`\mathrm{Im}(\mathrm{P}_\mathrm{q})=\overline{\mathrm{N}}_\mathrm{l}^\mathrm{q}`$ et évidemment définit par
$$\mathrm{P}_q=id\mathrm{Q}_q$$
nous obtenons alors par un calcul directe :
$`\mathrm{P}_q(U^q)=`$ $`q^2c^1\{q^2u[(q^3+q)(uW^qwU^q)+(1q^2)vV^q]`$
$`v(vU^qq^2uV^q)\}`$
$`\mathrm{P}_q(V^q)=`$ $`q^2c^1\{(q^3+q)u(q^2vW^q+wV^q)(q^3+q)(q^2uV^q`$
$`vU^q)(1q^2)v[(q^3+q)(uW^qwU^q)+(1q^2)vV^q]\}`$
$`\mathrm{P}_q(W^q)=`$ $`q^2c^1\{q^2v(q^2vW^qwV^q)w[(q^3+q)(uW^qwU^q)+`$
$`+(1q^2)vV^q]\}`$
Ainsi nous venons de prouver (comme dans le cas classique) que le $`𝒜_{0,q}^c`$-module $`T(\mathrm{H}_q)_l`$ est un $`𝒜_{0,q}^c`$-module projectif. (De même le $`𝒜_{0,q}^c`$-module $`T(\mathrm{H}_q)_r`$ est projectif).
## 6 Métrique et Connexion tressées
Nous définissons et montrons l’existence d’une (pseudo)métrique et d’une connexion (partiellement définie) tressées sur le module tangent $`T(\mathrm{H}_q)`$.
### 6.1 (Pseudo)métrique tressée
Rappelons que “pseudo” signifie que son analogue classique n’est pas définie positive. Par la suite nous omettrons cette précision.
###### Définition 6.1
L’opérateur
$$\begin{array}{cccc}<,>:& T(\mathrm{H}_q)_l_𝐊T(\mathrm{H}_q)_r& & 𝒜_{0,q}^c\hfill \\ & ab& & <ab>=<a,b>\hfill \end{array}$$
est une métrique tressée si : $`PT(\mathrm{H}_q)_l,QT(\mathrm{H}_q)_r,f𝒜_{0,q}^c`$
$$<fP,Q>=f<P,Q>,<P,Qf>=<P,Q>f$$
(30)
et $`<,>`$ est $`U_q(sl(2))`$-covariant, c’est-à-dire $`zU_q(sl(2))`$ :
$`z.<a,b>`$ $`=`$ $`<,>\mathrm{\Delta }(z).(ab)=<,>(z_{(1)}z_{(2)})(ab)`$
$`=`$ $`<z_{(1)}.a,z_{(2)}.b>\text{avec}(a,b)T(\mathrm{H}_q)_l\times T(\mathrm{H}_q)_r.`$
Si en outre on a :
$$<,>(𝖵^q\overline{𝖵}^q)_1=\mathrm{\hspace{0.17em}0},$$
(31)
la métrique tressée est dite $`q`$-symétrique. ($`\overline{𝖵}^q`$ coïncide avec $`𝖵^q`$ comme un espace vectoriel. Mais il engendre $`T(\mathrm{H}_q)`$ comme un $`𝒜_{0,q}^c`$-module droit.)
Soulignons que pour le moment nous introduisons la métrique tressée sur $`T(\mathrm{H}_q)_l_𝐊T(\mathrm{H}_q)_r`$. Nous pourrons la prolongée sur $`T(\mathrm{H}_q)_ϵ_𝐊T(\mathrm{H}_q)_ϵ`$ seulement après l’identification des modules tangents $`T(\mathrm{H}_q)_l`$ et $`T(\mathrm{H}_q)_r`$ (voir section 7.1).
###### Théorème 6.1
Il existe une unique (à un facteur près) métrique tressée et $`q`$-symétrique sur le module tangent $`T(\mathrm{H}_q)`$. Cette métrique tressée restreinte sur $`𝖵^q\overline{𝖵}^q`$ fournit la table suivante :
$$<U^q,\overline{U}^q>=uu,<U^q,\overline{V}^q>=uv,<V^q,\overline{U}^q>=vu,$$
$$<W^q,\overline{V}^q>=wv,<V^q,\overline{V}^q>=(1q^2)vvq^1(1+q^2)^2uw,$$
$$<W^q,\overline{W}^q>=ww,<U^q,\overline{W}^q>=q^1(1+q^2)^1vvq^2uw,$$
$$<V^q,\overline{W}>=vw,<W^q,\overline{U}^q>=q(1+q^2)^1vvq^2wu.$$
$`Preuve:`$ Décrivons d’abord tous les couplages
$$<,>:𝖵^q\overline{𝖵}^q𝒜_{0,q}^c$$
$`U_q(sl(2))`$-covariants. Pour cela, nous décomposons $`𝖵^{}\overline{𝖵}^{}`$ en $`U_q(sl(2))`$-modules irréductibles de dimension finie. La $`U_q(sl(2))`$-covariance impose les conditions :
$$<,>(𝖵^q\overline{𝖵}^q)_2=k𝖵_2^q,<,>(𝖵^q\overline{𝖵}^q)_0=\gamma $$
(32)
$`k`$ et $`\gamma `$ sont des constantes. (La $`q`$-symétrie exigera en outre la relation (31).)
La seconde étape de la démonstration consiste à déterminer les relations de dépendance entre les paramètres $`k`$ et $`\gamma `$. Nous la faisons par le biais des relations caractérisant le module tangent tressé gauche et le module tangent tressé droit. Cela revient à vérifier que la relation de dépendance définie par
$$<,>^{23}(𝖵𝖵^q)_0\overline{𝖵}^q=0$$
(33)
est compatible avec celle définie par
$$<,>^{12}𝖵^q(\overline{𝖵}^q𝖵)_0=0.$$
(34)
On étend enfin ce couplage à $`T(\mathrm{H}_q)_l_𝐊T(\mathrm{H}_q)_r`$ en utilisant la propriété (30). (Voir appendice A pour la forme explicite des équations symboliques).
Nous avons déterminés le couplage $`<,>`$ avec la condition (33). Mais on vérifie que la condition (34) est satisfaite pour ce couplage ainsi déterminé.
### 6.2 Connexion tressée
Rappelons d’abord que dans le cas classique ($`q=1`$), une connexion linéaire sur l’espace des champs de vecteurs $`E`$ d’une variété algébrique régulière (ou plus généralement d’une variété lisse) $``$ est l’application notée $``$
$$\begin{array}{cccc}:& EE& & E\hfill \\ & ab& & _ab\hfill \end{array}$$
$`𝐊`$-linéaire et satisfaisant aux propriétés suivantes :
1. $`_{fa}b=f_ab,(a,b)E^2,f\text{Fun}()`$
2. $`_afb=f_ab+(af)b`$.
Notons que la propriété 2. est la règle de dérivation de Leibniz pour la connexion $``$.
Lorsque la connexion $``$ est sans torsion (par exemple connexion de Levi-Civita), on a en plus
$$[a,b]=_ab_ba$$
(où $`[,]`$ désigne le crochet de Lie des champs de vecteurs $`a`$ et $`b`$), on en déduit alors :
$`_afb`$ $`=`$ $`_{fb}a[fb,a]`$
$`=`$ $`_{fb}a(f[b,a](af)b)`$
$`=`$ $`_{fb}a(_{fb}a_{fa}b(af)b)`$
$`=`$ $`f_ab+(af)b.`$
Ainsi si la connexion $``$ est sans torsion, la propriété 2. découle de la propriété 1. Ceci nous permet de nous passer de la règle de dérivation de Leibniz pour la connexion si cette dernière est sans torsion.
Nous généralisons la notion de connexion linéaire à notre cas tressé. Mais nous construisons dans ce cas plutôt une connexion partiellement définie sur le module tangent tressé. (“Partiellement définie” signifie qu’elle est définie sur un sous-ensemble de $`T(\mathrm{H}_q)T(\mathrm{H}_q)`$).
###### Définition 6.2
L’opérateur
$$\begin{array}{cccc}:& T(\mathrm{H}_q)_l_𝐊𝖵^q& & T(\mathrm{H}_q)_l\hfill \\ & ab& & _ab\hfill \end{array}$$
est une “connexion” tressée (sans torsion), s’il vérifie les propriét és :
1. $``$ est un morphisme dans la catégorie des $`U_q(sl(2))`$-modules, c’est-à-dire :
$$z._ab=_{z_{(1)}.a}z_{(2)}.b$$
$`aT(\mathrm{H}_q)_l,b𝖵^q,zU_q(sl(2))`$.
2. $`_{fa}b=f_ab,aT(\mathrm{H}_q)_l,b𝖵^q,f𝒜_{0,q}^c`$.
3. $$a^{ij}_{X_i}Y_j=[X_i,Y_j],X_i,Y_j𝖵^q\text{avec}a^{ij}X_iY_j𝖵_1^q.$$
(35)
Notons que (35) est la $`q`$-analogue de la notion de connexion sans torsion. Si nous arrivons en outre à étendre le $`q`$-crochet de Lie $`[,]_q`$ sur $`T(\mathrm{H}_q)`$ tout entier et à comprendre l’algèbre enveloppante de cette algèbre de Lie $`q`$-déformée (nous en avons besoin pour écrire la partie à gauche de la relation (35)), nous pourrons prolonger notre connexion (partiellement définie) sur $`T(\mathrm{H}_q)_ϵ_𝐊T(\mathrm{H}_q)_ϵ`$ en utilisant une analogue de (35). Malheureusement nous ne connaissons aucune façon de le faire.
###### Théorème 6.2
Il existe une connexion tressée $``$ sur le module tangent de l’hyperboloïde quantique (au sens de la définition (6.2). Elle est donnée sur $`𝖵^q`$ par
$`_{U^q}U^q=`$ $`\alpha (uvU^qq^2uuV^q),_{W^q}W^q=\alpha (wwV^qq^4vwW^q),`$
$`_{U^q}V^q=`$ $`\beta \{\{\alpha [(q^3+q)uwq^2vv]2q^2\}U^q\alpha (q^6+q^21)`$
$`uvV^q\}+\beta \alpha (q^3+q)uuW^q,`$
$`_{V^q}U^q=`$ $`\beta \{\{\alpha q^2[(q^3+q)uwq^2vv]+2\}U^q\alpha q^2(q^6+q^21)`$
$`uvV^q\}+\beta \alpha q^2(q^3+q)uuW^q,`$
$`_{V^q}W^q=`$ $`\beta \{\alpha q^3(1+q^2)wwU^q+\alpha (1+q^4q^6)vwV^q\}+\beta \{\alpha q^4`$
$`[(q+q^1)uwvv]2q^2\}W^q,`$
$`_{W^q}V^q=`$ $`\beta \{\alpha q^5(1+q^2)wwU+\alpha q^2(1+q^4q^6)vw^q\}+\beta \{\alpha q^6`$
$`[(q+q^1)uwvv]+2\}W^q,`$
$`_{W^q}U^q=`$ $`\beta \{\alpha q^6vwU^q+\{q^4(1q^2)\alpha [uw+[2]^1v^2]{\displaystyle \frac{2}{1+q^2}}\}`$
$`V^q\}q^6\beta \alpha uvW^q,`$
$`_{U^q}W^q=`$ $`\beta \{\alpha q^2vwU^q+\{(q^21)\alpha [uw+[2]^1v^2]+{\displaystyle \frac{2q}{1+q^2}}\}V^q`$
$`q^2\alpha uvW^q\},`$
$`_{V^q}V^q=`$ $`\beta \{\alpha q^3(1+q^2)^2vwU^q+\{q(1+q^2)(1q^4)[uw+[2]^1`$
$`v^2]+2(1q^2)\}V^q\}+q^3(1+q^2)^2\alpha uvW^q\},`$
$`\alpha ={\displaystyle \frac{2}{(1q^2+q^4)c}},\beta =(1+q^4)^1.`$
$`Preuve:`$ Comme dans le cas de la métrique tressée, nous décrivons d’abord toutes les applications
$$𝖵^q_𝐊𝖵^qT(\mathrm{H}_q)_l$$
$`U_q(sl(2))`$-covariantes. Cette $`U_q(sl(2))`$-covariance impose les conditions suivantes :
$$(𝖵^q𝖵^q)_2=\alpha (𝖵𝖵^q)_2,(𝖵^q𝖵^q)_0=0,\alpha 𝐊$$
(36)
complétées par les relations qui découlent de (35).
(Précisons que $`(𝖵^q𝖵^q)_2`$ désigne l’ensemble $`_ab`$$`ab`$ parcourent $`(𝖵^q𝖵^q)_2`$.)
Concrètement, $`U^qU^q`$ étant de poids deux et $`X._{U^q}U^q=0`$, considérons
$$_{U^q}U^q=\alpha (uvU^qq^2uuV^q)$$
(37)
$`uvU^qq^2uuV^q`$ est de poids 2. Notons que le choix $`_{U^q}U^q=uU^q`$ n’est pas compatible avec la relation
$$_K=(q^3+q)u_{W^q}+v_{V^q}+(q+q^1)w_{U^q}=0\text{dans}T(\mathrm{H}_q)_l.$$
Notons
$$J_i=\frac{1}{[i]_q}Y^iJ_{i1},i\{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}\},[i]_q=\frac{q^iq^i}{qq^1},\text{}$$
$$J_0=_{U^q}U^q=\alpha (uvU^qq^2uuV^q).$$
Par application de l’opérateur $`Y^i`$ à la relation (37) nous obtenons les quatre relations suivantes :
$$_{U^q}V^qq^2_{V^q}U^q=J_1,$$
$$_{U^q}W^q+q_{V^q}V^qq^4_{W^q}U^q=J_2,$$
$$_{V^q}W^q+q^2_{W^q}V^q=J_3,$$
$$_{W^q}W^q=J_4.$$
(Voir appendice B pour les formes explicites). Le choix de $`\alpha `$ est précisé par la deuxième équation du (36). (voir appendice B )
Remarque 6.2 Bien que nous ne savons pas définir la $`q`$-analogue de la notion de courbure (dans le cas classique elle est introduite localement) dans le cadre de notre approche globale, on peut déviner (à un facteur près) la forme de courbure correspondante en supposant que cette forme soit une déformation plate de son analogue classique. Sur la sphère, cette forme est donnée (à un facteur près) par
$$x(dydzdzdy)+y(dzdxdxdz)+z(dxdydydx).$$
Ainsi dans le cas tressé, pour obtenir la forme de courbure, il suffit de remplacer dans l’expression $`(q^3+q)uw+vv+(q+q^1)wu`$ les facteurs de droites respectivement par la $`q`$-analogue des formes $`dvdwdwdv,dwdududw,dudvdvdu`$.
## 7 Identification de $`T(\mathrm{H}_q)_l`$ et $`T(\mathrm{H}_q)_r`$
Nous nous intéressons à présent au problème qui consiste à identifier les modules tangents gauche $`T(\mathrm{H}_q)_l`$ et droit $`T(\mathrm{H}_q)_r`$. Cette identification est nécessaire, car elle permet d’étendre la métrique tressée sur $`T(\mathrm{H}_q)_l`$ (ou $`T(\mathrm{H}_q)_r`$).
Le fait que l’algèbre $`𝒜_{0,q}^c`$ ne soit munie d’aucune structure d’involution et que l’opérateur de tresse provenant de la $`R`$-matrice universelle du GQ $`U_q(sl(2))`$ ne soit pas involutif d’autre part, nous ont amenées à suggérer d’identifier les $`𝒜_{0,q}^c`$-modules $`T(\mathrm{H}_q)_l`$ et $`T(\mathrm{H}_q)_r`$ de la façon suivante – nous construisons une base de chacun des modules tangents, – puis nous construisons une application (dans la catégorie $`U_q(sl(2))Mod`$) entre $`T(\mathrm{H}_q)_l`$ et $`T(\mathrm{H}_q)_r`$ qui coïncide (pour $`q=1`$) avec l’application définie par la volte.
### 7.1 Base du $`𝒜_{0,q}^c`$-module $`T(\mathrm{H}_q)`$
Notons d’abord que la méthode de construction de cette base étant aussi bien valable dans le cas classique que dans notre cas tressé, nous considérons dans ce qui suit indifférement ces deux cas.
En considérant les modules
$$𝖵𝖵^{},𝖵_k𝖵^{},k=2,\mathrm{\hspace{0.17em}3},\mathrm{}$$
où les $`𝖵_k`$ (composantes de base de l’algèbre $`𝒜_{0,1}^c`$) sont des $`U(sl(2))`$-modules, nous déterminons les composantes qui “survivent” dans le module tangent. Il est évident que dans le produit $`𝖵𝖵^{}`$ seulement deux composantes survivent. À savoir :
$$(𝖵𝖵^{})_1,(𝖵𝖵^{})_2$$
puisque par construction la composante $`(𝖵𝖵^{})_0`$ est nulle dans le module tangent.
De façon analogue, dans le produit $`𝖵_2𝖵^{}`$, les composantes $`(𝖵_2𝖵^{})_3,(𝖵_2𝖵^{})_2`$ “survivent” et il n’est pas difficile de montrer que la composante $`(𝖵_2𝖵^{})_1`$ est nulle modulo les termes de $`𝐊𝖵^{}=𝖵^{}`$. En effet, par construction, les éléments de $`\mu ^{12}(𝖵(𝖵𝖵^{})_0)`$ sont nuls dans le module tangent. En réduisant en outre tout élément du produit $`𝖵𝖵`$ à sa forme canonique, il est la somme d’un élément de $`𝖵_2𝖵^2`$ et d’un élément de $`𝐊`$. D’où la preuve de ce dernier fait.
Ainsi de façon générale, dans le produit $`𝖵_k𝖵^{}`$ qui se décompose de la manière suivante
$$𝖵_k𝖵^{}=(𝖵_k𝖵^{})_{k+1}(𝖵_k𝖵^{})_k(𝖵_k𝖵^{})_{k1},$$
la composante $`(𝖵_k𝖵^{})_{k1}`$ est nulle modulo les termes appartenant à $`𝖵_j𝖵^{},j<k1`$.
Nous avons donc montré le fait suivant :
###### Proposition 7.1
Une base dans le module tangent gauche $`T(\mathrm{H}_q)_l`$ est formée par les $`U_q(sl(2))`$-modules
$$1.𝖵^q,2.(𝖵𝖵^q)_{1,2},3.(𝖵_2^q𝖵^q)_{2,3},4.(𝖵_3^q𝖵^q)_{3,4}\text{etc}.$$
On construit de façon similaire une base du module tangent droit $`T(\mathrm{H}_q)_r`$.
Identifions à présent les modules tangents $`T(\mathrm{H}_q)_l`$ et $`T(\mathrm{H}_q)_r`$ en définissant :
$$\alpha :T(\mathrm{H}_q)_lT(\mathrm{H}_q)_r$$
telle que l’application :
$$\alpha :(𝖵𝖵^q)_1(𝖵^q𝖵)_1,(𝖵𝖵^q)_2(𝖵^q𝖵)_2,$$
$$(𝖵_2^q𝖵^q)_2(𝖵^q𝖵_2^q)_2,\mathrm{}$$
soit un $`U_q(sl(2))`$-morphisme unique (à un facteur constant près sur chaque composante). Nous suggérons maintenant une façon “canonique” d’éliminer ce degré de liberté sur les composantes, de la manière suivante. On identifie les éléments de l’espace :
$$(𝖵𝖵^q)_2\text{et}(𝖵^q𝖵)_2,(𝖵_k^q𝖵^q)_{k+1}\text{et}(𝖵^q𝖵_k^q)_{k+1},k=2,3,\mathrm{}$$
qui coïncident quand on remplace $`𝖵^q`$ par $`𝖵`$. Quant aux composantes
$$(𝖵𝖵^q)_1\text{et}(𝖵^q𝖵)_1,(𝖵_k^q𝖵^q)_k\text{et}(𝖵^q𝖵_k^q)_k,$$
leurs éléments sont identifiés si la même opération amène aux images opposées. On peut facilement voir que dans le cas classique, cette identification et celle définie par la volte coïncident : c’est la motivation de notre méthode d’identification canonique.
Remarque 7.1 Lorsque l’algèbre $`𝒜_{0,q}^c`$ est munie d’un opérateur de tresse involutif ($`S^2=id`$), une telle identification est faite de façon similaire au cas classique en remplaçant la volte par $`S`$. Mais pour l’opérateur de tresse non involutif provenant du GQ $`U_q(sl(2))`$ ce n’est plus raisonnable. Considérons $`M_1`$ et $`M_2`$ deux $`𝒜_{0,q}^c`$-modules (gauche par exemple). Le problème réside dans le fait que dans le produit
$$m_1fm_2,m_1M_1,m_2M_2,f𝒜_{0,q}^c$$
il n’existe aucune façon raisonnable de transposer le facteur $`f`$ pour le mettre à gauche de telle manière que le produit tensoriel $`_{𝒜_{0,q}^c}`$ soit associatif et le module $`,M_1_{𝒜_{0,q}^c}M_2`$ soit une déformation plate de son analogue classique en supposant bien sûr que $`M_1`$ et $`M_2`$ soient des déformations plates de leurs analogues classiques respectifs.
En résumé, précisons une fois de plus que notre méthode qui consiste à introduire la métrique en deux étapes :
$``$ en définissant d’abord un couplage sur $`T(\mathrm{H}_q)_l_𝐊T(\mathrm{H}_q)_r`$,
$``$ et ensuite à identifier $`T(\mathrm{H}_q)_l`$ et $`T(\mathrm{H}_q)_r`$,
nous permet de contrôler le fait que notre construction ne soit pas contradictoire (i.e. la platitude a lieu).
Il existe dans la littérature plusieurs constructions de fibré quantique sur la sphère (\[BM\], \[S2\]) et de métrique quantique. Mais pour la plupart d’entre elles, la notion de platitude de ces constructions est carrément ignorée. C’est en cela que nos constructions sont différentes.
## Appendice
### A Hyperboloïde quantique et métrique tressée
Nous avons explicitement :
$`𝖵_0^q=Span((q^3+q)uw+vv+(q+q^1)wu),`$
$`𝖵_1^q=Span(q^2uvvu,(q^3+q)(uwwu)+(1q^2)vv,q^2vw+wv),`$
$`𝖵_2^q=Span(uu,uv+q^2vu,uwqvv+q^4wu,vw+q^2wv,ww).`$
(Dans ces expressions nous avons omis le symbole du produit tensoriel, ce que nous ferons lorsque cela n’apporte aucune confusion.)
$`𝒞_q`$ est le générateur de l’espace de dimension un ($`𝖵_0^q`$) :
$$𝒞_q=(q^3+q)uw+vv+(q+q^1)wu.$$
En posant :
$$x_0=u,x_1=q^2uvvu,x_2=uu$$
nous avons :
$`q^2uvvu=2u\mathrm{},(q^3+q)(uwwu)+(1q^2)vv=2v\mathrm{},`$
$`q^2vw+wv=2w\mathrm{},𝒞_q=(q^3+q)uw+vv+(q+q^1)wu=c,`$
Explicitons à présent les équations symboliques de la Section 6 :
$`<,>(𝖵^q\overline{𝖵}^q)_2=k𝖵_2^q`$.
Sur la composante de spin 2 de $`𝖵^q\overline{𝖵}^q,`$ l’élément $`U^q\overline{U}^q`$ est de poids 2 et nous avons nécessairement par la $`U_q(sl(2))`$-covariance :
$$<U^q,\overline{U}^q>=kuu$$
Par application de l’opérateur $`Y^n`$ (avec $`n\{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}\}`$) à $`<U^q,\overline{U}^q>`$ nous obtenons les quatre relations suivantes :
$$<U^q,\overline{V}^q>+q^2<V^q,\overline{U}^q>=k(q^2vu+uv),$$
$$<U^q,\overline{W}^q>+q<V^q,\overline{V}^q>q^4<W^q,\overline{U}^q>=k(qvvq^4wuuw),$$
$$<V^q,\overline{W}^q>+q^2<W^q,\overline{V}^q>=k(vw+q^2wv),$$
$$<W^q,\overline{W}^q>=kww.$$
$`<,>(𝖵^q\overline{𝖵}^q)_0=\gamma `$.
Sur la composante de spin zéro de $`𝖵^q\overline{𝖵}^q`$ nous avons
$$(q^3+q)<U^q,\overline{W}^q>+<V^q,\overline{V}^q>+(q+q^1)<W^q,\overline{U}^q>=\gamma .$$
$``$ Enfin la $`q`$-symétrie (31) se traduit par :
$$q^2<U^q,\overline{V}^q>=<V^q,\overline{U}^q>,<W^q,\overline{V}^q>=q^2<V^q,\overline{W}^q>,$$
$$(q^3+q)[<U^q,\overline{W}^q><W^q,\overline{U}^q>]=(q^21)<V^q,\overline{V}^q>.$$
Les relations de dépendance entre les paramètres $`k`$ et $`\gamma `$ découlent de (33) et (34) qui, en forme explicite se traduisent par :
$$<K,z>=0,z\overline{𝖵}^q,<z,\overline{K}>=0z𝖵^q\text{}$$
(38)
$$K=(q^3+q)uW^q+vV^q+(q+q^1)wU^q,$$
$$\overline{K}=(q^3+q)\overline{U}^qw+\overline{V}^qv+(q+q^1)\overline{W}^qu.$$
Déterminons cette relation de dépendance en utilisant la première équation de (38). Il suffit de considérer $`z=\overline{U}^q`$ dans $`<K,z>=0`$. On a alors
$$(q^3+q)u<W^q,\overline{U}^q>+q^2v<U^q,\overline{V}^q>+(q+q^1)w<U^q,\overline{U}^q>=0.$$
(39)
Or on obtient :
$`<U^q,\overline{U}^q>=kuu,<U^q,\overline{V}^q>=kuv,<V^q,\overline{W}^q>=kvw,`$ (40)
$`<V^q,\overline{V}^q>(q^3+q)<W^q,\overline{U}^q>=k(vv(q^3+q)wu),`$ (41)
$$<W^q,\overline{W}^q>=kww,q^2<V^q,\overline{V}^q>+(q^3+2q+q1)<W^q,\overline{U}^q>=\gamma .$$
(42)
Les relations (41) et (42) donnent :
$$<W^q,\overline{U}^q>=\alpha [\gamma kq^2(vv(q^3+q)wu)],$$
$$<V^q,\overline{V}^q>=\alpha [(q^3+q)\gamma +k(q^3+2q+q^1)vvk(q^3+q)(q^3+2q+q^1)wu]$$
$`\alpha =(2q^3+2q+q^1+q^5)^1`$.
Exprimons en outre la constante orbitale $`c`$ de l’hyperboloïde quantique en fonction de $`vv`$ et de $`wu`$. Dans l’algèbre $`𝒜_{0,q}^c`$ nous avons :
$$(q^3+q)(uwwu)+(1q^2)vv=0,(q^3+q)uw+vv+(q+q^1)wu=c\text{donc}$$
$$q^2vv+(q^3+2q+q^1)wu=c.$$
L’équation (39) s’écrit alors :
$$(q^3+q)\alpha u[\gamma kq^2(vv(q^3+q)wu]+kq^2vuv+k(q+q^1)wuu=0.$$
(43)
Faisons les changements nécessaires dans les deux derniers termes de l’équation (43) de façon à mettre $`u`$ à gauche dans ces termes. Nous pouvons alors réécrire l’équation (43) sous la forme :
$$\gamma \mathrm{\Gamma }_1+\mathrm{\Gamma }_2vv+\mathrm{\Gamma }_3wu=0$$
$$\mathrm{\Gamma }_1=(q^3+q)\alpha ,$$
$$\mathrm{\Gamma }_2=k[(q^3+q)\alpha q^2+q^4q^4(q+q^1)\frac{q^21}{q^3+q}],$$
$$\mathrm{\Gamma }_3=k[(q^3+q)^2\alpha q^2+(q+q^1)].$$
Nous en déduisons (pour $`q`$ générique) :
$`\gamma `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }_2}{\mathrm{\Gamma }_1}}vv{\displaystyle \frac{\mathrm{\Gamma }_3}{\mathrm{\Gamma }_1}}wu`$
$`=`$ $`{\displaystyle \frac{1}{q^2}}(1+q^4)k[q^2vv+(q^3+2q+q1)wu]`$
$`=`$ $`{\displaystyle \frac{1}{q^2}}(1+q^4)kc,c\text{est la contante orbitale non nulle}.`$
On vérifie qu’avec cette relation de dépendance trouvée, la deuxième équation de (38) est satisafaite.
### B Connexion tressée
Notons
$$J_i=\frac{1}{[i]_q}Y^iJ_{i1},i\{1,\mathrm{\hspace{0.17em}2},\mathrm{\hspace{0.17em}3},\mathrm{\hspace{0.17em}4}\},[i]_q=\frac{q^iq^i}{qq^1},\text{}$$
$$J_0=_{U^q}U^q=\alpha (uvU^qq^2uuV^q).$$
Par application de l’opérateur $`Y^i`$ à la relation (37) nous obtenons les quatre relations suivantes :
$$_{U^q}V^qq^2_{V^q}U^q=J_1,$$
$$_{U^q}W^q+q_{V^q}V^qq^4_{W^q}U^q=J_2,$$
$$_{V^q}W^q+q^2_{W^q}V^q=J_3,$$
$$_{W^q}W^q=J_4\text{avec :}$$
$$J_1=\alpha \{[(q^3+q)uwq^2vv)]U^q+(q^6+q^21)uvV^q(q^3+q)uuW^q\},$$
$$J_2=(1+q^2+q^4)\alpha \{q^2vwU^q+(1q^2)[uw+[2]^1v^2]V^q+q^2uvW^q]$$
$$J_3=\alpha \{q^3(q^2+1)wwU^q+(1+q^4q^6)vwV^q+q^4[(q+q^1uwvv]W^q\},$$
$$J_4=\alpha (wwV^qq^4vwW^q).$$
Des relations du (35), nous déduisons que :
$$_{U^q}W^q=_{W^q}U^q+\frac{q^21}{q^3+q}_{V^q}V^q+\frac{2}{q^3+q}V^q,$$
$$(q^3+q)J_2=(q^3+q)(1+q^4)_{W^q}U^q+(1+q^4)_{V^q}V^q2V^q,$$
$$(q+q^1)(q^21)_{W^q}U^q+q^2_{V^q}V^q+2V^q=0.$$
Par conséquent, nous savons exprimer les éléments suivants :
$$_{U^q}V^q,_{V^q}U^q,_{V^q}W^q,_{W^q}V^q,_{W^q}U^q,_{U^q}W^q,_{V^q}V^q$$
en fonction de $`\alpha `$. Plus précisément en fonction de ($`J_1`$) on a :
$`_{U^q}V^q=`$ $`{\displaystyle \frac{1}{1+q^4}}\{\{\alpha [(q^3+q)uwq^2vv]2q^2\}U^q\alpha (q^6+q^21)uvV^q`$
$`+\alpha (q^3+q)uuW^q\},`$
$`_{V^q}U^q=`$ $`{\displaystyle \frac{1}{1+q^4}}\{\{\alpha q^2[(q^3+q)uwq^2vv]+2\}U^q\alpha q^2(q^6+q^21)uvV^q`$
$`+\alpha q^2(q^3+q)uuW^q\}\}.`$
En fonction de ($`J_3`$) nous obtenons :
$`_{V^q}W^q=`$ $`{\displaystyle \frac{1}{1+q^4}}\{\alpha q^3(1+q^2)wwU^q+\alpha (1+q^4q^6)vwV^q+\{\alpha q^4`$
$`[(q+q^1)uwvv]2q^2\}W^q\},`$
$`_{W^q}V^q=`$ $`{\displaystyle \frac{1}{1+q^4}}\{\alpha q^5(1+q^2)wwU^q+\alpha q^2(1+q^4q^6)vwV^q+\{\alpha q^6`$
$`[(q+q^1)uwvv]+2\}W^q\}.`$
Nous avons également :
$`_{W^q}U^q=`$ $`{\displaystyle \frac{1}{1+q^4}}\{q^6\alpha vwU^q+\{q^4(1q^2)\alpha [uw+[2]^1v^2]{\displaystyle \frac{2q}{1+q^2}}\}V^q`$
$`q^6\alpha uvW^q\},`$
$`_{V^q}V^q=`$ $`{\displaystyle \frac{1}{1+q^4}}\{q^3(1+q^2)^2\alpha vwU^q+\{q(1+q^2)(1q^4)\alpha \{q(1+q^2)`$
$`(1q^4)\alpha [uw+[2]^1v^2]+2(1q^2)\}V^q+q^3(1+q^2)^2\alpha uvW^q\},`$
$`_{U^q}W^q=`$ $`{\displaystyle \frac{1}{1+q^4}}\{q^2\alpha vwU^q+\{(q^21)\alpha [uw+[2]^1v^2]+{\displaystyle \frac{2q}{1+q^2}}\}V^q`$
$`q^2\alpha uvW^q\}.`$
Il reste maintenant à préciser la constante $`\alpha `$. Pour cela, considérons la deuxième équation du (36) :
$$[(q^3+q)u_{W^q}+v_{V^q}+(q+q^1)w_{U^q}](z)=0z𝖵^q.$$
(44)
Par exemple pour $`z=U^q`$ dans (44) nous obtenons :
$$\beta [(q^4q^2+1)\alpha c+2][vU^qq^2uV^q]=0\text{et}vU^qq^2uV^q0\text{dans}T(\mathrm{H}_q)_l.$$
D’où la valeur de $`\alpha `$.
P. Akueson
I.S.T.V., Université de Valenciennes
B.P. 311 Valenciennes France
E-mail: akueson@univ-valenciennes.fr
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# HIGH-REDSHIFT GALAXIES IN COLD DARK MATTER MODELS
## 1 Introduction
The discovery and characterization of “Lyman-break” galaxies (LBGs) has opened a new window on the high-redshift universe, revealing a population of star-forming galaxies at $`z>3`$ whose comoving space density exceeds that of $`L_{}`$ galaxies today (Steidel et al., 1996; Lowenthal et al., 1997). These galaxies can be identified by their unusual colors in deep imaging surveys because the intrinsic continuum break at $`\lambda 912`$Å and the intergalactic absorption by the Ly$`\alpha `$ forest at $`\lambda <1216`$Å redshift into optical bands. Spectroscopic follow-up shows that photometry of Lyman-break objects yields robust approximate redshifts. From an optical imaging survey, one can therefore construct a sample of high-$`z`$ galaxies limited primarily by rest-frame ultraviolet (UV) luminosity, which, in the absence of dust extinction, is itself determined mainly by the instantaneous formation rate of massive stars. Application of this approach to the Hubble Deep Field (HDF; Williams et al., 1996) and other deep imaging surveys has yielded first attempts at one of the long-standing goals of observational cosmology, determination of the star formation history of the universe (e.g., Madau et al., 1996; Madau, 1997; Connolly et al., 1997; Steidel et al., 1999).
In this paper, we examine the ability of models based on inflation and cold dark matter (CDM) to account for the observed population of LBGs, using cosmological simulations that incorporate gravity, gas dynamics, and star formation. We consider five variants of the CDM scenario: three $`\mathrm{\Omega }_m=1`$ models, a spatially flat low density model with a cosmological constant, and an open universe low density model with $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$. The spatial clustering of the high-redshift galaxies in these simulations was discussed by Katz, Hernquist, & Weinberg (1999); here we focus on the masses and star formation properties of these galaxies.
Numerical simulations play two overlapping but distinct roles in cosmological studies. First, they provide quantitative predictions that can be compared to observations in order to test the underlying cosmological models. Second, they provide greater understanding of the observational phenomena themselves, by showing how observable structures might arise and evolve in a given cosmological scenario. In this paper we will emphasize the second of these roles, mainly because the numerical limitations of the simulations and our limited knowledge of the physics of star formation contribute uncertainties that are comparable to the differences between cosmological models. The examination of different cosmologies is still a useful exercise, however, because it shows how cosmological parameters and properties of primordial mass fluctuations affect the properties of the high-redshift galaxy population when other physical and numerical parameters are held fixed.
Hydrodynamic simulations complement the main alternative approach to the theoretical study of high-redshift galaxies, based on semi-analytic models of galaxy formation (e.g., Baugh et al., 1998; Kauffman et al., 1999; Somerville, Primack, & Faber, 2000). Semi-analytic models have the advantages of simplicity, flexibility, and speed. The price is a substantial number of approximations and tunable parameters; the values of some parameters are fixed by matching selected observations, leaving other observables as predictions of the model. Semi-analytic models incorporate simplified descriptions of gravitational collapse, mergers, and cooling of gas within dark halos. The strength of numerical simulations is their more realistic treatment of these processes. The only free parameters (apart from the physical parameters of the cosmological model being studied) are those related to the treatment of star formation and feedback. Given these parameters, simulations provide straightforward, untunable predictions. However, the simulation approach must contend with the numerical uncertainties caused by finite volume and finite resolution, and computational expense makes it a slow way to explore parameter space. Over the next few years, interactions between the numerical and semi-analytic approaches should strengthen both. Here we mainly present the numerical results on their own terms, with a brief comparison to interpretations based on semi-analytic models in §5.
We describe our numerical methods, treatment of star formation, and choice of cosmological model in §2. In §3 we present results for the LCDM model (CDM with a cosmological constant) at $`z=3`$, the redshift best probed by recent Lyman-break galaxy surveys. In §4 we broaden our scope, examining predictions of five different CDM models for the population of star-forming galaxies from $`z=6`$ to $`z=2`$. We discuss our results and prospects for future progress in §5.
## 2 Simulations
### 2.1 Numerical Parameters and Star Formation
All of our simulations use TreeSPH (Hernquist & Katz, 1989; Katz, Weinberg & Hernquist, 1996, hereafter KWH), a code that combines smoothed particle hydrodynamics (SPH; see Lucy, 1977; Gingold & Monaghan, 1977; Monaghan, 1992) with a hierarchical tree algorithm (Barnes & Hut, 1986) for computing gravitational forces. The method and illustrative cosmological applications are described in detail by KWH, so here we just specify the simulation parameters and recap the points that are most important to the present investigation.
Each of our simulations uses $`64^3`$ dark matter and $`64^3`$ SPH particles to model a triply periodic volume $`11.111h^1\mathrm{Mpc}`$ comoving Mpc on a side, where $`hH_0/(100\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1)`$. The simulations are evolved to $`z=2`$. For the three critical density ($`\mathrm{\Omega }_m=1`$) cosmological models, the dark matter particle mass is $`2.76\times 10^9M_{}`$ and the SPH particle mass is $`1.45\times 10^8M_{}`$. For the two low density ($`\mathrm{\Omega }_m=0.4`$) models, the dark matter particle mass is $`8.27\times 10^8M_{}`$ and the SPH particle mass is $`6.71\times 10^7M_{}`$. Gravitational forces are softened using a cubic spline kernel with a softening length $`ϵ=5h^1\mathrm{kpc}`$, equivalent to $`ϵ3.5h^1\mathrm{kpc}`$ for a Plummer softening law. The gravitational softening length is held fixed in comoving units, i.e., $`ϵ=1.25h^1`$ physical kpc at $`z=3`$. Particles have individual time steps that satisfy the conditions $`\mathrm{\Delta }t<0.4\mathrm{min}(ϵ/|𝐯|,\sqrt{ϵ/|𝐚|})`$, where $`𝐯`$ is the peculiar velocity and $`𝐚`$ is the acceleration. SPH particle time steps are also required to satisfy the Courant condition (see KWH). The maximum time step for any particle is $`\mathrm{\Delta }t_d=H_0^1/6000`$.
Radiative cooling is computed assuming primordial composition gas with helium abundance $`Y=0.24`$ by mass. All of the simulations incorporate a photoionizing UV background with the spectral shape and redshift history computed by Haardt & Madau (1996), but with intensity reduced by a factor of two in order to approximately match the mean opacity of the Ly$`\alpha `$ forest given our assumed baryon density (Croft et al., 1997). In practice, the photoionizing background has negligible effect on the Lyman-break galaxy population, at least in the mass range that our simulations can resolve (Weinberg, Hernquist & Katz, 1997a).
The gas that resides in collapsed dark matter halos exhibits a two-phase structure: hot gas at roughly the halo virial temperature with a density profile similar to that of the dark matter (but exhibiting a core at small radii), and radiatively cooled gas with $`T10^4\mathrm{K}`$ at much higher overdensity. In simulations that do not incorporate star formation, the clumps of radiatively cooled gas have masses and sizes comparable to the luminous regions of observed galaxies. Our star formation algorithm is essentially a prescription for turning this dense, cold gas into collisionless stars, returning energy from supernova feedback to the surrounding medium. We provide a brief synopsis of this algorithm here and refer the reader to KWH for details.
An SPH gas particle is “eligible” to form stars if it is Jeans unstable, resides in a region of converging flow, has an overdensity $`\rho _g/\overline{\rho _g}>55.7`$ (corresponding to the virial boundary of a singular isothermal sphere in the spherical collapse model), and has a hydrogen number density exceeding $`0.1\mathrm{cm}^3`$ (physical units). In practice, it is the physical density threshold that matters — gas with this density almost always satisfies the other criteria, except at very high redshift, where the overdensity threshold ensures that star formation does not occur in uncollapsed regions simply because the cosmic mean density is high. Once a gas particle is eligible to form stars, its star formation rate is given by
$$\frac{d\rho _{}}{dt}=\frac{d\rho _g}{dt}=\frac{c_{}ϵ_{}\rho _g}{t_g},$$
(1)
or
$$\frac{d\mathrm{ln}\rho _g}{dt}=\frac{c_{}ϵ_{}}{t_g},$$
(2)
where $`c_{}`$ is a dimensionless star formation rate parameter, $`ϵ_{}`$ is the fraction of the particle’s gas mass that will be converted to stellar mass in a single simulation timestep, and the gas flow timescale $`t_g`$ is the maximum of the local gas dynamical time, $`t_{\mathrm{dyn}}=(4\pi G\rho _g)^{1/2}`$, and the local cooling time. Each SPH particle has both a gas mass and a stellar mass (initially zero); the total gas$`+`$stellar mass contributes to gravitational forces, but only the gas mass is used in computing the SPH properties and forces. The probability $`p`$ that an eligible SPH particle undergoes a star formation event in an integration timestep of duration $`\mathrm{\Delta }t`$ is
$$p=1\mathrm{exp}\left(\frac{c_{}\mathrm{\Delta }t}{t_g}\right),$$
(3)
and if the particle does undergo such an event then $`ϵ_{}=1/3`$ of its remaining gas mass is converted into stars during that step. In the limit (nearly always satisfied in the simulations) that $`c_{}\mathrm{\Delta }t/t_g1`$, this algorithm yields the average star formation rate given by equation (1).<sup>1</sup><sup>1</sup>1In KWH, the description of the algorithm is accurate but their equations (44) and (45), which correspond to equations (1) and (2), are missing the factor of $`ϵ_{}`$. Once a particle’s gas mass falls below 5% of its original mass, it is converted into a collisionless, pure star particle, affected only by gravity, and its residual gas mass is redistributed to its SPH neighbors.
When an SPH particle undergoes star formation, recycled gas and supernova feedback energy are distributed to the particle and its neighbors, assuming a Miller-Scalo (1979) initial mass function truncated at $`0.1M_{}`$ and $`100M_{}`$ and $`10^{51}`$ ergs per supernova. This feedback energy is usually radiated away because it is released into a dense, gas rich medium with a short cooling time. Feedback therefore has only a modest impact in our simulations, and this is the physically appropriate result if the proto-galactic interstellar medium is fairly smooth and as dense as our simulations imply. It is possible that strong inhomogeneities in the interstellar medium (on scales well below our resolution limits) allow feedback to have a stronger effect in real proto-galaxies, and explicit modeling of this possibility is an important direction for future investigation. The scenario that we investigate here is a physically plausible limiting case.
In all of our simulations we set the star formation rate parameter $`c_{}`$ to 0.1 and $`ϵ_{}`$ to 1/3. As shown in KWH, the stellar masses of the simulated galaxies are insensitive to the value of $`c_{}`$. In the KWH tests, an order of magnitude increase to $`c_{}=1.0`$ changes the total stellar mass in the box at $`z=2`$ by only 15%, and the effect of a higher $`c_{}`$ is actually to decrease the stellar mass because star formation occurs in lower density gas where supernova feedback can have a stronger effect. Indeed, one obtains nearly the same galaxy population in simulations that do not include star formation at all, except that in this case the “galaxies” are the clumps of cold, dense gas instead of the clumps of cold, dense gas and stars (see KWH, figure 5). In our simulations, the rate at which a galaxy forms stars is governed mainly by the rate at which gas condenses from the hot halo into the cold clump; the regulation implied by equation (1) ensures that the gas condensation rate and the star formation rate cannot get too far out of step. The link between star formation rate and gas density is physically motivated, since denser gas is more gravitationally unstable and more easily able to radiate its energy. In the case where the cooling timescale is short and $`t_g=t_{\mathrm{dyn}}`$, equation (1) implies $`\dot{\rho }_{}\rho _g^{3/2}`$, similar to the Schmidt-law $`\dot{\rho }_{}\mathrm{\Sigma }_g^{3/2}`$ observed to hold over a large dynamic range in a wide variety of local galaxies (Schmidt, 1959; Kennicutt, 1998).
### 2.2 Cosmological Models
We consider five different cosmological models, all of which assume Gaussian primordial fluctuations and a universe dominated by cold, collisionless dark matter. In all cases we adopt a baryon density parameter $`\mathrm{\Omega }_b=0.0125h^2`$ based on Walker et al. (1991), though a higher $`\mathrm{\Omega }_b`$ is suggested by recent analyses of the Ly$`\alpha `$ forest opacity (Rauch et al., 1997; Weinberg et al., 1997b) and the deuterium abundance in high-redshift Lyman limit systems (Burles & Tytler 1997, 1998). The model we refer to as “standard” CDM (SCDM) assumes $`\mathrm{\Omega }_m=1`$, $`h=0.5`$, and an rms linear theory fluctuation in $`8h^1\mathrm{Mpc}`$ spheres of $`\sigma _8=0.7`$. For this model we use the parameterization of the CDM power spectrum given by Bardeen et al. (1986). The $`\sigma _8=0.7`$ normalization is roughly consistent with the observed abundance of rich galaxy clusters (White, Efstathiou, & Frenk, 1993), but the SCDM model does not reproduce the amplitude of cosmic microwave background anisotropies observed by the COBE-DMR experiment (Smoot et al., 1992; Bennett et al., 1996). Our second model, COBE-normalized CDM (CCDM), is the same as SCDM except that the normalization $`\sigma _8=1.2`$ is chosen to match the 4-year COBE data (Gorski et al. 1996; see Bunn & White 1997 for a discussion of the CDM normalization). With this value of $`\sigma _8`$ and $`\mathrm{\Omega }_m=1`$, the CCDM model produces galaxy clusters that are too massive to be consistent with observations.
One way to reconcile the COBE-DMR anisotropies and the observed cluster abundance within the context of $`\mathrm{\Omega }_m=1`$ CDM models is to assume that inflation generates a primeval power spectrum that is “tilted” instead of scale-invariant, $`P(k)k^n`$ with $`n<1`$. For our tilted CDM (TCDM) model, we adopt $`n=0.80`$ and the transfer function given by equation (D28) of Hu & Sugiyama (1996), which treats baryon damping effects more accurately than the original Bardeen et al. (1986) formulation. We normalize this analytic fit to the power spectrum to the amplitude $`\sigma _8=0.54`$ implied by COBE normalization, which we compute using the CMBFAST code of Seljak & Zaldarriaga (1996; Zaldarriaga, Seljak, & Bertschinger 1998), assuming the standard tensor mode contribution to microwave background anisotropies predicted by power law inflation models.
Another way to resolve the COBE/cluster conflict is to lower the value of $`\mathrm{\Omega }_m`$, reducing cluster masses for a given $`\sigma _8`$. We consider two different low-$`\mathrm{\Omega }_m`$ CDM models, one (LCDM) with a spatially flat universe and a cosmological constant $`\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m`$ and one (OCDM) with an open universe and $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$. For LCDM we adopt $`\mathrm{\Omega }_m=0.4`$, $`h=0.65`$, and a primeval spectral index $`n=0.93`$. With the tensor mode contribution, CMBFAST implies a normalization $`\sigma _8=0.8`$, which provides a good match to the cluster abundances for $`\mathrm{\Omega }_m=0.4`$ (White, Efstathiou, & Frenk, 1993). We again use the Hu & Sugiyama (1996) formulation of the transfer function. For OCDM, we adopt $`\mathrm{\Omega }_m=0.4`$, $`h=0.65`$, $`n=1.0`$, and a 2-year COBE-DMR normalization $`\sigma _8=0.75`$ (Ratra et al., 1997). Cluster masses in this model are lower than those in TCDM or LCDM, but they are consistent with current observations given their uncertainties (Cole et al., 1997). For OCDM we use the transfer function of Efstathiou, Bond, & White (1992) with $`\mathrm{\Gamma }=0.234`$; the Hu & Sugiyama (1996) formulation is more accurate, but we were unaware of it at the time we ran the OCDM simulation. In practice, the differences between different analytic or numerical formulations of transfer functions are of the same magnitude as the changes caused by slight shifts in the adopted values of $`h`$, $`\mathrm{\Omega }_b`$, or $`\mathrm{\Omega }_m`$.
Parameters of the five cosmological models are listed in Table 1. Figure 1 shows the linear theory power spectra of the five models, over the range of scales represented in the initial conditions of our simulations. Instead of $`P(k)`$ itself, we plot $`\mathrm{\Delta }^2(k)4\pi k^3P(k)`$, which (with our Fourier transform convention) is the contribution to the variance of linear mass fluctuations per unit interval of ln$`k`$ (see Peacock & Dodds, 1994). The differences in fluctuation amplitude among the three $`\mathrm{\Omega }_m=1`$ models are easy to see, as is the difference in $`P(k)`$ shape between the CCDM/SCDM models and the tilted (TCDM) and low density (OCDM/LCDM) models. Although LCDM has a slightly higher normalization than OCDM at $`z=0`$ ($`\sigma _8=0.80`$ vs. $`\sigma _8=0.75`$), OCDM has higher amplitude fluctuations at $`z=3`$ because of the smaller ratio of linear growth factors between $`z=3`$ and $`z=0`$ in an open universe.
## 3 Lyman-Break Galaxies in the LCDM Model
We will focus in this Section on the galaxy population of the LCDM model at $`z=3`$, then turn to other models and other redshifts in §4. Figure 2 shows particle distributions at $`z=3`$ from the full $`11.111h^1\mathrm{Mpc}`$ simulation cube (top panels) and a $`0.09h^1\mathrm{Mpc}`$ subcube containing the richest concentration of galaxies (bottom panels). The top panels show numerous concentrations of dense, shock heated gas, with typical temperatures $`T10^6\mathrm{K}`$ corresponding to the virial temperatures of the corresponding dark matter halos. However, the zoomed view at the lower left reveals the two phase structure that characterizes collapsed regions of the simulation. Extremely overdense clumps of $`10^4\mathrm{K}`$ gas, typically a few kpc in size, reside in a background of gas with $`T10^6\mathrm{K}`$. These dense concentrations of cold gas are, of course, the sites of star formation, as shown in the lower right panel. Because the knots of cold gas and stars stand out so clearly from the background, there is no ambiguity in identifying the galaxies in such a simulation; one only needs an algorithm that picks out these clumps. We use the SKID algorithm (Spline-Kernel Interpolated DENMAX; see KWH and Gelb & Bertschinger 1994), as implemented by Stadel, Katz, & Quinn<sup>2</sup><sup>2</sup>2see http://www-hpcc.astro.washington.edu/tools/SKID/, which identifies clumps of gravitationally bound particles associated with a common density maximum.
In collapsing dark matter halos that contain a small number of particles, the resolution of the SPH calculation becomes too low for the simulation to follow the cooling of the gas and subsequent star formation. Gardner et al. (1997, 2000) find that simulated halos with at least 60 dark matter particles nearly always contain a cold gas concentration, while halos with fewer particles often do not. We can therefore resolve the existence of galaxies in halos with mass $`M>60m_{\mathrm{dark}}`$, a quantity that we list in Table 1 as an indication of our dark matter mass resolution. At $`z=3`$, the halo circular velocity corresponding to this limiting mass (plus the associated baryon mass) is $`v_c140\mathrm{km}\mathrm{s}^1`$ in the $`\mathrm{\Omega }_m=1`$ models and $`v_c95\mathrm{km}\mathrm{s}^1`$ in the low density models (see Gardner et al. 1997, equation 3). Comparison of LCDM simulations with $`64^3`$ and $`128^3`$ SPH particles in the same $`11.111h^1\mathrm{Mpc}`$ volume (Gardner et al., in preparation; Aguirre et al., in preparation) suggests that the requirement for accurate estimation of a galaxy’s star formation rate via equation (1) is somewhat more stringent, corresponding to 60 or more particles in the condensed baryon phase (cold gas plus stars). We list $`60m_{\mathrm{SPH}}`$ as an approximate baryon mass resolution limit for star formation calculations in Table 1.
For the galaxies at $`z=3`$, the asterisks in Figure 3 compare the star formation rate averaged over the previous 20 Myr to the star formation rate averaged over the previous 200 Myr. The two rates are usually within a factor of two of each other and are often much closer, indicating that star formation in our simulated galaxies is fairly steady over intervals of 200 Myr. Circles in Figure 3 show the “instantaneous” star formation rate calculated by applying the prescription of §2 to the gas particle distribution. This is the rate that would be used to calculate star formation in the simulation’s next system timestep (of duration $`\mathrm{\Delta }t=2.5`$ Myr). The tight correlation between this estimate of the star formation rate and the rates averaged over longer intervals demonstrates that the instantaneous estimate is not sensitive to numerical “noise” in particle densities and positions, at least for systems with star formation rate $`>1M_{}/`$yr. We will henceforth use the instantaneous estimate as our standard measure of the star formation rate (hereafter denoted SFR), with Figure 3 as evidence that our results are insensitive to this choice.
Figure 4 plots galaxies’ instantaneous star formation rates against their stellar masses. The correlation between SFR and $`M_{}`$ is reasonably strong, but there is enough scatter that a sample selected by a threshold in SFR (above a horizontal line in Figure 4) excludes some fairly massive galaxies and includes others that are substantially further down the stellar mass function. Nearly all galaxies are forming stars at a rate faster than the rate $`\mathrm{SFR}=M_{}/t`$ that would build them up steadily over the age of the universe; this result is not surprising, since the galaxies do not start to form stars at $`t=0`$. The ratio SFR$`/\mathrm{SFR}`$ is substantially higher for low mass galaxies than for high mass galaxies. Since this ratio is correlated with the overall shape of a galaxy’s spectral energy distribution, Figure 4 implies that less massive galaxies should be bluer than more massive galaxies. This trend could be caused partly by our limited numerical resolution, since the simulated galaxies do not form stars at the correct efficiency until they are sufficiently massive. However, the trend appears to be present even in the fairly well resolved systems.
Figure 5 plots the ratio of stellar mass to baryon mass (stars plus cold gas) as a function of baryon mass. The low mass galaxies are gas rich, while the most massive systems are predominantly stellar. This trend is physically plausible, but it is almost certainly exaggerated by our underestimate of star formation rates in poorly resolved systems. It should therefore be treated as a tentative prediction, awaiting confirmation with higher resolution simulations.
All of the simulated galaxies reside in dark matter halos, and the more massive halos frequently contain several galaxies (see Gardner et al. 2000, figure 5). As a characteristic of these halos, we calculate the circular velocity $`v_c=(GM/r)^{1/2}`$ from the total mass (dark plus baryonic) within a physical radius $`r=40h^1\mathrm{kpc}`$ around each simulated galaxy. Figure 6 plots these circular velocities against the baryon masses (left) and star formation rates (right) of the LCDM galaxies at $`z=3`$. There is a well defined lower ridge line in the $`v_c`$ vs. $`M_b`$ plot, but there are also galaxies with low $`M_b`$ and high $`v_c`$, which are usually “satellite” galaxies in halos with several distinct baryonic subunits. The galaxies with high SFR tend to reside in relatively massive halos, but halos above a given $`v_c`$ threshold host galaxies with a wide range of SFR, and the correlation between SFR and halo $`v_c`$ becomes increasingly weak below SFR$`=10M_{}\mathrm{yr}^1`$. Changing the radius for $`v_c`$ definition from $`r=40h^1\mathrm{kpc}`$ to $`r=20h^1\mathrm{kpc}`$ or $`r=80h^1\mathrm{kpc}`$ makes little difference to the appearance of Figure 6. For the galaxies with high SFR, the circular velocities in Figure 6 are large compared to the typical nebular emission line widths $`\sigma 70\mathrm{km}\mathrm{s}^1`$ measured in LBGs (Pettini et al., 1998), but it is unclear that emission line widths have much to do with the dark matter potential well depth even in local starburst galaxies (Lehnert & Heckman, 1996).
Figure 7 illustrates the build-up of the galaxy population in the LCDM model from $`z=6`$ to $`z=2`$. Each panel shows a projection of the box at the indicated redshift, with each galaxy represented by a circle whose area is proportional to the instantaneous star formation rate. By $`z=6`$ there are already ten galaxies in the box with SFR$`>1M_{}\mathrm{yr}^1`$ and two with SFR$`>10M_{}\mathrm{yr}^1`$. As time goes on, the star formation rates of the most massive galaxies tends to increase, though this trend saturates after $`z=4`$. The total number of star-forming galaxies in the box increases steadily, with more than 200 galaxies in the box at $`z=2`$. Most strikingly, the locations of newly forming galaxies are strongly correlated with the locations of pre-existing galaxies, so that the backbone of structure traced by the galaxy population remains similar from $`z=6`$ to $`z=2`$, although it becomes better defined as more galaxies form. We show in KHW that the galaxy correlation function $`\xi (r)`$ remains nearly constant in comoving coordinates from $`z=4`$ to $`z=2`$ and displays little dependence on the cosmological model, although the amplitude of $`\xi (r)`$ does increase at a given redshift if one selects only the most massive systems. Comparisons of other models in forms similar to Figure 7 appear in KHW (figure 1).
## 4 High-redshift Star Formation in CDM Cosmologies
We now turn to the main quantitative results of the paper, predictions of the star formation rates of high-redshift galaxies in our simulations of the five CDM models listed in Table 1. Each simulation represents a theoretical model that combines cosmological assumptions with assumptions about galactic scale star formation. While the predictions are not sensitive to the value of $`c_{}`$, the one free parameter in our star formation algorithm, they do depend on the general features of the algorithm itself: the star formation rate is an increasing function of the local density of cold gas, and supernova feedback energy is deposited mainly in the dense interstellar medium of forming galaxies and is therefore radiated away rather efficiently. We return to this point in §5.
Figure 8 shows the cumulative distribution of the simulated galaxies as a function of instantaneous SFR, at redshifts $`z=6`$, 5, 4, 3, and 2. The vertical axis represents the comoving space density of all simulated galaxies with star formation rate exceeding the indicated SFR, in $`h^3\mathrm{Mpc}^3`$. The amplitudes and redshift evolution of these distributions depend strongly on the amplitude of mass fluctuations in the cosmological model (see Figure 1). The CCDM simulation, with the highest fluctuation amplitude, already has 15 galaxies with SFR$`>60M_{}\mathrm{yr}^1`$ by $`z=6`$. The number of rapidly star-forming galaxies declines slowly from $`z=5`$ to $`z=3`$ and drops substantially between $`z=3`$ and $`z=2`$, though even at this redshift the CCDM model has the highest overall star formation rate of any of the models. The TCDM model, with the lowest fluctuation amplitude, has no star formation in galaxies resolved by our simulation until $`z=4`$. It exhibits a rapid rise in the number of star-forming galaxies between $`z=4`$ and $`z=3`$, and a steepening of the distribution function from $`z=3`$ to $`z=2`$. The flatness of the $`z=3`$ distribution function may be in part a numerical artifact, since the lower mass systems in this low amplitude model are barely resolved, and their star formation rates may be correspondingly underestimated. The other three models, with intermediate fluctuation amplitude, show intermediate behavior. For example, the LCDM model displays a steady rise in the star formation rate from $`z=6`$ to $`z=4`$, then little change from $`z=4`$ to $`z=2`$.
To facilitate comparison between Figure 8 and existing or future observational data, we list in Table 2 the conversions from SFR to apparent magnitude and from comoving $`h^3\mathrm{Mpc}^3`$ to number per arcmin<sup>2</sup> per unit redshift. Specifically, $`C_V`$ is the volume conversion factor and $`m_{10}`$ is the apparent magnitude on the AB system at observed wavelength $`\lambda _{\mathrm{obs}}=1500\times (1+z)`$Å that corresponds to a star formation rate SFR$`=10M_{}\mathrm{yr}^1`$. The value of $`m_{10}`$ is similar in our high- and low-density models because the effect of low-$`\mathrm{\Omega }_m`$ is approximately cancelled by the increase in $`H_0`$. We adopt the conversion from SFR to UV continuum luminosity quoted by Pettini et al. (1998), $`L_{1500}=10^{29}(\mathrm{SFR}/10M_{}\mathrm{yr}^1)\mathrm{erg}\mathrm{s}^1\mathrm{Hz}^1`$, which in turn is based on Bruzual & Charlot (1993) population synthesis models assuming continuous star formation and a Salpeter initial mass function extending from $`0.1M_{}`$ to $`100M_{}`$. For example, the LCDM model predicts $`10^{2.66}`$ galaxies per $`h^3\mathrm{Mpc}^3`$ with SFR above $`20M_{}\mathrm{yr}^1`$ at $`z=5`$, corresponding to $`C_V\times n=664\times 10^{2.66}=1.46`$ galaxies per arcmin<sup>2</sup> per unit redshift with apparent AB magnitude less than $`25.452.5\times \mathrm{log}(20/10)=24.70`$ at $`\lambda _{\mathrm{obs}}=9000`$Å. The conversions in Table 2 can be calculated using the formulas in Hogg (1999) together with the definition $`m_{AB}=2.5\mathrm{log}f_\nu 48.60`$. Note, however, that these magnitudes assume no dust extinction, while the UV continuum slopes of observed Lyman-break galaxies imply typical UV extinctions of $`0.52.5`$ magnitudes (Adelberger & Steidel, 2000).
The horizontal lines in Figure 8 mark the comoving space densities of spectroscopically confirmed objects in the LBG samples of Steidel et al. (1996) and Lowenthal et al. (1997), which have a mean redshift $`z3`$. Assuming that these surveys pick out the galaxies with the highest star formation rates, the intersections of the dotted ($`z=3`$) simulation curves with these horizontal lines yield the predicted star formation rates for galaxies near the sample magnitude limits. The Steidel et al. (1996) magnitude limit of $`25.5`$ corresponds to a star formation rate of approximately $`4M_{}\mathrm{yr}^1`$ with the assumptions described above (see Table 2), but $`12.5`$ magnitudes of dust extinction would increase the implied SFR by factors of $`2.510`$. The CCDM simulation predicts a star formation rate of $`90M_{}\mathrm{yr}^1`$ for objects with this space density, which is clearly too high unless the true dust extinction is much larger than Adelberger & Steidel (2000) estimate. The other simulations predict star formation rates of $`3050M_{}\mathrm{yr}^1`$, which is consistent with the Steidel et al. (1996) results if the typical extinction is $`22.5`$ magnitudes. The Lowenthal et al. (1997) survey of the HDF probes one magnitude deeper than the Steidel et al. (1996) sample, and thus a factor of 2.5 lower in SFR. The comoving space density of the Lowenthal et al. (1997) galaxies is higher by a factor of 3.5, and since Lowenthal et al. (1997) only observed $`2/3`$ of their Lyman-break candidates (with a 44% success rate), the true space density at this magnitude limit could be higher by a factor of $`1.53`$, as indicated by the arrow in Figure 8. The predicted star formation rates in the CCDM simulation again appear too high, while the predictions of the SCDM, OCDM, and LCDM models appear roughly compatible with the Lowenthal et al. (1997) results if the dust extinction correction is fairly large and the true space density is $`23`$ times the lower limit. The TCDM model predicts very low star formation rates at the Lowenthal et al. (1997) space density. The numerical prediction is clearly ruled out by the data, though this failure of the TCDM model should still be viewed with some caution until it is confirmed at higher numerical resolution. We also note that the small size of the simulation volumes leads to significant uncertainties in the predicted star formation rates at these low comoving densities: the Steidel et al. (1996) space density corresponds to 3.5, 1.5, and 1.2 galaxies in the $`(11.111h^1\mathrm{Mpc})^3`$ simulation volume for the critical density, open, and flat-$`\mathrm{\Lambda }`$ models, respectively.
Figure 9 presents a more detailed comparison between numerical predictions and observational data at $`z=3`$. The upper panel shows the differential distribution of star formation rates. The LCDM result is shown as a solid histogram, but to preserve visual clarity we show distributions for other models as connected lines. Here we omit galaxies with baryonic mass $`M_b<60m_{\mathrm{SPH}}`$ because limited numerical resolution would cause us to underestimate their star formation rates. The distributions therefore decline at low SFR because of the absence of low mass galaxies.
In the lower panel we convert the predictions to observable units using the conversions in Table 2. Points with $`1\sigma `$ statistical error bars show the luminosity function of Lyman-break galaxies at $`z3`$ with (squares) and without (circles) correction for dust extinction, from Adelberger & Steidel (2000; based on data from Steidel et al. 1999). As Adelberger & Steidel (2000) emphasize, the extinction-corrected points are quite uncertain: the corrections assume that the correlation between UV continuum slope and extinction observed in local starburst galaxies (Meurer, Heckman, & Calzetti, 1999) also holds at $`z=3`$, and the points at faint magnitudes rely on an extrapolation of the luminosity function. However, these extinction corrections are probably the best that can be made with current data, and they yield plausible consistency between the population of UV-selected Lyman-break galaxies and constraints on dust emission from sub-mm counts and the infrared background (Adelberger & Steidel, 2000).
Given the theoretical and observational uncertainties (which include numerical limitations, IMFs, the value of $`\mathrm{\Omega }_b`$, extinction corrections, incompleteness, and $`H_0`$), we do not wish to draw strong conclusions from Figure 9. The CCDM model appears to predict excessively luminous galaxies, as expected from the discrepancies already noted, and this discrepancy would be more severe if the simulations had used the Burles & Tytler (1997, 1998) estimate of $`\mathrm{\Omega }_b`$ instead of the Walker et al. (1991) estimate. (Gardner et al. \[in preparation\] find that the SFR scales approximately as $`\mathrm{\Omega }_b^{1.5}`$.) The galaxy population in the TCDM model is probably too faint, unless the true extinction corrections are surprisingly small or the numerical resolution effects are more severe than we think. However, a higher $`\mathrm{\Omega }_b`$ would improve the performance of this model. The other three models appear roughly compatible with the data. The limited dynamic range of the resolved galaxy populations in the simulations prevents a detailed comparison to the shape of the observed luminosity function. In future work, we will combine simulations of the LCDM model with different resolutions and box sizes to represent the galaxy population over a wider mass range.
Figure 10 shows the globally averaged density of star formation as a function of redshift, a representation of the cosmic star formation history made famous by Madau et al. (1996). The simulation results accord with the impressions from Figure 8. In particular, comparison of Figure 10 to Figure 1 shows that the cosmic star formation history depends strongly on the amplitude of mass fluctuations. The high-amplitude, CCDM model predicts a high-amplitude SFR curve that peaks at high redshifts. The globally averaged SFR in this model declines slowly at $`z<4`$. In the other models, the SFR appears to be reaching a plateau at $`z2`$, when the simulations stop. The low-amplitude, TCDM model predicts the lowest SFR, especially at high redshift. Data points in Figure 10 are taken from the analysis of Steidel et al. (1999), based on their own data and the data of Lilly et al. (1996) and Connolly et al. (1997). At each redshift, Steidel et al. (1999) estimate the global SFR by integrating the UV continuum luminosity function for galaxies with $`L>0.1L_{}`$, so the points do not include the contribution of the faintest galaxies (which are often below the survey magnitude limits). The open circles show estimates with no extinction correction, while the filled circles incorporate extinction corrections of 0.44 magnitudes at $`z<2`$ and 0.67 magnitudes at $`z>2`$. Open squares show the extinction-corrected estimates converted to the cosmological parameters of our LCDM model; points for the OCDM model parameters would lie between the open squares and filled circles.
Unfortunately, the global SFR is a difficult quantity to predict robustly from numerical simulations with a limited dynamic range, since they miss the contribution from galaxies below the resolution limit and underestimate the contribution from rare, massive galaxies, which are unlikely to form in a small simulation volume. Figure 2a of Weinberg et al. (1999) illustrates these resolution and box size effects using several simulations of an LCDM model (one with a higher baryon density). Because of these missing contributions, one should regard the curves in Figure 10 as lower limits to the true predictions of these models. The theoretical and observational uncertainties make it difficult to draw strong conclusions from Figure 10. The CCDM model appears to predict too much star formation. The SCDM and OCDM predictions agree fairly well with the extinction-corrected estimates of Steidel et al. (1999), but contributions from lower mass galaxies or an increase in $`\mathrm{\Omega }_b`$ would make this agreement worse. The LCDM predictions are somewhat low, but they might plausibly be boosted towards reasonable agreement with higher resolution and/or higher $`\mathrm{\Omega }_b`$. The TCDM predictions are far below the observational estimates. Of course, most of the action in the observational data takes place at $`z<2`$, after these simulations stop. Figure 2b of Weinberg et al. (1999) shows preliminary results from simulations of a similar set of models (from Davé et al. 1999), extended to $`z=0`$. The global SFR declines in all of the models at $`z<1`$, though not as sharply as the data points in Figure 10.
In Table 3 we list a more robust prediction of the simulations, the contribution to the globally averaged SFR from galaxies above our estimated resolution limit, those with $`M_b>60m_{\mathrm{SPH}}`$. We also list the number density $`n_{\mathrm{res}}`$ of such galaxies at each redshift. To the extent that galaxy absolute-magnitude correlates with baryonic mass, the corresponding observational quantity could be computed from a galaxy survey by choosing a limiting magnitude at which the galaxy number density is $`n_{\mathrm{res}}`$ and summing the contribution to the global SFR from galaxies above this magnitude limit.
## 5 Discussion
Our main result is that inflationary CDM models, combined with straightforward assumptions about galactic scale star formation, predict a substantial population of star-forming galaxies at $`z2`$. The stellar masses and star formation rates of these high-redshift systems are sensitive to the amplitude of the underlying mass power spectrum (compare, e.g., Figure 1 and Figure 10.) The results of the LCDM, OCDM, and SCDM simulations appear roughly consistent with the observed properties of Lyman-break galaxies, given the theoretical and observational uncertainties. The low-amplitude, TCDM model predicts an anemic LBG population that is probably inconsistent with current observations, though this conclusion may be sensitive to our finite numerical resolution and our adopted value of $`\mathrm{\Omega }_b`$. The high-amplitude, CCDM model appears to predict too much high-redshift star formation, by a significant factor.
The Ly$`\alpha `$ forest offers a more direct probe of the amplitude of mass fluctuations at high redshift (Croft et al., 1998). Recent observational analyses (Croft et al. 1999; McDonald et al. 2000; Croft et al., in preparation) imply that the matter power spectrum at $`z23`$ is similar to that in our LCDM, OCDM, and SCDM models but incompatible with the CCDM or TCDM models. It is reassuring that the models supported by the Ly$`\alpha `$ forest data appear to be the most compatible with the star formation properties of observed LBGs, though an increase in $`\mathrm{\Omega }_b`$ to the values supported by recent D/H studies (Burles & Tytler 1997, 1998) and other Ly$`\alpha `$ forest analyses (Rauch et al., 1997; Weinberg et al., 1997b) might spoil this agreement to some extent. We will examine the influence of $`\mathrm{\Omega }_b`$ on the high-redshift galaxy population elsewhere (Gardner et al., in preparation); our initial results imply that galaxy star formation rates in the SCDM model scale roughly as $`\mathrm{\Omega }_b^{1.5}`$. In KHW, we showed that the clustering of high-redshift galaxies in these simulations is consistent with observed LBG clustering (Adelberger et al., 1998; Giavalisco et al., 1998), and that the clustering is insensitive to the cosmological model because galaxies form at the same “biased” locations of the dark matter distribution in all five simulations.
In Gardner et al. (2000), we examine the predictions of these simulations for damped Ly$`\alpha `$ absorption, which is the other main observational probe of the high-redshift galaxy population. The galaxies resolved in these simulations account for only a fraction of the observed damped Ly$`\alpha `$ absorption at $`z3`$, ranging from $`3\%`$ in TCDM to $`30\%`$ in LCDM, SCDM, and CCDM to $`50\%`$ in OCDM. Since the simulations already go to higher space densities than existing spectroscopic samples of LBGs, our results imply that these LBG samples are not yet deep enough to include the galaxies responsible for most damped Ly$`\alpha `$ absorption. Haehnelt, Steinmetz, & Rauch (2000) reach a similar conclusion using analytic arguments. By extrapolating the simulation results with the aid of the Press-Schechter (1974) mass function, Gardner et al. (2000) conclude that absorption in lower mass systems is sufficient to account for observed damped Ly$`\alpha `$ absorption in any of these cosmological models, with the possible exception of TCDM. However, the uncertainties in the extrapolation are large, and definitive examination of the compatibility between LBG and damped Ly$`\alpha `$ constraints will require higher resolution simulations.
Semi-analytic methods, sometimes combined with N-body simulations of the dark matter distribution, are the main alternative to hydrodynamic simulations for theoretical investigation of high-redshift galaxy formation. Using these methods, several groups have found that CDM models like the ones studied here can reproduce the numbers, luminosities, colors, and clustering properties of observed LBGs (e.g., Baugh et al., 1998; Governato et al., 1998; Kauffman et al., 1999; Somerville, Primack, & Faber, 2000). The semi-analytic papers have led to three rather different suggestions about the nature of Lyman-break galaxies. In the first view, observed LBGs are the most massive galactic systems present at high redshift, forming stars at a fairly steady rate (Baugh et al., 1998). In the second view, interactions play a crucial role in triggering bursts of star formation, and many LBGs are low mass systems boosted temporarily, and briefly, to prominence (Kolatt et al., 1999; Somerville, Primack, & Faber, 2000). A third, intermediate case is one in which LBGs are massive galaxies experiencing bursts of star formation stimulated by minor or major mergers (Somerville, Primack, & Faber, 2000). This variety of views is mirrored to some extent in the observational literature on LBGs (compare, for example, Steidel et al. 1996 to Lowenthal et al. 1997 or Trager et al. 1997).
Our simulations suggest a picture that is intermediate between the extremes of this debate, but closest to the first point of view. Star formation in the simulated galaxies is steady on timescales of 200 Myr (Figure 3), and the instantaneous star formation rate is fairly well correlated with stellar mass (Figure 4). However, there is enough scatter in galaxy star formation rates that a sample of galaxies selected above a SFR threshold includes objects with a substantial range of stellar masses (Figure 4), and these may reside in halos with a wide range of circular velocities (Figure 6). The simulations automatically include interactions and mergers, but they do not resolve the existence of the low mass systems envisioned to play an important role in the extreme version of the collision-induced starburst model.
The properties of the simulated LBG population depend on the cosmological initial conditions and on the basic physics of gravity and gas dynamics, but they also depend on our adopted prescription for galactic scale star formation. The crucial features of this prescription are (1) that the local star formation timescale decreases with increasing gas density, as implied by studies of local galaxies (Kennicutt, 1998), and (2) that supernova feedback energy is deposited mainly in the dense interstellar medium, where it is usually radiated away before it has a large dynamical impact. Since we do not require any external triggers for star formation, an isolated galaxy that steadily accretes cold gas will convert that gas into stars, albeit over many orbital times. Interactions and mergers can enhance star formation by driving gas to higher density, but galaxies do not build up large reservoirs of dense gas that wait to be ignited. Limited resolution may reduce the influence of interactions and mergers in these simulations, since they do not resolve low mass satellites and do not resolve the nuclear star formation that is prominent in high resolution simulations of starbursts induced by minor (Mihos & Hernquist 1994a; 1996) or major (Mihos & Hernquist, 1994b; Hernquist & Mihos, 1995) mergers. A more efficient feedback mechanism could also lead to more episodic star formation histories, by producing cycles of starbursts followed by suppressed accretion and cooling. Since we have not investigated scenarios in which interactions or feedback play a larger role, we cannot draw conclusions about their viability. However, our results suggest that the straightforward treatment of star formation described in §2.1 is sufficient to explain at least the basic properties of the observed LBG population within the CDM cosmological framework.
The clearest prediction of the simulations is that the Lyman-break galaxies studied by Steidel et al. (1996, 1999) and Lowenthal et al. (1997) represent the tip of an iceberg. The cumulative distribution curves in Figure 8 should be taken as lower limits to the predicted galaxy number densities, especially at high redshifts, since limited resolution causes these simulations to underestimate the star formation rates in low mass systems, or to miss the systems entirely. Nonetheless, the curves for, e.g., the LCDM model show that there should be large numbers of $`z=3`$ galaxies below the magnitude limits of current LBG samples, and significant numbers of galaxies with SFR$`10M_{}\mathrm{yr}^1`$ even at $`z5`$. Recent searches have already yielded a number of spectroscopically confirmed galaxies at $`z=57`$ (Spinrad et al., 1998; Weymann et al., 1998; Chen, Lanzetta, & Pascarelle, 1999; Hu, Cowie, & McMahon, 1999), and analyses of deep HST/NICMOS images show candidate objects to redshifts $`z10`$ (Yahata et al., 2000; Dickinson, 2000). Systematic characterization of this faint galaxy population will be challenging, so it will be some time before simulations and data can be compared quantitatively in the very high redshift regime. But the existence of a significant population of star-forming galaxies at $`z>5`$ is a natural prediction of the CDM scenario.
Figures 4 and 5 also imply some correlations between observable properties of $`z=3`$ galaxies. While less massive galaxies tend to have lower star formation rates, they usually have higher ratios of instantaneous SFR to time-averaged SFR, and they should therefore have bluer spectral energy distributions unless they are more heavily reddened by dust. Less massive galaxies also tend to be more gas rich. Unfortunately, both of these predicted trends could be exaggerated by numerical resolution effects, so we do not regard them as very robust.
This paper presents a first attempt to characterize the star formation properties of high-redshift galaxies using hydrodynamic simulations, but there is much room for progress with future simulations. The emerging consensus on cosmological parameters, if it survives the tightening of observational constraints, makes the task easier by focusing effort on a preferred background model. Within such a framework, simulations with different box sizes and resolutions can be combined to model the galaxy population over a wider dynamic range of mass and redshift, improving the comparison to the observed luminosity function (as in Figure 9) and global star formation history (as in Figure 10). Weinberg et al. (1999) take a first step along this path, using multiple simulations of the LCDM model (with higher $`\mathrm{\Omega }_b`$) to predict cumulative SFR distributions from $`z=0.5`$ to $`z=10`$. Since the present simulations resolve galaxies far below the limits of current LBG spectroscopic samples, simulations of larger volumes at lower resolution will improve the comparison between predicted and observed LBG clustering. Higher resolution simulations, on the other hand, can probe the connection between LBGs and damped Ly$`\alpha `$ systems and test the robustness of some of the trends found in §3. They can also provide predictions of the structural properties of high-redshift galaxies, such as size and morphology, though these may be best investigated with simulations that zoom in to follow the formation of individual objects (e.g., Haehnelt, Steinmetz, & Rauch, 1998; Contardo, Steinmetz, & Fritze-von Alvensleben, 1998). In the long run, we also hope to examine different formulations of star formation and feedback, to determine what descriptions of these physical processes match the observed properties of galaxies and the intergalactic medium over the full range of accessible redshifts.
Where are the Lyman-break galaxies today? Because our present simulations stop at $`z=2`$, we will save a detailed examination of this question for a future paper on the assembly history of galaxies, using simulations (like those of Davé et al. 1999) that continue to $`z=0`$. A first look at these simulations suggests that there is no simple, one-sentence answer. Between $`z=3`$ and $`z=0`$, galaxies accrete fresh material and merge with each other, and many new galaxies form that had no $`z=3`$ progenitors at all (at least above the numerical resolution limits). The particles that lie in galaxies at $`z=3`$ are distributed at $`z=0`$ among galaxies with a wide range of environments and masses, though the most massive $`z=0`$ galaxies always contain some of these particles and the least massive often do not. Any link between LBGs and present-day ellipticals, or bulges, or halos, or clusters, can at best capture a general trend, one that is likely to be violated nearly as often as it is obeyed. Fortunately, cosmological simulations are an ideal tool for characterizing the full range of possible histories, providing the theoretical thread that can tie snapshots of the galaxy population at different redshifts into a coherent picture of galaxy formation and evolution.
We thank Kurt Adelberger and Chuck Steidel for helpful discussions and for providing the data plotted in Figures 9 and 10. We also thank James Lowenthal for helpful discussions. This work was supported by NASA Astrophysical Theory Grants NAG5-3922, NAG5-3820, and NAG5-3111, by NASA Long-Term Space Astrophysics Grant NAG5-3525, and by the NSF under grants ASC93-18185, ACI96-19019, and AST-9802568. The simulations were performed at the San Diego Supercomputer Center.
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# Distinguished properties of the gamma process, and related topics
## 1 Introduction
The purpose of this work is to link various questions from the probability theory, combinatorics, and representation theory, which relate to the gamma process and the Poisson–Dirichlet measures, and which were not regarded as a single whole. Since in spite of many aspects of our work its main object is the gamma process, first we would like to present some fundamental properties of the classical gamma process and its role for further considerations.
The standard gamma distribution (e.g. see ) with parameter $`\alpha >0`$ is a distribution on the positive half-line with the following density:
$$\frac{u^{\alpha 1}e^u}{\mathrm{\Gamma }(\alpha )},u>0.$$
Random variable with the gamma distribution is called a gamma variable. The classical gamma process ($`\gamma _t`$, $`t0`$) is an increasing stochastic process with independent homogeneous increments such that $`\gamma _t`$ has the gamma distribution with parameter $`t`$ for each $`t>0`$.
The following fundamental properties of this process have led us to more general developments and definitions discussed below, and we suggest that the reader comes back to them as a “toy model” when reading more general later parts of our paper.
If two variables $`\xi `$, $`\varphi `$ are independent and have standard gamma distributions, then $`\xi +\varphi `$ and $`\frac{\xi }{\xi +\varphi }`$ are also independent. Moreover, this is a characteristic property of the gamma distribution (Lukacs’ theorem, see Section 4 and ). Namely, if $`x_1`$, $`x_2`$ are independent positive random variables, and $`x_1+x_2`$ and $`\frac{x_1}{x_1+x_2}`$ are also independent, then $`x_1`$ and $`x_2`$ are common multiples of standard gamma variables with some parameters.
The first property easily implies
For each $`t>0`$, the process ($`\gamma _u/\gamma _t`$, $`ut`$) is independent on ($`\gamma _v`$, $`vt`$).
From property 2 one can deduce the following
For each $`t>0`$ and $`a>0`$ the law of the process $`(\frac{\gamma _u}{1+a},ut)`$ is equivalent to that of $`(\gamma _u,ut)`$ with Radon–Nikodym density $`(1+a)^t\mathrm{exp}(a\gamma _t)`$.
This property will lead us to a very important quasi-invariance of the law of the gamma process which was discovered in in a different way and was used for the representation theory (see Section 5).
It was noticed that for a small index $`\alpha `$ the stable distribution could be renormalized in such a way that the limit when $`\alpha `$ tends to zero is the gamma distribution . This fact leads us to an important conclusion (obtained also in another way in ) that the law of the gamma process is a limit of renormalized stable processes when $`\alpha 0`$ (see Section 9).
The last property allows us to consider the gamma process as an antithesis to the Brownian motion: both correspond to the extreme values of the parameter $`\alpha [0,2]`$, and in between we have stable processes which can be considered as a deformation from the Brownian motion to the gamma process. Moreover, we can see many strikingly similar properties, for example:
* Classical Bernstein’s characteristic property of Gaussian measures (components of both pairs of variables $`(X,Y)`$ and $`(X+Y,XY)`$ are independent if and only if the distributions of $`X`$ and $`Y`$ are Gaussian) is similar to our property 1.
* Let $`(B_t,t0)`$ be a Brownian motion. For every $`t>0`$ the Brownian bridge $`(B_u\frac{u}{t}B_t,ut)`$ is independent on $`(B_v,vt)`$, which is similar to our property 2.
* For every $`b`$ and every $`t>0`$ the law (a measure in the space of realizations) of the process $`(B_u+bu,ut)`$ is equivalent to that of $`(B_u,ut)`$ with Radon–Nikodym density $`\mathrm{exp}(bB_tb^2t/2)`$. The large group of symmetries of the gamma process is a multiplicative group of hyperbolic rotations (see Section 5). The law of the Brownian motion has a large orthogonal group of symmetries. Stable processes perhaps also have such large groups of symmetries, but they consist of non-linear transformations, and can be considered as a homotopy between hyperbolic and orthogonal rotations.
In view of this parallelism, it would be quite interesting to connect better the Brownian motion and the gamma process, and one such attempt is recorded in the Appendix, where we have gathered further properties of the gamma process which are not so close to the main context of the paper.
In view of the previous property stating that the gamma process is a limit of stable processes it is natural to ask what kind of stability do the gamma distribution and the gamma process have? In Section 10 we introduce a notion of generalized stability for sigma-finite measures, and show that the Lebesgue measure on the real line and the quasi-Lebesgue measure which is a sigma-finite measure equivalent to the law of the gamma process are also stable in this new sense.
Now we present our general framework and the list of topics touched on in the paper. The first one is the theory of Poisson–Dirichlet measures from the point of view of the gamma process. More general, we consider the theory of Lévy processes (a class of processes with independent values) without Gaussian component. The laws of such processes admit a canonical decomposition into a so-called conic part, i.e. a measure on the cone of positive series, and a standard product measure on sequences of points of the base space. This decomposition is of the most general character. It appeared first in , but the original proof was rather complicated and shaded the key relations with measures on positive series and a special role of stable and gamma processes, while our proof is based on very general and simple considerations.
This decomposition leads us to the theory of measures on positive series, i.e. the theory of random positive series, and their projections on the simplex of series with unit sum. Particularly, we deal with the Poisson–Dirichlet measures $`\mathrm{PD}(\theta )`$ and their generalizations which were studied by Kingman , Vershik–Shmidt , Pitman–Yor and their followers, and have numerous applications in combinatorics (e.g. ), number theory , mathematical biology (e.g. , see a detailed survey in ), etc. See also some interesting discussion on the Poisson–Dirichlet measures in most recently published book . These random series may be roughly characterized as “the most random convergent series” — a kind of white noise on convergent series.
The key feature in this work is the quasi-invariance of the gamma measure (the law of the gamma process) with respect to an infinite-dimensional group of multiplicators. It was first discovered and used in the works of Gelfand–Graev–Vershik on the representation theory of current groups, more exactly, of the group $`\mathrm{SL}(2,F)`$, where $`F`$ is an algebra of functions on a manifold. This property of the gamma measure followed from rather indirect considerations. The same considerations prompted the existence of an equivalent $`\sigma `$-finite measure which is invariant under multiplications by non-negative functions with zero integral of logarithm. We call this measure quasi-Lebesgue, because of its key property which is an infinite-dimensional generalization of the well-known property of the finite-dimensional Lebesgue measure. We mean invariance of the Lebesgue measure under the action of diagonal matrices with determinant $`1`$ (a Cartan subgroup). In this work we prove this quasi-invariance directly, starting from the characteristic functional (the Laplace transform) of the gamma measure, and construct explicitly the corresponding $`\sigma `$-finite measure. The same quasi-invariance implies new symmetry properties of the Poisson–Dirichlet measures with respect to Markovian transformations.
However, the most important is the following link outlined in . Both the gamma measure and the quasi-Lebesgue measure are weak limits of $`\alpha `$-stable measures (the laws of $`\alpha `$-stable processes) when $`\alpha `$ goes to zero. In terms of the Poisson–Dirichlet measures this fact was proved in . More exactly, Pitman–Yor define a two-parameter family of measures $`\mathrm{PD}(\alpha ,\theta )`$ on the simplex of positive series with unit sum and show that they converge to the Poisson–Dirichlet measures $`\mathrm{PD}(\theta )`$ when $`\alpha 0`$. It turns out that this convergence follows from the convergence of renormalized $`\alpha `$-stable measures to the gamma measure when $`\alpha 0`$. This convergence is a “commutative” analogue of the key fact discovered in which deals with the limit (more exactly, the derivative with respect to the parameter in $`0`$ which corresponds to a renormalization before taking the limit) of positive definite spherical functions of the complementary series of $`\mathrm{SL}(2,R)`$ when the parameter tends to a critical value (the so-called canonical state). Thus the Poisson–Dirichlet measures are directly related to the representation theory.
Another corollary of the quasi-invariance of the gamma measure allows to obtain easily the Markov–Krein identity which in our context relates the distribution of a linear functional on the gamma process and the distribution of the same functional on the normalized gamma process.
The paper is organized as follows.
In Section 2 we define a general class of Lévy processes. The main properties of these processes are studied in Section 3. Though in the sequel we deal only with the stable and gamma processes, the Decomposition Theorem 1 is proved by so general and natural considerations that we give it in the most general form that does not complicate the argument. This theorem states that the law of a Lévy process is the product of the conic part (the measure on the cone of positive convergent series) and a product measure on sequences of points of the base space. In fact, our proof applies to even more general processes. Theorem 2 is a characterization of measures on the cone of positive convergent series which are obtained as conic parts of Lévy processes (so-called measures of product type). In some sense, these measures enjoy the greatest possible independence of coordinates. Note that passing from the simplex to the cone (a kind of poissonization) simplifies many questions. For example, characterization of the conic parts of Lévy processes is simpler than characterization of the simplicial parts.
Section 4 contains definition and basic properties of the gamma process and the Poisson–Dirichlet measures which are the simplicial parts of the gamma process.
In Section 5 we prove the key property of the gamma measure, namely its quasi-invariance with respect to a group of multiplicators (Theorem 3). This group is rather wide and consists of all non-negative measurable functions with summable logarithm. Thus the gamma measure is a multiplicative analogue of the Wiener measure which is quasi-invariant with respect to a wide group of additive shifts. This parallel should certainly be considered more thoroughly. It is not known whether the measure with such supply of preserving transformations is unique. We give a partial counter-example for a smaller group of transformations. As to the quasi-invariance under the whole group, the gamma measure seems to be the unique measure enjoying this property.
In Section 6 we apply quasi-invariance of the gamma process to obtain the corresponding property of the Poisson–Dirichlet measures.
In Section 7 we introduce a $`\sigma `$-finite measure which is already invariant (projective invariant) under the same group of multiplicators. It is called quasi-Lebesgue since it generalizes a well-known property of the Lebesgue measure in $`^n`$.
Section 8 is devoted to definitions and basic properties of the stable processes. The simplicial parts of these properties are related to the two-parameter Poisson–Dirichlet measures $`\mathrm{PD}(\alpha ,\theta )`$.
In Section 9 we prove the statement suggested in that the gamma process is a weak limit of renormalized stable processes. As was noted above, this fact is related to the representation theory .
In Section 10 we give a new definition of stability which can be applied to $`\sigma `$-finite measures as well (unlike the classical definition). According to this definition, the quasi-Lebesgue measure is zero-stable.
Finally, in Sections 11 and 12 we deal with the Markov–Krein transform and its generalizations. We present a new probabilistic interpretation of this transform: formulae of this kind relate the distribution of a linear functional on the process with the distribution of the same functional on the normalized process. This interpretation sets the same question for general Lévy processes.
In the Appendix we try to trace some further connection of our topics, namely, some new and unexpected links with Brownian motion.
The topics touched upon in this paper stimulate many new problems, only a small part of which is mentioned above.
## 2 Definition of Lévy processes on arbitrary spaces
Let $`(X,\nu )`$ be a standard Borel space with a non-atomic finite non-negative measure $`\nu `$, and let $`\nu (X)=\theta `$ be the total charge of $`\nu `$. We denote by
$$D=\{z_i\delta _{x_i},x_iX,z_i,|z_i|<\mathrm{}\}$$
a real linear space of all finite real discrete measures on $`X`$.
Consider a class of measures $`\mathrm{\Lambda }`$ on the half-line $`_+`$ satisfying the following conditions,
$`\mathrm{\Lambda }(0,\mathrm{})=\mathrm{},`$ (1)
$`\mathrm{\Lambda }(1,\mathrm{})<\mathrm{},`$ (2)
$`{\displaystyle _0^1}s𝑑\mathrm{\Lambda }(s)<\mathrm{},`$ (3)
$`\mathrm{\Lambda }(\{0\})=0.`$ (4)
Let $`F_\mathrm{\Lambda }`$ be the infinitely divisible distribution with Lévy measure $`\mathrm{\Lambda }`$, i.e. the Laplace transform $`\psi _\mathrm{\Lambda }`$ of $`F_\mathrm{\Lambda }`$ is given by
$$\psi _\mathrm{\Lambda }(t)=\mathrm{exp}\left(_0^{\mathrm{}}(1e^{ts})𝑑\mathrm{\Lambda }(s)\right).$$
###### Definition 1
A homogeneous Lévy process on the space $`(X,\nu )`$ with Lévy measure $`\mathrm{\Lambda }`$ satisfying (1)–(4) is a generalized process on $`D`$ whose law $`P_\mathrm{\Lambda }=P_\mathrm{\Lambda }(\nu )`$ has the Laplace transform
$$𝔼\left[\mathrm{exp}\left(_Xa(x)𝑑\eta (x)\right)\right]=\mathrm{exp}\left(_X\mathrm{log}\psi _\mathrm{\Lambda }(a(x))𝑑\nu (x)\right),$$
(5)
where $`a`$ is an arbitrary non-negative bounded Borel function on $`X`$.
The correctedness of this definition is guaranteed by the following explicit construction (see \[20, chapter 8\]). Consider a Poisson point process on the space $`X\times _+`$ with mean measure $`\nu \times \mathrm{\Lambda }`$. We associate with a realization $`\mathrm{\Pi }=\{(X_i,Z_i)\}`$ of this process an element
$$\eta =\underset{(X_i,Z_i)\mathrm{\Pi }}{}Z_i\delta _{X_i}D.$$
(6)
Then $`\eta `$ is a random discrete measure obeying the law $`P_\mathrm{\Lambda }`$. Note that if $`\mathrm{\Lambda }`$ is a $`\delta `$-measure $`\delta _z`$ for some $`z_+`$, then $`\mathrm{\Pi }`$ is the Poisson process on the set $`X\times \{z\}`$ (which we identify with $`X`$) with mean measure $`\nu `$, and the corresponding random element $`\eta `$ is a measure that has equal charges $`z`$ at the points of this process. Thus a Lévy process with an arbitrary measure $`\mathrm{\Lambda }`$ is a continual convolution of independent Poisson processes on $`X`$ corresponding to different levels (charges).
It follows that the law $`P_\mathrm{\Lambda }`$ of the Lévy process is concentrated on the cone $`D^+=\{z_i\delta _{x_i}D:z_i>0\}D`$ consisting of all finite positive discrete measures on X. The conditions (1)–(4) imposed on the measure $`\mathrm{\Lambda }`$ have the following meaning: (1) implies that the random measure $`\eta `$ has an infinite number of atoms; (2) together with (3) guarantees that $`\eta `$ is a finite measure, i.e. the sum of charges is finite; finally, (4) means that our Lévy process has no Gaussian component.
Remarks. 1. Our definition of the Lévy process is closely related to the notion of completely random measure, see , \[20, chapter 8\].
2. If $`X=_+`$ and $`\nu `$ is the Lebesgue measure on $`_+`$, we recover an ordinary definition of a subordinator (a process with stationary positive independent increments) corresponding to the Lévy measure $`\mathrm{\Lambda }`$.
3. It is easy to see that $`P_\mathrm{\Lambda }(\nu )=P_{\theta \mathrm{\Lambda }}(\nu /\theta )`$, i.e. we may consider only normalized parameter measures $`\nu `$. Thus in the sequel we assume $`\nu (X)=1`$.
## 3 Decomposition theorem for Lévy processes and measures of product type on the cone
Consider the cone
$$C=\{z=(z_1,z_2,\mathrm{}):z_1z_2\mathrm{}0,z_i<\mathrm{}\}l^1.$$
We now define a special class of measures on $`C`$ indexed by infinitely divisible distributions on the half-line. Fix an integer $`n`$ and a probability vector $`p=(p_1,\mathrm{},p_n)`$ (i.e. a vector $`p`$ with $`p_1,\mathrm{},p_n>0`$ and $`p_1+\mathrm{}+p_n=1`$). Consider a sequence $`\xi _i`$ of i.i.d. variables such that $`P(\xi _i=k)=p_k`$ for $`k=1,\mathrm{}n`$. For $`Q=(Q_1,Q_2,\mathrm{})C`$, denote by $`\mathrm{\Sigma }_k^{(p)}=\mathrm{\Sigma }_k^{(p)}(Q)`$ the random sum $`\mathrm{\Sigma }_k^{(p)}=_{i:\xi _i=k}Q_i`$. Let $`Q`$ be a random series with distribution $`\varkappa `$ on $`C`$ such that the distribution $`F`$ of the sum $`Q_i`$ is infinitely divisible.
###### Definition 2
We say that a series $`Q`$ (and its distribution $`\varkappa `$) is of product type, if for each $`n`$ and each probability vector $`p`$ the sums $`\mathrm{\Sigma }_1^{(p)},\mathrm{}\mathrm{\Sigma }_n^{(p)}`$ are independent and $`\mathrm{\Sigma }_k^{(p)}`$ obeys the law $`F^{p_k}`$.
We define a map $`T:D^+C\times X^{\mathrm{}}`$ by
$$T\eta =((Q_1,Q_2,\mathrm{}),(X_1,X_2,\mathrm{})),\text{ if }\eta =Q_i\delta _{X_i}.$$
(7)
###### Definition 3
Let $`P`$ be a distribution on the space $`D^+`$, and let $`\eta `$ be a random process obeying the law $`P`$. The random sequence of charges $`Q_1,Q_2,\mathrm{}`$ is called the conic part of the process $`\eta `$, and its distribution on the cone $`C`$ is called the conic part of the law $`P`$.
Note that in view of representation (6) the conic part of the Lévy process with Lévy measure $`\mathrm{\Lambda }`$ is just the ordered sequence of points of the Poisson process on $`_+`$ with mean measure $`\mathrm{\Lambda }`$. Thus the conic part depends only on $`\mathrm{\Lambda }`$ and not on the $`(X,\nu )`$. In fact, the following theorem shows that studying the Lévy process may be essentially reduced to studying its conic part, since the construction of the process includes the parameter measure in a trivial way. This fundamental property of homogeneous Lévy processes is a particular case of the representation theorem first proved in . We present here a simpler proof of this fact. The key point of our proof is the following lemma. Let $`(X,\nu )`$ be a standard Borel space with continuous probability measure $`\nu `$. Denote $`X^k=X\times \mathrm{}\times X`$ ($`k`$ factors), $`\nu ^k=\nu \times \mathrm{}\times \nu `$ ($`k`$ factors) and let $`\nu _{diag}`$ be the image of $`\nu `$ under the diagonal map $`x(x,\mathrm{},x)`$.
###### Lemma 1
Let $`\tau `$ be some continuous probability measure on $`X^k`$. If for each measure preserving transformation $`L`$ of $`(X,\nu )`$, the transformation $`L^k=L\times \mathrm{}\times L`$ ($`k`$ factors) preserves $`\tau `$, then $`\tau `$ is a convex combination of $`\nu ^k`$ and $`\nu _{diag}`$.
* For simplicity, assume $`k=2`$, the general case being quite similar. The diagonal $`\mathrm{\Delta }=\{(x,x),xX\}`$ is obviously an invariant subset for the group $`G=\{L\times L\}`$, where $`L`$ runs over the set of all $`\nu `$-preserving transformations of the space $`X`$. Thus it suffices to show that if $`\tau `$ is concentrated on the set $`(X\times X)\mathrm{\Delta }`$, then $`\tau =\nu \times \nu `$, and if $`\tau `$ is concentrated on $`\mathrm{\Delta }`$, then $`\tau =\nu _{diag}`$.
In the first case let $`\xi _n=\{A_i\}_{i=1}^{2^n}`$ be an arbitrary partition of the space $`X`$ into $`2^n`$ sets of equal $`\nu `$-measure $`1/2^n`$. Denote by $`\xi _n^2`$ the corresponding partition of the space $`X\times X`$, i.e. $`\stackrel{~}{\xi }_n=\{A_{ij}\}`$, where $`A_{ij}=A_i\times A_j`$. The group $`G`$ acts transitively on the set of non-diagonal elements of $`\xi _n`$. Thus all non-diagonal elements have equal $`\tau `$-measure. Denote $`Y_n=(X\times X)(A_i\times A_i)`$ and $`\epsilon _n=\tau (Y_n)`$. Since $`\tau `$ is concentrated on $`(X\times X)\mathrm{\Delta }`$, we have $`\epsilon _n1`$ as $`n\mathrm{}`$. Considering finer partitions and using the above argument, we obtain that for each $`k`$, if a rectangle $`A\times BY_n`$ and $`\nu (A)=\nu (B)=\nu (Y_n)/2^k`$, then $`\tau (A\times B)=\epsilon _n/4^k`$. But then the restriction of $`\tau `$ on the set $`Y_n`$ equals $`\epsilon _n(\nu \times \nu )`$. Letting $`n\mathrm{}`$, we obtain $`\tau =\nu \times \nu `$.
In the second case, identifying the diagonal $`\mathrm{\Delta }`$ with $`X`$, we obtain that $`\tau `$ is a measure on $`X`$ which is invariant under all $`\nu `$-preserving transformations, hence obviously $`\tau =\nu `$. $``$
###### Theorem 1
Let $`\eta =Q_i\delta _{X_i}`$ be a homogeneous Lévy process on the space $`(X,\nu )`$ with Lévy measure $`\mathrm{\Lambda }`$. Then $`TP_\mathrm{\Lambda }=\varkappa _\mathrm{\Lambda }\times \nu ^{\mathrm{}}`$, i.e. $`X_1,X_2,\mathrm{}`$ is a sequence of i.i.d. random variables with common distribution $`\nu `$, and this sequence is independent of the conic part $`\{Q_i\}_i`$.
* Let $`L:XX`$ be a $`\nu `$-preserving transformation of $`X`$. This transformation acts on the space $`D`$ by substituting coordinates, i.e. $`z_i\delta _{x_i}z_i\delta _{Lx_i}`$, and it is clear that the law $`P_\mathrm{\Lambda }`$ of a homogeneous Lévy process on $`(X,\nu )`$ is invariant under $`L`$. Denote by $`P_\mathrm{\Lambda }^z`$ the conditional measure of $`P_\mathrm{\Lambda }`$ given the conic part equal to $`zC`$. The transformation $`L`$ acts “fibre-wise”, i.e. it does not change the conic part, hence $`L`$ preserves almost all conditional measures $`P_\mathrm{\Lambda }^z`$. In particular, if we denote by $`(P_\mathrm{\Lambda }^z)_k`$ the conditional distribution of the first $`k`$ points $`X_1,\mathrm{},X_k`$ on the space $`X^k`$, then the transformation $`L^k`$ preserves $`(P_\mathrm{\Lambda }^z)_k`$.
Now it follows from Lemma 1 that for almost all $`z`$ $`(P_\mathrm{\Lambda }^z)_k=\nu ^k`$ for all $`k`$, i.e. $`P_\mathrm{\Lambda }^z=\nu ^{\mathrm{}}`$, and Theorem 1 follows. $``$
###### Theorem 2
The measure $`\varkappa `$ on the cone $`C`$ is the conic part of some Lévy process $`P_\mathrm{\Lambda }`$ with Lévy measure $`\mathrm{\Lambda }`$ satisfying (1)–(4) if and only if it is of product type with $`F=F_\mathrm{\Lambda }`$.
* Given fixed $`n`$ and a probability vector $`p=(p_1,\mathrm{},p_n)`$, consider a partition $`X=A_1\mathrm{}A_n`$ of the space $`X`$ such that $`\nu (A_k)=p_k`$, $`k=1,\mathrm{},n`$. Let $`X_1,X_2,\mathrm{}`$ be the sequence of i.i.d. variables with common distribution $`\nu `$ and assume $`\xi _i=k`$, if $`X_iA_k`$. Then the random variables $`\xi _i`$ form a sequence of i.i.d. variables, and $`P(\xi _i=k)=p_k`$. Consider a random process
$$\eta =Q_i\delta _{X_i},$$
(8)
where the sequence $`Q_1,Q_2,\mathrm{}`$ is independent of $`\{X_i\}`$ and obeys the law $`\varkappa `$. Let $`\mathrm{\Sigma }_k^{(p)}=_{i:\xi _i=k}Q_i`$. It is easy to see that for arbitrary $`t_1,\mathrm{},t_n>0`$
$$𝔼\left[\mathrm{exp}\left(\underset{k=1}{\overset{n}{}}t_k\mathrm{\Sigma }_k^{(p)}\right)\right]=𝔼\left[\mathrm{exp}\left(_Xa(x)𝑑\eta (x)\right)\right],$$
(9)
where $`a`$ is a step function such that $`a(x)=t_i`$, if $`xA_i`$.
Now let $`\varkappa `$ be the conic part of the law $`P_\mathrm{\Lambda }`$ of some Lévy process. Then, by Theorem 1, the process $`\eta `$ defined by (8) obeys $`P_\mathrm{\Lambda }`$, and it follows from the Laplace transform formula (5) that the right-hand side of (9) equals
$$\underset{i=1}{\overset{n}{}}\psi _\mathrm{\Lambda }(t_i)^{\nu (A_i)}=\underset{i=1}{\overset{n}{}}\psi _\mathrm{\Lambda }(t_i)^{p_i}.$$
Since the left-hand side of (9) is the Laplace transform of the common distribution of the variables $`\mathrm{\Sigma }_1^{(p)},\mathrm{},\mathrm{\Sigma }_n^{(p)}`$, we obtain that $`\mathrm{\Sigma }_1^{(p)},\mathrm{},\mathrm{\Sigma }_n^{(p)}`$ are independent, and $`\mathrm{\Sigma }_k^{(p)}`$ obeys the law $`F^{p_k}`$, i.e. $`\varkappa `$ is of product type.
Conversely, let $`\varkappa `$ be a measure of product type corresponding to an infinitely divisible law $`F`$. Define a random process $`\eta `$ on an arbitrary measurable space $`(X,\nu )`$ satisfying the conditions of Definition 1 by (8). The above argument shows that $`\eta `$ satisfies (5) with $`\mathrm{\Lambda }`$ equal to the Lévy measure of $`F`$ for all positive step functions $`a`$, and one can easily extend this to all bounded positive Borel functions by continuity. Thus $`\eta `$ obeys $`P_\mathrm{\Lambda }`$, and $`\varkappa `$ is the conic part of $`P_\mathrm{\Lambda }`$. $``$
The strong law of large numbers for Poisson processes (see \[20, 4.55\]) implies the following result on the asymptotic behavior of the vectors $`Z=(Z_1,Z_2,\mathrm{})C`$ obeying the law $`\varkappa _\mathrm{\Lambda }`$.
###### Proposition 1 ()
Let $`m_\mathrm{\Lambda }(t)=\mathrm{\Lambda }(t,\mathrm{})`$, $`t>0`$. Then
$$\underset{n\mathrm{}}{lim}\frac{m_\mathrm{\Lambda }(Z_n)}{n}=1$$
(10)
for almost all with respect to $`\varkappa _\mathrm{\Lambda }`$ vectors $`ZC`$.
Note that we may rewrite (10) as
$$\underset{z0}{lim}\frac{m_\mathrm{\Lambda }(z)}{\mathrm{\#}\{i:Z_i>y\}}=1.$$
(11)
In other words, if $`\eta `$ is a random discrete measure obeying the law $`P_\mathrm{\Lambda }`$, then the number of charges of $`\eta `$ which are greater than $`z`$ has the same asymptotics when $`z0`$ as the tail $`\mathrm{\Lambda }(z,\mathrm{})`$ of the Lévy measure $`\mathrm{\Lambda }`$.
Denote by $`D_1^+D^+`$ the simplex of all normalized atomic measures. Then $`D^+=D_1^+\times [0,\mathrm{})`$, i.e. each $`\eta D^+`$ can be represented as
$$\eta =(\eta /\eta (X),\eta (X)).$$
(12)
The second coordinate is the total charge of the measure $`\eta `$. It follows from the definition of a Lévy process that $`\eta (X)`$ obeys the infinite divisible law $`F_\mathrm{\Lambda }`$ corresponding to the Lévy measure $`\mathrm{\Lambda }`$. The first coordinate is called the normalization of the measure $`\eta `$. Note that, in general, the law of a Lévy process is not a product measure in this decomposition (see Lemmas 2 and 3 below).
Using this decomposition, consider a map $`T^{}:D^+_+\times \mathrm{\Sigma }\times X^{\mathrm{}}`$, where $`\mathrm{\Sigma }=\{y=(y_1,y_2,\mathrm{}):y_1y_2\mathrm{}0,y_1+y_2+\mathrm{}=1\}`$ is the infinite-dimensional simplex, and
$$T^{}\eta =(\eta (X),(Q_1/\eta (X),Q_2/\eta (X),\mathrm{}),(X_1,X_2,\mathrm{})),\text{if}\eta =Q_i\delta _{X_i}.$$
###### Definition 4
The normalized sequence of charges $`Q_1/\eta (X),Q_2/\eta (X),\mathrm{}`$ is called the simplicial part of the process and its distribution $`\sigma _\mathrm{\Lambda }`$ on $`\mathrm{\Sigma }`$ is called the simplicial part of the law $`P_\mathrm{\Lambda }`$.
## 4 The gamma process and Poisson–Dirichlet distributions
Let $`(X,\nu )`$ be a standard Borel space with a non-atomic finite non-negative measure $`\nu `$, and let $`\nu (X)=\theta `$ be the total charge of $`\nu `$.
###### Definition 5
The gamma process with scale parameter $`\lambda >0`$ on the space $`(X,\nu )`$ is a Lévy process on $`(X,\nu )`$ corresponding to the Lévy measure with density $`d\mathrm{\Lambda }_\mathrm{\Gamma }^\lambda (z)=z^1e^{\lambda z}dz`$, $`z>0`$.
As in general case, instead of considering a non-normalized parameter measure $`\nu `$, we may take its normalization $`\overline{\nu }=\nu /\nu (X)`$ and the Lévy measure with density $`\theta z^1e^{\lambda z}dz`$. But in this particular case it is often more convenient to use the above definition with a non-normalized measure. In the sequel we shall use both variants without additional mention.
The corresponding infinitely divisible law is the gamma distribution $`𝒢_{\theta ,\lambda }`$ on $`_+`$ with shape parameter $`\theta `$ and scale parameter $`\lambda `$, i.e.
$$d𝒢_{\theta ,\lambda }=\frac{\lambda ^\theta }{\mathrm{\Gamma }(\theta )}t^{\theta 1}e^{\lambda t}dt,t>0.$$
Note that $`\lambda `$ is a trivial scale parameter. Namely, if $`\eta `$ is the gamma process with scale parameter $`1`$, then the gamma process $`\eta ^\lambda `$ with scale parameter $`\lambda `$ is obtained from $`\eta `$ by multiplying by $`\lambda `$, i.e. $`\eta ^\lambda =\lambda \eta `$. Thus we will consider only gamma processes with scale parameter $`1`$.
The law $`P_\mathrm{\Gamma }=P_\mathrm{\Gamma }(\nu )`$ of the gamma process (called the gamma measure on the space $`(X,\nu )`$) is thus given by the Laplace transform
$$𝔼_\mathrm{\Gamma }\left[\mathrm{exp}\left(_Xa(x)𝑑\eta (x)\right)\right]=\mathrm{exp}\left(_X\mathrm{log}\left(1+a(x)\right)𝑑\nu (x)\right),$$
(13)
where $`a`$ is an arbitrary non-negative bounded Borel function on $`X`$.
Let $`=(X,\nu )`$ be the set of (classes$`mod0`$ of) non-negative measurable functions on the space $`X`$ with $`\nu `$-summable logarithm,
$$=\{a:X_+:_X|\mathrm{log}a(x)|𝑑\nu (x)<\mathrm{}\}.$$
It follows from the Poisson construction (6) and Campbell’s theorem on sums over Poisson processes (see \[20, 3.2\]) that each function $`a`$ correctly defines a measurable linear functional $`\eta f_a(\eta )=_Xa(x)𝑑\eta (x)`$ on $`D`$, and formula (13) holds for all $`a`$.
It is well known that the gamma distribution enjoys the following property. If $`Y`$ and $`Z`$ are independent gamma variables with the same scale parameter, then the variables $`Y+Z`$ and $`\frac{Y}{Y+Z}`$ are independent. Moreover, a remarkable result of Lukacs (similar to the famous Bernstein’s characterization of normal distributions) states that this property is characteristic of the gamma distribution, i.e. if $`Y`$ and $`Z`$ are independent non-degenerate positive random variables, and the variables $`Y+Z`$ and $`\frac{Y}{Y+Z}`$ are independent, then both $`Y`$ and $`Z`$ have gamma distributions with the same scale parameter. In other words, the independence condition may be formulated as follows. Let us describe a point $`x=(x_1,x_2)`$ in the first quadrant $`_+\times _+`$ by the sum $`x_1+x_2`$ of its coordinates and its projection onto the unit 2-simplex (i.e. interval) $`\{y=(y_1,y_2):y_1,y_20,y_1+y_2=1\}`$. Then the distribution $`𝒢_{\theta _1,\lambda }\times 𝒢_{\theta _2,\lambda }`$ on $`_+\times _+`$ is a product measure in these coordinates. These results imply the corresponding statements for the gamma process which are a key point for many important properties of $`P_\mathrm{\Gamma }`$.
###### Lemma 2
In representation (12) the gamma measure is a product measure $`P_\mathrm{\Gamma }=𝒢_\theta \times \overline{P}_\mathrm{\Gamma }`$, i.e. the total charge $`\gamma (X)`$ of the gamma process and the normalized gamma process $`\overline{\gamma }=\gamma /\gamma (X)`$ are independent. The distribution of the total charge is the gamma distribution $`𝒢_{\theta ,1}`$ with shape parameter $`\theta `$ and scale parameter 1.
###### Lemma 3
If the law $`P_\mathrm{\Lambda }`$ of some Lévy process is a product measure in representation (12), then $`P_\mathrm{\Lambda }`$ is a gamma process, i.e. $`d\mathrm{\Lambda }(z)=z^1e^{\lambda z}dz`$, $`z>0`$, for some $`\lambda >0`$.
###### Definition 6 ()
The simplicial part of the gamma measure $`P_\mathrm{\Gamma }(\nu )`$ with $`\nu (X)=\theta `$ is called the Poisson–Dirichlet distribution with parameter $`\theta `$ and denoted by $`\mathrm{PD}(\theta )`$.
The above definition is just one of many other possible definitions of the Poisson–Dirichlet distributions. These distributions arise in many fields of pure and applied mathematics. They play an important role in statistics because of their connection with Dirichlet distributions and Dirichlet random measures . In number theory Poisson–Dirichlet measures arise in the problem of distribution of prime divisors of a random integer . There are many asymptotic combinatorial problems leading to the measures $`\mathrm{PD}(\theta )`$, such as the distribution of the cycle lengths of a random permutation or the distribution of the degrees of the irreducible factors of a random monic polynomial over a finite field , etc. The Poisson–Dirichlet distributions also play an important role in applications to population genetics, ecology and physics. A (non-complete) survey of different aspects of Poisson–Dirichlet measures can be found in .
It follows from Lemma 2 that $`T^{}P_\mathrm{\Gamma }=𝒢_\theta \times \mathrm{PD}(\theta )\times \overline{\nu }^{\mathrm{}}`$, i.e. the conic part of the gamma measure is a product measure $`𝒢_{\theta ,1}\times \mathrm{PD}(\theta )`$. We call this measure the conic Poisson–Dirichlet distribution with parameter $`\theta `$ and denote it by $`\mathrm{CPD}(\theta )`$.
Note that in case of the gamma process $`m(t)=\theta _t^{\mathrm{}}s^1e^s𝑑s\theta \mathrm{log}t`$. Thus, by Proposition 1,
$$\underset{n\mathrm{}}{lim}\frac{\mathrm{log}Z_n}{n}=\frac{1}{\theta }$$
(14)
almost surely with respect to $`\mathrm{CPD}(\theta )`$. It follows that the same asymptotics holds for $`\mathrm{PD}(\theta )`$, i.e. $`lim_n\mathrm{}\frac{\mathrm{log}Y_n}{n}=\frac{1}{\theta }`$ for almost all vectors $`Y\mathrm{\Sigma }`$ with respect to $`\mathrm{PD}(\theta )`$. In particular, we see that the measures $`\mathrm{PD}(\theta )`$ (as well as $`\mathrm{CPD}(\theta )`$) are mutually singular for different $`\theta `$.
Many properties of ordinary Poisson–Dirichlet distributions have their natural analogues for conic distributions. For example, it is well known that the measure $`\mathrm{PD}(\theta )`$ may be obtained by the following stick breaking process. Let $`Y_1`$ be a random variable on the interval $`[0,1]`$ obeying the law $`\theta (1t)^{\theta 1}dt`$, $`t[0,1]`$. If we have already constructed variables $`Y_1,\mathrm{},Y_n`$, then $`Y_{n+1}`$ has the same distribution scaled on the interval $`[Y_n,1]`$. Thus we obtain a random sequence $`0=Y_0<Y_1<Y_2<\mathrm{}<1`$. Let $`Z_k=Y_kY_{k1}`$, $`k=1,2,\mathrm{}`$. The Poisson–Dirichlet measure $`\mathrm{PD}(\theta )`$ is the distribution of the order statistics $`Z_{(1)}Z_{(2)}\mathrm{}`$ of the sequence $`Z_1,Z_2,\mathrm{}`$. It follows that the conic Poisson–Dirichlet measure $`\mathrm{CPD}(\theta )`$ may be obtained by the randomized version of this procedure. Namely, for the first step we choose the random length $`L`$ of the interval with gamma distribution $`𝒢_{\theta ,1}`$, and then proceed as before starting with the interval $`[0,L]`$.
As shown in , the Poisson–Dirichlet measures $`\mathrm{PD}(\theta )`$ can be naturally included into a two-parameter family $`\mathrm{PD}(\alpha ,\theta )`$ of distributions on the simplex $`\mathrm{\Sigma }`$. This family is obtained by the following non-stationary version of the stick breaking process. Let $`Y_1`$ be a random variable on the interval $`[0,1]`$ obeying the beta distribution $`B(1\alpha ,\theta +\alpha )`$. If we have already constructed variables $`Y_1,\mathrm{},Y_n`$, then $`Y_{n+1}`$ has the beta distribution $`B(1\alpha ,\theta +(n+1)\alpha )`$ scaled on the interval $`[Y_n,1]`$. Thus we obtain a random sequence $`0=Y_0<Y_1<Y_2<\mathrm{}<1`$. Let $`Z_k=Y_kY_{k1}`$, $`k=1,2,\mathrm{}`$. The two-parameter Poisson–Dirichlet measure $`\mathrm{PD}(\alpha ,\theta )`$ is the distribution of the order statistics $`Z_{(1)}Z_{(2)}\mathrm{}`$ of the sequence $`Z_1,Z_2,\mathrm{}`$.
The range of admissible parameters is the union of the sets
$$\{(\alpha ,\theta ):\mathrm{\hspace{0.33em}0}\alpha <1,\theta >\alpha \}\text{and}\{(\alpha ,m\alpha ):\alpha <0,m\}.$$
The first case $`\alpha (0,1)`$ is the most interesting, the second one being a sort of degenerate case. The ordinary Poisson–Dirichlet distributions $`\mathrm{PD}(\theta )`$ correspond to $`\alpha =0`$, i.e. $`\mathrm{PD}(\theta )=\mathrm{PD}(0,\theta )`$. See for various properties of the measures $`\mathrm{PD}(\alpha ,\theta )`$. In particular, the distributions $`\mathrm{PD}(\alpha ,\theta )`$ with fixed $`\alpha 0`$ and different $`\theta `$ are absolutely continuous (unlike the case $`\alpha =0`$). We discuss some questions related to the two-parameter Poisson–Dirichlet distributions in Sections 8 and 12.
## 5 Multiplicative quasi-invariance of the gamma process
As was mentioned above, the law $`P_\mathrm{\Lambda }`$ of each Lévy process is invariant under all $`\nu `$-preserving transformations of the space $`(X,\nu )`$ which act on $`D`$ by substituting the coordinates. However, the gamma measure $`P_\mathrm{\Gamma }`$ enjoys additional invariance properties. We present now a large group of linear transformations of the space $`D`$ (preserving the cone $`D^+`$) for which $`P_\mathrm{\Gamma }`$ is a quasi-invariant measure.
Consider the above defined class $``$ of non-negative functions on $`X`$ with $`\nu `$-summable logarithm. Each function $`a`$ defines not only a linear functional $`f_a`$ on $`D`$ but also a multiplicator $`M_a:DD`$ by $`(M_a\eta )(x)=a(x)\eta (x)`$, that is $`M_a\eta =a(x_i)z_i\delta _{x_i}`$ for $`\eta =z_i\delta _{x_i}`$. Note that the set $``$ is a commutative group with respect to pointwise multiplication of functions, and $`M_a`$ is a group action of $``$. Denote by $`\stackrel{~}{a}`$ the function $`\stackrel{~}{a}(x)=(1/a(x))1`$.
The following property of the gamma process was first discovered in in quite different terms; it plays an important role in the representation theory of the current group $`\mathrm{SL}(2,F)`$, where $`F`$ is the space of functions on a manifold.
###### Theorem 3
For each $`a`$, the measure $`P_\mathrm{\Gamma }`$ is quasi-invariant under $`M_a`$, and the corresponding density is given by the following formula,
$$\frac{d(M_aP_\mathrm{\Gamma })}{dP_\mathrm{\Gamma }}(\eta )=\mathrm{exp}\left(_X\mathrm{log}a(x)𝑑\nu (x)\right)\mathrm{exp}\left(_X\stackrel{~}{a}(x)𝑑\eta (x)\right).$$
(15)
* Fix $`a`$ and let $`\xi =L_a\eta `$. Consider an arbitrary function $`b`$. Then $`_Xb(x)𝑑\xi (x)=_Xa(x)b(x)𝑑\eta (x)`$. Thus, in view of (13), the Laplace transform $`𝔼\left[\mathrm{exp}\left(_Xb(x)𝑑\xi (x)\right)\right]`$ equals
$$\begin{array}{c}𝔼\left[\mathrm{exp}\left(_Xa(x)b(x)𝑑\eta (x)\right)\right]=\mathrm{exp}\left(_X\mathrm{log}\left(1+a(x)b(x)\right)𝑑\nu (x)\right)=\hfill \\ \hfill =\mathrm{exp}\left(_X\mathrm{log}a(x)𝑑\nu (x)\right)\mathrm{exp}\left(_X\mathrm{log}\left(\frac{1}{a(x)}+b(x)\right)𝑑\nu (x)\right).\end{array}$$
Using (13) once more, we may consider the last factor as the Laplace transform of $`P_\mathrm{\Gamma }`$ calculated on the function $`(\frac{1}{a(x)}1)+b(x)=\stackrel{~}{a}(x)+b(x)`$. Denote $`I(a)=_X\mathrm{log}a(x)𝑑\nu (x)`$. Then we have
$$\begin{array}{c}𝔼\left[\mathrm{exp}\left(_Xb(x)𝑑\xi (x)\right)\right]=I(a)𝔼\left[\mathrm{exp}\left(_X\left(\stackrel{~}{a}(x)+b(x)\right)𝑑\eta (x)\right)\right]=\hfill \\ \hfill =𝔼\left[I(a)\mathrm{exp}\left(_X\stackrel{~}{a}(x)𝑑\eta (x)\right)\mathrm{exp}\left(_Xb(x)𝑑\eta (x)\right)\right],\end{array}$$
and Theorem 3 follows. $``$
In particular, if we consider multiplication by constant $`c>0`$, then the corresponding density depends only on the total charge, namely
$$\frac{d(M_cP_\mathrm{\Gamma })}{dP_\mathrm{\Gamma }}(\eta )=\frac{1}{c^\theta }\mathrm{exp}\left(\left(1\frac{1}{c}\right)\eta (X)\right).$$
(16)
###### Theorem 4
The action of the group $``$ on the space $`(D^+,P_\mathrm{\Gamma })`$ is ergodic.
* Let $`G:D^+`$ be a $`P_\mathrm{\Gamma }`$-measurable functional on $`D^+`$ which is invariant under all $`M_a`$ i.e. $`G(M_a\eta )=G(\eta )`$ a.e. with respect to $`P_\mathrm{\Gamma }`$. Consider an arbitrary Borel function $`k:`$. Then for each $`a`$
$$\begin{array}{c}𝔼\left[k(G(\eta ))\right]=𝔼\left[k(G(M_a\eta ))\right]=\hfill \\ \hfill 𝔼\left[k(G(\eta ))\mathrm{exp}\left(_X\stackrel{~}{a}(x)𝑑\eta (x)\right)\right]\mathrm{exp}\left(_X\mathrm{log}a(x)𝑑\nu (x)\right),\end{array}$$
where $`𝔼`$ denotes the expectation with respect to $`P_\mathrm{\Gamma }`$. But in view of (13) the last factor equals
$$\left(𝔼\left[\mathrm{exp}\left(_X\stackrel{~}{a}(x)𝑑\eta (x)\right)\right]\right)^1=\left(𝔼\left[\mathrm{exp}\left(f_{\stackrel{~}{a}}(\eta )\right)\right]\right)^1,$$
hence we have
$$𝔼\left[k(G(\eta ))\mathrm{exp}\left(f_{\stackrel{~}{a}}(\eta )\right)\right]=𝔼\left[k(G(\eta ))\right]𝔼\left[\mathrm{exp}\left(f_{\stackrel{~}{a}}(\eta )\right)\right].$$
Thus $`G`$ is independent of every functional $`f_a`$, and Theorem 4 follows. $``$
It is natural to ask whether the quasi-invariance property stated in Theorem 3 is characteristic of the gamma process. If we fix the density, then the answer is positive (), i.e. the gamma measure is the only measure on $`D^+`$ satisfying (15). On the other hand, for a smaller group of multiplications, the answer may be negative, as the following example shows. Let us call a process quasi-multiplicative if its law is invariant under all transformations $`M_a`$ with constant $`a>0`$. For simplicity, consider subordinators (i.e. Lévy processes for $`X=_+`$ and $`\nu `$ equal to the Lebesgue measure).
###### Proposition 2 ()
Let $`\eta `$ be a subordinator with Lévy measure $`\mathrm{\Lambda }(dx)=k(x)dx`$, where $`k(x)>0`$, and denote $`g(x)=xk(x)`$. Then $`\eta `$ is quasi-multiplicative if and only if for all $`a>0`$
$$_0^1\left(\sqrt{g(x/a)}\sqrt{g(x)}\right)^2\frac{dx}{x}<\mathrm{}.$$
(17)
It is shown in that for each $`m<1/2`$ any function $`k_m(x)`$ that satisfies
$$k_m(x)=\frac{1}{x}\left(\mathrm{log}\frac{1}{x}\right)^{2m}\text{ for }x<1/2\text{ and }_{1/2}^{\mathrm{}}k_m(x)𝑑x<\mathrm{}$$
(18)
provides an example of a quasi-multiplicative subordinator that is not equivalent to any scaled gamma process. It is clear that if $`\eta `$ is quasi-multiplicative, then its law is quasi-invariant under all step functions with finitely many steps, but the quasi-invariance under the whole group $``$ does not take place in this example. It is not known if there exists a measure different from $`P_\mathrm{\Gamma }`$ which has this property.
Computation of the law of hitting time of the drifted gamma process is coherent with the quasi-invariance of the gamma process. This fact is also parallel to some property of Brownian motion and will be considered in more details elsewhere.
## 6 Quasi-invariance of the Poisson–Dirichlet distributions
Let $`a`$. According to the general theory of polymorphisms (see ), the transformation $`M_a`$ induces a Markovian operator $`R_a`$ on the cone $`C`$. Namely, let $`z=(z_1,z_2,\mathrm{})C`$. Consider the conditional distribution $`P_\mathrm{\Gamma }^z`$ of the gamma process on $`(X,\nu )`$, given the conic part equal to $`z`$. Then the random image of the point $`z`$ under $`R_a`$ is the conic part of the process $`M_a\eta `$, where $`\eta `$ obeys the law $`P_\mathrm{\Gamma }^z`$. It follows from Theorem 1 that
$$R_az=V(a(X_1)z_1,a(X_2)z_2,\mathrm{}),$$
where $`(X_1,X_2,\mathrm{})`$ is a sequence of i.i.d. random variables on $`X`$ with common distribution $`\nu `$, and $`V`$ denotes a map that arranges the coordinates in non-increasing order.
In a similar way, the transformation $`M_a`$ induces a Markovian operator $`S_a`$ on the simplex $`\mathrm{\Sigma }`$,
$$S_ay=V(\frac{a(X_1)y_1}{\sigma },\frac{a(X_2)y_2}{\sigma },\mathrm{}),$$
where the sequence $`(X_1,X_2,\mathrm{})`$ is as before, and $`\sigma =a(X_1)y_1+a(X_2)y_2+\mathrm{}`$.
Note that the definitions of the operators $`S_a`$ and $`R_a`$ depend only on the distribution of the function $`a`$. Thus we may assume that $`X=[0,1]`$ and $`\nu =\theta \lambda `$, where $`\lambda `$ is the Lebesgue measure on the interval.
Theorems 34 immediately imply
###### Theorem 5
1) The Poisson–Dirichlet distribution $`\mathrm{PD}(\theta )`$ is quasi-invariant under the Markovian operator $`S_a`$ for all $`a`$, and
$$\frac{dS_a\mathrm{PD}(\theta )}{d\mathrm{PD}(\theta )}(y)=\mathrm{exp}\left(\theta _0^1\mathrm{log}a(s)𝑑s\right)_0^{\mathrm{}}\frac{\sigma ^{\theta 1}}{\mathrm{\Gamma }(\theta )}\left(\underset{i=1}{\overset{\mathrm{}}{}}L_{1/a}(\sigma y_i)\right)𝑑\sigma ,$$
where $`L_{1/a}()`$ is the Laplace transform of the distribution of the function $`1/a(t)`$ with respect to the uniform distribution on the interval $`[0,1]`$.
2) The Poisson–Dirichlet distribution $`\mathrm{PD}(\theta )`$ is ergodic with respect to $`\{S_a\}_a`$.
## 7 Quasi-Lebesgue measure and a representation of the current group
In this section we define, following , a $`\sigma `$-finite measure on $`D^+`$ which is equivalent to $`P_\mathrm{\Gamma }`$ and invariant under a subgroup $`_0`$ of $``$ consisting of all functions $`a`$ such that $`_X\mathrm{log}a(x)𝑑\nu (x)=0`$.
###### Definition 7
Consider a $`\sigma `$-finite measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ on $`D^+`$ defined by
$$\frac{d\stackrel{~}{P}_\mathrm{\Gamma }}{dP_\mathrm{\Gamma }}(\eta )=\mathrm{exp}(\eta (X)).$$
(19)
It is called the quasi-Lebesgue measure.
Theorem 3 implies
###### Theorem 6
For each $`a`$, the quasi-Lebesgue measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ is quasi-invariant under $`M_a`$ with a constant density
$$\frac{dM_a(\stackrel{~}{P}_\mathrm{\Gamma })}{d\stackrel{~}{P}_\mathrm{\Gamma }}=\mathrm{exp}\left(_X\mathrm{log}a(x)𝑑\nu (x)\right).$$
###### Corollary 1
If $`_X\mathrm{log}a(x)𝑑\nu (x)=0`$, then the quasi-Lebesgue measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ is invariant with respect to $`M_a`$.
We see that the measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ is invariant with respect to an infinite-dimensional multiplicative group whose action generalizes the action of the group of diagonal matrices with determinant $`1`$ in a finite-dimensional vector space. Thus we may consider the measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ as an infinite-dimensional analogue of the Lebesgue measure. This property was much used in for the representation theory of the group $`\mathrm{SL}(2,F)`$.
Let us consider the group of triangular matrices of order $`2`$
$$T_{a,b}=\left(\begin{array}{cc}a()& b()\\ 0& a()^1\end{array}\right)$$
with $`b,\mathrm{log}aL^1(X,m)`$ for some measurable space $`(X,m)`$ (a current group in the terminology of physicists).
###### Theorem 7
The formula
$$U(T_{a,b})F(\eta )=\mathrm{exp}\left(_X\mathrm{log}a(x)𝑑\nu (x)+i_Xa(x)b(x)𝑑\eta (x)\right)F(M_a\eta ).$$
defines a unitary irreducible representation of this group in the space $`L^2(\stackrel{~}{P}_\mathrm{\Gamma })`$.
* The representation is correctly defined and its unitarity follows from the invariance property of $`\stackrel{~}{P}_\mathrm{\Gamma }`$. The irreducibility follows from the ergodicity of the action of the group $``$ of multiplicators. $``$
Remarks. 1. This representation may be extended to the group $`\mathrm{SL}(2,F)`$, for this we need to define only one operator, namely, the image of the matrix $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$.
2. This representation of $`\mathrm{SL}(2,F)`$ was firstly introduced in in a completely different way, but later was interpreted with the space $`L^2(\widehat{P}_\mathrm{\Gamma })`$.
3. This representation is a continual analogue of the classical representations of the group $`\mathrm{SL}(2,)`$ in $`L^2(_+)`$.
Note that passing from $`P_\mathrm{\Gamma }`$ to $`\stackrel{~}{P}_\mathrm{\Gamma }`$ we do not change the conditional measures $`P_\mathrm{\Gamma }^s`$ of $`P_\mathrm{\Gamma }`$, given the conic part equal to $`s`$, and modify only the factor measure on the half-line, i.e. $`T^{}\stackrel{~}{P}_\mathrm{\Gamma }=m_\theta \times \mathrm{PD}(\theta )\times \overline{\nu }^{\mathrm{}}`$ and $`T\stackrel{~}{P}_\mathrm{\Gamma }=\stackrel{~}{\mathrm{PD}}(\theta )\times \overline{\nu }`$, where $`m_\theta `$ has density $`t^{\theta 1}/\mathrm{\Gamma }(\theta )`$, $`t>0`$ (in particular, $`m_1`$ is just the Lebesgue measure on the half-line), and $`\stackrel{~}{\mathrm{PD}}(\theta )=m_\theta \times \mathrm{PD}(\theta )`$.
Recall that the transformation $`M_a`$ induces a Markovian operator $`R_a`$ on the cone $`C`$ (see Section 6).
###### Corollary 2
The $`\sigma `$-finite measure $`\stackrel{~}{\mathrm{PD}}(\theta )`$ on the cone $`C`$ is invariant under the Markovian operator $`R_a`$ for all $`a`$.
## 8 Stable processes and general Poisson–Dirichlet measures
###### Definition 8
Let $`\alpha (0,1)`$. The standard $`\alpha `$-stable process on the space $`(X,\nu )`$ is a Lévy process with Lévy measure
$$d\mathrm{\Lambda }_\alpha =\frac{c\alpha }{\mathrm{\Gamma }(1\alpha )}s^{\alpha 1}ds,s>0,$$
(20)
where $`c>0`$ is an arbitrary fixed positive number.
The corresponding infinitely divisible law, i.e. the distribution of the sum of charges with respect to $`P_\alpha `$, is the $`\alpha `$-stable law $`F_\alpha `$ on $`_+`$.
Denote by $`P_\alpha `$ the law of the $`\alpha `$-stable process. The Laplace transform of $`P_\alpha `$ equals
$$𝔼_\alpha \left[\mathrm{exp}\left(_Xa(x)𝑑\eta (x)\right)\right]=\mathrm{exp}\left(c_Xa(x)^\alpha 𝑑\nu (x)\right),$$
(21)
for an arbitrary measurable function $`a:X_+`$ with $`_Xa(x)^\alpha 𝑑\nu (x)<\mathrm{}`$.
###### Proposition 3 ()
The simplicial part of an $`\alpha `$-stable process with $`\alpha (0,1)`$ is the Poisson–Dirichlet distribution $`\mathrm{PD}(\alpha ,0)`$.
###### Proposition 4
The conic part of the law of the $`\alpha `$-stable process is concentrated on the set
$$\{zC:\underset{n\mathrm{}}{lim}z_nn^{1/\alpha }=\left(\frac{c}{\mathrm{\Gamma }(1\alpha )}\right)^{1/\alpha }\}.$$
(22)
* In case of an $`\alpha `$-stable process, we have $`m(t)=\frac{c}{\mathrm{\Gamma }(1\alpha )}t^\alpha `$, and (22) follows immediately from Proposition 1. $``$
For $`s>0`$, denote by $`\varkappa _\alpha ^s`$ the conditional distribution of the conic part $`\varkappa _\alpha `$ of the stable process on the simplex $`\mathrm{\Sigma }_s=\{z=(z_1,z_2,\mathrm{}):z_i=s\}`$ of monotone sequences with sum $`s`$ (i.e. the conic part of the conditional distribution of the law $`P_\alpha `$ on the set $`D_s=\{\eta =z_i\delta _{x_i}D^+:z_i=s\}`$ of positive discrete measures with total charge $`s`$).
###### Corollary 3
The homothetic projection of the conditional measure $`\varkappa _\alpha ^s`$ on the unit simplex $`\mathrm{\Sigma }`$ is concentrated on the set
$$\{y\mathrm{\Sigma }:\underset{n\mathrm{}}{lim}y_nn^{1/\alpha }=\frac{1}{s}\left(\frac{c}{\mathrm{\Gamma }(1\alpha )}\right)^{1/\alpha }\}$$
(23)
###### Corollary 4
The homothetic projections of the measures $`\varkappa _\alpha ^s`$ and $`\varkappa _\alpha ^t`$ on the unit simplex $`\mathrm{\Sigma }`$ are singular for all pairs $`s,t>0`$, $`st`$. Thus the distribution $`\mathrm{PD}(\alpha ,0)`$ is a continual sum of a family of singular distributions.
The following statement shows how one may recover the conic part of the stable process starting with its simplicial part $`\mathrm{PD}(\alpha ,0)`$.
###### Corollary 5 ()
Let the vector $`Q=(Q_1,Q_2,\mathrm{})\mathrm{\Sigma }`$ have the distribution $`\mathrm{PD}(\alpha ,0)`$. The limit $`L(Q)=lim_n\mathrm{}n^{1/\alpha }Q_n`$ exists almost surely. Let $`S(Q)=\frac{1}{L(Q)}\frac{c}{\mathrm{\Gamma }(1\alpha )}`$. Then the $`\alpha `$-stable process $`\eta `$ on the space $`(X,\nu )`$ may be represented as
$$\eta =S(Q)\underset{i=1}{\overset{\mathrm{}}{}}Q_i\delta _{X_i},$$
where $`Q`$ obeys $`\mathrm{PD}(\alpha ,0)`$ and $`X_1,X_2,\mathrm{}`$ is a sequence of i.i.d. variables on $`X`$ with common distribution $`\nu `$.
The Poisson–Dirichlet distribution $`\mathrm{PD}(\alpha ,\theta )`$ with $`\alpha ,\theta 0`$ is not the simplicial part of any Lévy process. However, one may obtain it as the simplicial part of the process that has density with respect to a stable process. Namely, let $`\theta >\alpha `$ and consider the law $`P_{\alpha ,\theta }`$ on $`D`$ which has the density
$$\frac{dP_{\alpha ,\theta }}{dP_\alpha }(\eta )=\frac{c_{\alpha ,\theta }}{\eta (X)^\theta }$$
(24)
with respect to the $`\alpha `$-stable law $`P_\alpha `$. Here $`c_{\alpha ,\theta }=c^{\theta /\alpha }\frac{\mathrm{\Gamma }(\theta +1)}{\mathrm{\Gamma }(\theta /\alpha +1)}`$ is a normalizing constant.
###### Proposition 5 ()
The simplicial part of the law $`P_{\alpha ,\theta }`$ is the Poisson–Dirichlet distribution $`\mathrm{PD}(\alpha ,\theta )`$.
## 9 The gamma measure as a weak limit of laws of $`\alpha `$-stable processes when $`\alpha `$ tends to zero
The purpose of this section is to show that it is natural to consider the gamma process as a weak limit of the $`\alpha `$-stable processes when $`\alpha 0`$. We present several settings of this statement.
Let $`k>0`$ and consider the measure $`P_{\alpha ,k}`$ on $`D`$ given by
$$\frac{dP_{\alpha ,k}}{dP_\alpha }(\eta )=\frac{\mathrm{exp}\left(\gamma \eta (X)\right)}{𝔼_\alpha \left[\mathrm{exp}\left(\gamma \eta (X)\right)\right]}=e^{\gamma ^\alpha }e^{\gamma \eta (X)},$$
(25)
where $`\gamma =\frac{k}{\alpha ^{1/\alpha }}`$. Denote by $`\widehat{P}_{\alpha ,k}`$ the law of the process $`\gamma \eta `$, where $`\eta `$ obeys $`P_{\alpha ,k}`$. The following theorem was formulated in .
###### Theorem 8 ()
The measures $`\widehat{P}_{\alpha ,k}`$ converge weakly to $`P_\mathrm{\Gamma }`$ when $`\alpha 0`$.
* It follows from (25), (21) that the Laplace transform of the measure $`\widehat{P}_{\alpha ,k}`$ equals
$$𝔼_{\widehat{P}_{\alpha ,k}}\mathrm{exp}(f_a(\eta ))==\mathrm{exp}(\gamma ^\alpha _X((a(x)+1)^\alpha 1)d\nu (x)).$$
But
$$\gamma ^\alpha \left((a(x)+1)^\alpha 1\right)=\frac{k^\alpha }{\alpha }\left(\alpha \mathrm{log}(a(x)+1)+o(\alpha )\right)\mathrm{log}(a(x)+1),$$
as $`\alpha 0`$, hence
$$𝔼_{\widehat{P}_{\alpha ,k}}\mathrm{exp}(f_a(\eta ))\mathrm{exp}\left(_X\mathrm{log}(a(x)+1)𝑑\nu (x)\right)=𝔼_{P_\mathrm{\Gamma }}\mathrm{exp}(f_a(\eta )),$$
and Theorem 7 follows. $``$
This important result is a key point of the construction of two-parameter Poisson–Dirichlet distributions $`\mathrm{PD}(\alpha ,\theta )`$ (see Section 4). In particular, one obtains the following corollary.
###### Corollary 6 ()
For a fixed $`\theta 0`$, the distributions $`\mathrm{PD}(\alpha ,\theta )`$ converge to $`\mathrm{PD}(0,\theta )=\mathrm{PD}(\theta )`$ when $`\alpha 0`$.
###### Corollary 7
Let $`\stackrel{~}{P}_{\alpha ,k}`$ be a measure with constant density $`e^{\gamma ^\alpha }`$ with respect to the $`\alpha `$-stable law $`P_\alpha `$. Denote by $`\stackrel{ˇ}{P}_{\alpha ,k}`$ the law of the process $`\gamma \eta `$, where $`\eta `$ obeys $`\stackrel{~}{P}_{\alpha ,k}`$. Then the measures $`\stackrel{ˇ}{P}_{\alpha ,k}`$ converge weakly to the quasi-Lebesgue measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$.
That the quasi-Lebesgue measure is a kind of the limit case of stable processes was suggested in . See Section 10 for another understanding of this statement.
There is the following useful way to formalize the transition to a $`\sigma `$-finite limit. Consider the convolution of $`n`$ copies of the measures $`P_{1/n}`$ and multiply it by a function of $`n`$. Then the limit in $`n`$ will be the $`\sigma `$-finite measure under consideration. In terms of the Laplace transform this procedure is equivalent to the following relation:
$$\underset{n\mathrm{}}{lim}\mathrm{exp}(n(x^{1/n}1))=\frac{d}{d\alpha }\mathrm{exp}(x^\alpha )|{}_{\alpha =0}{}^{}=\frac{1}{x}.$$
Differentiating by $`\alpha `$ at the point $`\alpha =0`$ is just a “commutative” analogue of the main technique used in for constructing so-called canonical states and the representation theory of semi-simple currents for the group $`\mathrm{SL}(2,R)`$.
## 10 Equivalent definition of stable laws, $`\sigma `$–finite stable measures and zero–stable laws
In this section we give another definition of stable laws which is valid for $`\sigma `$-finite measures. We describe a piece of theory (to be exposed in details elsewhere) of $`\sigma `$-finite stable measures showing that it is natural to consider the Lebesgue measure as a zero-stable law.
Let $`F`$ be a distribution on $``$. Consider the distribution $`F\times F`$ on $`\times `$. Let $`_\alpha `$ be the $`\alpha `$-norm in the space $`(\times )^{}`$ of linear functionals on $`\times `$, i.e. if $`f(x_1,x_2)=a_1x_1+a_2x_2`$, then $`f_\alpha =(|a_1|^\alpha +|a_2|^\alpha )^{1/\alpha }`$. Let us consider only stable laws depending on one parameter $`\alpha (0,2]`$. Then the ordinary definition of an $`\alpha `$-stable law on $``$ is equivalent to the following one.
###### Definition 9
The law $`F`$ is $`\alpha `$-stable, if the following condition holds. If two linear functionals $`f_1`$ and $`f_2`$ on $`\times `$ have the same $`\alpha `$-norm, i.e. $`f_1_\alpha =f_2_\alpha `$, then $`f_1`$ and $`f_2`$ have the same distribution with respect to $`F\times F`$.
Note that in case of linear functionals $`f_1`$ and $`f_2`$ the equality of distributions is equivalent to the existence of a $`F\times F`$-preserving transformation $`L:\times \times `$ such that $`f_2=f_1L`$. Thus we obtain the following definition of a stable law which applies to $`\sigma `$-finite measures.
###### Definition 10
The measure $`F`$ (may be $`\sigma `$-finite) is called $`\alpha `$-stable, if the following condition holds. If two linear functionals $`f_1`$ and $`f_2`$ on $`\times `$ have the same $`\alpha `$-norm, i.e. $`f_1_\alpha =f_2_\alpha `$, then there exists a $`F\times F`$-preserving transformation $`L:\times \times `$ such that $`f_2=f_1L`$.
Let $`a_1a_20`$. If $`\alpha 0`$, then
$`2^{1/\alpha }f_\alpha `$ $`=`$ $`\left({\displaystyle \frac{|a_1|^\alpha +|a_2|^\alpha }{2}}\right)^{1/\alpha }=\left({\displaystyle \frac{1}{2}}|a_1|^\alpha +{\displaystyle \frac{1}{2}}|a_2|^\alpha \right)^{1/\alpha }=`$
$`=`$ $`\left({\displaystyle \frac{1}{2}}+\alpha \mathrm{log}|a_1|+O(\alpha ^2)+{\displaystyle \frac{1}{2}}+\alpha \mathrm{log}|a_2|+O(\alpha ^2)\right)^{1/\alpha }=`$
$`=`$ $`\left(1+\alpha \mathrm{log}|a_1a_2|+O(\alpha ^2)\right)^{1/\alpha }\mathrm{log}|a_1a_2|,`$
thus it is natural to consider the quasi-norm $`f_0=|a_1a_2|`$ as a natural limit of $`\alpha `$-norms when $`\alpha `$ tends to $`0`$. We obtain the following definition of a zero-stable measure on $``$.
###### Definition 11
The measure $`F`$ on $``$ is called zero-stable, if for each two linear functionals $`f_1`$ and $`f_2`$ on $`\times `$ with $`f_1_0=f_2_0<\mathrm{}`$, there exists a transformation $`L:\times \times `$ preserving $`F\times F`$ such that $`f_2=f_1L`$.
Remark. Note that in case $`\alpha =2`$, corresponding to the normal law $`F`$, the $`F\times F`$-preserving transformation $`L`$, which connects functionals of equal 2-norm, is a rotation. In the “opposite” case $`\alpha =0`$ the corresponding transformation is hyperbolic, of the form $`(a_1,a_2)(ca_1,\frac{a_2}{c})`$. In both cases we have a linear mapping. Unlike these two extreme cases, in intermediate cases $`\alpha (0,2)`$ the transformation $`L`$ is not linear.
###### Proposition 6
For any $`\beta <1`$, the measure on $``$ with density $`|x|^\beta dx`$ is zero-stable. In particular, the Lebesgue measure on $``$ is zero-stable.
* Easy calculation. $``$
Now we may construct a theory of $`\sigma `$-finite stable processes on arbitrary spaces. Definition 8 of an $`\alpha `$-stable process on the space $`X`$ is equivalent to the following one. We recall that each bounded Borel function $`a`$ on $`X`$ defines a linear functional $`f_a`$ on the space $`D`$ of finite discrete measures on $`X`$ by $`f_a(\eta )=_Xa(x)𝑑\eta (x)`$. Let $`_\alpha `$ denote the $`\alpha `$-norm $`a_\alpha =(|a(x)|^\alpha 𝑑\nu (x))^{1/\alpha }`$, and let $`a_0=\mathrm{exp}\left(_X\mathrm{log}|a(x)|d\nu (x)\right)`$.
###### Definition 12
The law $`P_\alpha `$ on $`D`$ is $`\alpha `$-stable, where $`\alpha [0,2]`$, if for each two linear functionals $`f_{a_1}`$ and $`f_{a_2}`$ with $`a_1_\alpha =a_2_\alpha <\mathrm{}`$, there exists a transformation $`L:DD`$ preserving the law $`P_\alpha `$ such that $`f_{a_2}=f_{a_1}L`$.
It is easy to check that the quasi-Lebesgue measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ satisfies this condition. Indeed, if $`a_1_0=a_2_0`$, then $`_X\mathrm{log}(a_2/a_1)(x)𝑑\nu (x)=0`$. Hence the multiplicator $`M_{a_2/a_1}`$ preserves the measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ by Corollary 1, and it is obvious that $`f_{a_2}=f_{a_1}M_{a_2/a_1}`$. Thus we obtain the following proposition.
###### Proposition 7
The quasi-Lebesgue measure $`\stackrel{~}{P}_\mathrm{\Gamma }`$ is zero-stable.
Note that formulae (13) and (19) imply that for all $`a`$
$`𝔼_{\stackrel{~}{P}_\mathrm{\Gamma }}\left[\mathrm{exp}\left({\displaystyle _X}a(x)𝑑\eta (x)\right)\right]=\mathrm{exp}\left({\displaystyle _X}\mathrm{log}a(x)𝑑\nu (x)\right)`$
$`=\mathrm{exp}\left({\displaystyle _X}\mathrm{log}\varphi (a(x))𝑑\nu (x)\right),`$
where $`\varphi (t)=1/t`$ is the Laplace transform of the Lebesgue measure on $`_+`$ which is zero-stable. Comparing with (5), we see that we could define a zero-stable process as a Lévy process with Lévy measure corresponding to a zero-stable law.
## 11 Distributions of linear functionals of the gamma processes and the Markov–Krein transform
In this section we show that the Markov–Krein identity known in the context of Dirichlet processes may be interpreted as a formula relating the distribution of a linear functional with respect to the gamma process and the distribution of the same functional with respect to the normalized gamma process. This interpretation allows to prove it immediately, using only a formula for the Laplace transform of the gamma process.
Given a function $`a`$, denote by $`\mu _a`$ the distribution of the linear functional $`\eta f_a(\eta )=_Xa(x)𝑑\eta (x)`$ on $`D`$ with respect to the law $`\overline{P}_\mathrm{\Gamma }`$ of the normalized gamma process, and let $`\nu _a`$ be the distribution of the function $`a`$ with respect to the (normalized) parameter measure $`\nu `$. The following property is characteristic for the gamma process.
###### Theorem 9
The measures $`\mu _a`$ and $`\nu _a`$ are related by the following integral identity,
$$_{}\frac{1}{(1+zu)^\theta }d\mu _a(u)=\mathrm{exp}(_X\mathrm{log}(1+zu)^\theta d\nu _a(u)).$$
(26)
This formula was first obtained in in the context of Dirichlet processes by hard analytic arguments (see also simpler proofs in and ). But the relation with the gamma process, which is a key point of our simple proof, has been overlooked. Note that the left-hand side of identity (26) is the generalized Cauchy–Stieltjes transform of the distribution $`\mu _a`$. It is natural to call the right-hand side the multiplicative version of the generalized Cauchy–Stieltjes transform of the distribution $`\nu _a`$. In view of (13), it is equal to the Laplace transform of the gamma measure $`P_\mathrm{\Gamma }`$ calculated on the function $`a`$. Thus one may regard formula (26) as relating an integral transform (namely, the Cauchy–Stieltjes transform) of the distribution $`\mu _a`$ of the functional $`f_a`$ with respect to the normalized gamma process and an integral transform (namely, the Laplace transform) of its distribution with respect to the non-normalized gamma process.
In case of $`\theta =1`$, the identity (26) means that the distribution $`\mu _a`$ is the Markov–Krein transform of the measure $`\nu _a`$. This transform arises in many contexts, such as the Markov moment problem, continued fractions theory, exponential representations of functions of negative imaginary type, the Plancherel growth of Young diagrams, etc. (see for a detailed survey).
* Using (13), Lemma 2 and the Fubini theorem we obtain that the right-hand side of (26) equals
$`\mathrm{exp}\left({\displaystyle _X}\mathrm{log}(1+za(x))𝑑\nu (x)\right)=𝔼_{P_\mathrm{\Gamma }}\left[\mathrm{exp}\left(z{\displaystyle _X}a(x)𝑑\gamma (x)\right)\right]`$
$`=`$ $`𝔼_{P_\mathrm{\Gamma }}\left[\mathrm{exp}\left(z\gamma (X){\displaystyle _X}a(x)𝑑\overline{\gamma }(x)\right)\right]`$
$`=`$ $`𝔼_{P_{\overline{\mathrm{\Gamma }}}}\left[{\displaystyle \frac{1}{\mathrm{\Gamma }(\theta )}}{\displaystyle _0^{\mathrm{}}}t^{\theta 1}\mathrm{exp}\left(tzt{\displaystyle _X}a(x)𝑑\overline{\gamma }(x)\right)\right]`$
$`=`$ $`𝔼_{P_{\overline{\mathrm{\Gamma }}}}\left[{\displaystyle \frac{1}{(1+z_Xa(x)𝑑\overline{\gamma }(x))^\theta }}\right],`$
and Theorem follows. $``$
Remarks. 1. According to a personal communication of P. Diaconis, the idea of proving formula (26) using the Laplace transform formula for the gamma process was used by F. Huffer in case of discrete parameter measure $`\nu `$ (in this case the gamma process is just a sum of independent gamma variables and the normalized gamma process is a random point of a finite-dimensional simplex obeying a Dirichlet distribution). But the fact that this argument works for continuous parameter measures as well, which simplifies the proof in essential manner, seems to have been overlooked.
2. It follows from the known results on the Markov–Krein transform that the distribution $`\mu _a`$ of a linear functional $`f_a`$ is absolutely continuous (see for an explicit formula for the density).
3. See for similar results on the common distributions of several linear functionals of the Dirichlet process. It is easy to extend the proof of Theorem 9 to obtain these results.
## 12 The two-parameter generalization of the Markov–Krein formula
In an analogue of the Markov–Krein identity is obtained for the distribution of a linear functional with respect to the generalized Dirichlet process associated with the two-parameter Poisson–Dirichlet distributions $`\mathrm{PD}(\alpha ,\theta )`$. Using the key idea of Section 11, we present here a new proof of this identity based on relation of the two-parameter Poisson–Dirichlet family to the stable processes.
Let $`\alpha (0,1)`$ and $`\theta >\alpha `$. Denote by $`\overline{P}_{\alpha ,\theta }`$ the normalization (in sense of (12)) of the law $`P_{\alpha ,\theta }`$. (Recall that the simplicial part of this law is $`\mathrm{PD}(\alpha ,\theta )`$, see Section 8.) Given an arbitrary measurable function $`a:X_+`$ such that $`_Xa(x)^\alpha 𝑑\nu (x)<\mathrm{}`$, let $`\mu _a`$ be the distribution of the functional $`f_a`$ with respect to $`\overline{P}_{\alpha ,\theta }`$, and let $`\nu _a`$ be the distribution of the function $`a`$ with respect to the (normalized) measure $`\overline{\nu }`$.
###### Theorem 10
The measures $`\mu _a`$ and $`\nu _a`$ are related by the following integral identity:
1) if $`\theta 0`$,
$$\left(_{}(1+zu)^\theta 𝑑\mu _a(u)\right)^{\frac{1}{\theta }}=\left(_{}(1+zu)^\alpha 𝑑\nu _a(u)\right)^{\frac{1}{\alpha }};$$
(27)
2) if $`\theta =0`$,
$$\mathrm{exp}\left(_{}\mathrm{log}(1+zu)^\alpha d\mu _a(u)\right)=_{}(1+zu)^\alpha d\nu _a(u).$$
(28)
* 1) Denote the left-hand side of the desired identity by $`A^{1/\theta }`$ and the right-hand side by $`B^{1/\alpha }`$. Using the identity
$$\frac{1}{r^\theta }=\frac{1}{\mathrm{\Gamma }(\theta )}_0^{\mathrm{}}t^{\theta 1}e^{rt}𝑑t,$$
we obtain
$`A`$ $`=`$ $`c_{\alpha ,\theta }𝔼^\alpha \left[\left(\eta (X)+z{\displaystyle _X}a(x)𝑑\eta (x)\right)^\theta \right]`$
$`=`$ $`{\displaystyle \frac{c_{\alpha ,\theta }}{\mathrm{\Gamma }(\theta )}}𝔼^\alpha \left[{\displaystyle _0^{\mathrm{}}}t^{\theta 1}\mathrm{exp}\left(t\left(\eta (X)+z{\displaystyle _X}a(x)𝑑\eta (x)\right)\right)𝑑t\right]`$
$`=`$ $`{\displaystyle \frac{c_{\alpha ,\theta }}{\mathrm{\Gamma }(\theta )}}{\displaystyle _0^{\mathrm{}}}t^{\theta 1}𝔼^\alpha \left[\mathrm{exp}\left({\displaystyle _X}t(1+za(x))𝑑\eta (x)\right)\right]𝑑t.`$
By the Laplace transform formula (21), the expectation equals precisely $`e^{t^\alpha B}`$, thus
$$A=\frac{c_{\alpha ,\theta }}{\mathrm{\Gamma }(\theta )}_0^{\mathrm{}}t^{\theta 1}e^{t^\alpha B}𝑑t,$$
and (27) follows by changing variables.
2) Follows from (27) by letting $`\theta 0`$. $``$
## Appendix
This is a continuation of the discussion on properties of the gamma process initiated in the Introduction. We will start with a new statement which is called the subordinating property:
There is the identity in law
$$(\beta _{\gamma _t},t0)\stackrel{\mathrm{law}}{=}\left(\frac{1}{\sqrt{2}}(\gamma _t^{(1)}\gamma _t^{(2)}),t0\right),$$
where the left-hand side, $`(\beta _u,u0)`$, is a Brownian motion independent from a gamma process $`(\gamma _t,t0)`$, and in the right-hand side $`\gamma ^{(1)}`$ and $`\gamma ^{(2)}`$ are two independent gamma processes.
Combining the quasi-invariance property 3 (see Introduction) and the subordinating property 6 of the gamma process, one easily deduces the following proposition, which expresses the conditional law of $`(\gamma _t,t0)`$ given $`(X_t\beta _{\gamma _t},t0)`$, where $`(\beta _u,u0)`$ is a Brownian motion independent on $`(\gamma _t,t0)`$.
###### Proposition 8 ()
Conditionally on X, the process $`(\gamma _t,t0)`$ is distributed as $`\frac{1}{2}T_{V_t}^{(1)}`$, $`t0`$, where $`(T_a^{(1)},a0)`$ is the process of the first hitting times of levels $`a0`$ of a Brownian motion with drift 1 independent of $`X`$, and $`V_t=_0^t|dX_s|`$ is a gamma process taken at $`(2t)`$.
The above Proposition suggests the existence of some deep links between the Brownian motion and the gamma process. We now present two occurrences of the gamma process related to the Brownian motion and to the Bessel processes.
###### Theorem 11 ()
Let $`(C_t,t0)`$ denote the Cauchy process, and for each $`x`$ define $`\mu _x=sup\{s1:C_s+xs=\mathrm{max}_{u1}(C_u+xu)\}`$. Then
$$(\mu _x,x)\stackrel{\mathrm{law}}{=}(\gamma _{\rho (x)}/\gamma _1,x),$$
where $`(\gamma _u,u0)`$ is the standard gamma process, and $`\rho (x)=\frac{1}{2}+\frac{1}{\pi }arctg(x)P(C_1x)`$.
See also \[1, Proposition 1, p. 1535\] for some occurrences of $`(\gamma _u/\gamma _1,u1)`$ in relation with embedded regenerative sets.
To present the second instance we have in mind, we recall that the law $`Q_0^\delta `$ on $`C(_+,_+)`$ of the square of a Bessel process of (fractional) dimension $`\delta `$, i.e. of $`_+`$-valued diffusion with infinitesimal generator $`2x\frac{d^2}{dx^2}+\delta \frac{d}{dx}`$, is infinitely divisible, more precisely, $`Q_0^\delta Q_0^\delta ^{}=Q_0^{\delta +\delta ^{}}`$ for every $`\delta ,\delta ^{}0`$.
In the derivation of the Lévy–Khinchin representation of $`Q_0^\delta `$ in was shown the existence of a two-parameter process $`(X_t^\delta ,\delta 0,t0)`$ which may be described as follows:
* it has homogeneous independent increments in $`\delta `$, with values in the space $`C(_+,_+)`$;
* for each $`\delta >0`$, $`(X_t^\delta ,t0)`$ obeys the law $`Q_0^\delta `$;
* for each $`t>0`$, $`(X_t^\delta ,\delta 0)`$ is distributed as $`(2t\gamma _{\delta /2},\delta 0)`$.
The limit of the Bessel processes at the critical value of the parameter is similar to the situation with stable and gamma processes when $`\alpha `$ tends to zero (see Section 9).
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# Neutrino oscillations in high energy cosmic neutrino flux
## Abstract
I discuss the effects of neutrino oscillations on high energy cosmic neutrinos which come from cosmologically distant astrophysical sources. I incorporate all the up-to-date constraints from the solar, atmospheric, reactor, accelerator data and give the possible pattern for the ratio of the high energy cosmic neutrinos in the cases of three and four neutrino schemes.
## 1. Introduction
A lot of attention has been focused on high energy neutrinos ($`E10^6`$ GeV) which come from cosmologically distant astrophysical sources such as Active Galactic Nuclei and Gamma Ray Burst fireballs (typical distance is 100 Mpc), since they can be distinguished from atmospheric neutrinos for such high energies and identification of flavors of neutrinos may be possible in new km<sup>2</sup> surface area neutrino telescopes . The effects of neutrino oscillations on the high energy cosmic neutrinos have been discussed in the past , and the purpose of this talk is to update the analysis by taking into account all the constraints from the solar, atmospheric, reactor and accelerator data in the three or four neutrino framework (This talk is based on the work ).
## 2. Analysis
Since the path length of neutrinos is much larger than any possible neutrino oscillation length which is suggested from the solar, atmospheric or LSND data for the energy $`E_\nu `$ 10<sup>6</sup> GeV, I will average over rapid oscillations throughout this talk. Then the oscillation probability is given by
$`P(\nu _\alpha \nu _\beta ;L=\mathrm{})=\delta _{\alpha \beta }{\displaystyle \underset{jk}{}}U_{\alpha j}^{}U_{\beta j}U_{\alpha k}U_{\beta k}^{}`$ $`=`$ $`{\displaystyle \underset{j}{}}|U_{\alpha j}|^2|U_{\beta j}|^2,`$ (1)
where $`U_{\alpha k}`$ stands for the MNS mixing matrix and I have ignored odd functions in $`\mathrm{sin}(\mathrm{\Delta }m_{ij}^2L/4E)`$ which oscillate rapidly as $`L\mathrm{}`$.
The electron and muon neutrinos are mainly produced in the decay chain of charged pions whereas the tau neutrinos are mainly produced in the decay chain of charmed mesons at a suppressed level . The ratio of the intrinsic high energy cosmic neutrinos flux is typically $`F^0(\nu _e)`$ : $`F^0(\nu _\mu )`$ : $`F^0(\nu _\tau )=\mathrm{\hspace{0.17em}1}`$ : $`2`$ : $`<10^5`$. For simplicity I assume that the ratio is $`F^0(\nu _e)`$ : $`F^0(\nu _\mu )`$ : $`F^0(\nu _\tau )=\mathrm{\hspace{0.17em}1}`$ : $`2`$ : 0. Thus the ratio of flux of neutrinos in the far distance is given by
$`\left(\begin{array}{c}F(\nu _e)\\ F(\nu _\mu )\\ F(\nu _\tau )\end{array}\right)=P\left(\begin{array}{c}F^0(\nu _e)\\ F^0(\nu _\mu )\\ F^0(\nu _\tau )\end{array}\right)=P\left(\begin{array}{c}1\\ 2\\ 0\end{array}\right)F^0(\nu _e),`$ (11)
where a matrix $`P`$ has components $`\left(P\right)_{\alpha \beta }=P(\nu _\alpha \nu _\beta ;L=\mathrm{})`$ (See (1)).
Since currently we do not know the precise total cosmic neutrino flux I will mainly focus my discussion on the ratio of different flavors of neutrinos. To plot the ratio of the three neutrino flavors, I introduce a triangle representation. Fig. 1 is a unit regular triangle and the position of the point gives the ratio of the high energy neutrino flux, where $`F_\alpha F(\nu _\alpha )`$ is given by (11).
In the three flavor framework, because of the constraint of the CHOOZ data and the atmospheric neutrino data of Superkamiokande and Kamiokande, it has been known that $`|U_{e3}|^2`$ is small and $`|U_{\mu j}|^2|U_{\tau j}|^2|`$ (See, e.g., ). Using the allowed region for $`|U_{\alpha j}|^2`$ ($`\alpha =e,\mu ,\tau `$) in the atmospheric neutrino analysis of and in solar neutrino analysis of , the possible ratio of the high energy neutrino flux is calculated numerically and is given in Fig. 2. The allowed region is a small area around the midpoint $`F(\nu _e)=F(\nu _\mu )=F(\nu _\tau )=1/3`$.
In the four neutrino scheme one needs in principle tetrahedron to express the ratio of the four neutrino flux, but since we do not observe the cosmic sterile neutrino I normalize the flux of each active neutrino by the total flux of active ones:
$`\left(\begin{array}{c}\stackrel{~}{F}(\nu _e)\\ \stackrel{~}{F}(\nu _\mu )\\ \stackrel{~}{F}(\nu _\tau )\end{array}\right){\displaystyle \frac{1}{F(\nu _e)+F(\nu _\mu )+F(\nu _\tau )}}\left(\begin{array}{c}F(\nu _e)\\ F(\nu _\mu )\\ F(\nu _\tau )\end{array}\right).`$ (18)
After redefining the flux this way ($`\stackrel{~}{F}F`$), we can plot the ratio of each active neutrino with the same triangle graph as in the three neutrino case.
If one demands that the number $`N_\nu `$ of effective neutrinos in Big Bang Nucleosynthesis (BBN) be less than 4, then it can be shown that the $`4\times 4`$ MNS mixing matrix splits into two $`2\times 2`$ block diagonal matrices. In this case the ratio is given by a small region depicted in Fig. 3.
On the other hand, some people give conservative bound for $`N_\nu `$, and without the BBN constraint $`N_\nu <4`$ the only restrictions come from the solar and atmospheric neutrino data. The analysis of the solar neutrino data in the four neutrino scheme with ansatz $`U_{e3}=U_{e4}=0`$ has been done recently in . The analysis of the atmospheric neutrino data in the four neutrino framework has been done in again with ansatz $`U_{e3}=U_{e4}=0`$. Using these results, the ratio of the high energy cosmic neutrinos is evaluated and is given in Fig. 4 which has much wider allowed region than any other case. This scheme may be distinguished from others if one has good precision in future experiments.
I would like to thank hospitality of Summer Institute 99 at Yamanashi, Japan where this work was started. This research was supported in part by a Grant-in-Aid for Scientific Research of the Ministry of Education, Science and Culture, #12047222, #10640280.
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# 1 Introduction
## 1 Introduction
The general form of the spectrum of the universal infrared background radiation was predicted by Primack et al. , Malkan and Stecker , Dwek et al. and Fall et al. . The origin is currently thought to be due to star formation in the early universe producing starlight predominantly in the 1 micron range which was then re-processed into the 100 micron range due to heating of surrounding dust. Early data from the DIRBE and FIRAS instruments on COBE (see ref. for a review) and from ISO were mainly upper or lower limits which did not severely constrain the models. New results from these experiments are becoming available due to a more sophisticated treatment of the galactic and solar-system radiation foreground subtraction. In particular, new determinations of the infrared background intensity have very recently become available at 1.25, 2.2 and 3.5 microns , 15.0 microns , 60 and 100 microns , and 140 and 240 microns . In Fig. 1 we show the predictions of Malkan & Stecker together with the new observations and recent lower limits based on infrared galaxy source counts and a model including the effects of an infrared burst phase of ultraluminous infrared galaxies. The observations are in excellent agreement with prediction inasmuch as the expected trend (two peaks and valley) is found, the agreement being particularly good between 1.25 and 15 microns. However, the level of the observed flux in the 60–240 micron range is higher, and in the case that this is a truly universal background it has important consequences both for star formation in the early universe, and for intergalactic absorption of TeV gamma-rays which leads to a crisis in our understanding of high energy phenomena.
Immediately following the discovery of Markarian 421 as a TeV source it became clear that interactions with a universal infrared background may have a strong influence on the propagation of TeV gamma-rays through intergalactic space , in the same way that the microwave background affects the propagation of 1000 TeV gamma-rays . The observation by several telescopes of a very high level of emission by Markarian 501 during 1997 (see ref. for a review) showed that the spectrum extended to 10 TeV and beyond. At that time, it was pointed out by Stanev and Franceschini that if the spectrum continued much beyond 10 TeV that we would “have to revise our concepts about the propagation of TeV gamma-rays in the intergalactic space, and that something complicates the process.” This has now come about as the analysis of HEGRA data for the whole of 1997 shows that the spectrum of Markarian 501 extends well beyond 10 TeV .
## 2 Implications of the new IR observations for the propagation of Gamma-rays through the universe
Photon-photon pair production is a very elementary process in QED, and can safely be calculated at energies relevant to TeV astronomy. For propagation of energetic particles (gamma-rays) of energy $`E_\gamma `$, mass $`m_\gamma `$(=0) and velocity $`\beta _\gamma c`$(=$`c)`$, through isotropic radiation the reciprocal of the mean free path for collisions with photons is given by
$$x_{\gamma \gamma }(E_\gamma )^1=\frac{1}{8E_{\gamma }^{}{}_{}{}^{2}\beta _\gamma }_{\epsilon _{\mathrm{min}}}^{\mathrm{}}𝑑\epsilon \frac{n(\epsilon )}{\epsilon ^2}_{s_{\mathrm{min}}}^{s_{\mathrm{max}}(\epsilon ,E_\gamma )}𝑑s(sm_\gamma ^2c^4)\sigma (s),$$
(1)
where $`n(\epsilon )`$ is the differential photon number density, and $`\sigma (s)`$ is the total cross section for a centre of momentum frame energy squared given by $`s=m_\gamma ^2c^4+2\epsilon E_\gamma (1\beta _\gamma \mathrm{cos}\theta )`$ where $`\theta `$ is the angle between the directions of the energetic particle (gamma-ray) and soft photon, $`s_{\mathrm{min}}=(2m_ec^2)^2`$, $`\epsilon _{\mathrm{min}}=(s_{\mathrm{min}}m_\gamma ^2c^4)/[2E_\gamma (1+\beta _\gamma )]`$, and $`s_{\mathrm{max}}(\epsilon ,E_\gamma )=m_\gamma ^2c^4+2\epsilon E_\gamma (1+\beta _\gamma )`$. For photon-photon pair production by gamma-rays, we take $`m_\gamma =0`$ and $`\beta _\gamma =1`$, and the photon-photon pair production cross section from .
We have calculated the mean free path for this process for interactions of gamma-rays with the infrared and cosmic microwave backgrounds, and this is shown in Fig. 2. Also shown are results based on the upper and lower models for the IR background due to Stecker and Malkan , extended to shorter wavelengths. We note that above 10 TeV the mean free path using the new IR intensity (thick solid curve) is almost an order of magnitude lower, and this has very important consequences for the HEGRA observations of Markarian 501.
We show in Fig. 3 the time-averaged flux from Markarian 501 observed by HEGRA using the stereoscopic system of telescopes . We note that the spectrum smoothly extends to at least 20 TeV. This observation was confirmed by using the stand-alone telescopes of HEGRA, although with less statistical significance and with rather different systematics . It is important to note that the Crab Nebula, as a steady galactic source, has approximately the same emission level as Markarian 501 in the high state, and has been continuously monitored by the HEGRA telescopes . The Crab spectrum is also shown in Fig. 3 for comparison (shifted down by a factor of 100 for clarity). The Crab Nebula spectrum determined extends again up to 20 TeV as a straight power-law with a very similar spectral index, and this suggests reliable detection of photons up to 20 TeV. For the following discussion, an understanding of the absolute energy calibration is of considerable importance. This was investigated in great detail by Monte Carlo methods, and in addition, checked by absolute measurement of the cosmic ray proton flux using the HEGRA telescope system in about the same energy range as in the Crab and Markarian 501 observations. In this energy range the cosmic ray proton flux is rather well measured by direct methods, and a compilation of data can be found in ref. . The absolute energy scale of HEGRA was then estimated to have an uncertainty of 15% . In addition, there is data available on Markarian 501 from the Whipple and CAT telescopes, and this is in very good agreement with the HEGRA data but does not reach as high in energy, mainly because of their lower exposure.
We have applied the mean free path in the infrared background radiation plotted in Fig. 2 to “correct” the observed flux from Markarian 501 assuming a Hubble constant of $`H_0=65`$ km s<sup>-1</sup> Mpc<sup>-1</sup>, and this is also shown in Fig. 3. We note that this implies a dramatic (by several orders of magnitude) increase in luminosity at the source (see right hand scale). Even taking into account relativistic beaming, with a typical Doppler factor of 20, this would imply an extraordinarily high luminosity for an active galaxy of the BL Lac class of which Markarian 501 is a member. The large correction factor is a direct consequence of the new IR data at 60 and 100 microns and the implied shift of the “dust peak” to somewhat shorter wavelengths. It is worthwhile pointing out that the correction factor depends on the value of the Hubble constant. If it were larger than that we have assumed, then the distance to the source would be smaller, and the absorption less, but unacceptably high values of $`H_0`$ would be required to give a reasonable source luminosity. In fact, the problem may be even more severe if future infrared measurements would fill in the “valley” between 5 and 40 microns giving an IR background such as shown by the short dashed lines in Figure 1. In this case, the mean free path would be as shown by the short dashed curves in Figure 2, with the result that correction for absorption would cause an even more dramatic rise of the source spectrum above a few TeV.
## 3 Immediate consequences
The source spectrum of Markarian 501, i.e. the observed spectrum corrected for propagation from the source, is in gross disagreement with expectations of the Synchrotron self-Compton (SSC) model routinely used to interpret the observed spectra of active galaxies. In this model, accelerated electrons produce most of the observed radio through X-ray emission as synchrotron radiation, and the same population Compton scatters synchrotron photons to gamma-ray energies. The apparent success of the SSC model depends critically on which intergalactic absorption is used, as the observed sources are at the gamma-ray horizon . Now that the most recent determinations of the diffuse IR background are high, this makes it impossible to fit the gamma-ray observations using the SSC model. For example, Guy et al. took a relatively high level for the IR background, but still a factor of 2 short of recent determination. Petry et al. and Bednarek and Protheroe used the IR background at a level of that predicted by Malkan & Stecker , while Sambruna et al. assume a specific shape of the IR background in order to keep the TeV data compatible with the SSC model.
The corrected TeV data (Fig. 3) have a very prominent upturn above 10 TeV which cannot be explained by SSC models. Such an upturn may occur naturally in some proton blazar models , particularly those which predict a very flat or nearly monoenergetic spectrum of synchrotron radiating protons . However, even with these models it appears impossible to obtain the extreme turn-up in the TeV source spectrum inferred if the new IR data are used to correct for propagation.
In principle the turn up could be a pile-up caused in a pair–Compton cascade during propagation through the IR background . However, for this to be the cause the intergalactic magnetic field would have to be extremely small otherwise the pile-up would be distributed as a halo , and any time-variability, such as the observed flaring on time scales of hours would be washed out.
## 4 Discussion: What is the solution?
Many possible solutions have been mentioned by Finkbeiner et al. , and here we provide a quantitative discussion of some of them. The simplest solution to the problem would be if the recent determinations of the universal infrared background were an overestimate. This is not completely ruled out, and they are not incompatible with star formation and dust-heating in the early Universe. In a similar manner, a shift in the energy scale of the HEGRA data downwards by a factor of two would have a similar effect. This appears rather extreme, and not very likely. If both the infrared data and the HEGRA TeV data are confirmed, then we may be forced to consider more radical possibilities.
An interesting idea based on known physics was proposed by Harwit, Protheroe and Biermann . It considers a TeV Cherenkov event as being due to a coherent superposition of a number of lower energy photons simultaneously arriving at the top of the atmosphere, and masquerading as a single TeV photon with the sum of the energies. Such Bose-Einstein condensates would suffer losses on a distance scale comparable to that of the mean free path in the infrared radiation of the individual photons of the condensate. This idea can be tested on the basis of extensive simulations, and by detailed comparison of Cherenkov images of gamma ray showers from the Crab nebula and Markarian 501. Such analysis is underway. <sup>1</sup><sup>1</sup>1After submission of this paper a preprint by the HEGRA Collaboration appeared (astro-ph/0006092) which renders the possibility of a Bose-Einstein condensate explanation as being rather unlikely. Even for a Bose-Einstein condensate with, on average, as few as two arriving photons, the effect on the intergalactic absorption of 20 TeV photons from Markarian 501 would be dramatic. We show in Fig. 4 the result of correcting for intergalactic absorption for this case in an approximate way (plotted as asterisks). In calculating this, we assumed for simplicity that each event recorded by the telescope consisted of a Bose-Einstein condensate of two photons, and we then calculated the average number $`N_0`$ of photons of energy $`E_\gamma /2`$ that would have been emitted for each pair of photons arriving with total energy $`E_\gamma `$. The energy of the emitted bunch was then on average $`N_0E_\gamma /2`$, and the correction factor used was therefore $`N_0/2`$. Because of the change in occupation number, and hence energy, of the bunch of photons during propagation, the precise meaning of the corrected spectrum ($`E_\gamma ^2F(E_\gamma )`$) is ambiguous, and caution should be applied when interpreting it. Nevertheless, integrating the corrected luminosity ($`E_\gamma ^2L(E_\gamma )`$) over $`\mathrm{ln}(E_\gamma )`$ will give the correct total emitted power. Thus, we note that this process is a viable mechanism for solving the problem, and has not yet been ruled out by the observations.
Another possibility which has been suggested is violation of Lorentz invariance . Solutions of that kind have also been invoked to explain the observed high energy particles above the GZK cut-off. Amelino-Camelia et al. suggested that the velocity of propagation of energetic particles is modified in vacuo due to microscopic quantum fluctuations occurring on scales of the order of the Planck length. Their dispersion relation was used by Kifune to show that intergalactic space will be much more transparent to TeV photons. We follow his treatment, in which energetic photons have velocity $`\beta _\gamma c=(1\xi _\gamma E_\gamma /E_0)c`$, and momentum $`p`$ given by $`p^2c^2=(E^2+\xi _\gamma E_\gamma ^3/E_0)`$ from which we can define an effective mass $`m_\gamma `$ for high energy photons by $`m_\gamma ^2c^4=\xi _\gamma E_\gamma ^3/E_0`$. Taking $`\xi _\gamma =1`$ and $`E_0=1.2\times 10^{19}`$ GeV (Planck mass) , and using $`\beta _\gamma `$ and $`m_\gamma ^2`$ in Equation 1, we get the result shown by the thick chain line in Figure 2. One should be cautious, however, as Equation 1 was derived assuming $`s`$ to be Lorentz invariant. Nevertheless, our result gives an indication of the likely effect of Lorentz Invariance violation. Applying our result to correct the HEGRA 1997 Markarian 501 data we obtain the spectrum shown by the diamonds in Fig. 4, and find that indeed the absorption in intergalactic space is sufficiently reduced to solve the problem, resulting in a spectrum that is largely consistent with both leptonic and hadronic models.
However, it is remarkable that an acceptable source spectrum for Markarian 501 is obtained simply by choosing the Planck mass as the energy scale for violation of Lorentz Invariance. Limits for this energy scale have been obtained previously by looking for time delays in variable sources such as Gamma Ray Bursts and Markarian 421 . By observing instead the energy spectrum of distant sources like Markarian 501 above 10 TeV, we may also have a sensitive tool for exploring the energy scale of Lorentz Invariance violation.
New data on Markarian 501 and Markarian 421 relevant to the issue discussed in this paper are soon to come as both sources are observed to be in a high flaring state in the ongoing observation period (May 2000). Detecting TeV gamma-rays from a more distant BL Lac object (e.g. $`z`$ =0.03–0.09) would provide for a much needed additional constraint on the phenomenon observed at Markarian 501.
Acknowledgements
The research of RJP is funded by the Australian Research Council. HM thanks the University of Adelaide for hospitality while this work was carried out.
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# LOCAL AUTOMORPHISMS OF THE UNITARY GROUP AND THE GENERAL LINEAR GROUP ON A HILBERT SPACE
## 1 Introduction
The study of automorphism groups of algebraic structures is of great importance in every field of mathematics. In a series of papers (see and the references therein) we investigated these groups from the point of view of how they are determined by their local actions. Our investigations were motivated by the paper of Kadison on local derivations and by a problem of Larson in initiating the study of local automorphisms of Banach algebras. The structures that we treated so far were mainly $`C^{}`$-algebras and we considered the following question: When is it true that any local automorphism, that is, any linear transformation which pointwise equals an automorphism (this automorphism may, of course, differ from point to point) is an automorphism?
It is easy to see that if we drop the assumption of the linearity of the transformations in question, then the corresponding statements are no longer true. However, if instead of linearity plus locality we assume the so-called 2-locality, then we can obtain positive results (for the first such result see ). 2-locality means that our transformation (linearity is not assumed any more) is supposed to be equal to an automorphism at every pair of points. Notice that in this way we arrive at a question that can be raised in any algebraic structure. For example, this observation motivated us to consider the problem for the orthomodular poset of all projections on a Hilbert space whose structure plays a fundamental role in the mathematical foundations of quantum mechanics (see ).
In the present paper we study the analogous problem in the case of two important groups appearing in pure algebra and in the theory of operator algebras. They are the unitary group and the general linear group.
We begin with some notation and definitions. Throughout the paper $`H`$ denotes a complex infinite-dimensional separable Hilbert space. By $`B(H)`$, $`U(H)`$, $`GL(H)`$, and $`A(H)`$ we denote the algebra of all bounded linear operators on $`H`$, the multiplicative group of all unitary operators on $`H`$, the general linear group on $`H`$ consisting of all bounded invertible linear operators on $`H`$, and the set of all additive maps on $`H`$ (that is, the set of all maps $`A:HH`$ satisfying $`A(x+y)=Ax+Ay`$ $`(x,yH)`$), respectively. For $`TB(H)`$ we denote its spectrum by $`\sigma (T)`$.
A mapping $`\varphi :GL(H)GL(H)`$ is called a 2-local automorphism of the general linear group if for every $`X,YGL(H)`$ there is an automorphism $`\varphi _{X,Y}`$ of the group $`GL(H)`$, depending on $`X`$ and $`Y`$, such that $`\varphi _{X,Y}(X)=\varphi (X)`$ and $`\varphi _{X,Y}(Y)=\varphi (Y)`$. In the case of the unitary group the situation is somewhat different. Because of certain reasons (see, for example, Theorem 2.1), $`U(H)`$ is considered here as a topological group and by an automorphism of $`U(H)`$ we mean a uniformly continuous group-automorphism. Now, not surprisingly, a mapping $`\varphi :U(H)U(H)`$ is called a 2-local automorphism of the unitary group if for every $`X,YU(H)`$ there is an automorphism $`\varphi _{X,Y}`$ of $`U(H)`$ (in the above sense), such that $`\varphi _{X,Y}(X)=\varphi (X)`$ and $`\varphi _{X,Y}(Y)=\varphi (Y)`$.
## 2 Local automorphisms of the unitary group
It was proved in that any uniformly continuous group isomorphism between the unitary groups of two $`AW^{}`$-factors is implemented by a linear or conjugate-linear \*-isomorphism of the factors themselves. As a particular case, concerning $`B(H)`$ we have the following result.
###### Theorem 2.1
(Sakai) Let $`\varphi :U(H)U(H)`$ be a uniformly continuous automorphism. Then there exists either a unitary or an antiunitary operator $`U`$ on $`H`$ such that
$$\varphi (V)=UVU^{}(VU(H)).$$
(1)
As for the 2-local automorphisms of the unitary group $`U(H)`$ we have the following statement. In the proof we use the notation $`xy`$ which denotes the operator on $`H`$ defined by
$$(xy)(z)=z,yx(zH)$$
for any $`x,yH`$.
###### Theorem 2.2
Every 2-local automorphism of $`U(H)`$ is an automorphism.
Proof. Let $`\varphi :U(H)U(H)`$ be a 2-local automorphism. For every projection $`PB(H)`$, the operator $`I2P`$ is unitary. Since $`\varphi `$ is locally of the form $`(\text{1})`$, we obtain that $`\varphi (I2P)=I2P^{}`$ for some projection $`P^{}`$. Consider the transformation
$$P(I\varphi (I2P))/2.$$
This is a 2-local automorphism of the orthomodular poset of all projections on $`H`$. By \[5, Proposition\], it is an automorphism and there is either a unitary or an antiunitary operator $`U`$ on $`H`$ such that our transformation is of the form
$$P(I\varphi (I2P))/2=UPU^{}.$$
We have $`\varphi (I2P)=U(I2P)U^{}`$ for every projection $`P`$. Transforming the original map $`\varphi `$ by this operator $`U`$, we can suppose without loss of generality that $`\varphi (I2P)=I2P`$ for every projection $`P`$. Let $`VU(H)`$ and pick an arbitrary unit vector $`xH`$. Let $`P`$ be the orthogonal projection onto the subspace spanned by $`x`$, that is, let $`P=xx`$. By the local property of $`\varphi `$ we have an either unitary or antiunitary operator $`U_{V,P}`$ such that $`\varphi (V)=U_{V,P}VU_{V,P}^{}`$ and $`\varphi (I2P)=U_{V,P}(I2P)U_{V,P}^{}`$. Since $`\varphi (I2P)=I2P`$, it follows that $`P=U_{V,P}PU_{V,P}^{}`$. We compute
$$\varphi (V)x,xxx=xx\varphi (V)xx=U_{V,P}PU_{V,P}^{}U_{V,P}VU_{V,P}^{}U_{V,P}PU_{V,P}^{}=$$
$$U_{V,P}PVPU_{V,P}^{}=U_{V,P}Vx,xxxU_{V,P}^{}.$$
Since $`U_{V,P}`$ is either linear or conjugate-linear, we have either
$$\varphi (V)x,xxx=Vx,xU_{V,P}xU_{V,P}x$$
or
$$\varphi (V)x,xxx=\overline{Vx,x}U_{V,P}xU_{V,P}x.$$
We deduce that for every $`xH`$ either
$$\varphi (V)x,x=Vx,x$$
or
$$\varphi (V)x,x=\overline{Vx,x}$$
holds true. It is rather elementary to verify (see \[5, Lemma\]) that this implies that either $`\varphi (V)=V`$ or $`\varphi (V)=V^{}`$. We show that either $`\varphi (V)=V`$ for every $`VU(H)`$ or $`\varphi (V)=V^{}`$ for every $`VU(H)`$. To see this, observe that $`\varphi (iI)`$ is either $`iI`$ or $`iI`$. First assume that $`\varphi (iI)=iI`$. We assert that in this case we have $`\varphi (V)=V`$ $`(VU(H))`$. Suppose on the contrary that there is a non-selfadjoint unitary operator $`V`$ for which $`\varphi (V)=V^{}`$. Let $`\lambda \sigma (V)`$. By the spectral theorem of normal operators we can choose a sequence $`(x_n)`$ of pairwise orthogonal unit vectors in $`H`$ such that $`Vx_n,x_n\lambda `$. We extend $`(x_n)`$ to a complete orthonormal sequence $`(x_n^{})`$ in $`H`$. Pick pairwise different complex numbers $`\lambda _n`$ of modulus 1 from the open upper half-plane and consider the unitary operator $`W=_n\lambda _nx_n^{}x_n^{}`$. By the local property of $`\varphi `$ we have an either unitary or antiunitary operator $`U_{i,W}`$ such that $`\varphi (iI)=U_{i,W}iIU_{i,W}^{}`$ and $`\varphi (W)=U_{i,W}WU_{i,W}^{}`$. Since we have supposed that $`\varphi (iI)=iI`$, it follows that $`U_{i,W}`$ is unitary. So, $`\varphi (W)=_n\lambda _nU_{i,W}x_n^{}U_{i,W}x_n^{}`$ and, on the other hand, we know that $`\varphi (W)=W`$ or $`\varphi (W)=W^{}`$. These result in $`\varphi (W)=W`$. Once again, by the local property of $`\varphi `$ we have an either unitary or antiunitary operator $`U_{W,V}`$ such that
$$\varphi (W)=U_{W,V}WU_{W,V}^{}\mathrm{and}\varphi (V)=U_{W,V}VU_{W,V}^{}.$$
Since $`\varphi (W)=W`$, it follows that $`U_{W,V}`$ is necessarily linear and from the equalities
$$\underset{n}{}\lambda _nx_n^{}x_n^{}=W=\varphi (W)=U_{W,V}WU_{W,V}^{}=\underset{n}{}\lambda _nU_{W,V}x_n^{}U_{W,V}x_n^{}$$
we conclude that $`U_{W,V}`$ is diagonizable with respect to $`(x_n^{})`$. Therefore, we can compute
$$x_n,Vx_n=V^{}x_n,x_n=\varphi (V)x_n,x_n=$$
$$U_{W,V}VU_{W,V}^{}x_n,x_n=VU_{W,V}^{}x_n,U_{W,V}^{}x_n=Vx_n,x_n.$$
If $`n`$ goes to infinity, we obtain that $`\lambda =\overline{\lambda }`$. Since $`\lambda `$ was an arbitrary element of $`\sigma (V)`$, we infer that $`V=V^{}`$ which is a contradiction. So, we have $`\varphi (V)=V`$ for every $`VU(H)`$ and this shows that $`\varphi `$ is an automorphism of $`U(H)`$. We now consider the case when $`\varphi (iI)=iI`$. Similarly as above one can verify that then we have $`\varphi (V)=V^{}`$ $`(VU(H))`$. By the local property of $`\varphi `$ it follows that for every $`V,V^{}U(H)`$ there exists an either unitary or antiunitary operator $`U`$ such that $`V^{}=UVU^{}`$ and $`V_{}^{}{}_{}{}^{}=UV^{}U^{}`$. If $`V`$ is diagonal with respect to an orthonormal basis defined in the same way as $`W`$ and $`V^{}`$ permutes the same basis, then one can easily arrive at a contradiction.
The proof of the theorem is now complete.
## 3 Local automorphisms of the general linear <br>group
We first remark that we were unable to find the description of the general form of automorphisms of $`GL(H)`$ in the literature. However, applying a result of Radjavi on a factorization of invertible operators into a product of involutions and some automatic continuity techniques it is possible to obtain the general form of the automorphisms of $`GL(H)`$ as a consequence of results of Rickart on the isomorphisms of some analogues of the classical groups . So, we begin with a statement on the form of the automorphisms of the general linear group $`GL(H)`$.
###### Theorem 3.1
Let $`\varphi :GL(H)GL(H)`$ be an automorphism of the general linear group. Then $`\varphi `$ is of one of the following forms:
* there exists a bounded linear invertible operator $`T:HH`$ such that
$$\varphi (X)=TXT^1(XGL(H)),$$
* there exists a bounded conjugate-linear invertible operator $`T:HH`$ such that
$$\varphi (X)=TXT^1(XGL(H)),$$
* there exists a bounded linear invertible operator $`T:HH`$ such that
$$\varphi (X)=\left(TX^1T^1\right)^{}(XGL(H)),$$
* there exists a bounded conjugate-linear invertible operator $`T:HH`$ such that
$$\varphi (X)=\left(TX^1T^1\right)^{}(XGL(H)).$$
Remark. If $`\varphi `$ is of type (i), (ii), (iii), (iv), then for every $`XGL(H)`$ we have $`\sigma (\varphi (X))=\sigma (X)`$, $`\sigma (\varphi (X))=\overline{\sigma (X)}`$, $`\sigma (\varphi (X))=(\overline{\sigma (X)})^1`$, $`\sigma (\varphi (X))=\sigma (X)^1`$, respectively.
Proof. Let $`\varphi :GL(H)GL(H)`$ be an automorphism. By a result of Rickart \[8, Theorem 5.1\], \[9, Theorem I\], $`\varphi `$ must be either of the form $`\varphi (X)=\tau (X)TXT^1`$ or of the form $`\varphi (X)=\tau (X)(TX^1T^1)^{}`$, where $`\tau :GL(H)𝐂`$ is a multiplicative map and $`T:HH`$ is a bijective additive map satisfying $`T(\lambda x)=f(\lambda )Tx`$, $`\lambda 𝐂`$, $`xH`$, for some ring automorphism $`f`$ of C (one has to be careful when applying the result of Rickart since $`S^{}`$ in his papers denotes the adjoint of an operator $`S`$ defined as for Banach space operators, while here, of course, $`S^{}`$ denotes the adjoint in the Hilbert space sense). Let us first show that $`\tau (X)1`$. Since $`\tau (I)=1`$ we have $`\tau (S)\{1,1\}`$ for any involution $`SGL(H)`$, that is, for any $`S`$ satisfying $`S^2=I`$. According to every element of $`GL(H)`$ can be written as a product of involutions, and consequently, the range of $`\tau `$ is contained in $`\{1,1\}`$. An arbitrary involution $`S`$ can be expressed as $`S=(IP)P`$ where $`P`$ is an idempotent. For $`R=(IP)+iP`$ we have $`S=R^2`$, and therefore, $`\tau (S)=1`$. Applying once again we conclude that $`\tau `$ is identically equal to 1.
Therefore, either $`\varphi (X)=TXT^1`$ or $`\varphi (X)=(TX^1T^1)^{}`$. In the second case we can compose $`\varphi `$ by $`X(X^{})^1`$ to conclude that in both cases $`XTXT^1`$ is an automorphism of $`GL(H)`$. We have to prove that $`f:𝐂𝐂`$ is either the identity or the complex conjugation and that $`T`$ is bounded. To prove this one can apply automatic continuity techniques (note that for this part of the proof the assumption that $`H`$ is infinite-dimensional is indispensable). However, we shall use a shorter way of reducing the problem to a known result.
So, assume that $`\varphi (X)=TXT^1`$ is an automorphism of $`GL(H)`$. Clearly, $`\varphi ^1(X)=T^1XT`$. Define additive mappings $`\psi ,\phi :B(H)A(H)`$ by $`\psi (X)=TXT^1`$ and $`\phi (X)=T^1XT`$, $`XB(H)`$. If $`|\lambda |>X`$, then $`\psi (X)=\psi (X\lambda I)+\psi (\lambda I)GL(H)+GL(H)B(H)`$. So, $`\psi `$ is a multiplicative map from $`B(H)`$ into $`B(H)`$ which is also bijective since $`\phi :B(H)B(H)`$ is its inverse. By there exists a bijective bounded linear or conjugate-linear map $`S:HH`$ such that $`TXT^1=SXS^1`$ for every $`XB(H)`$, or equivalently, the additive map $`S^1T`$ commutes with every $`XB(H)`$. It follows that $`T=\lambda S`$ for some nonzero scalar $`\lambda `$. This completes the proof of the statement that every automorphism $`\varphi `$ of $`GL(H)`$ has one of the forms (i), (ii), (iii) or (iv).
In the proof of the main result of this section we shall need the following lemma. Let $`K`$ be a nonempty subset of the complex plane and let $`\lambda `$ be a complex number. We use the following notation: $`K\lambda =\{\mu \lambda :\mu K\}`$, $`\overline{K}=\{\overline{\mu }:\mu K\}`$, and $`r(K)=sup\{|\mu |:\mu K\}[0,\mathrm{}]`$. If also $`0K`$, then $`K^1=\{\mu ^1:\mu K\}`$.
Lemma. Let $`K`$ be a nonempty compact subset of C and $`\lambda `$ a complex number such that $`\mathrm{Im}\lambda >r(K)+1`$. Then $`0K\lambda `$, $`K\lambda (K\lambda )^1`$, $`K\lambda \overline{K\lambda }`$, and $`K\lambda (\overline{K\lambda })^1`$.
Proof. Clearly, $`0K\lambda `$ and $`|\lambda |>r(K)+1`$. It follows that $`r(K\lambda )|\lambda |r(K)>1`$. On the other hand, $`r((K\lambda )^1)=\frac{1}{|\lambda \lambda _0|}`$ for some $`\lambda _0K`$. Since
$$\frac{1}{|\lambda \lambda _0|}\frac{1}{|\lambda ||\lambda _0|}\frac{1}{|\lambda |r(K)}<1,$$
we have $`r(K\lambda )>r((K\lambda )^1)=r((\overline{K\lambda })^1)`$, which further yields that $`K\lambda (K\lambda )^1`$ and $`K\lambda (\overline{K\lambda })^1`$. Furthermore, our assumption implies that $`K\lambda `$ belongs to the open lower half-plane, and consequently, $`K\lambda \overline{K\lambda }`$. This completes the proof.
###### Theorem 3.2
Every 2-local automorphism of $`GL(H)`$ is an automorphism.
Proof. Assume that $`\varphi `$ is a 2-local automorphism of $`GL(H)`$. Composing it with an appropriate automorphism of $`GL(H)`$ we can assume with no loss of generality that $`\varphi (2iI)=2iI`$. It follows from the Remark that $`\varphi _{X,2iI}`$ has to be of type (i) for every $`XGL(H)`$. In particular, we have $`\sigma (\varphi (X))=\sigma (X)`$ $`(XGL(H))`$, and $`\varphi (\lambda I)=\lambda I`$ $`(\lambda 𝐂)`$. Denote by $`𝒮`$ the set of all operators $`XGL(H)`$ satisfying $`\sigma (X)\overline{\sigma (X)}`$ and $`\sigma (X)\sigma (X)^1`$, $`\sigma (X)(\overline{\sigma (X)})^1`$. If $`X𝒮`$ and $`Y`$ is an arbitrary element of $`GL(H)`$ then $`\varphi _{X,Y}`$ has to be of type (i). In particular, $`\varphi _{X,Y}(\lambda I)=\lambda I`$ $`(\lambda 𝐂)`$.
Now, let $`XB(H)`$ be any bounded linear operator on $`H`$. Denote by $`L_X`$ the set of all complex numbers $`\lambda `$ such that $`X\lambda I`$ is an invertible operator contained in $`𝒮`$. By Lemma, this set is always nonempty. We define $`\psi :B(H)B(H)`$ by $`\psi (X)=\varphi (X\lambda I)+\lambda I`$ where $`\lambda L_X`$. First we have to show that $`\psi `$ is well-defined. So, assume that $`\mu `$ also belongs to $`L_X`$. Then we already know that $`\varphi _{X\lambda I,X\mu I}`$ is of type (i), and consequently,
$$\varphi (X\lambda I)+\lambda I=\varphi _{X\lambda I,X\mu I}(X\lambda I)+\lambda I=$$
$$\varphi _{X\lambda I,X\mu I}(X)=\varphi (X\mu I)+\mu I.$$
Next, we show that the restriction of $`\psi `$ to $`GL(H)`$ coincides with $`\varphi `$. In order to do this we first observe that if $`\phi :GL(H)GL(H)`$ is an automorphism of type (i) and if $`X,Y`$ are arbitrary elements of $`GL(H)`$ such that $`X+Y`$ is also invertible, then $`\phi (X+Y)=\phi (X)+\phi (Y)`$. So, for any $`XGL(H)`$ and $`\lambda L_X\{0\}`$ we have
$$\psi (X)=\varphi (X\lambda I)+\lambda I=\varphi _{X,X\lambda I}(X\lambda I)+\lambda I=$$
$$\varphi _{X,X\lambda I}(X)+\varphi _{X,X\lambda I}(\lambda I)+\lambda I=\varphi _{X,X\lambda I}(X)=\varphi (X).$$
It is well-known that every algebra automorphism of $`B(H)`$ is inner. We show that $`\psi `$ is a 2-local automorphism of $`B(H)`$, that is, for every pair $`X,YB(H)`$ there exists a bounded linear invertible $`S:HH`$ such that $`\psi (X)=SXS^1`$ and $`\psi (Y)=SYS^1`$. By Lemma we know that there exists $`\lambda 𝐂`$ such that $`\lambda L_XL_Y`$. The automorphism $`\varphi _{X\lambda I,Y\lambda I}`$ is of type (i), and therefore spatially implemented by a bounded linear invertible operator, say $`S`$. Then
$$\psi (X)=\varphi (X\lambda I)+\lambda I=\varphi _{X\lambda I,Y\lambda I}(X\lambda I)+\lambda I=$$
$$S(X\lambda I)S^1+\lambda I=SXS^1,$$
and similarly,
$$\psi (Y)=SYS^1.$$
Applying the result of the second author on 2-local automorphisms of $`B(H)`$ we conclude that there exists a bounded linear invertible operator $`T:HH`$ such that $`\psi (X)=TXT^1`$ for every $`XB(H)`$. Consequently, $`\varphi (X)=TXT^1`$ for every $`XGL(H)`$. This completes the proof.
## Acknowledgement
We express our gratitude to Professor R.V. Kadison for providing information on the mentioned papers of C.E. Rickart.
This research was supported by a Hungarian-Slovene grant and also by a grant from the Ministry of Science of Slovenia.
Lajos Molnár
Institute of Mathematics
Lajos Kossuth University
P.O. Box 12
4010 Debrecen, Hungary
e-mail:molnarl@math.klte.hu
Peter Šemrl
Department of Mathematics
University of Ljubljana
Jadranska 19
1000 Ljubljana, Slovenia
e-mail:peter.semrl@fmf.uni-lj.si
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# SELF-SIMILAR PERFECT FLUIDS
## 1 Basic facts about homotheties
Throughout this paper $`(M,g)`$ will denote a space-time: $`M`$ then being a Hausdorff, simply connected, four-dimensional manifold, and $`g`$ a Lorentz metric of signature (-,+,+,+). All the structures will be assumed smooth.
A global vector field $`X`$ on $`M`$ is called homothetic if either one of the following conditions holds on a local chart
$$_Xg_{ab}=2\lambda g_{ab},X_{a;b}=\lambda g_{ab}+F_{ab},$$
(1)
where $`\lambda `$ is a constant on $`M`$ , $``$ stands for the Lie derivative operator, a semi-colon denotes a covariant derivative with respect to the metric connection, and $`F_{ab}=F_{ba}`$ is the so-called homothetic bivector. If $`\lambda 0`$, $`X`$ is called proper homothetic and if $`\lambda =0`$, $`X`$ is a Killing vector field (KV) on $`M`$. For a geometrical interpretation of (1) we refer the reader to$`^{\mathrm{?},\mathrm{?}}`$.
A necessary condition that $`X`$ be homothetic is
$$X_{}^{a}{}_{;bc}{}^{}=R_{}^{a}{}_{bcd}{}^{}X^d,$$
(2)
where $`R_{}^{a}{}_{bcd}{}^{}`$ are the components of the Riemann tensor in a coordinate chart; thus, a homothetic vector field (HVF) is a particular case of affine collineation$`^\mathrm{?}`$ and therefore it will satisfy
$$_XR_{}^{a}{}_{bcd}{}^{}=_XR_{ab}=_XC_{}^{a}{}_{bcd}{}^{}=0,$$
(3)
where $`R_{ab}`$ ($`R_{}^{c}{}_{acb}{}^{}`$) and $`C_{}^{a}{}_{bcd}{}^{}`$ stand, respectively for the components of the Ricci and the conformal Weyl tensor.
The set of all HVFs on $`M`$ forms a finite dimensional Lie algebra under the usual bracket operation and will be referred to as the homothetic algebra, $`_r`$, $`r`$ being its dimension. The set of all Killing vectors fields on $`M`$ forms a finite dimensional Lie algebra (dimension $`s`$) under the same bracket operation, and will be referred to here as the Lie algebra of isometries, $`𝒢_s`$ which is contained in (i.e., is a subalgebra of) $`_r`$. Furthermore, it is immediate to see by direct computation that the Lie bracket of an HVF with a KV is always a KV. From these considerations it immediately follows that the highest possible dimension of $`_r`$ in a four-dimensional manifold is $`r=11`$.
If $`rs`$ then $`s=r1`$ necessarily, and one may choose a basis $`\{X_A\}_{A=1\mathrm{}r}\{X_1,\mathrm{},X_{r1},X\}`$ for $`_r`$, in such a way that $`X`$ is proper homothetic and $`X_1,\mathrm{},X_{r1}`$ are Killing vector fields spanning $`𝒢_{r1}`$. If these vector fields in the basis of $`_r`$ are all complete vector fields, then each member of $`_r`$ is complete and Palais’ theorem $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ guarantees the existence of an $`r`$-dimensional Lie group of homothetic transformations of $`M`$ ($`H_r`$) in a well-known way; otherwise, it gives rise to a local group of local homothetic transformations of $`M`$ and, although the usual concepts of isotropy and orbits still hold, a little more care is required $`^\mathrm{?}`$.
The following result $`^{\mathrm{?},\mathrm{?}}`$ will be useful: The orbits associated with $`_r`$ and $`𝒢_{r1}`$ can only coincide if either they are four-dimensional or three-dimensional and null. (This result still holds if $`_r`$ is replaced by the conformal Lie algebra $`𝒞_r`$ and does not depend on the maximality of $`_r`$ or $`𝒞_r`$).
The set of zeroes of a proper HVF, i.e., $`\{pM:X(p)=0\}`$ (fixed points of the homothety), either consists of topologically isolated points, or else is part of a null geodesic. The latter case corresponds to the well-known (conformally flat or Petrov type N) plane waves $`^{\mathrm{?},\mathrm{?}}`$.
At any zero of a proper HVF on $`M`$ all Ricci and Weyl eigenvalues must necessarily vanish and thus the Ricci tensor is either zero or has Segre type $`\{(2,11)\}`$ or $`\{(3,1)\}`$ (both with zero eigenvalue), whereas the Weyl tensor is of the Petrov type $`O`$, $`N`$ or $`III`$ $`^\mathrm{?}`$ (see also $`^\mathrm{?}`$ for vacuum space-times).
## 2 Basic facts about perfect fluids admitting HVFs
The energy-momentum tensor for a perfect fluid is given by
$$T_{ab}=(\mu +p)u_au_b+pg_{ab},$$
(4)
where $`\mu `$ and $`p`$ are, respectively, the energy density and the pressure as measured by an observer comoving with the fluid, and $`u^a`$ ($`u^au_a=1`$) is the four-velocity of the fluid. If $`X`$ is an HVF then, from Einstein’s Field Equations (EFE) it follows that
$$_XT_{ab}=0,$$
(5)
and this implies in turn $`^\mathrm{?}`$
$$_Xu_a=\lambda u_a,_Xp=2\lambda p,_X\mu =2\lambda \mu .$$
(6)
Thus, the Lie derivatives of $`u_a`$, $`p`$ and $`\mu `$ with respect to a KV vanish identically.
If a barotropic equation of state exists, $`p=p(\mu )`$, and the space-time admits a proper HVF $`X`$ then $`^\mathrm{?}`$
$$p=(\gamma 1)\mu ,$$
(7)
where $`\gamma `$ is a constant ($`0\gamma 2`$ in order to comply with the weak and dominant energy conditions). Of particular interest are the values $`\gamma =1`$ (pressure-free matter, “dust”) and $`\gamma =4/3`$ (radiation fluid). In addition, the value $`\gamma =2`$ (stiff-matter) has been considered in connection with the early Universe. Furthermore, values of $`\gamma `$ satisfying $`0\gamma <2/3`$, while physically unrealistic as regards a classical fluid, are of interest in connection with inflationary models of the Universe. In particular, the value $`\gamma =0`$, for which the fluid can be interpreted as a positive cosmological constant, corresponds to exponential inflation, while the values $`0<\gamma <2/3`$ correspond to power law inflation in FRW models $`^\mathrm{?}`$, but it is customary to restrict $`\gamma `$ to the range $`1\gamma 2`$.
If the proper HVF $`X`$ and the four-velocity $`u`$ are mutually orthogonal (i.e., $`u^aX_a=0`$) and a barotropic equation of state is assumed, it follows that $`\gamma =2`$, i.e., $`p=\mu `$ stiff-matter $`^\mathrm{?}`$, on the other hand, if $`X^a=\alpha u^a`$ the fluid is then shear-free. Further information on this topic can be found in $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$.
## 3 The “dimensional count-down”
In this section, the maximal Lie algebra of global HVF on $`M`$ will be denoted as $`_r`$ ($`r`$ being its dimension), and it will be assumed that at least one member of it is proper homothetic.
The case of multiply transitive action is thoroughly studied in $`^\mathrm{?}`$. We summarize in the following table the results given there, which follow invariably from considerations on the associated Killing subalgebra and the fixed point structure of the proper HVF.
The first entry in the table gives the dimension of the group of homotheties, the second and third entries stand for the nature and dimension of the homothetic and Killing orbits respectively (e.g.: $`N_2`$, $`T_2`$ and $`S_2`$ denote Null, Timelike and Spacelike two-dimensional orbits respectively, $`O_3`$ stands for three-dimensional orbits of either nature, timelike, spacelike or null), the fourth and fifth entries give the Petrov and Segre type(s) of the associated Weyl and Ricci tensors. Finally, the last two entries give respectively the possible interpretation whenever it is in some sense unique, and the existence or non-existence of perfect fluid solutions for that particular case, along with some supplementary information; thus FRW stands for Friedmann-Robertson-Walker, LRS for Locally Rotationally Symmetric, and Bianchi refers to that family of perfect fluid solutions. The cases that cannot arise are labeled as “Not Possible”, and wherever no information is given on the Petrov and Segre types, it is to be understood that all types are possible in principle. The Segre type of the Ricci tensor of the case described in the last row, is unrestricted except in that it must necessarily have two equal (spacelike) eigenvalues; perfect fluid solutions of these characteristics constitute special cases of spherically, plane or hyperbolically symmetric perfect fluid space-times. For further information on LRS spacetimes, see $`^{\mathrm{?},\mathrm{?}}`$; for the case $`r=4`$ transitive and null three-dimensional Killing orbits, see $`^{\mathrm{?},\mathrm{?}}`$. Regarding spatially homogeneous Bianchi models, see $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$; and for the last three cases occurring in the table, see $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$.
| $`r`$ | $`O_m`$ | $`K_n`$ | Petrov | Segre | Interpretation | PF info. |
| --- | --- | --- | --- | --- | --- | --- |
| $`11`$ | $`M`$ | $`M`$ | $`O`$ | $`0`$ | Flat | $``$ |
| $`10`$ | $`M`$ | $`M`$ | - | - | Not Possible | $``$ |
| $`9`$ | $`M`$ | $`M`$ | - | - | Not Possible | $``$ |
| $`8`$ | $`M`$ | $`M`$ | $`O`$ | $`\left\{(2,11)\right\}`$ | Gen. P. wave | $``$ |
| $`7`$ | $`M`$ | $`M`$ | $`N`$ | $`0,\left\{(2,11)\right\}`$ | Gen. P. wave | $``$ |
| $`7`$ | $`M`$ | $`T_3`$ | $`O`$ | $`\left\{(1,11)1\right\}`$ | Tachyonic Fl. | $``$ |
| $`7`$ | $`M`$ | $`N_3`$ | - | - | Not Possible | $``$ |
| $`7`$ | $`N_3`$ | $`N_3`$ | $`O`$ | $`\left\{(2,11)\right\}`$ | Gen. P. wave | $``$ |
| $`7`$ | $`M`$ | $`S_3`$ | $`O`$ | $`\{1,\left(111\right)\}`$ | Perfect Fluid | FRW |
| $`6`$ | $`M`$ | $`M`$ | - | - | Not Possible | $``$ |
| $`6`$ | $`N_3`$ | $`N_3`$ | $`N`$ | $`\left\{(2,11)\right\}`$ | Gen. P. wave | $``$ |
| $`5`$ | $`M`$ | $`M`$ | - | - | Not Possible | $``$ |
| $`5`$ | $`M`$ | $`N_3`$ | - | - | - | $``$ |
| $`5`$ | $`N_3`$ | $`N_3`$ | - | - | Not Possible | $``$ |
| $`5`$ | $`M`$ | $`T_3`$ | $`D,N,O`$ | $`\{1,1\left(11\right)\},\{2,\left(11\right)\}`$ | | LRS |
| $`5`$ | $`M`$ | $`S_3`$ | $`D,O`$ | $`\left\{(1,1)11\right\},\left\{(2,1)1\right\}`$ | | LRS |
| $`4`$ | $`M`$ | $`N_3`$ | $`II,III,D,N,O`$ | $`\left\{(1,1)\left(11\right)\right\},\left\{(2,11)\right\}`$ | Plane waves | $``$ |
| $`4`$ | $`N_3`$ | $`N_3`$ | - | - | Not Possible | $``$ |
| $`4`$ | $`M`$ | $`T_3`$ | | | | Bianchi |
| $`4`$ | $`M`$ | $`S_3`$ | | | | Bianchi |
| $`4`$ | $`O_3`$ | $`N_2`$ | $`N,O`$ | $`\{3,1\},\{2,11\},\left\{(1,1)11\right\}`$ | | $``$ |
| $`4`$ | $`O_3`$ | $`T_2`$ | $`D,O`$ | $`\left\{(1,1)11\right\}`$ | | $``$ |
| $`4`$ | $`O_3`$ | $`S_2`$ | $`D,O`$ | $`\left\{\left(11\right)\right\}`$ | | $``$ |
The case $`r=3`$ has an associated Killing subalgebra $`𝒢_2`$ and the respective dimensions of their orbits are 3 and 2 (see for instance $`^{\mathrm{?},\mathrm{?},\mathrm{?},\mathrm{?}}`$ and references cited therein). When the Killing subalgebra has null orbits, the metric is of Kundt’s class $`^\mathrm{?}`$ and perfect fluids are excluded. If the Killing orbits are timelike, the solutions can then be interpreted as special cases of axisymmetric stationary space-times $`^{\mathrm{?},\mathrm{?}}`$, and if they are spacelike as special cases of inhomogeneous cosmological solutions or cylindrically symmetric space-times. In both cases, perfect fluid solutions have been found.
## References
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warning/0005/math0005251.html
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ar5iv
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text
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# The domain algebra of a 𝐶𝑃-semigroup
## 1. Basic properties of $`𝒜`$
Let $`\varphi =\{\varphi _t:t0\}`$ be a $`CP`$-semigroup as defined in the abstract. We first recall four characterizations of the domain of the generator of $`\varphi `$.
###### Lemma 1
Let $`A(H)`$. The following are equivalent.
###### Demonstration proof
The implications (i) $``$ (ii) and (iii) $``$ (iv) are trivial, and (ii) $``$ (iii) is a straightforward consequence of the Banach-Steinhaus theorem.
###### Demonstration proof of (iv) $``$ (i)
Since the unit ball of $`(H)`$ is weakly sequentially compact, the hypothesis (iv) implies that there is a sequence $`t_n0+`$ such that
$$\frac{1}{t_n}(\varphi _{t_n}(A)A)T(H)$$
in the weak operator topology. We claim: for every $`s>0`$,
$$_0^s\varphi _\lambda (T)𝑑\lambda =\varphi _s(A)A.$$
$`1.1`$
The integral on the left is interpreted as a weak integral; that is, for $`\xi ,\eta H`$,
$$_0^s\varphi _\lambda (T)\xi ,\eta 𝑑\lambda =\varphi _s(A)\xi ,\eta A\xi ,\eta .$$
To see that, fix $`\lambda >0`$. Since $`\varphi _\lambda `$ is weakly continuous on bounded sets in $`(H)`$ we have
$$\frac{1}{t_n}(\varphi _{\lambda +t_n}(A)\varphi _\lambda (A))=\varphi _\lambda (\frac{1}{t_n}(\varphi _{t_n}(A)A))\varphi _\lambda (T)$$
in the weak operator topology, as $`n\mathrm{}`$. By the bounded convergence theorem, we find that for fixed $`\xi ,\eta H`$,
$$\underset{n\mathrm{}}{lim}\frac{1}{t_n}(_0^s\varphi _{\lambda +t_n}(A)\xi ,\eta 𝑑\lambda _0^s\varphi _\lambda (A)\xi ,\eta 𝑑\lambda )=_0^s\varphi _\lambda (T)\xi ,\eta 𝑑\lambda .$$
Writing
$$_0^sf(\lambda +t_n)𝑑\lambda _0^sf(\lambda )𝑑\lambda =_s^{s+t_n}f(\lambda )𝑑\lambda _0^{t_n}f(\lambda )𝑑\lambda ,$$
the left side of the preceding formula becomes
$$\underset{n\mathrm{}}{lim}(\frac{1}{t_n}_s^{s+t_n}\varphi _\lambda (A)\xi ,\eta 𝑑\lambda \frac{1}{t_n}_0^{t_n}\varphi _\lambda (A)\xi ,\eta 𝑑\lambda )$$
which, because of continuity of $`\varphi `$ in the time parameter, is $`\varphi _s(A)\xi ,\eta A\xi ,\eta `$, as asserted in (1.1).
To prove the strong-$``$ convergence asserted in (i), fix $`\xi H`$ and use (1.1) to write
$`{\displaystyle \frac{1}{s}}(\varphi _s(A)\xi A\xi )T\xi ={\displaystyle \frac{1}{s}}{\displaystyle _0^s}\varphi _\lambda (T)\xi 𝑑\lambda {\displaystyle _0^s}T\xi 𝑑\lambda `$
$`{\displaystyle \frac{1}{s}}{\displaystyle _0^s}\varphi _\lambda (T)\xi T\xi 𝑑\lambda ({\displaystyle \frac{1}{s}}{\displaystyle _0^s}\varphi _\lambda (T)\xi T\xi ^2𝑑\lambda )^{1/2}.`$
The integrand of the last term expands as follows
$`\varphi _\lambda (T)\xi T\xi ^2`$ $`=\varphi _\lambda (T)^{}\varphi _\lambda (T)\xi ,\xi 2\mathrm{}\varphi _\lambda (T)\xi ,T\xi +T\xi ^2`$
$`\varphi _\lambda (T^{}T)\xi ,\xi 2\mathrm{}\varphi _\lambda (T)\xi ,T\xi +T\xi ^2,`$
the last inequality by the Schwarz inequality for unital CP maps. Since $`\varphi _\lambda (T^{}T)`$ (resp. $`\varphi _\lambda (T)`$) tends weakly to $`T^{}T`$ (resp. $`T`$) as $`\lambda 0+`$, it follows that
$$\underset{s0+}{lim\; sup}\frac{1}{s}_0^s\varphi _\lambda (T)\xi T\xi ^2𝑑\lambda T^{}T\xi ,\xi 2T\xi ,T\xi +T\xi ^2=0,$$
and we conclude that $`\frac{1}{s}(\varphi _s(A)A)`$ tends strongly to $`T`$ as $`s0+`$.
Similarly, $`\frac{1}{s}(\varphi _s(A)A)^{}=\frac{1}{s}(\varphi _s(A^{})A^{})`$ tends strongly to $`T^{}`$.
###### Definition
Let $`𝒟`$ be the set of all operators $`A(H)`$ for which the four conditions of Lemma 1 are satisfied. $`L:𝒟(H)`$ denotes the generator of $`\varphi `$,
$$L(A)=\underset{t0+}{lim}\frac{1}{t}(\varphi _t(A)A),A𝒟.$$
It is obvious that $`𝒟`$ is a self-adjoint linear subspace of $`(H)`$, that $`L(A^{})=L(A)^{}`$ for $`A𝒟`$, and a standard argument shows that $`𝒟`$ is dense in $`(H)`$ in the $`\sigma `$-strong operator topology.
###### Lemma 2
For every operator $`A𝒟`$ we have
$$L(A)=\underset{t>0}{sup}\frac{1}{t}\varphi _t(A)A.$$
###### Demonstration proof
The inequality $``$ is clear from the fact that $`L(A)`$ is the weak limit of operators $`\frac{1}{t}(\varphi _t(A)A)`$ near $`t=0+`$, i.e.,
$$L(A)\underset{t0+}{lim\; sup}\frac{1}{t}\varphi _t(A)A\underset{t>0}{sup}\frac{1}{t}\varphi _t(A)A.$$
For $``$, set $`T=L(A)`$. Using (1.1), we can write for every $`t>0`$
$$\frac{1}{t}\varphi _t(A)A=\frac{1}{t}_0^t\varphi _\lambda (T)𝑑\lambda \frac{1}{t}_0^t\varphi _\lambda (T)𝑑\lambda T,$$
since $`\varphi _\lambda 1`$ for every $`\lambda 0`$.
###### Theorem A
$`𝒜=\{A𝒟:A^{}A𝒟,AA^{}𝒟\}`$ is a $``$-subalgebra of $`(H)`$.
###### Demonstration proof
$`𝒜`$ is obviously a self-adjoint set of operators. We have to show that $`𝒜`$ is a vector space satisfying $`𝒜𝒜𝒜`$.
Fix $`t>0`$. By Stinespring’s theorem we can write
$$\varphi _t(X)=V_t^{}\pi _t(X)V_t,X(H)$$
$`1.2`$
where $`V_t`$ is an isometry from $`H`$ into some other Hilbert space $`H_t`$ and where $`\pi _t:(H)(H_t)`$ is a normal $``$-homomorphism of von Neumann algebras. $`P_t=V_tV_t^{}`$ is a self-ajoint projection in $`(H_t)`$.
For $`t>0`$ we will consider the seminorms $`p_t`$, $`q_t`$ defined on $`(H)`$ as follows
$`p_t(X)`$ $`=t^1\varphi _t(X)X,`$
$`q_t(X)`$ $`=t^{1/2}P_t\pi _t(X)\pi _t(X)P_t,X(H).`$
###### Lemma 3
For every operator $`X(H)`$ we have the following characterizations.
###### Remark Remark
The second assertion of Lemma 3 requires clarification. By definition, an operator $`X`$ belongs to $`𝒜`$ iff all four operators $`X,X^{},X^{}X,XX^{}`$ belong to the domain of the generator $`L`$ of $`\varphi =\{\varphi _t:t0\}`$. In that case both operators $`\sigma _L(dX^{}dX)`$ and $`\sigma _L(dXdX^{})`$ are well defined by the above formulas. The “symbol” map $`\sigma _L`$ will be discussed more fully in section 3.
###### Demonstration proof of Lemma 3
The assertion (i) follows from Lemmas 1 and 2 above. In order to prove (ii) we require the following more concrete expression for the seminorm $`q_t`$,
$`q_t(X)=`$ $`1.3`$
$`\mathrm{max}({\displaystyle \frac{1}{t}}(\varphi _t(X^{}X)\varphi _t(X)^{}\varphi _t(X))^{1/2},{\displaystyle \frac{1}{t}}(\varphi _t(XX^{})\varphi _t(X)^{}\varphi _t(X^{}))^{1/2}).`$
To prove (1.3) we decompose the commutator $`\pi _t(X)P_tP_t\pi _t(X)`$ into a sum
$$\pi _t(X)P_tP_t\pi _t(X)=(\mathrm{𝟙}P_t)\pi _t(X)P_tP_t\pi _t(X)(\mathrm{𝟙}P_t).$$
Since the first term $`(\mathrm{𝟙}P_t)\pi _t(X)P_t`$ has initial space in $`P_tH_t`$ and final space in $`(\mathrm{𝟙}P_t)`$, and the second term has the opposite property, it follows that
$$\pi _t(X)P_tP_t\pi _t(X)=\mathrm{max}((\mathrm{𝟙}P_t)\pi _t(X)P_t,P_t\pi _t(X)(\mathrm{𝟙}P_t)).$$
We have
$`(\mathrm{𝟙}P_t)\pi _t(X)P_t^2`$ $`=V_t^{}\pi _t(X^{})(\mathrm{𝟙}P_t)\pi _t(X)V_t`$
$`=V_t^{}\pi _t(X^{}X)V_tV_t^{}\pi _t(X^{})V_tV_t^{}\pi _t(X)V_t`$
$`=\varphi _t(X^{}X)\varphi _t(X)^{}\varphi _t(X).`$
Similarly,
$$P_t\pi _t(X)(\mathrm{𝟙}P_t)^2=V_t^{}\pi _t(X)(\mathrm{𝟙}P_t)\pi _t(X^{})V_t=\varphi _t(XX^{})\varphi _t(X)^{}\varphi _t(X^{}),$$
and formula (1.3) follows from these two expressions.
Now if $`X𝒜`$ then all four operators $`X,X^{},X^{}X,XX^{}`$ belong to $`𝒟`$, hence all four limits
$`\underset{t0+}{lim}{\displaystyle \frac{1}{t}}(\varphi _t(X^{}X)X^{}X)=L(X^{}X),`$
$`\underset{t0+}{lim}{\displaystyle \frac{1}{t}}(\varphi _t(XX^{})X^{}X)=L(XX^{}),`$
$`\underset{t0+}{lim}{\displaystyle \frac{1}{t}}(\varphi _t(X)X)=L(X),`$
$`\underset{t0+}{lim}{\displaystyle \frac{1}{t}}(\varphi _t(X^{})X^{})=L(X^{})`$
exist relative to the strong operator topology. Writing
$`\varphi _t(X^{}X)\varphi _t(X)^{}\varphi _t(X)=`$ $`1.4`$
$`(\varphi _t(X^{}X)`$ $`X^{}X)X^{}(\varphi _t(X)X)(\varphi _t(X^{})X^{})\varphi _t(X)`$
and using strong continuity of multiplication on bounded sets, we find that the limit
$$\underset{t0+}{lim}\frac{1}{t}(\varphi _t(X^{}X)\varphi _t(X^{})\varphi _t(X))=L(X^{}X)X^{}L(X)L(X^{})X=\sigma _L(dX^{}dX)$$
exists relative to the strong operator topology.
In the same way we deduce the existence of the strong limit
$$\underset{t0+}{lim}\frac{1}{t}(\varphi _t(XX^{})\varphi _t(X)\varphi _t(X^{}))=L(XX^{})XL(X^{})L(X)X^{}=\sigma _L(dXdX^{}).$$
It follows that for every $`X𝒜`$ the seminorms $`q_t(X)`$ are bounded for $`t>0`$, and for such $`X`$ we have
$$\mathrm{max}(\sigma _L(dX^{}dX)^{1/2},\sigma _L(dXdX^{})^{1/2})\underset{t0+}{lim\; sup}q_t(X).$$
Conversely, suppose we are given an operator $`X𝒟`$ for which the seminorms $`q_t(X)`$ are bounded for $`t>0`$. We have to show that $`X^{}X`$ and $`XX^{}`$ belong to $`𝒟`$; since $`𝒟`$ is self-adjoint and the seminorms $`q_t`$ are symmetric in that $`q_t(X^{})=q_t(X)`$, it is enough to show that $`X^{}X`$ belong to $`𝒟`$. (1.4) implies that for fixed $`t>0`$,
$`\varphi _t(X^{}X)X^{}X=`$ $`1.5`$
$`(\varphi _t(X^{}X)`$ $`\varphi _t(X^{})\varphi _t(X))+X^{}(\varphi _t(X)X)+(\varphi _t(X^{})X^{})\varphi _t(X)`$
Because of (1.3), the first term on the right of (1.5) is bounded in norm by $`M_1t`$ where $`M_1`$ is a positive constant. Similarly, since $`X`$ and $`X^{}`$ belong to $`𝒟`$ the second and third terms are bounded in norm by terms of the form $`M_2t`$ and $`M_3t`$ respectively, hence
$$\varphi _t(X^{}X)X^{}X(M_1+M_2+M_3)t.$$
By Lemma 1, $`X^{}X`$ must belong to $`𝒟`$.
Turning now to the proof of Theorem A, (or more properly, to the proof that $`𝒜`$ is an algebra), Lemma 3 tells us that $`𝒜`$ consists of all operators $`X(H)`$ for which
$$\underset{t>0}{sup}p_t(X)<\mathrm{},\text{and }\underset{t>0}{sup}q_t(X)<\mathrm{}.$$
Since $`p_t`$ and $`q_t`$ are both seminorms, it follows that $`𝒜`$ is a complex vector space which is obviously closed under the $``$-operation.
To see that $`𝒜`$ is closed under multiplication, pick $`X,Y𝒜`$. According to Lemma 3, it is enough to show
$$\underset{t>0}{sup}q_t(XY)<\mathrm{}$$
$`1.6`$
and
$$\underset{t>0}{sup}p_t(XY)<\mathrm{}$$
$`1.7`$
To prove (1.6) we claim that
$$q_t(XY)q_t(X)Y+Xq_t(Y).$$
$`1.8`$
Indeed, writing $`[A,B]`$ for the commutator $`ABBA`$ we have
$$[P_t,\pi _t(XY)]=[P_t,\pi _t(X)]\pi _t(Y)+\pi _t(X)[P_t,\pi _t(Y)],$$
and hence
$`q_t(XY)`$ $`=t^{1/2}[P_t,\pi _t(XY)]`$
$`t^{1/2}[P_t,\pi _t(X)\pi _t(Y)+\pi _t(X)t^{1/2}[P_t,\pi _t(Y)],`$
from which (1.8) is evident.
Finally, consider the condition (1.7). By definition of $`𝒜`$, $`A𝒜`$ implies $`A^{}A𝒟`$. Since $`𝒜`$ is now known to be a linear space we can assert that if $`X,Y𝒜`$ then for every $`k=0,1,2,3`$ we have $`Y+i^kX𝒜`$, hence $`(Y+i^kX)^{}(Y+i^kX)𝒟`$ and by the polarization formula
$$X^{}Y=\frac{1}{4}\underset{k=0}{\overset{3}{}}i^k(Y+i^kX)^{}(Y+i^kX),$$
$`X^{}Y`$ must also belong to $`𝒟`$. Since $`𝒜^{}=𝒜`$, we can replace $`X^{}`$ with $`X`$ to conclude that $`XY𝒟`$. (1.7) now follows from Lemma 3 (i).
###### Corollary
Let $`𝒟`$ be the domain of the generator of a $`CP`$-semigroup acting on $`(H)`$ and let $`A`$ be a self-adjoint operator such that $`A𝒟`$ and $`A^2𝒟`$. Then $`p(A)𝒟`$ for every polynomial $`p(x)=a_0+a_1x+\mathrm{}+a_nx^n`$.
## 2. Examples
In this section we describe two classes of examples which are in a sense at opposite extremes. In the first class of examples of $`CP`$-semigroups $`\varphi =\{\varphi _t:t0\}`$, each $`\varphi _t`$ leaves the $`C^{}`$-algebra $`𝒦`$ of all compact operators invariant, $`\varphi _t(𝒦)𝒦`$, its domain algebra $`𝒜`$ is strongly dense in $`(H)`$, and its generator restricts to a “second order” differential operator on $`𝒜`$ (see formula (1.1) of ). In the second class of examples, the individual maps satisfy $`\varphi _t(𝒦)𝒦=\{0\}`$ for $`t>0`$, $`𝒜`$ is not strongly dense in $`(H)`$, and its generator is degenerate in the sense that it restricts to a derivation on $`𝒜`$.
We first recall the class of examples of $`CP`$-semigroups of , including the heat flow of the $`CCR`$ algebra. While for simplicity we confine the discussion to the case of one degree of freedom, the reader will note that everything carries over verbatim to the case of $`n`$ degrees of freedom, $`n=1,2,\mathrm{}`$.
Let $`\{W_z:z^2\}`$ be an irreducible Weyl system acting on a Hilbert space $`H`$. Thus, $`z^2W_z`$ is a strongly continuous mapping from $`^2`$ into the unitary operators on $`H`$ which satisfies the canonical commutation relations in Weyl’s form
$$W_{z_1}W_{z_2}=e^{i\omega (z_1,z_2)}W_{z_1+z_2},z_1,z_2^2,$$
$`\omega `$ denoting the symplectic form on $`^2`$ given by
$$\omega ((x,y),(x^{},y^{}))=\frac{1}{2}(x^{}yxy^{}).$$
Let $`\{\mu _t:t0\}`$ be a one-parameter family of probability measures on $`^2`$ which is a semigroup under the natural convolution of measures
$$\mu \nu (S)=_{^2\times ^2}\chi _S(z+w)𝑑\mu (z)𝑑\nu (w),$$
which satisfies $`\mu _0=\delta _{(0,0)}`$, and which is measurable in $`t`$ in the natural sense. It is convenient to define the Fourier transform of a measure $`\mu `$ in terms of the symplectic form $`\omega `$ as follows,
$$\widehat{\mu }(z)=_^2e^{i\omega (z,\zeta )}𝑑\mu (\zeta ),z^2.$$
Given such a semigroup of probability measures $`\{\mu _t:t0\}`$ there is a unique $`CP`$ semigroup $`\varphi =\{\varphi _t:t0\}`$ acting on $`(H)`$ which satisfies
$$\varphi _t(W_z)=\widehat{\mu }_t(z)W_z,z^2,t0$$
see , Proposition 1.7. Two cases of particular interest are
$$\varphi _t(W_z)=e^{t|z|^2}W_z,t0$$
$`CCRheatflow`$
where $`|(x,y)|`$ denotes the Euclidean norm $`(x^2+y^2)^{1/2}`$, and
$$\varphi _t(W_z)=e^{t|z|}W_z,t0.$$
$`Cauchyflow`$
For both of these examples a straightforward estimate shows that for fixed $`z^2`$ there is a constant $`M>0`$ such that
$$\varphi _t(W_z)W_z=|\widehat{\mu }_t(z)1|Mt,t>0$$
and hence $`W_z𝒟`$. Since $`W_z`$ is unitary, $`\mathrm{𝟙}=W_z^{}W_z=W_zW_z^{}`$ belongs to $`𝒟`$ , and hence $`W_z`$ belongs to the domain algebra $`𝒜`$ of $`\varphi `$ for every $`z^2`$. We conclude that for these examples, the domain algebra is strongly dense in $`(H)`$.
Indeed, it is not hard to show that $`𝒜`$ contains a $``$-algebra of compact operators that is norm-dense in the algebra $`𝒦`$ of all compact operators. Unlike the examples to follow, these flows leave $`𝒦`$ invariant in the sense that $`\varphi _t(𝒦)𝒦`$ for all $`t0`$, and can therefore be considered as $`CP`$-semigroups which act on the separable $`C^{}`$-algebra $`𝒦`$, rather than than as $`CP`$-semigroups acting on $`(H)`$.
We now describe a class of examples of $`CP`$ semigroups whose domain algebras are not strongly dense in $`(H)`$. These examples are inspired by a class of $`CP`$ semigroups that have emerged in recent work of Robert Powers, to whom we are indebted for useful discussions.
Let $`H=L^2(0,\mathrm{})`$ and let $`U=\{U_t:t0\}`$ be the semigroup of isometries $`U_t\xi (x)=\xi (xt)`$ for $`xt`$, $`U_t\xi (x)=0`$ for $`0x<t`$. Fix a real number $`\alpha >0`$, and let $`f`$ be the unit vector in $`L^2(0,\mathrm{})`$ obtained by normalizing the exponential function $`u(x)=e^{\alpha x}`$, $`x0`$. One has $`U_t^{}f=e^{\alpha t}f`$ for every $`t0`$, hence the vector state $`\omega (A)=Af,f`$ satisfies $`\omega (U_tAU_t^{})=e^{2\alpha t}\omega (A)`$, $`A(H)`$.
We consider the family of unit-preserving normal completely positive maps $`\varphi =\{\varphi _t:t0\}`$ defined on $`(H)`$ by
$$\varphi _t(A)=\omega (A)E_t+U_tAU_t^{},t0.$$
where $`E_t=\mathrm{𝟙}U_tU_t^{}`$ is the projection on the subspace $`L^2(0,t)L^2(0,\mathrm{})`$. Since
$$\omega (E_t)=\omega (\mathrm{𝟙})\omega (U_tU_t^{})=1e^{2\alpha t},$$
it follows that $`\omega (\varphi _t(A))=\omega (A)`$ for every $`A`$. A routine computation now shows that $`\varphi `$ satisfies the semigroup property $`\varphi _s\varphi _t=\varphi _{s+t}`$, hence $`\varphi `$ is a $`CP`$ semigroup.
Let $`𝒟`$ be the domain of the generator of $`\varphi `$ and let $`𝒜`$ be the domain algebra
$$𝒜=\{A𝒟:A^{}A𝒟,AA^{}𝒟\}.$$
Theorem A implies that $`𝒜`$ is a unital $``$-algebra and we calculate its strong closure.
###### Proposition
The strong closure of $`𝒜`$ consists of all operators $`B(H)`$ such that $`B`$ commutes with the rank-one projection $`f\overline{f}`$.
Thus the strong closure of $`𝒜`$ consists of all operators $`B`$ such that both $`B`$ and $`B^{}`$ have $`f`$ as an eigenvector.
###### Demonstration proof
By Lemma 1, the domain $`𝒟`$ of the generator of $`\varphi `$ consists of all operators $`A`$ with the property
$$\varphi _t(A)AMt,\text{for all }t0,$$
$`2.1`$
where $`M`$ is a positive constant depending on $`A`$.
First, we show that $`ff`$ commutes with $`𝒜`$. Choose $`A𝒜`$. In order to show that $`A`$ commutes with $`f\overline{f}`$, it is enough to show that
$$\omega (A^{}A)=\omega (AA^{})=|\omega (A)|^2,$$
$`2.2`$
since (2.2) implies
$$Af\omega (A)f^2=\omega (A^{}A)2|\omega (A)|^2+|\omega (A)|^2=0,$$
and similarly $`A^{}f\omega (A^{})f=0`$. Multiplying $`\varphi _t(A)A`$ on the right by $`E_t`$ and using the fact that $`\varphi _t(A)E_t=\omega (A)E_t`$ we conclude from (2.1) that
$$\underset{t0}{lim}\omega (A)E_tAE_t=0.$$
Replacing $`A`$ with $`A^{}A`$ and $`AA^{}`$ one also finds
$$\underset{t0}{lim}\omega (A^{}A)E_tA^{}AE_t=\underset{t0}{lim}\omega (AA^{})E_tAA^{}E_t=0.$$
Taken together, these three limits imply that $`\omega (A^{}A)=\omega (AA^{})=|\omega (A)|^2`$, as required.
To prove the opposite inclusion it is enough to show that for every self-adjoint operator $`A(H)`$ satisfying $`Af=0`$ there is a sequence $`A_n`$ of self-adjoint operators in $`𝒜`$ which converges weakly to $`A`$ (recall that $`𝒜`$ is a self-adjoint algebra containing the identity). Fix such an $`A`$ and, for every $`ϵ>0`$, set
$$A_ϵ=\varphi _ϵ(A)=\omega (A)E_ϵ+U_ϵAU_ϵ^{}=U_ϵAU_ϵ^{}.$$
$`A_ϵ`$ converges weakly to $`A`$ as $`ϵ0`$. Moreover, $`A_ϵ`$ is supported in the interval $`(ϵ,\mathrm{})`$ in the sense that $`A_ϵE_ϵ=E_ϵA_ϵ=0`$, and in addition we have $`A_ϵf=0`$ since
$$A_ϵf=U_ϵAU_ϵ^{}f=e^{\alpha ϵ}U_ϵAf=0.$$
We show that each $`A_ϵ`$ can be weakly approximated by self-adjoint elements of the domain algebra.
###### Lemma
Suppose $`ϵ>0`$ and let $`A`$ be a self-adjoint operator in $`(H)`$ such that (i) $`Af=0`$ and (ii) $`A`$ is supported in $`(ϵ,\mathrm{})`$ in the sense that $`AE_ϵ=E_ϵA=0`$. Let $`u`$ be a $`C^{\mathrm{}}`$ function having compact support in $`[0,ϵ]`$ and consider
$$B=_0^{\mathrm{}}u(s)U_sAU_s^{}𝑑s=_0^{\mathrm{}}u(s)\varphi _s(A)𝑑s.$$
Then $`B^n𝒟`$ for every $`n=1,2,\mathrm{}`$, and in particular $`B𝒜`$.
###### Demonstration proof
Observe first that $`B`$ has both properties (i) and (ii), hence so does $`B^n`$ for every $`n`$. Thus for $`t<ϵ`$ we have
$$\varphi _t(B^n)B^n=U_tB^nU_t^{}B^n=U_tB^nU_t^{}B^nU_tU_t^{}=(U_tB^nB^nU_t)U_t^{}.$$
This implies that for sufficiently small $`t`$
$$\varphi _t(B^n)B^n=U_tB^nB^nU_t.$$
We conclude that $`B^n𝒟`$ iff there is a constant $`K>0`$ such that
$$U_tB^nB^nU_tKt,\text{for all }t>0.$$
$`2.3`$
To prove (2.3), one uses the Leibniz rule for the derivation $`D(X)=U_tXXU_t`$ to estimate $`U_tB^nB^nU_t`$ in terms of $`U_tBBU_t`$,
$$D(B^n)nB^{n1}D(B)=nB^{n1}U_tBBU_t.$$
Since $`B`$ has been smoothed it belongs to the domain $`𝒟`$, hence there is a constant $`M`$ such that $`U_tBBU_tMt`$, hence $`U_tB^nB^nU_tnMB^{n1}t`$.
The proof of the Proposition is completed by choosing $`A=A_ϵ`$ in the hypothesis of the Lemma and by choosing a sequence $`u_k`$ of nonnegative $`C^{\mathrm{}}`$ functions, each of which has integral $`1`$, such that $`u_k(x)=0`$ outside the interval $`0x1/k`$. A standard argument shows that the sequence of self-adjoint operators
$$B_k=_0^{\mathrm{}}u_k(s)\varphi _s(A_ϵ)𝑑s$$
converges weakly to $`A_ϵ`$, and the Lemma implies that $`B_k𝒜`$ for $`k>1/ϵ`$.
Thus the strong closure $`𝒜^{}`$ of $`𝒜`$ has the form $`(H_0)`$ where $`H_0H`$ is a subspace of codimension one in $`H`$, and the following implies that these examples are “almost” $`E_0`$-semigroups in the sense that there is an $`E_0`$-semigroup $`\alpha =\{\alpha _t:t0\}`$ acting on $`(H_0)`$ such that $`\varphi _t`$ acts as follows on $`𝒜^{}`$,
$$\varphi _t(B\lambda )=\alpha _t(B)\lambda ,B(H_0),\lambda .$$
###### Corollary
Let $`\overline{𝒜}`$ be the strong closure of $`𝒜`$. Then $`\varphi _t(\overline{𝒜})\overline{𝒜}`$ for every $`t0`$ and $`\{\varphi _t_{\overline{𝒜}}:t0\}`$ is a semigroup of endomorphisms of this von Neumann algebra.
###### Demonstration proof
We show that $`\varphi _t(𝒜)\overline{𝒜}`$, and for $`A,B𝒜`$ one has $`\varphi _t(AB)=\varphi _t(A)\varphi _t(B)`$.
Choose $`A𝒜`$, and let $`f`$ and $`\omega (A)=Af,f`$ be as in the definition of $`\varphi _t`$,
$$\varphi _t(A)=\omega (A)E_t+U_tAU_t^{},A𝒜,t0.$$
Since $`f`$ is an eigenvector for both $`A`$ and $`A^{}`$ and $`U_t^{}f=e^{\alpha t}f`$, one can verify directly that $`\varphi _t(A)f=\omega (A)f`$ and $`\varphi _t(A)^{}f=\omega (A^{})f`$, and the Proposition implies that $`\varphi _t(𝒜)\overline{𝒜}`$. Finally, for $`A`$, $`B𝒜`$ one has $`\omega (AB)=\omega (A)\omega (B)`$, and $`\varphi _t(AB)=\omega (AB)E_t+U_tABU_t^{}=\omega (A)\omega (B)E_t+U_tAU_t^{}U_tBU_t^{}=\varphi _t(A)\varphi _t(B)`$. By normality of $`\varphi _t`$, the formula $`\varphi _t(AB)=\varphi _t(A)\varphi _t(B)`$ persists for operators $`A`$,$`B`$ in the strong closure of $`𝒜`$.
## 3. The symbol of the generator: properties and structure
There are two useful characterizations of the generators of uniformly continuous $`CP`$-semigroups, i.e., those whose generators are everywhere defined bounded linear maps on $`(H)`$. The first is due to Lindblad and independently to Gorini et al (also see , Theorem 4.2). The second characterization is due to Evans and Lewis , based on work of Evans . These two results can be paraphrased as follows.
###### Theorem
Let $`L:(H)(H)`$ be a bounded linear map and let $`\varphi =\{\varphi _t:t0\}`$ be the semigroup defined on $`(H)`$ by $`\varphi _t=\mathrm{exp}(tL)`$. The following are equivalent.
A linear map $`L:(H)(H)`$ satisfying property (3) of Theorem 3.1 is called conditionally completely positive . While the characterization (2) tells us exactly which bounded linear maps generate $`CP`$ semigroups, the cited decomposition of $`L`$ into a sum of more familiar mappings is unfortunately not unique.
The purpose of this section is to make two observations. First, we point out that the notion of a conditionally completely positive linear map defined on a $``$-algebra is more properly formulated in terms of the bimodule of noncommutative 2-forms over that algebra; and once that is done the “symbol” of the map becomes analogous to a Riemannian metric. Second, we show that by making use of the domain algebra of section 1, this notion becomes appropriate for the generators of arbitrary $`CP`$-semigroups.
Let $`𝒜`$ be the domain algebra of a $`CP`$ semigroup $`\varphi =\{\varphi _t:t0\}`$ acting on $`(H)`$
$$𝒜=\{A𝒟:A^{}A𝒟,AA^{}𝒟\},$$
where $`𝒟`$ is the natural domain of the generator $`L`$ of $`\varphi `$. We first recall the definition of the module of noncommutative 1-forms $`\mathrm{\Omega }^1(𝒜)`$, and 2-forms $`\mathrm{\Omega }^2(𝒜)`$. The algebraic tensor product of vector spaces $`𝒜𝒜`$ can be considered an involutive bimodule over $`𝒜`$, with
$`a(xy)b`$ $`=axyb,`$
$`(xy)^{}`$ $`=y^{}x^{}.`$
The map $`d:𝒜𝒜𝒜`$ defined by $`dx=\mathrm{𝟙}xx\mathrm{𝟙}`$ is a derivation for which $`(dx)^{}=d(x^{})`$, and it is a universal derivation of $`𝒜`$ in the sense that if $`E`$ is any $`𝒜`$-bimodule and $`D:𝒜E`$ is a linear map satisfying $`D(xy)=xD(y)+D(x)y`$ for all $`x,y𝒜`$, then there is a unique homomorphism of $`𝒜`$-modules $`\theta :\mathrm{\Omega }^1(𝒜)E`$ such that $`\theta d=D`$. Every element of $`\mathrm{\Omega }^1(𝒜)`$ is a finite sum of the form
$$\omega =\underset{k=1}{\overset{r}{}}a_kdx_k,$$
and the involution in $`\mathrm{\Omega }^1(𝒜)`$ satisfies
$$(adx)^{}=d(x^{})a^{}=d(x^{}a^{})+x^{}d(a^{}).$$
Finally, $`\mathrm{\Omega }^1(𝒜)`$ is the kernel of the multiplication map $`\mu :𝒜𝒜𝒜`$ defined by $`\mu (xy)=xy`$, and thus we have an exact sequence of $`𝒜`$-modules
$$0\mathrm{\Omega }^1(𝒜)𝒜𝒜\underset{𝜇}{}𝒜0.$$
$`3.1`$
$`\mathrm{\Omega }^2(𝒜)`$ is defined by
$$\mathrm{\Omega }^2(𝒜)=\mathrm{\Omega }^1(𝒜)_𝒜\mathrm{\Omega }^1(𝒜),$$
and any element of $`\mathrm{\Omega }^2(𝒜)`$ can be written as a sum
$$\omega =\underset{k=1}{\overset{r}{}}a_kdx_kdy_k.$$
The involution in $`\mathrm{\Omega }^2(𝒜)`$ satisfies
$$(adxdy)^{}=d(y^{})d(x^{})a^{}=d(y^{})d(x^{}a^{})d(y^{}x^{})d(a^{})+y^{}d(x^{})d(a^{}).$$
Since $`𝒜`$ is a $``$-subalgebra of $`(H)`$, we may also think of $`(H)`$ as an $`𝒜`$-bimodule. Now a straightforward argument shows that for every linear mapping $`L:𝒜(H)`$ there is a unique homomorphism of bimodules $`\sigma _L:\mathrm{\Omega }^2(𝒜)(H)`$ which satisfies
$$\sigma _L(dxdy)=L(xy)xL(y)L(x)y+xL(\mathrm{𝟙})y,x,y𝒜$$
$`3.2`$
(see section 2 of for more detail). $`\sigma _L\mathrm{hom}(\mathrm{\Omega }^2,(H))`$ is called the symbol of the linear map $`L`$.
Consider now the special case in which $`𝒜=(H)`$, $`L:(H)(H)`$ is a bounded linear mapping, and let $`\varphi =\{\varphi _t=\mathrm{exp}tL:t0\}`$ is the semigroup of bounded operators on $`(H)`$ generated by $`L`$. The preceding theorem gives two characterizations of the maps $`L`$ for which each $`\varphi _t=\mathrm{exp}tL`$ is completely positive; however, the following characterization is perhaps more in spirit with the theory of differential operators on manifolds.
###### Theorem
Let $`L:(H)(H)`$ be a bounded linear map. To the two characterizations (2) (3) above, one can append the following equivalent condition
This characterization is Proposition 1.6 of ; a fuller discussion of these issues can be found in . Notice that the sense of the inequality $``$ is determined by the fact that the involution in $`\mathrm{\Omega }^1`$ satisfies $`(dx)^{}=d(x^{})`$, and hence for $`\omega =dx`$ we have $`\omega ^{}\omega =d(x^{})dx`$. In particular, for $`\omega =dx`$ where $`x`$ is a self-adjoint element we have $`\sigma _L(\omega ^2)0`$ while $`\sigma _L(\omega ^{}\omega )0`$.
###### Remark Remarks
There is a rather compelling analogy between this characterization of the generators of CP semigroups and the generator of the heat flow of a Riemannian manifold, namely the Laplacian. More precisely, let $`M`$ be a complete (but not necessarily compact) Riemannian manifold and consider its natural Hilbert space $`L^2(M)`$. The Laplacian $`\mathrm{\Delta }`$ acts naturally as a densely defined operator on $`L^2(M)`$ and generates a semigroup of bounded operators $`\mathrm{exp}t\mathrm{\Delta }`$, $`t0`$, acting on $`L^2(M)`$ (the book of Davies is a good reference). This semigroup maps bounded functions in $`L^2(M)`$ to bounded functions in $`L^2(M)`$, and the latter determines a semigroup of normal linear maps on the abelian von Neumann algebra $`L^{\mathrm{}}(M)`$ which carries nonnegative functions to nonnegative functions and fixes the constant functions.
In order to discuss the symbol of $`\mathrm{\Delta }`$ we introduce local coordinates in some open set $`UM`$ to identify $`U`$ with an open region in $`^n`$. For clarity, we will be explicit with notation. At each point $`xU`$ the tangent space $`T_x(M)`$ is identified with $`^n`$, and for a smooth function $`f`$ on $`M`$ the differential $`df`$ takes the following form
$$df(x,v)=\frac{d}{dt}f(x+tv)|_{t=0}=\underset{k=1}{\overset{n}{}}\frac{f}{x_k}(x)v_k.$$
The metric gives rise to a an operator function $`xUG(x)`$ on $`^n`$ by way of
$$v,w_{T_x(M)}=G(x)v,w_^n,v,wT_x(M),xU,$$
where $`,_^n`$ denotes the Euclidean inner product on $`^n`$. $`G(x)`$ is an invertible positive operator on $`^n`$ for every $`xU`$. For two vector fields $`\xi ,\eta `$ on $`M`$ we have
$$\xi (x),\eta (x)_{T_x(M)}=G(x)\xi (x),\eta (x)_^n=\underset{i,j=1}{\overset{n}{}}g_{ij}(x)\xi _j(x)\eta _i(x),$$
for $`xU`$, $`(g_{ij}(x))`$ being the matrix of $`G(x)`$ relative to the usual orthonormal basis for $`^n`$.
The inner product on the tangent space $`T_x(M)`$ promotes naturally to an inner product on the cotangent space $`T_x^{}(M)`$. Indeed, the Riesz lemma implies that every linear functional $`\rho `$ on $`T_x(M)`$ is associated with a unique vector $`\rho _{}T_x(M)`$ via
$$\rho (v)=v,\rho _{}_{T_x(M)},$$
and the inner product in $`T_x^{}(M)`$ is defined by
$$\rho ,\sigma _{T_x^{}(M)}=\rho _{},\sigma _{}_{T_x(M)}.$$
With these conventions one finds that for a smooth function $`f`$ and a point $`xU`$, $`df(x,)_{}`$ becomes the vector in $`^n`$ with components $`v_1,\mathrm{},v_n`$,
$$v_i=\underset{j=1}{\overset{n}{}}g^{ij}(x)\frac{f}{x_j}(x),$$
$`(g^{ij}(x))=(g_{ij}(x))^1`$ being the matrix of the inverse operator $`G(x)^1`$. For points $`xU`$ one has
$$(df)_{},(dg)_{}_{T_x(M)}=\underset{i,j=1}{\overset{n}{}}g^{ij}(x)\frac{f}{x_j}\frac{f}{x_i}.$$
$`3.3`$
We first recall that the dualized Riemannian metric (whose values are inner products on the cotangent spaces $`T_x^{}(M)`$) can be linearized naturally so that it becomes a $`C^{\mathrm{}}(M)`$-linear map of the the module $`\mathrm{\Omega }^{(2)}(M)`$ of symmetric 2-forms. More explicitly, let $`\mathrm{\Omega }^1(M)`$ be the usual module of $`1`$-forms and let $`\mathrm{\Omega }^{(2)}(M)`$ be the submodule of $`\mathrm{\Omega }^1(M)_{C^{\mathrm{}}(M)}\mathrm{\Omega }^1(M)`$ consisting of all elements that are fixed under the action of the reflection $`R`$ defined by $`R:\omega _1\omega _2\omega _2\omega _1`$. For $`\omega _1,\omega _2\mathrm{\Omega }^1(M)`$ we write $`\omega _1\omega _2`$ for the symmetrized product
$$\omega _1\omega _2=\frac{1}{2}(\omega _1\omega _2+\omega _2\omega _1)\mathrm{\Omega }^{(2)}(M).$$
There is a unique homomorphism of $`C^{\mathrm{}}(M)`$-modules $`G^{}:\mathrm{\Omega }^{(2)}(M)C^{\mathrm{}}(M)`$ satisfying $`G^{}(dfdg)(x)=df,dg_{T_x^{}(M)}`$ for all $`xM`$, and in local coordinates (3.3) implies that $`G^{}`$ has the form
$$G^{}(dfdg)(x)=\underset{i,j=1}{\overset{n}{}}g^{ij}(x)\frac{f}{x_j}\frac{g}{x_i},xU.$$
$`3.4`$
If one knows $`G^{}`$ as a homomorphism of $`C^{\mathrm{}}(M)`$-modules then one also knows the inner product in each cotangent space $`T_x^{}(M)`$, and hence one can recover the original metric as an inner product on tangent spaces by duality and the Riesz lemma as above.
We now relate these remarks to the symbol of the Laplacian $`\mathrm{\Delta }`$ of $`M`$. The symbol of any differential operator $`L:C^{\mathrm{}}(M)C^{\mathrm{}}(M)`$ of order at most 2 is associated with the bilinear form defined on $`C^{\mathrm{}}(M)`$ by
$$f,gC^{\mathrm{}}(M)L(fg)fL(g)gL(f)+fgL(\mathrm{𝟙}).$$
A straightforward argument shows that there is a (necessarily unique) homomorphism of $`C^{\mathrm{}}(M)`$-modules $`\sigma _L:\mathrm{\Omega }^{(2)}(M)C^{\mathrm{}}(M)`$ satisfying
$$\sigma _L(dfdg)=L(fg)fL(g)gL(f)+fgL(1).$$
In particular, this defines the symbol of any second order differential operator on $`C^{\mathrm{}}(M)`$), as an element of $`\mathrm{hom}(\mathrm{\Omega }^2(M),C^{\mathrm{}}(M))`$.
Restricting attention to the operator $`L=\mathrm{\Delta }`$, one sees that for each $`fC^{\mathrm{}}(M)`$ the restriction of $`\mathrm{\Delta }(f)`$ to $`U`$ has the form
$$\mathrm{\Delta }(f)(x)=\underset{i,j=1}{\overset{n}{}}g^{ij}(x)\frac{^2f}{x_ix_j}+\underset{k=1}{\overset{n}{}}u_k(x)\frac{f}{x_k},$$
$`3.5`$
where $`u_1,\mathrm{},u_n`$ are appropriate smooth functions (see p. 147 of ). Using (3.5) one easily computes the symbol of $`\mathrm{\Delta }`$, and because of the local formula (3.4) for $`G^{}`$ one obtains $`\sigma _\mathrm{\Delta }=2G^{}`$. From these remarks we conclude that the symbol of the Laplacian (considered as an element of $`\mathrm{hom}(\mathrm{\Omega }^{(2)}(M),C^{\mathrm{}}(M))`$) is precisely the Riemannian metric of $`M`$ in its dualized form.
Returning now to the case of a general CP semigroup $`\varphi =\{\varphi _t:t0\}`$ acting on $`(H)`$, let $`𝒜`$ be the domain algebra of the generator of $`\varphi `$. Letting $`L`$ be the restriction of the generator to $`𝒜`$, it is natural to ask the extent to which the generator can be identified with something analogous to a Riemannian metric (more precisely, to the homomorphism of $`C^{\mathrm{}}(M)`$-modules $`G^{}:\mathrm{\Omega }^{(2)}(M)C^{\mathrm{}}(M)`$ that the dualized Riemannian metric determines). We have already defined the symbol $`\sigma _L:\mathrm{\Omega }^2(𝒜)(H)`$ as a homomorphism of $`𝒜`$-modules, and the following asserts that $`\sigma _L`$ does behave as if it were a (perhaps degenerate) Riemannian metric.
###### Proposition
Let $`\varphi =\{\varphi _t:t0\}`$ be a $`CP`$ semigroup acting on $`(H)`$ and consider the restriction $`L`$ of the generator to the domain algebra $`L:𝒜(H)`$. Then the symbol of $`L`$ satisfies
$$\sigma _L(\omega ^{}\omega )0,\omega \mathrm{\Omega }^1(𝒜);$$
and more generally for all $`\xi _1,\mathrm{},\xi _nH`$ and $`\omega _1,\mathrm{},\omega _n\mathrm{\Omega }^1(𝒜)`$ we have
$$\underset{i,j=1}{\overset{n}{}}\sigma _L(\omega _j^{}\omega _i)\xi _i,\xi _j0.$$
The proof is a computation, facilitated by the following formula.
###### Lemma
Let $`\omega _1`$, $`\omega _2`$ be elements of $`\mathrm{\Omega }^1(𝒜)`$ having the form
$$\omega _k=\underset{p=1}{\overset{s}{}}A_{kp}B_{kp},k=1,2,$$
where $`A_{k1}B_{k1}+\mathrm{}+A_{ks}B_{ks}=0`$ for $`k=1,2`$. Then
$$\sigma _L(\omega _1\omega _2)=\underset{p,q=1}{\overset{s}{}}A_{1p}L(B_{1p}A_{2q})B_{2q}.$$
###### Demonstration proof of Lemma
Since $`\mathrm{\Omega }^1(𝒜)`$ is spanned by elements of the form $`AdX`$, as well as by elements of the form $`dYB`$, $`A,B,X,Y𝒜`$, and since $`\sigma _L`$ is a bimodule homomorphism, it suffices to check the formula for $`\omega _1`$, $`\omega _2`$ of the particular form $`\omega _1=dX`$, $`\omega _2=dY`$.
Writing
$`dXdY`$ $`=(X\mathrm{𝟙}\mathrm{𝟙}X)(Y\mathrm{𝟙}\mathrm{𝟙}Y)`$
$`=(X\mathrm{𝟙})(Y\mathrm{𝟙})(\mathrm{𝟙}X)(Y\mathrm{𝟙})(X\mathrm{𝟙})(\mathrm{𝟙}Y)+(\mathrm{𝟙}X)(\mathrm{𝟙}Y)`$
the right side of the asserted formula for $`\sigma _L(dXdY)`$ has the form
$``$ $`(XL(Y)L(XY)XL(\mathrm{𝟙})Y+L(X)Y)=`$
$`L(XY)XL(Y)L(X)Y+XL(\mathrm{𝟙})Y=\sigma _L(dXdY),`$
as required.
###### Demonstration proof of Proposition
Because of the exact sequence (3.1), every element $`\omega \mathrm{\Omega }^1(𝒜)`$ can be written
$$\omega =A_1B_1+\mathrm{}+A_sB_s,$$
where $`A_k,B_k`$ are elements of $`𝒜`$ satisfying $`A_1B_1+\mathrm{}+A_sB_s=0`$. Choose elements $`\omega _1,\mathrm{},\omega _n\mathrm{\Omega }^1(𝒜)`$ of the form
$$\omega _k=\underset{p=1}{\overset{s}{}}A_{kp}B_{kp},k=1,\mathrm{},n$$
where $`_pA_{kp}B_{kp}=0`$ for $`k=1,\mathrm{},n`$. We have
$$\omega _k^{}=\underset{p=1}{\overset{n}{}}B_{kp}^{}A_{kp}^{}$$
so that the product $`\omega _k^{}\omega _j\mathrm{\Omega }^2(𝒜)`$ is given by
$$\omega _k^{}\omega _j=\underset{p,q=1}{\overset{s}{}}(B_{kp}^{}A_{kp}^{})(A_{jq}B_{jq}).$$
The Lemma implies that
$$\sigma _L(\omega _k^{}\omega _j)=\underset{p,q=1}{\overset{s}{}}B_{kp}^{}L(A_{kp}^{}A_{jq})B_{jq}.$$
$`3.6`$
If we now choose vectors $`\xi _kH`$, $`k=1,\mathrm{},n`$ then we find that
$`{\displaystyle \underset{k,j}{}}\sigma _L(\omega _k^{}\omega _j)\xi _j,\xi _k`$ $`={\displaystyle \underset{i,j,p,q}{}}L(A_{kp}^{}A_{jq})B_{jq}\xi _j,B_{kp}\xi _k`$ $`3.7`$
$`={\displaystyle \underset{\alpha ,\beta }{}}L(A_\beta ^{}A_\alpha )\eta _\alpha ,\eta _\beta ,`$
where in the third term, $`\alpha `$ and $`\beta `$ run over all pairs $`\{(k,p):1kn,1ps\}`$ and where the $`\eta _\alpha `$ are defined by $`\eta _{(k,p)}=B_{kp}\xi _k`$. Finally, the last term on the right of (3.7) can be rewritten in terms of the $`ns\times ns`$ operator matrix $`\stackrel{~}{A}`$ having the entries $`A_\alpha `$ along a single row and zeros along all the other rows as follows
$$L^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\stackrel{~}{\eta },\stackrel{~}{\eta },$$
where $`\stackrel{~}{\eta }`$ is the column vector with components $`\eta _\alpha `$, and where $`L^{(ns)}`$ is the natural map induced by $`L`$ on matrices over $`𝒜`$ be applying $`L`$ to the elements of the matrix term-by-term.
Thus we have to show that $`L^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\stackrel{~}{\eta },\stackrel{~}{\eta }0`$. Recalling that the definition of $`L`$ on elements of $`𝒜`$ is
$$L(X)=\underset{t0+}{lim}\frac{1}{t}(\varphi _t(X)X)$$
and the fact that $`𝒜`$ is a $``$-subalgebra of the domain of $`L`$, it follows that
$$L^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\stackrel{~}{\eta },\stackrel{~}{\eta }=\underset{t0+}{lim}\frac{1}{t}(\varphi _t^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\stackrel{~}{\eta },\stackrel{~}{\eta }\stackrel{~}{A}^{}\stackrel{~}{A}\stackrel{~}{\eta },\stackrel{~}{\eta }).$$
Notice that $`\stackrel{~}{A}^{}\stackrel{~}{A}\stackrel{~}{\eta },\stackrel{~}{\eta }=0`$. Indeed, by inspection of the components of the column vector $`\stackrel{~}{A}\stackrel{~}{\eta }`$ we find that it is the column vector having a single (possibly) nonzero component and that component is
$$\underset{\alpha }{}A_\alpha \eta _\alpha =\underset{k,p}{}A_{kp}B_{kp}\xi _k=\underset{k=1}{\overset{n}{}}(\underset{p=1}{\overset{s}{}}A_{kp}B_{kp})\xi _k=0,$$
since $`_pA_{kp}B_{kp}=0`$ for every $`k`$. Thus we have to show that
$$\underset{t0+}{lim}\frac{1}{t}\varphi _t^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\stackrel{~}{\eta },\stackrel{~}{\eta }0.$$
$`3.8`$
Now since for each $`t>0`$ the map $`\varphi _t`$ is unital and completely positive, the Schwarz inequality for completely positive maps implies
$$\varphi _t^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\varphi _t^{(ns)}(\stackrel{~}{A})^{}\varphi _t^{(ns)}(\stackrel{~}{A}),$$
and hence for positive $`t`$ we have
$$\frac{1}{t}\varphi _t^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\stackrel{~}{\eta },\eta \frac{1}{t}\varphi _t^{(ns)}(\stackrel{~}{A})\eta ,\varphi _t^{(ns)}(\stackrel{~}{A})\eta =\frac{1}{t}\varphi _t^{(ns)}(\stackrel{~}{A})\eta ^2.$$
We claim that the term on the right tends to zero as $`t0+`$. Indeed, since the operator matrix $`\stackrel{~}{A}`$ belongs to the domain of the generator of the CP semigroup $`\varphi ^{(ns)}=\{\varphi _t^{(ns)}:t0\}`$, Lemma 1 implies that there is a constant $`M>0`$ such that for every positive $`t`$, $`\varphi _t^{(ns)}(\stackrel{~}{A})\stackrel{~}{A}Mt`$. It follows that
$$\varphi _t^{(ns)}(\stackrel{~}{A})\stackrel{~}{\eta }=\varphi _t^{(ns)}(\stackrel{~}{A})\stackrel{~}{\eta }\stackrel{~}{A}\eta Mt\stackrel{~}{\eta }$$
and hence
$$\underset{t0+}{lim\; sup}\frac{1}{t}\varphi _t^{(ns)}(\stackrel{~}{A})\stackrel{~}{\eta }^2\underset{t0+}{lim}\frac{1}{t}(M^2t^2\stackrel{~}{\eta }^2)=0.$$
It follows that
$$\underset{t0+}{lim}\frac{1}{t}\varphi _t^{(ns)}(\stackrel{~}{A}^{}\stackrel{~}{A})\stackrel{~}{\eta },\stackrel{~}{\eta }\underset{t0+}{lim}\frac{1}{t}\varphi _t^{(ns)}(\stackrel{~}{A})\eta ^2=0,$$
and the inequality (3.8) follows.
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# Skewness as a probe of non-Gaussian initial conditions
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## Abstract
We compute the skewness of the matter distribution arising from non–linear evolution and from non–Gaussian initial perturbations. We apply our result to a very generic class of models with non–Gaussian initial conditions and we estimate analytically the ratio between the skewness due to non-linear clustering and the part due to the intrinsic non-Gaussianity of the models. We finally extend our estimates to higher moments.
preprint: UGVA-DPT 1999/11-zzzz astro-ph/0005087
\] The source of the initial density fluctuations which have led to the formation of structure, observed in the Universe today is unknown. Determining its nature will certainly be of utmost importance for the fruitful relation between high energy physics and cosmology.
In models which presently attract most attention, initial density fluctuations are generated during an inflationary phase. In the simplest inflationary models, the initial fluctuations obey Gaussian statistics. If this picture is correct, the deviations from Gaussianity we observe today were induced by nonlinear gravitational instability . However, it is also conceivable that the present deviations from Gaussianity have two components: gravitationally induced and intrinsic, coming from the initial conditions rather than nonlinear dynamics . Here, we investigate to what extent an intrinsic component can be ‘washed out’ by nonlinear dynamics and on which scales it could be either detected or constrained from above in future galaxy surveys.
We start by deriving a general expression for the so-called skewness parameter, $`S_3`$, including the effect of an initial non–Gaussianity, non–linear evolution and smoothing. We then estimate the normalized $`N`$–point cumulant, $`S_N`$, for a wide class of models and compare it with the result obtained in Gaussian models due to mild non–linearities.
If the galaxies trace the spatial mass distribution, galaxy surveys can be used to estimate the cumulants of the mass density contrast field, given by
$$M_N(R)(\delta _R)^N(𝐱,\eta _0)_c$$
(1)
of the smoothed density field $`\delta _R(𝐱,\eta )d^3𝐱^{}W_R(|𝐱𝐱^{}|)\delta (𝐱^{},\eta )`$, where $`\delta (𝐱,\eta )`$ is the density field, $`\eta `$ and $`\eta _0`$ the conformal time and its value today, and $`W_R`$ is a window function (e.g. Gaussian or top–hat) of width $`R`$. The brackets in (1) denote an ensemble average and the subscript $`c`$ indicates that we deal with the connected part of the $`N`$–point function. For a Gaussian field, all cumulants of order $`N>2`$ vanish: $`M_N=0`$. $`M_2`$ is the variance while $`M_3`$ is a measure of the asymmetry of the distribution, known as skewness. We will also use the more common normalized cumulant,
$`S_N(R)=M_N(R)/(M_2(R))^{(N1)}.`$
This ratio is constant (independent of $`R`$) in the weakly non-linear regime . To calculate the general expression for $`M_3(R)`$ in the weakly nonlinear regime, we follow the method developed in . Expanding $`\delta (𝐱,\eta )`$ in a perturbative series, $`\delta _1+\delta _2+𝒪(3)`$ and solving the system of coupled Euler, Poisson and continuity equations at second order leads, in Fourier space, to $`\mathrm{\Delta }_1(\eta ,𝐤)=D(\eta ,𝐤)`$ and
$$\mathrm{\Delta }_2(\eta ,𝐤)=(2\pi )^{3/2}d^3𝐪J(𝐪,𝐤𝐪)D(\eta ,𝐪)D(\eta ,𝐤𝐪)$$
where we consider only the fastest growing modes and we use the convention
$$\mathrm{\Delta }_N(\eta ,𝐤)=(2\pi )^{3/2}\delta _N(\eta ,𝐱)\mathrm{e}^{i𝐤𝐱}d^3𝐱.$$
At late times where a possible source term or seed has decayed, the time and space dependence of the function $`D`$ can be factorized, $`D(\eta ,𝐤)=D_+(\eta )\epsilon (𝐤)`$, where $`D_+`$ is the standard linear growing mode . Perturbation theory gives
$$J(𝐤,𝐪)=\frac{2}{3}(1+\kappa )+(q/k)P_1(\mu )+\frac{2}{3}\left(\frac{1}{2}\kappa \right)P_2(\mu ),$$
(2)
where the $`P_{\mathrm{}}`$ is the Legendre polynomial of order $`\mathrm{}`$, $`\mu 𝐤𝐪/kq`$. The quantity $`\kappa `$ is a weak function of $`\mathrm{\Omega }`$; for $`\mathrm{\Omega }>0.01`$, $`\kappa (3/14)\mathrm{\Omega }^{0.03}`$ . The smoothing applies order by order. In Fourier space, we have $`\mathrm{\Delta }_R(\eta ,𝐤)=D(\eta ,𝐤)W_k,`$ $`W_k`$ being the Fourier transform of the window function. To fifth order, the skewness is
$`M_3`$ $`=`$ $`\delta _{R,1}^3+3\delta _{R,1}^2\delta _{R,2}+𝒪(5).`$ (3)
We introduce the two–,three– and four–point power spectra as $`12𝒫_2(k_1)\delta (𝐤_1+𝐤_2)`$,
$`123𝒫_3(𝐤_1,𝐤_2)\delta (𝐤_1+𝐤_2+𝐤_3),`$ (4)
$`1234_c𝒫_4(𝐤_1,𝐤_2,𝐤_3)\delta (𝐤_1+𝐤_2+𝐤_3+𝐤_4).`$ (5)
(The Dirac $`\delta `$ is a simple consequence of statistical homogeneity which we assume throughout.) Here $`12\mathrm{}ND(\eta ,𝐤_1)D(\eta ,𝐤_2)\mathrm{}D(\eta ,𝐤_N)`$. The functions $`𝒫_2`$ and $`𝒫_3`$ are also known as the power spectrum and the bispectrum, respectively. Inserting the Fourier transforms of $`\delta _1`$ and $`\delta _2`$ after smoothing in (3), expressing the correlators of $`D`$ in terms of the power spectra (5) and performing one integration using the Dirac function in (5), we obtain
$`M_3(R)={\displaystyle \frac{d^3𝐤d^3𝐪}{(2\pi )^6}𝒫_3(𝐤,𝐪)W_kW_qW_{|𝐤+𝐪|}}`$ (8)
$`+{\displaystyle \frac{d^3𝐤d^3𝐪}{(2\pi )^6}𝒫_2(k)𝒫_2(q)W_kW_qW_{|𝐤+𝐪|}J(𝐤,𝐪)}+`$
$`{\displaystyle \frac{d^3𝐤d^3𝐪d^3𝐩}{(2\pi )^6}𝒫_4(𝐤,𝐪𝐤,𝐩)W_qW_pW_{|𝐪+𝐩|}J(𝐤,𝐪𝐤)}`$
For a Gaussian field, $`𝒫_4=𝒫_3=0`$ and the only non-vanishing contribution comes from the second term in the above expression. For a top hat window, this term gives $`M_3=(34/7\gamma )M_2^2`$, with $`\gamma =d\mathrm{log}M_2(R)/d\mathrm{log}R`$ . Note also that $`\gamma (R)`$ is the logarithmic slope of the two-point correlation function of the density fluctuations - the Fourier transform of $`𝒫_2(k)`$. It is usually assumed that $`\gamma >0`$ (condition of hierarchical clustering, see e.g. ).
The class of models we want to analyze are those where fluctuations in the dark matter are induced by the energy and momentum of an inhomogeneously distributed component which contributes only a small fraction to the total energy momentum tensor and which interacts only gravitationally with the cosmic fluid. Such a component is denoted as ‘seed’ . As stressed above, we need to compute the $`N`$–point power spectra of the density field at the end of the linear regime. The comoving linear density fluctuation $`D`$ of the cosmic matter–radiation fluid evolves according to
$`\ddot{D}+H\left(16w+3c_s^2\right)\dot{D}+k^2c_s^2D`$ (9)
$`\frac{3}{2}\left(1+8w3w^26c_s^2\right)H^2D=\mathrm{SS}(𝐤,\eta ),`$ (10)
with $`\mathrm{SS}(1+w)4\pi G(f_\rho +3f_P)`$, $`f_\rho `$ and $`f_P`$ being the inhomogeneous energy density and pressure of the seeds. When the seed is a scalar field $`\varphi `$ with vanishing potential, $`f_\rho +3f_P=\dot{\varphi }^2`$. $`G`$ is Newton’s constant, $`a`$ denotes the cosmic scale factor, a dot refers to the derivative with respect to conformal time, $`H\dot{a}/a`$, $`wP/\rho `$ and $`c_s^2\dot{P}/\dot{\rho }`$ are respectively the enthalpy and the adiabatic sound speed of the cosmic fluid.
Equation (9) can be solved by a Green’s function, $`𝒢`$,
$$D(𝐤,\eta )=_{\eta _i}^\eta 𝒢(𝐤,\eta ,\eta ^{})\mathrm{SS}(𝐤,\eta ^{})𝑑\eta ^{},$$
(11)
where $`\eta _i`$ is some early initial time deep in the radiation era. For the linear part of the reduced $`N`$–point function we then obtain
$`D(𝐤_1,\eta )\mathrm{}D(𝐤_N,\eta )_c={\displaystyle _{\eta _i}^\eta }𝑑\eta _1\mathrm{}𝑑\eta _N`$ (12)
$`𝒢(𝐤_1,\eta ,\eta _1)\mathrm{}𝒢(𝐤_N,\eta ,\eta _N)\mathrm{SS}(𝐤_1,\eta _1)\mathrm{}\mathrm{SS}(𝐤_N,\eta _N)_c.`$ (13)
We define the connected $`N`$–point function of the source by
$`\mathrm{SS}(1)\mathrm{}\mathrm{SS}(N)_cF_N(𝐤_1,\mathrm{}𝐤_N;\eta _1\mathrm{}\eta _N)\delta \left({\displaystyle 𝐤_i}\right),`$
where $`(i)(𝐤_i,\eta _i)`$. Again, the $`\delta `$ function of the sum of all momenta is a consequence of the statistical homogeneity.
We now assume that the reduced $`N`$-point function of the source can be replaced by its ‘perfectly coherent approximation’ given by
$`F_N(𝐤_1,\mathrm{},𝐤_{N1};\eta _1,\mathrm{},\eta _N)\mathrm{sign}(F_N)\times `$ (14)
$`\sqrt[N]{|F_N(𝐤_1,\mathrm{},𝐤_{N1};\eta _1,\mathrm{},\eta _1)\mathrm{}F_N(𝐤_1,\mathrm{},𝐤_{N1};\eta _N,\mathrm{},\eta _N)|}`$ (15)
(here and below, $`𝐤_N`$ is always given by $`𝐤_N=(𝐤_1+\mathrm{}+𝐤_{N1})`$). This approximation is exact if the evolution equation for $`\mathrm{SS}`$ is linear and the randomness is entirely due to initial conditions. Then the source term is of the form $`\mathrm{SS}(𝐤,\eta )=R(𝐤)s(k,\eta )`$, where only $`R`$ is a random variable and $`s`$ is a deterministic solution to the linear evolution equation of $`\mathrm{SS}`$ which can be taken out of the average $``$. This is the key property which renders the $`N`$-point function decoherent. Then $`F_N`$ can be written as
$`F_N(𝐤_1,\mathrm{},𝐤_{N1};\eta _1,\mathrm{},\eta _N)`$ (16)
$`s(1)\mathrm{}s(N)R(𝐤_1)\mathrm{}R(𝐤_N)_c`$ (17)
which is clearly of the form (15).
An important example are models with no sources but with non-Gaussian initial conditions for $`D`$. Such models, like e.g. the recent $`\chi ^2`$ Peebles model , are always perfectly coherent and therefore included in our analysis: In this case $`D(𝐤,\eta )=R(𝐤)d(k,\eta )`$, where $`R`$ is a non-Gaussian random variable given by the initial condition and $`d`$ is a deterministic homogeneous solution of Eq. (9). Clearly, if we choose $`\mathrm{SS}(𝐤,\eta )=R(𝐤)\delta (\eta \eta _{in})`$ and $`𝒢(k,\eta ,\eta ^{})=d(k,\eta )`$, $`D`$ is of the form (11). Therefore, models where the non-Gaussianity is purely due to initial conditions are always perfectly coherent. As the equation of motion for $`D`$ is second order, the homogeneous solution has in principle two modes, $`D=R_1(𝐤)d_1(k,\eta )+R_2(𝐤)d_2(k,\eta )`$, but since we shall evaluate the $`N`$-point functions deeply in the matter era, the decaying mode will have disappeared and may thus be neglected in our analysis.
Models where the source term is due to a scalar field which evolves linearly in time are not perfectly coherent, since $`\mathrm{SS}`$ is given by the components of the energy momentum tensor which are quadratic in the fields. Numerical calculations, however, have shown that this non-linearity is not severe and perfect coherence is a relatively good approximation . One example of this kind are axionic seeds in pre-big bang cosmology for which decoherence has been tested and is found to be on the level of less than 5% for the CMB power spectrum. In Fig. 1 the functions $`D_2(k,\eta )`$ and $`D_3(k,k,\eta )`$ as obtained by a full numerical calculation are compared to their coherent approximation (15) for the large-N limit of global $`O(N)`$ symmetric scalar fields. This is another example where the scalar field evolution is linear and the only non-linearity in the source term is due to the energy momentum tensor being quadratic in the field .
For topological defects, especially for cosmic strings, the perfectly coherent approximation misses several important features (like the ‘smearing out’ of secondary acoustic peaks). However, we believe that our generic scaling result holds also in this case, as is indicated by numerical simulations of global texture: even though global texture show considerable decoherence , the same scaling law for higher moments which we derive here has been discovered numerically .
Under the perfectly coherent approximation Eq. (13) can be factorized as the product of the $`N`$ solutions, $`D_{Nj}(𝐤_1,\mathrm{},𝐤_{N1},\eta )`$ of the equations (9) with source term $`[F_N(𝐤_1,\mathrm{},𝐤_{N1};\eta ,\mathrm{},\eta )]^{1/N}`$, where $`𝐤_j`$ is the wave number $`𝐤`$ appearing in the term $`c_s^2k^2`$ on the left hand side of (9) and the other wave numbers have to be considered like parameters of the source term,
$`D(𝐤_1,\eta )\mathrm{}D(𝐤_N,\eta )_c`$ (18)
$`\left[{\displaystyle \underset{j=1}{\overset{N}{}}}D_{Nj}(𝐤_1,\mathrm{},𝐤_{N1},\eta )\right]\delta ({\displaystyle 𝐤_i})`$ (19)
$`𝒫_N(𝐤_1,\mathrm{},𝐤_{N1},\eta )\delta ({\displaystyle 𝐤_i}).`$ (20)
To continue, we assume that $`F_N`$ is a simple power law in the $`k_i`$ on super-Hubble scales and that it decays after Hubble crossing. This behavior is certainly correct for all examples discussed in the literature so far. We can then make the following ansatz
$$F_N\{\begin{array}{cc}_{n=1}^N\frac{k_n^\alpha }{k_0^\alpha }(f(\eta )\eta )^N\eta ^3\hfill & \text{if }k_i\eta 1,i\{1,\mathrm{},N\}\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$
(21)
Here $`f`$ is a dimensionless function and $`k_0`$ is an arbitrary scale. For scale invariant seeds (e.g. topological defects) $`f`$ is just a constant and $`\alpha =0`$. For axion seeds generated during a pre-big bang phase, $`\alpha `$ depends on the spectral index of the axion field, which in turn is determined by the evolution law of the extra dimension . For the Peebles model $`\alpha `$ is given by the power spectrum of the scalar field $`\varphi `$ and $`f`$ is a delta-function. Since $`F_N`$ is symmetrical in the variables $`𝐤_j`$ we can order them such that $`k_1k_2\mathrm{}k_N`$.
Let us discuss the temporal behavior of the variables $`D_{Nj}`$. As long as $`k_1\eta <1`$, the term $`c_s^2k_j^2D`$ can be neglected in Eq. (9) and the Green’s function is a power law. At $`k_1\eta 1`$ the source term decays and as long as a perturbation remains super horizon, it just grows like $`\eta ^2`$, so that for $`k_j\eta <1<k_1\eta `$,
$`D_{Nj}g(1/k_1)k_1^{2+3/N}(\eta k_1)^2\mathrm{\Pi }_{n=1}^N(k_n/k_0)^\alpha `$
where
$`g(\eta )={\displaystyle \frac{4\pi G}{\eta ^{23/N}}}{\displaystyle _{\eta _{in}}^\eta }G(\eta ,\eta ^{})f(\eta ^{})\eta ^{(23/N)}{\displaystyle \frac{d\eta ^{}}{\eta ^{}}},`$
and we have to take the part of the integral above which converges when $`\eta _{in}0`$.
Once the perturbation enters the horizon it either starts oscillating with roughly constant amplitude or continues to grow $`\eta ^2`$, depending on whether $`k_j`$ enters during the radiation or matter dominated era. At late time, $`\eta \eta _{eq}`$ and $`k\eta 1`$, we therefore obtain
$`D_{Nj}`$ $``$ $`g(1/k_1)k_1^{2+3/N}(k_1/k_j)^2`$
$`{\displaystyle \underset{n=1}{\overset{N}{}}}(k_n/k_0)^\alpha \{\begin{array}{cc}\left(\frac{\eta }{\eta _{eq}}\right)^2\hfill & \text{if }k_j\eta _{eq}>1\hfill \\ (\eta k_j)^2\hfill & \text{if }k_j\eta _{eq}<1\hfill \end{array}`$
where $`\eta _{eq}`$ is the time of equality between the matter and radiation densities. Defining $`0j_{eq}N`$ so that $`k_j\eta _{eq}>1`$ for all $`jj_{eq}`$ we obtain for the connected $`N`$-point function
$`𝒫_N(𝐤_1,\mathrm{},𝐤_{N1},\eta )`$ $``$ $`g(1/k_1)^Nk_1^3\eta ^{2N}`$ (24)
$`{\displaystyle \underset{n=1}{\overset{N}{}}}\left({\displaystyle \frac{k_n}{k_0}}\right)^\alpha {\displaystyle \underset{j=1}{\overset{j_{eq}}{}}}\left({\displaystyle \frac{1}{k_j\eta _{eq}}}\right)^2`$
Using this result for the ordinary power spectrum, $`𝒫_2`$, we can express $`𝒫_N`$ is terms of products of $`𝒫_2`$ as
$`𝒫_N(𝐤_1,\mathrm{},𝐤_{N1},\eta )`$ (25)
$`k_1^{3(1N/2)}{\displaystyle \underset{j=1}{\overset{N}{}}}\left(\sqrt{𝒫_2(k_j,\eta )}{\displaystyle \frac{g(1/k_1)k_j^{3/2}}{g(1/k_j)k_1^{3/2}}}\right).`$ (26)
For the class of models considered and under the assumption of perfect coherence, we have determined the connected $`N`$–point power spectra in the linear regime which are the input of the skewness (8).
$`M_3`$ has two contributions: A linear one due to the initial non–Gaussianity (contained in $`𝒫_3`$) and one due to non–linear clustering which induces skewness even in an originally Gaussian distribution of perturbations; it contains a Gaussian part ($`𝒫_2^2`$) and a non-Gaussian term ($`𝒫_4`$). We decompose the skewness as
$$M_3=M_3^{(L)}+M_3^{(NL)}$$
We want to estimate the ratio of these two contributions. Under our approximation (26), the first term of (8) reduces to
$`M_3^{(L)}={\displaystyle \frac{d^3𝐤d^3𝐤^{}}{(2\pi )^6}W_kW_qW_{|𝐤+𝐤^{}|}\sqrt{𝒫_2(k)𝒫_2(q)𝒫_2(|𝐤+𝐪|)}}`$ (27)
$`k_{\mathrm{max}}^{3/2}\left[{\displaystyle \frac{g(1/k_{\mathrm{max}})^3(kq|𝐤+𝐪|)^{3/2}}{g(1/k)g(1/q)g(1/|𝐤+𝐪|)k_{\mathrm{max}}^{9/2}}}\right],`$ (28)
where $`k_{\mathrm{max}}\mathrm{max}\{k,q,|𝐤+𝐪|\}`$. $`M_3^{(NL)}`$ is given by the second and third terms in (8).
To estimate analytically the ratio $`M_3^{(L)}/M_3^{(NL)}=S_3^{(L)}/S_3^{(NL)}`$, we make the following approximations:
* We assume that $`𝒫_2`$ is a simple power law within the range of scales of interest, namely all the modes which enter the horizon during the radiation era, this is $`0.1h^1`$Mpc$`2\pi /k20h^2`$Mpc, name. $`𝒫_2(k)=k^3(k/k_{})^\gamma .`$
* We also assume that $`g(\eta )\eta ^r`$.
* We replace the window function by a simple cut–off at $`k=1/R`$.
* For symmetry reasons we may integrate over the triangle $`qkR`$ and then multiply the result by 2.
* Since in our integration region, $`qk`$, we replace $`|𝐤+𝐪|`$ by $`k`$.
With these approximations the angular dependence of the integrand disappears and the integrals over $`𝐤`$ and $`𝐪`$ in (8) can be trivially performed leading to
$`M_3^{(L)}(R)`$ $``$ $`{\displaystyle \frac{4(k_{}R)^{3\gamma /2}}{(2\pi )^43\gamma (3+\gamma /2+r)}}`$ (29)
for $`\gamma >0\text{ and }3+\gamma /2+r>0`$ (30)
$`M_3^{(NL)}(R)`$ $``$ $`{\displaystyle \frac{(k_{}R)^{2\gamma }}{(2\pi )^4\gamma ^2}}\text{for}\gamma >0,`$ (31)
where we have just considered the Gaussian contribution, $`𝒫_2^2`$ to $`M_3^{(NL)}`$.
Since $`k_{}`$ is just the scale beyond which the density contrast $`D(𝐱)^2_{R=1/k}𝒫_2(k)k^3`$ is larger than unity and non-linearities become important, we define the non–linearity scale $`R_{\mathrm{lin}}=1/k_{}`$. The ratio between the skewness due to the non–Gaussianity in the linear perturbation and the one due to dynamical nonlinearities is then
$$\frac{S_3^{(L)}}{S_3^{(NL)}}\frac{4\gamma }{3(3+\gamma /2r)}\left(\frac{R}{R_{\mathrm{lin}}}\right)^{\gamma /2}.$$
(32)
This is our main result. It is readily checked that the non-Gaussian contribution, $`𝒫_4`$, to $`M_3^{(NL)}`$ behaves just like the contribution $`M_3^{(NL)}`$ and thus only modifies the pre-factor in (32), which should not be taken too seriously in view of the relatively crude approximations which we have employed to obtain our result.
This computation of the skewness is easily generalized to higher moments. As our computation shows, linear non–Gaussianities scale like
$$M_N^{(L)}(R)(R/R_{\mathrm{lin}})^{N\gamma /2}.$$
(33)
The dominant non–linear contribution to the connected $`N`$-point function which is also present in Gaussian theories contains $`N2`$ second order terms $`D_2`$ and therefore scales like
$$M_N^{(NL,\mathrm{Gauss})}(R)(R/R_{\mathrm{lin}})^{(N1)\gamma }.$$
(34)
The lowest order non-linearity for a generic non-Gaussian model, however just comes from the non-Gaussian term with $`N+1`$ factors of $`D`$. The non-Gaussian non-linear corrections therefore generically scale like
$$M_N^{(NL,\mathrm{noGauss})}(R)(R/R_{\mathrm{lin}})^{(N+1)\gamma /2}.$$
(35)
Only for $`N=3`$ the two terms (34) and (35) scale in the same way. For all higher $`N`$’s the non-Gaussian contribution dominates in the mildly non-linear regime, $`RR_{\mathrm{lin}}`$. From Eq. (35) we infer that in on large scales the ratios for all reduced $`N`$-point functions very generically scale like
$$\frac{S_N^{(L)}(R)}{S_N^{(NL)}(R)}\left(\frac{R}{R_{\mathrm{lin}}}\right)^{\gamma /2}.$$
(36)
This expression agrees with other analytic predictions as well as numerical simulations in a global texture model . The agreement with the texture simulations which are decoherent suggests that the validity of our result extends beyond the conditions under which Eq. (36) was derived. More important than decoherence is that the source term decays at late times and therefore the density perturbations just evolve according to the homogeneous solution. This implies that at late times the $`N`$-point functions behave like the homogeneous growing mode to the $`N`$th power, while the reduced $`N`$-point function induced by non-linear clustering from Gaussian perturbations scales like the growing mode to the $`2(N1)`$th power. Since topological defect sources decay on sub-horizon scales, we conclude that the derived scaling behavior is also valid for them (this argument will be expanded in our follow up publication ).
Our result implies that on small scales ($`RR_{\mathrm{lin}}`$), the dominant contribution to the cumulants comes from nonlinear Newtonian gravitational clustering, and the Gaussian term actually dominates. Intrinsic deviations from Gaussianity are difficult to detect on small scales. Hence, we should look for signs of intrinsic non-Gaussianity at large scales ($`R>R_{\mathrm{lin}}`$). This suggestion was expressed earlier based on qualitative physical arguments ; however, our present result is derived from first principles for a specific class of initial conditions – coherent seeds.
If galaxies trace mass, the measurements of the two-point correlation function suggest $`R_{\mathrm{lin}}10h^1`$Mpc and $`\gamma (R)1.8`$ for $`10\mathrm{k}\mathrm{p}\mathrm{c}hR15`$ Mpc (here $`h`$ is the usual parameterization for the Hubble constant in units of 100 km s<sup>-1</sup>Mpc<sup>-1</sup>); the slope $`\gamma `$ becomes steeper at larger separations $`R`$ . A frequently considered theoretical possibility for long-wave tail of the initial $`𝒫_2(k)`$, called the Zel’dovich-Harrison spectrum, would give $`\gamma =4`$ at large separations. Hence, we can expect all $`S_N`$s to “blow up” with increasing scale for the class of non-Gaussian models considered here, in contrast with models with Gaussian initial conditions. The available measurements of $`S_3(R)`$ and $`S_4(R)`$ do not show such a rise with scale and have already been used to constrain texture models . Likewise, there are indications that the existent data from the APM Galaxy Survey may may already extend to sufficiently large scales to constrain the $`\chi ^2`$ Peebles model . With surveys presently underway like the Sloan Digital Sky Survey , the prospects for using the approach outlined here to probe the statistics of the cosmological initial conditions will become even better.
In this work we derived a scaling law for the “intrinsic to induced” skewness ratio (32) for coherent seeds. We also showed how to generalize this law to higher cumulants. We plan to follow these calculations with more detailed predictions for coherent seed models and to confront our analytic results with numerical simulations as well as observational data from galaxy surveys . Let us also repeat that the derived scaling laws seem to be more general than their derivation as they have been obtained numerically for global texture which are decoherent seeds. We actually believe that the origin of the scaling laws is not coherence but mainly the decay of the sources at late time and we therefore conjecture that they hold also for topological defects.
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# Existential Contextuality and the Models of Meyer, Kent and Clifton
## I Introduction
Meyer Meyer , Kent KentA and Clifton and Kent KentB (MKC in the sequel) have recently proposed a new class of hidden variables models in which values are only assigned to a restricted subset of the set of all observables. MKC claim that their models “nullify” the Kochen-Specker theorem Bell ; Koch ; MerminB ; Peres ; Bub . In Appleby self2 we showed that this claim is unfounded: the MKC models do not, in any way, invalidate the essential physical point of the Kochen-Specker theorem (for other critical discussions see refs. Cabel ; Havli ; Mermin ; self0 ; MKCextra ). Nevertheless, the MKC models are still of much interest. Together with the models proposed by Pitowsky Pitowsky they show that the physical interpretation of the Kochen-Specker theorem involves some important subtleties which, in the past, have not been sufficiently appreciated.
The purpose of this paper is to show that the MKC models exhibit a novel kind of contextuality, which has not previously been remarked in the literature, and which is even more strikingly at variance with classical assumptions than the usual kind of contextuality, featuring in the Kochen-Specker theorem. In the usual kind of contextuality it is only the *value* assigned to an observable which is context-dependent. In the MKC models, however, it is the very *existence* of an observable which is context-dependent (its existence, that is, as a physical property whose value can be revealed by measurement). This phenomenon may be described as existential contextuality <sup>1</sup><sup>1</sup>1It should be noted that the concept of existential contextuality introduced here is completely unrelated to the concept of ontological contextuality discussed by Heywood and Redhead Hey , Redhead Red and Pagonis and Clifton Pag . . It confirms the point made in ref. self2 , that the MKC models do not, as MKC claim, provide a classical explanation for non-relativistic quantum mechanics.
This paper was originally motivated by a seeming inconsistency in MKC’s statement KentA ; KentB , that their models are both non-contextual *and* non-local. There do, of course, exist theories which have both these properties (Newtonian gravity, for example). However, in the framework of quantum mechanics the phenomena of contextuality and non-locality are closely connected, as has been stressed by Mermin MerminB (also see Heywood and Redhead Hey and Basu *et al* Basu ). The discussion in Mermin MerminB suggests that, if one were to examine the predictions the MKC models make regarding (for example) the GHZ set up GHZ , then one might expect to find evidence that these models are not only non-local (as MKC state), but also contextual (which they deny). As we will see, this is in fact the case.
## II GHZ set-up, with a locality assumption
Consider Mermin’s variant MerminB ; MerminC of the GHZ set-up GHZ , as illustrated in Fig. 1. This arrangement is usually regarded as a way of demonstrating non-locality. As discussed in the Introduction, we will show that it can also be used to demonstrate a form of contextuality.
The system consists of three spin-$`1/2`$ particles. Let $`\widehat{𝝈}^{(r)}`$ denote the Pauli spin vector for particle $`r`$, and let $`^{(r)}`$ be the 2-dimensional Hilbert space on which it acts. The spin state is
$$|\psi =\frac{1}{\sqrt{2}}\left(|1,1,1|1,1,1\right)$$
(1)
where $`|s_1,s_2,s_3`$ denotes the joint eigenstate of $`\widehat{\sigma }_z^{(1)}`$, $`\widehat{\sigma }_z^{(2)}`$, $`\widehat{\sigma }_z^{(3)}`$ with eigenvalues $`s_1`$, $`s_2`$, $`s_3`$. The particles emerge from a source and pass through three space-like separated detectors (see Fig. 1). For each $`r`$ the corresponding detector measures one of the two observables $`\widehat{\sigma }_x^{(r)}`$ or $`\widehat{\sigma }_y^{(r)}`$. One has MerminC
$$\widehat{\sigma }_x^{(1)}\widehat{\sigma }_x^{(2)}\widehat{\sigma }_x^{(3)}|\psi =|\psi $$
(2)
and
$$\widehat{\sigma }_x^{(1)}\widehat{\sigma }_y^{(2)}\widehat{\sigma }_y^{(3)}|\psi =\widehat{\sigma }_y^{(1)}\widehat{\sigma }_x^{(2)}\widehat{\sigma }_y^{(3)}|\psi =\widehat{\sigma }_y^{(1)}\widehat{\sigma }_y^{(2)}\widehat{\sigma }_x^{(3)}|\psi =|\psi $$
(3)
Consequently, if the detectors are strictly ideal, and if they are set precisely at the combination $`xxx`$, then the product of measured values must necessarily be $`1`$. Similarly, if the detectors are strictly ideal, and if they are set precisely at one of the combinations $`xyy`$, $`yxy`$, $`yyx`$, then the product of measured values must necessarily be $`+1`$.
MKC argue that it would not, in practice, be possible to align the detectors with infinite precision, implying that the detectors, instead of performing ideal measurements <sup>2</sup><sup>2</sup>2It should be stressed that the detectors still perform *non-ideal* measurements of $`\widehat{\sigma }_{j_1}^{(1)}`$, $`\widehat{\sigma }_{j_2}^{(2)}`$, $`\widehat{\sigma }_{j_3}^{(3)}`$ (see Appleby self2 ). of the observables $`\widehat{\sigma }_{j_1}^{(1)}`$, $`\widehat{\sigma }_{j_2}^{(2)}`$, $`\widehat{\sigma }_{j_3}^{(3)}`$ (with $`j_r=x`$ or $`y`$), may actually perform ideal measurements of a slightly different set of commuting observables $`\widehat{\tau }^{(1)}`$, $`\widehat{\tau }^{(2)}`$, $`\widehat{\tau }^{(3)}`$. They postulate that the observables $`\widehat{\tau }^{(1)}`$, $`\widehat{\tau }^{(2)}`$, $`\widehat{\tau }^{(3)}`$ are always such that their joint spectral resolution is a subset of a countable set $`𝒫_\mathrm{d}`$ (in the notation of Clifton and Kent KentB ), which is dense in the space of all projections on $`_1_2_3`$. For each $`r`$ detector $`r`$ reveals the pre-existing value of the observable $`\widehat{\tau }^{(r)}`$ which it does in fact ideally measure.
The fact that the observables $`\widehat{\tau }^{(r)}`$ may not precisely coincide with the observables $`\widehat{\sigma }_{j_r}^{(r)}`$ means that there may be a small, non-zero probability of obtaining the “wrong” measurement outcome (*i.e.* 1 for the combination $`xxx`$, and $`1`$ for the combinations $`xyy`$, $`yxy`$, $`yyx`$). This is consistent with the unavoidable imprecision of real, laboratory measurements.
In this paper we are, for simplicity, confining ourselves to the kind of measurement envisaged by MKC, in which the imprecision is entirely due to the detectors not being aligned precisely in the directions specified. It should be stressed that such measurements are still highly idealised. MKC assume that there is always *some* observable which a detector ideally measures. They overlook the fact that a real, laboratory instrument does not, typically, perform an ideal measurement of anything: neither the nominal observable, which the experimenter records as having been measured, nor any other observable either. We discuss this point further in Appleby self2 .
In their published papers MKC take the view that the difference between $`\widehat{\tau }^{(r)}`$ and $`\widehat{\sigma }_{j_r}^{(r)}`$ is due to detector $`r`$ not being aligned with infinite precision. On this view $`\widehat{\tau }^{(r)}`$ must be a local observable of the form $`\widehat{\tau }^{(r)}=𝐧_r\widehat{𝝈}^{(r)}`$, where $`𝐧_r`$ is a unit vector close to the unit vector in the $`j_r`$ direction, representing the actual alignment <sup>3</sup><sup>3</sup>3We are assuming ideal detectors, so the concept “actual alignment of detector $`r`$” is unambiguous. In the case of non-ideal detectors this concept may not be sharply defined (see Appleby self2 ). of detector $`r`$. Kent’s KentC subsequent suggestion, that $`\widehat{\tau }^{(r)}`$ may be a non-local admixture of observables pertaining to more than one particle, will be discussed in the next section.
Given an arbitrary triplet of unit vectors $`(𝐧_1,𝐧_2,𝐧_3)`$, define projections
$$\widehat{P}_{s_1s_2s_3}=\frac{1}{8}(1+s_1𝐧_1\widehat{𝝈}^{(1)})(1+s_2𝐧_2\widehat{𝝈}^{(2)})(1+s_3𝐧_3\widehat{𝝈}^{(3)})$$
(4)
where, for each $`r`$, $`s_r=\pm 1`$. These projections constitute the joint spectral resolution for the operators $`𝐧_1\widehat{𝝈}^{(1)}`$, $`𝐧_2\widehat{𝝈}^{(2)}`$, $`𝐧_3\widehat{𝝈}^{(3)}`$. Define $`S_6^{}`$ to be the set of vector triplets $`(𝐧_1,𝐧_2,𝐧_3)`$ for which the corresponding projections $`\widehat{P}_{s_1s_2s_3}`$ all $`𝒫_\mathrm{d}`$ (where, in the notation of Clifton and Kent KentB , $`𝒫_\mathrm{d}`$ is the countable set of projections on which the MKC valuations are defined). $`S_6^{}`$ is a countable, dense subset of $`S_2\times S_2\times S_2`$ (where $`S_2`$ is the unit 2-sphere). Its significance is that $`(𝐧_1,𝐧_2,𝐧_3)`$ represents a possible set of alignments for the three detectors if and only if $`(𝐧_1,𝐧_2,𝐧_3)S_6^{}`$.
We will now show that $`S_6^{}`$ cannot be a Cartesian product of the form $`S_{}^{}{}_{2}{}^{(1)}\times S_{}^{}{}_{2}{}^{(2)}\times S_{}^{}{}_{2}{}^{(3)}`$, with $`S_{}^{}{}_{2}{}^{(1)},S_{}^{}{}_{2}{}^{(2)},S_{}^{}{}_{2}{}^{(3)}S_2^{}`$. We will then use this to show that the MKC models exhibit a novel kind of contextuality.
In order to establish this result suppose that $`S_6^{}`$ *is* of the form $`S_{}^{}{}_{2}{}^{(1)}\times S_{}^{}{}_{2}{}^{(2)}\times S_{}^{}{}_{2}{}^{(3)}`$. We will show that this assumption leads to a contradiction.
For each $`r`$ let $`𝐧_{rx}`$, $`𝐧_{ry}`$ be a fixed pair of vectors $`S_{}^{}{}_{2}{}^{(r)}`$ such that $`𝐧_{rx}`$ (respectively $`𝐧_{ry}`$) is close to $`𝐞_x`$ (respectvely $`𝐞_y`$), the unit vector in the $`x`$ (respectively $`y`$) direction. Then
$`\psi |(𝐧_{1x}\widehat{𝝈}^{(1)})(𝐧_{2x}\widehat{𝝈}^{(2)})(𝐧_{3x}\widehat{𝝈}^{(3)})|\psi `$ $`=(1ϵ_0)`$ (5a)
$`\psi |(𝐧_{1x}\widehat{𝝈}^{(1)})(𝐧_{2y}\widehat{𝝈}^{(2)})(𝐧_{3y}\widehat{𝝈}^{(3)})|\psi `$ $`=(1ϵ_1)`$ (5b)
$`\psi |(𝐧_{1y}\widehat{𝝈}^{(1)})(𝐧_{2x}\widehat{𝝈}^{(2)})(𝐧_{3y}\widehat{𝝈}^{(3)})|\psi `$ $`=(1ϵ_2)`$ (5c)
$`\psi |(𝐧_{1y}\widehat{𝝈}^{(1)})(𝐧_{2y}\widehat{𝝈}^{(2)})(𝐧_{3x}\widehat{𝝈}^{(3)})|\psi `$ $`=(1ϵ_3)`$ (5d)
where $`ϵ_a0`$ for each $`a`$. Let $`ϵ=\mathrm{max}(ϵ_0,ϵ_1,ϵ_2,ϵ_3)`$. It follows from Eqs. (2) and (3) and the continuity of the expectation values that $`ϵ0`$ as $`𝐧_{rx}𝐞_x`$, $`𝐧_{ry}𝐞_y`$ for $`r=1,2,3`$. The fact that $`S_6^{}`$ is dense in $`S_2\times S_2\times S_2`$ means that $`S_{}^{}{}_{2}{}^{(r)}`$ is dense in $`S_2`$ for $`r=1,2,3`$. It follows that the vectors $`𝐧_{rj}`$ can be chosen so as to make $`ϵ`$ arbitrarily small.
Let $`\mathrm{\Lambda }`$ be the hidden state space, and for each $`\lambda \mathrm{\Lambda }`$ let $`s_{rj}(\lambda )`$ be the corresponding valuation of $`𝐧_{rj}\widehat{𝝈}^{(r)}`$. We have $`s_{rj}(\lambda )=\pm 1`$ for all $`r,j`$. Define
$`f_0(\lambda )`$ $`=s_{1x}(\lambda )s_{2x}(\lambda )s_{3x}(\lambda )`$ (6a)
$`f_1(\lambda )`$ $`=s_{1x}(\lambda )s_{2y}(\lambda )s_{3y}(\lambda )`$ (6b)
$`f_2(\lambda )`$ $`=s_{1y}(\lambda )s_{2x}(\lambda )s_{3y}(\lambda )`$ (6c)
$`f_3(\lambda )`$ $`=s_{1y}(\lambda )s_{2y}(\lambda )s_{3x}(\lambda )`$ (6d)
Then $`f_a(\lambda )=\pm 1`$ for all $`a,\lambda `$. Also
$$f_0(\lambda )f_1(\lambda )f_2(\lambda )f_3(\lambda )=\left(s_{1x}(\lambda )s_{2x}(\lambda )s_{3x}(\lambda )s_{1y}(\lambda )s_{2y}(\lambda )s_{3y}(\lambda )\right)^2=1$$
(7)
for all $`\lambda `$.
Let $`\mu `$ be the probability measure on $`\mathrm{\Lambda }`$ corresponding to the state $`|\psi `$. The assumption that $`S_6^{}=S_{}^{}{}_{2}{}^{(1)}\times S_{}^{}{}_{2}{}^{(2)}\times S_{}^{}{}_{2}{}^{(3)}`$ implies that
$$1ϵ1ϵ_a=f_a(\lambda )𝑑\mu 1$$
(8)
for all $`a`$. For each $`a`$ let $`A_a`$ be the set
$$A_a=\{\lambda \mathrm{\Lambda }:f_a(\lambda )=1\}$$
(9)
Then it follows from Inequality (8) that
$$1ϵf_a(\lambda )𝑑\mu =2\mu (A_a)1$$
(10)
for all $`a`$. It follows that $`\mu (A_a)1ϵ/2`$ for all $`a`$ and, consequently, that $`\mu (A_0A_1A_2A_3)12ϵ`$. We noted above that, with a suitable choice of the vectors $`𝐧_{rj}`$, $`ϵ`$ can be made arbitrarily small. It follows that there exist vectors $`𝐧_{rj}`$ such that $`\mu (A_0A_1A_2A_3)>0`$ (in fact, there exist vectors $`𝐧_{rj}`$ such that $`\mu (A_0A_1A_2A_3)1`$). On the other hand, it follows from Eq. (7) that $`\mu (A_0A_1A_2A_3)=0`$ for every choice of $`𝐧_{rj}`$—which is a contradiction.
We have thus shown that the set $`S_6^{}`$ does not have the form of a Cartesian product for any model of MKC type. This has important consequences: for it implies that it must, in general, happen that a change in the alignment of one detector forces a change in the alignment of at least one of the other two detectors. This represents a form of non-locality. However, the point which concerns us here is that it also represents a form of contextuality.
It is a particularly striking form of contextuality. Let $`S_{}^{}{}_{2}{}^{(r)}`$ be the set of possible alignments for detector $`r`$. In the usual kind of contextuality $`S_{}^{}{}_{2}{}^{(r)}`$ is fixed, and it is only the values assigned to the members of this set which depend on the measurement context. However, in the MKC models it is the set $`S_{}^{}{}_{2}{}^{(r)}`$ itself which depends on the measurement context. In other words, it is not simply the *value*, but the very *existence* of an observable which is context-dependent (its existence, that is, as a physical property whose value can be revealed by measurement).
## III GHZ set-up, with non-local detectors
In the last section we assumed that detector $`r`$ reveals the value of a local observable, defined on the state space of particle $`r`$. Kent KentC has objected to this assumption. He suggests, instead, that, on the level of the hidden variables, the detectors may function as non-local devices, which reveal the values of non-local admixtures of observables pertaining to more than one particle.
Let us begin by noting that this suggestion involves a significant departure from the view taken in MKC’s published papers. In their published papers MKC argue that the observable $`\widehat{\tau }^{(r)}`$, whose value is revealed by detector $`r`$, is also the observable which detector $`r`$ ideally measures <sup>4</sup><sup>4</sup>4MKC only consider ideal detectors. Their assumption, that a measurement reveals the pre-existing value of the observable which is ideally measured, is obviously not applicable in the case of non-ideal detectors (see Appleby self2 ). , the discrepancy between $`\widehat{\tau }^{(r)}`$ and $`\widehat{\sigma }_{j_r}^{(r)}`$ being entirely attributable to an inaccuracy in the alignment of detector $`r`$. Clearly, detector $`r`$ can only perform ideal quantum measurements of local observables pertaining to particle $`r`$. Consequently, the position adopted in MKC’s published papers implies that $`\widehat{\tau }^{(r)}`$ must be a local observable pertaining to particle $`r`$—as we assumed in the last section.
If a detector reveals the pre-existing value of some non-local observable then, on the level of the hidden variables, it must be interacting non-locally with more than one particle. This interaction would represent a further element of non-locality in the theory, additional to the non-locality required by the standard arguments (Bell, GHZ, *etc*.). A model of this kind would thus be even more strongly non-classical than the models originally proposed in MKC’s published papers.
Nevertheless, the fact that the observables $`\widehat{\tau }^{(r)}`$ may be assumed to be arbitrarily close to local observables of the form $`𝐧_r\widehat{𝝈}^{(r)}`$ means that a model of the kind indicated will still be consistent with the empirical predictions of conventional quantum mechanics. The question consequently arises, whether the phenomenon of existential contextuality, discussed in the last section, also occurs in models of this more general kind. It is easily seen that the answer to this question is in the affirmative.
Let $`𝒫`$ be the set of all projection operators on $`_1_2_3`$, and let $`𝒫_\mathrm{d}`$ be the countable, dense subset of $`𝒫`$ on which the MKC truth-functions are defined (where we are employing the notation of Clifton and Kent KentB , as before). Let $`\overline{𝒫}_\mathrm{d}`$ be the set of self-adjoint operators on $`_1_2_3`$ whose spectral resolutions are contained in $`𝒫_\mathrm{d}`$. The triplet $`(\widehat{\tau }^{(1)},\widehat{\tau }^{(2)},\widehat{\tau }^{(3)})`$ of commuting observables whose values are revealed by the three detectors must $`\overline{𝒫}_\mathrm{d}\times \overline{𝒫}_\mathrm{d}\times \overline{𝒫}_\mathrm{d}`$. It is determined by the hidden state of the three detectors. Let $`T\overline{𝒫}_\mathrm{d}\times \overline{𝒫}_\mathrm{d}\times \overline{𝒫}_\mathrm{d}`$ be the set of all possible triplets $`(\widehat{\tau }^{(1)},\widehat{\tau }^{(2)},\widehat{\tau }^{(3)})`$, as determined by the set of all possible hidden detector states. The set $`T`$ is the analogue, in the more general setting of this section, of the set $`S_6^{}`$ defined in the last section.
We may now show, using a straightforward modification of the argument in the last section, that $`T`$ is not a Cartesian product. In fact, suppose that $`T`$ *was* of the form $`T^{(1)}\times T^{(2)}\times T^{(3)}`$ (with $`T^{(r)}\overline{𝒫}_\mathrm{d}`$ for $`r=1,2,3`$). We could then choose, for each $`r=1,2,3`$ and $`j=x,y`$, operators $`\widehat{\tau }_j^{(r)}T^{(r)}`$ such that $`\widehat{\tau }_j^{(r)}\widehat{\sigma }_j^{(r)}`$ for all $`r`$, $`j`$. This would imply
$`1ϵ`$ $`\psi |\widehat{\tau }_x^{(1)}\widehat{\tau }_x^{(2)}\widehat{\tau }_x^{(3)}|\psi 1`$ (11a)
$`1ϵ`$ $`\psi |\widehat{\tau }_x^{(1)}\widehat{\tau }_y^{(2)}\widehat{\tau }_y^{(3)}|\psi 1`$ (11b)
$`1ϵ`$ $`\psi |\widehat{\tau }_y^{(1)}\widehat{\tau }_x^{(2)}\widehat{\tau }_y^{(3)}|\psi 1`$ (11c)
$`1ϵ`$ $`\psi |\widehat{\tau }_y^{(1)}\widehat{\tau }_y^{(2)}\widehat{\tau }_x^{(3)}|\psi 1`$ (11d)
where the positive constant $`ϵ`$ can be chosen arbitrarily small (compare Eqs. (5) in the last section). If $`ϵ<1/2`$ we can show that these inequalities lead to a contradiction, by an argument which is essentially the same as the argument following Eqs. (5) in the last section. It follows that $`T`$ is not a Cartesian product.
We conclude that the MKC models still exhibit the phenomenon of existential contextuality described in the last section, even on the assumption that the detectors may reveal the pre-existing values of non-local observables.
## IV Conclusion
In this paper we have argued that the MKC models are contextual. It follows that they do not, as MKC claim, provide a classical explanation for the empirically verifiable predictions of non-relativistic quantum mechanics. This confirms the conclusion reached in Appleby self2 , on the basis of a different, completely independent argument.
We would, however, stress that, notwithstanding these criticisms, it appears to us that the work of MKC is deeply interesting, and important. We have argued that MKC’s attempt to explain non-relativistic quantum mechanics in classical terms is misconceived. Nevertheless their work is still valuable because, together with the earlier work of Pitowsky Pitowsky , it shows that the physical interpretation of the Kochen-Specker theorem is a great deal more subtle than may superficially appear. It consequently leads to a deeper understanding of the conceptual implications of quantum mechanics.
The work of MKC and Pitowsky is also interesting because it enlarges the scope of the hidden variables concept in a new and imaginative way. In the past hidden variables theories have primarily been motivated by purely philosophical considerations. The emphasis has been largely (though not entirely—see Valentini Val1 ; Val2 ) on constructing alternative interpretations of conventional quantum mechanics. Recently, however, ’t Hooft tHooft1 ; tHooft2 (in one way) and Faraggi and Matone Far1 and Bertoldi *et al* Far2 (in another way) have speculated that Planck scale physics may most appropriately be described in terms of a hidden variables theory which is *not* equivalent to conventional quantum mechanics. A theory of this kind, if it could be constructed, would be *empirically* significant.
In this connection it may be worth noting that ’t Hooft tHooft2 has argued that, on the level of Planck scale variables, it may not be possible to rotate a detector at will so as to measure either the $`x`$ or $`y`$ components of a particle’s spin, and that this may provide a way of circumventing the Bell theorem. This proposal is similar to MKC’s attempt to circumvent the Kochen-Specker theorem. Our analysis of the MKC models would consequently seem to indicate that one cannot restore classicality in the manner ’t Hooft suggests.
On the other hand, there is no evident reason why one should demand non-contextuality and locality in respect of a theory of the kind proposed by ’t Hooft. Such a theory must, by definition, restore the concept of a world of objective facts. However, this concept is by no means exclusive to classical physics. In other respects the theory might be highly non-classical. Indeed, it might be even more highly non-classical than conventional quantum mechanics. The aim is to understand the actual constitution of the physical universe. There is no clear reason to exclude, at the outset, the possibility that the world actually is contextual and non-local.
The ideas of ’t Hooft and Faraggi *et al* are admittedly speculative. They do, however, provide an additional motive for investigating new and more imaginative implementations of the hidden variables hypothesis.
###### Acknowledgements.
It is a pleasure to thank Prof. A. Zeilinger for his hospitality during the programme “Quantum Measurement and Information” held at ESI in Vienna. It is also a pleasure to thank N.D. Mermin, A. Peres, A. Kent, P. Busch, D. Home, K. Svozil, I. Pitowsky, G. Mahler, and an anonymous referee for their stimulating and helpful comments.
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# Schrödinger revisited: an algebraic approach
## I Introduction
In a previous paper, Monk and Hiley have suggested that instead of using the traditional Hilbert space description of quantum phenomena, one should give primary consideration to the algebraic structure, not only because it has a number of mathematical advantages that have already been pointed out by Dirac , but because it offers the possibility of a radically different interpretation of the quantum formalism. By algebraic structures, we mean exploiting the rich possibilities contained in, for example, C\* and W\* (von Neumann) algebras, which play an important role in field theory as well as in equilibrium and non-equilibrium statistical mechanics . In spite of the potential richness of these methods, they have not been used in the general debate on the foundations of quantum theory mainly because of the abstract nature of the mathematics.
In their paper, Monk and Hiley outlined how this mathematics can be simplified so that the approach becomes much more transparent. Furthermore, one can develop an interpretation for this formalism provided one is willing to give up the basic ideas of particles- and/or fields-in-interaction and instead, think in terms of process. Indeed, recent work has shown how the algebraic approach does have a potential for taking the discussions of the meaning of the quantum formalisms into new domains. It is this background that provides the motivation for the present paper. However, we will not assume any detailed prior knowledge of this abstract algebraic structure. Our purpose here is to show how we can re-write the equations of elementary quantum mechanics in a purely algebraic way, which, as we show, yields some interesting new insights.
Traditionally, the algebraic approach has implied the use of the Heisenberg picture. Here the operators (or elements of the algebra, in our case) become time dependent and carry the dynamics of the quantum system. One clear advantage of this approach is that all the elements of the algebra are representation-independent and so we are not tied to any one particular representation. A further advantage is that the equations of motion have a close structural similarity to the classical equations of motion, viz, commutator brackets directly replace Poisson brackets .
In contrast to this, the Schrödinger picture has the time development entirely tied into the wave function, and as such is a representation dependent object which appears to exists only in Hilbert space and does not appear in the algebra. Thus, the Schrödinger picture does not seem to have the generality of the representation-independent Heisenberg picture.
In this paper, we will show that this representation dependence of the Schrödinger equation is only apparent and in section II we will show that it can be written in a representation independent form. This means writing the wave function as a “wave operator” (in the left ideal of the algebra) so that the Schrödinger equation becomes purely algebraic and independent of any representation in a Hilbert space.
The way to do this has been known for a long time . Indeed, Monk and Hiley have already shown in simple terms that the key step involves expressing the wave operator in the left ideal. This means that the wave function must be replaced by the density operator, $`\rho `$, even in the case of pure states. Here the density operator plays the role of an idempotent and it is this idempotency that is central in the whole approach.
Starting with the quantum Liouville equation and writing the density operator as the product of elements respectively in the left and right ideals in the algebra, we find that the Schrödinger picture may be expressed through two algebraic forms which are representation independent. These forms are the two equations (19) and (20) for the density and phase operators. The first, expressed in terms of the commutator of the Hamiltonian and the density operator, is an alternative form of the quantum Liouville equation which, as is well known, describes the conservation of probability. The second describes the time dependence of the phase and is expressed in terms of the anti-commutator of the Hamiltonian and the density operator. This second equation, which does not appear in the literature as far as we are aware, becomes an equation for the conservation of energy in systems when energy is well defined.
In section III, the time dependent phase operator equation (20) is shown to be gauge invariant and reproduces some well known results of gauge theory in a very direct and simple way. For example, one can immediately derive the Aharonov-Bohm phase for a particle travelling in a vector potential while a trivial extension incorporates the Aharonov-Casher phase. This latter phase arises when a neutral particle with a magnetic moment passes a line-charge. In addition to these examples, the Berry phase and its associated energy follow almost trivially from the same equation.
In sections IV and V we use these equations to explore in more detail the Bohm interpretation \[BI\] . Because the equations (19) and (20) are representation independent, we can construct a BI based on trajectories in any representation. To do this, we introduce a generalisation of the current density operator and demonstrate its use in both ordinary space and momentum space. In section V and through the examples in section VI we show in detail how one can construct a consistent BI in the $`p`$-representation. This is contrary to the assertion that this is not possible even in “the simplest case to construct an acceptable causal interpretation” and we discuss the significance of this statement in the light of our examples.
Our examples not only remove one of serious criticisms of this interpretation, namely that it does not use the full symplectic symmetry of the quantum formalism, but also provide us with new insights into the meaning of BI and its relation to standard quantum mechanics. In all of this work no appeal is made to any classical formalism whatsoever, showing that the BI is quantum through-and-through.
Perhaps the most important conclusion of this work is to show the BI arises directly from the non-commutative structure of the quantum mechanical phase space. Non-commutative geometries are not built on any form of well defined continuous manifolds. We are forced to construct “shadow manifolds” . As we show in section V, these shadow manifolds have the structure of a phase space. One is constructed using the $`x`$-representation and the other uses the $`p`$-representation. These spaces are different but converge to the same phase space in the classical limit.
In the final section VII we discuss the consequences of this construction for the BI. In fact we show that our approach more clearly illustrates the ideas that Bohm and Hiley presented in the final chapter of their book . There it was argued that a new way of exploring the meaning of the quantum formalism required a new order, the implicate order, this order having its origins in the mathematics of non-commutative geometry . The shadow phase spaces are examples of explicate orders. Finally, we briefly discuss how the BI fits into this general scheme.
## II The algebraic approach to the Schrödinger equation
We begin by writing the Schrödinger equation in a general representation
$$i\frac{\psi (a_i,t)}{t}=H\psi (a_i,t)(\mathrm{}=1)$$
(1)
where the $`a_i`$ are the eigenvalues of $`A`$ an algebraic element or operator in the algebra<sup>*</sup><sup>*</sup>*Throughout this paper, we set $`\mathrm{}=1`$ and use the convention of representing operators by capitals and eigenvalues by lower case letters.. Equation (1) introduces the state vectors $`\psi (a_i,t)`$ that are not elements of the algebra; rather, they are elements of a separate vector (Hilbert) space and as such depend on a representation. In this form, the Schrödinger equation does not appear to be part of the algebra even though it uses elements of the operator algebra.
On the other hand, in the Heisenberg approach we write the Hamiltonian flow of the operator $`A`$ A as
$$A(0)A(t)=M(t)^1A(0)M(t)$$
(2)
where $`M(t)=\mathrm{exp}[iHt]`$, so giving rise to the Heisenberg equation of motion
$$\frac{dA}{dt}=\frac{1}{i}[A,H]_{}.$$
(3)
This means that the Heisenberg time evolution can be regarded as an inner automorphism in the algebra $`M(t):AAA`$ A. The equation of motion can be generalised to include the explicit time dependence of $`A`$ giving
$$\frac{dA}{dt}=\frac{A}{t}+\frac{1}{i}[A,H]_{}.$$
(4)
Note that this equation is representation-free and the time evolution is discussed entirely within the algebra itself.
In the algebraic approach , the state function is introduced through a density operator, $`\rho `$. This operator is actually in the quantum algebra so that the Schrödinger time development must be implicit within the algebra itself. How, then, can the time evolution be algebraically expressed within the Schrödinger picture, without reference to either Hilbert space or a particular representation in it?
In the usual approach to quantum mechanics, the density operator is, unfortunately, not introduced as a primitive notion in the theory. Rather, it is introduced almost as an after-thought when it is found necessary to deal with mixed states. But using the density operator as a starting point has the advantage of including both pure states and mixed states together and of satisfying the idempotent condition $`\rho =\rho ^2`$. Moreover, if we adopt the further defining condition that the density operator must satisfy the Liouville theorem
$$\frac{d\rho }{dt}=0,$$
(5)
then, since the density operator is also an element of the algebra, equation (4) immediately leads to the quantum Liouville equation
$$\frac{\rho }{t}=\frac{1}{i}[H,\rho ]_{}.$$
(6)
If $`\rho =\psi \psi `$, is an element of the algebra, then it must also be possible to identify the ket and the bra with particular elements of the algebra. A ket is an element of a vector space, which when multiplied from the left it must remain in that space. The algebraic equivalent of this vector space is a left ideal I<sub>L</sub>. An element, $`B\epsilon `$, of a left ideal, where $`B`$ is a wave operator and $`\epsilon `$ a primitive idempotent, corresponds to a ket. Similarly an element, $`\epsilon C`$, of a right ideal, I<sub>R</sub>, corresponds to a bra. This suggests we write $`\rho =B\epsilon \epsilon C=B\epsilon C`$, which means that a pure state density operator corresponds to a two sided ideal, subject to the conditions $`\rho =\rho ^2`$ and tr$`\rho =1`$.
To see how to write the algebraic equivalent to the Schrödinger equation, let us substitute $`\rho =B\epsilon C`$ into the equation of motion (6), so that
$$i\mathrm{}\left(\frac{B}{t}\right)\epsilon C+i\mathrm{}B\epsilon \left(\frac{C}{t}\right)=HB\epsilon CB\epsilon CH.$$
(7)
Since $`B`$ and $`C`$ are operator elements outside the ideals of the algebra, it can be assumed that there exist $`B^{}:B^{}B=\mathrm{𝟏}`$ and $`C^{}:CC^{}=\mathrm{𝟏}`$. Multiplying the above equation from the left by $`B^{}`$ and from the right by $`C^{}`$, we find after re-arrangement, that
$$B^{}\left(i\mathrm{}\frac{B}{t}HB\right)\epsilon =\epsilon \left(i\mathrm{}\frac{C}{t}+CH\right)C^{}.$$
(8)
Since $`B(C)`$ is any non-null element of the algebra, we can write
$$i\left(\frac{B}{t}\right)\epsilon =HB\epsilon $$
(9)
and
$$i\epsilon \left(\frac{C}{t}\right)=\epsilon CH.$$
(10)
We see immediately that equations (9) and (10), which are respectively in the left and right ideals of the algebra, have the same general form as the Schrödinger equation and its conjugate counterpart ($`H`$ is assumed to be Hermitian). We stress here again that $`B`$ and $`C`$ are elements of the algebra and not elements of a Hilbert space.
To see exactly how these two equations are related to the usual Hilbert space formalism, we specifically choose the wave operator $`B`$ to be a function of the position operator $`X`$ so that $`B\epsilon =B(X,t)`$ I<sub>L</sub> and then project $`B(X,t)`$ into a complex function belonging to $`L^2(x,\mu )`$ viz,$`L^2(x,\mu )`$ means square integrable complex functions with measure $`\mu `$.
$$\eta :B(X,t)B(x,t)$$
(11)
so that
$$B(X,t)(x)=B(x,t).$$
(12)
This is the usual wave functionNot all elements of a left ideal produce state functions that are physically meaningful. We will not discuss these restriction here. (See Ballentine .), conventionally written as $`\mathrm{\Psi }(x,t)`$. It is now straight forward to show that equation (9) becomes the Schrödinger equation
$$i\frac{\mathrm{\Psi }(x,t)}{t}=H(x)\mathrm{\Psi }(x,t).$$
(13)
The conjugate equation can be derived by first assuming the dual projection
$$\eta :C(X,t)C^{}(x,t)$$
(14)
so that
$$C(X,t)(x)=C^{}(x,t).$$
(15)
Again it is straight forward to show that equation (10) leads to the conjugate Schrödinger equation. Thus equations (9) and (10) are the algebraic, representation independent equivalents of the Schrödinger equation.
Now let us continue developing the general structure. We write the wave operators $`B`$ and $`C`$ in the mutually conjugate forms $`B=\mathrm{exp}[iS_Q(t)]`$ and $`C=\mathrm{exp}[iS_Q^{}(t)]`$ where $`S_Q=Si\mathrm{ln}R`$. In this case equations (9) and (10) become the dual pair
$$\frac{S_Q}{t}B\epsilon =HB\epsilon $$
(16)
and
$$\epsilon C\frac{S_Q^{}}{t}=\epsilon CH,$$
(17)
which are quantum algebraic equivalents<sup>§</sup><sup>§</sup>§Equations (16) and (17) cannot be simplified further, since $`B\epsilon I_L`$ and $`\epsilon CI_R`$ do not have inverses. of the Hamilton-Jacobi equation of classical mechanics
$$\frac{S_{cl}}{t}+H=0,$$
(18)
where $`S_{cl}`$ is the classical action. It is natural, therefore, to call $`S_Q`$ the quantum action.
Equations (16) and (17) respectively evolve in the left and right ideals, which are mutually dual spaces, so reflecting the essential duality between the Schrödinger equation and its complex conjugate. This duality can be lifted out of the left and right ideals of the algebra and reflected in another pair of algebraic equations. Post- and pre-multiplying equations (16) and (17) by $`\epsilon C`$ and $`B\epsilon `$ respectively and then adding and subtracting the resulting equations, we find
$$\left(\frac{S_Q}{t}\rho +\rho \frac{S_Q^{}}{t}\right)=[H,\rho ]_{}$$
(19)
and
$$\left(\frac{S_Q}{t}\rho \rho \frac{S_Q^{}}{t}\right)=[H,\rho ]_+$$
(20)
Of note is the appearance of the commutator and the anti-commutator with the Hamiltonian on the RHS of these two equations. They are mathematically equivalent to the Schrödinger equation and its conjugate and are general in the sense that they are independent of a specific representation. Expressing equation (19) in Hermitian form, we immediately recover (6), the quantum Liouville equation. Equation (20) cannot in general be reduced to a simpler algebraic form but may be recognised as a symmetrised operator form of the Hamilton-Jacobi equation.
The meaning of equations (19) and (20) can be further clarified by looking at the diagonal elements in the $`a`$-representation. Thus,
$$\frac{\rho _R(a)}{t}+\frac{1}{i}[\rho ,H]_{}(a)=0$$
(21)
and
$$\rho _R(a)\frac{S(a)}{t}+\frac{1}{2}[\rho ,H]_+(a)=0.$$
(22)
Here we have written \[….\](a) = $`a|[\mathrm{}.]|a`$ and $`\rho _R=R^2`$. If we now choose $`A`$ to be the position operator, then equation (21) takes the form
$$\frac{𝒫}{t}+𝐣=\mathrm{𝟎}$$
(23)
where $`𝒫=𝒫(𝐱)=\rho _𝐑(𝐱)=𝐱|\rho |𝐱`$ and $`𝐣`$ is a probability current. Thus, the Liouville equation (21) is identified with the conservation of probability as expected. Equation (22) describes the time variation of the quantum phase and so we will call it (and the more general form (20)) the quantum phase equation. In a state in which the energy is well defined, this equation becomes
$$\frac{S}{t}=E.$$
(24)
In this case, equation (22) expresses the conservation of energy in Hamilton-Jacobi form.
The Liouville equation (21) is well known and plays a prominent role in quantum statistical mechanics. The quantum phase equation (22) does not usually appear in the literature, although something similar has been used by George et al. in their discussions of irreversible quantum processes. In their case, the anti-commutator is simply introduced by defining it to be the energy super-operator. What we show here is that this operator comes directly from the Schrödinger equation and although the extension to super operator status is possible, this generalisation is not necessary for the purposes of this paper.
In summary then, equations (19) and (20) are simply the algebraic equivalents of the Schrödinger equation when it is written in a way that does not depend on a specific representation. It is now easy to confirm that these two equations, when expressed in a particular representation, are simply the real and imaginary parts of the Schrödinger equation under a polar decomposition of the wave function written in that particular representation.
## III The quantum phase equation
### A Gauge invariance
We will first examine equation (22) in some detail. Let us begin by looking at the form of this equation in the $`x`$-representation when we choose the Hamiltonian $`H=p^2/2m+V(x)`$. Here equation (22) becomes
$$\frac{S}{t}+\frac{1}{2m}\left(\frac{S}{x}\right)^2+V(x)\frac{1}{2mR}\left(\frac{^2R}{x^2}\right)=0$$
(25)
This equation is the real part of the Schrödinger equation in the $`x`$-representation.
Before examining this equation in detail we must ensure that it is gauge invariant. To show that this is the case, let us first introduce the gauge transformation $`V^{}(x)=V(x)+V_0`$. This must be accompanied by the phase transformation $`\psi ^{}(x,t)=\varphi (x)\mathrm{exp}[i(E+V_0)t]`$. Since we are considering a well defined energy state, the transformed equation (20) will read
$$((\frac{S_Q}{t}+V_0)\rho \rho (\frac{S_Q^{}}{t}+V_0))=[H,\rho ]_++[V_0,\rho ]_+$$
(26)
because $`\rho ^{}=\rho `$, $`H^{}=H+V_0`$ and $`S_Q^{}=S_Q+V_0`$. Then, since $`[V_0,\rho ]_+=2V_0\rho `$, we immediately recover equation (20), so establishing its gauge invariance and that of equation (22).
Gauge invariance in this case involves the phase change $`S^{}=S+S_0`$ then since
$$V^{}(x,t)=V(x)+V_0(t)$$
(27)
our equation gives
$$\frac{S_0}{t}=V_0(t)$$
(28)
so that
$$S_0=_{t_0}^tV_0(t^{})𝑑t^{}$$
(29)
It will immediately be recognised that this is the expression for the scalar part of the Aharonov-Bohm effect .
We can also obtain the magnetic phase shift from equation (22) by starting from the Hamiltonian
$$H=\frac{1}{2m}(𝐏e𝐀)^2(c=1)$$
(30)
On expanding this Hamiltonian, we find
$`H`$ $`={\displaystyle \frac{1}{2m}}𝐏^2{\displaystyle \frac{e}{2m}}\left(𝐏𝐀+𝐀𝐏\right)+{\displaystyle \frac{e^2}{2m}}A^2`$ (31)
$`=H_{freeparticle}+H_{int}+H_{freefield}`$ (32)
The corresponding phase will then consist of three terms
$$S=S_{freeparticle}+S_{int}+S_{freefield}$$
(33)
so that the diagonal form
$$\rho _R\left(\frac{S_{int}}{t}\right)+\frac{1}{2}[H_{int},\rho ]_+=0$$
(34)
gives
$$\rho _R\frac{S_{int}}{t}=e𝐀𝐣_x$$
(35)
in the $`x`$-representation. If we then write $`𝐣_x=\rho _R𝐯`$, we find
$$S_{int}=e_{t_0}^t𝐀𝐯𝑑t=e_{x_0}^x𝐀𝑑𝐱$$
(36)
We immediately recognise this equation as the expression for the Aharonov-Bohm phase for the vector potential This effect was derived in the $`x`$-representation from the ‘guidance’ condition by Philippidis et al . A more recent discussion using this approach rather than the method we use can be found in Sjöqvist and Carlsen .
The phase for the Aharonov-Casher effect , which involves a neutral particle with a magnetic moment passing a line of electric charges, also follows trivially once the Hamiltonian
$$H=\frac{1}{2m}\left(𝐏𝐄\times \mu \right)^2\frac{\mu E^2}{m}$$
(37)
is assumed. The additional phase change also follows trivially from the same procedure used for the vector potential.
It should also be noted that the Berry phase emerges directly from equation (22). In this case the behaviour of the quantum system depends on some additional cyclic parameter $`(t)`$. The phase now becomes a function of this parameter. Thus equation (22) becomes
$$\rho _R\left(\frac{S}{t}+\dot{}\frac{S}{}\right)+\frac{1}{2}[\rho ,H]_+=0$$
(38)
giving an extra phase factor $`\dot{}\frac{S}{}`$. Thus the contribution to the phase from this extra degree of freedom is
$$S_{Berryphase}=_{t_0}^t\dot{}\frac{S}{}𝑑t$$
(39)
To evaluate this term, we need to consider specific problems which means going to a specific Hamiltonian in a specific representation. This representation is generally the $`x`$-representation. Berry considered the case of the precession of nuclear spin in a magnetic field in his original paper, and showed that
$$\frac{S}{}=\mathrm{}(t),t|_{}(t),t$$
(40)
so that
$$S_{Berryphase}=\mathrm{}__0^{}(t),t|_{}|(t),t𝑑$$
(41)
which is exactly the result obtained by Berry .
### B The $`x`$\- and $`p`$-representations
Having seen how the additional phase changes arise for these simple gauge fields, we now return to examine the details of equation (25). To do this, let us consider the case of the harmonic oscillator with Hamiltonian $`H=p^2/2m+Kx^2/2`$. In this case equation (25) reads
$$\frac{S_x}{t}+\frac{1}{2m}\left(\frac{S_x}{x}\right)^2+\frac{Kx^2}{2}\frac{1}{2mR_x}\left(\frac{^2R_x}{x^2}\right)=0$$
(42)
where we have inserted the suffix $`x`$ to emphasise that this is equation (22) expressed in the $`x`$-representation.
Now let us write down the corresponding equation in the $`p`$-representation. This takes the form
$$\frac{S_p}{t}+\frac{p^2}{2m}+\frac{K}{2}\left(\frac{S_p}{p}\right)^2\frac{K}{2R_p}\left(\frac{^2R_p}{p^2}\right)=0$$
(43)
It should be noted that although the functional forms of these two equations are clearly different, they nevertheless have the same energy content. This can be very easily checked for the ground state of the harmonic oscillator. One can quickly show that both equations give the ground state energy to be $`\omega /2`$, the zero-point energy.
In spite of the differences in functional form, there are structural similarities between these two equations. These arise essentially because we have chosen a symmetric Hamiltonian. For example, instead of the $`p`$ that appears in what looks like a kinetic energy term in equation (43), we have $`(S_x/x)`$ in equation (42), and instead of $`x`$ in the potential energy term in equation (43), we have $`(S_p/p)`$. The last term in each equation has the same general form except with the roles of $`x`$ and $`p`$ interchanged.
Since equation (42) is the real part of the Schrödinger equation, we can identifyHolland has called expressions of this type ‘local expectation values’.
$$p_r=\frac{\mathrm{}\left[\psi ^{}(x)P\psi (x)\right]}{|\psi (x)|^2}=\left(\frac{S_x}{x}\right)$$
(44)
Substituting this into equation (42) we find
$$\frac{S_x}{t}+\frac{p_r^2}{2m}+\frac{K}{2}x^2\frac{1}{2mR_x}\left(\frac{^2R_x}{x^2}\right)=0$$
(45)
In the Bohm interpretation, $`p_r`$ was identified with the “beable” momentum. With this identification equation (44) makes it now quite clear why the beable momentum is a function of $`x`$, in contrast to the classical momentum which is always an independent variable.
If $`p_r`$ is a momentum, then clearly equation (45) looks like an equation for the total energy of the quantum system. If we make the assumption that this is an expression for the conservation of energy then, in quantum theory, we must have an additional quality of energy represented by the last term on the RHS. This term is, of course, the quantum potential energy. As has been shown elsewhere this new quality of energy offers an explanation of quantum processes like interference, barrier penetration and quantum non-separability, all of which are quantum phenomena .
Notice that equation (44) allows the possibility of approaching the classical limit smoothly. In this limit $`S_xS_{cl}`$, $`p_rp_{cl}=(S_{cl}/x)`$ and the quantum potential energy becomes negligible so that equation (45) becomes the classical Hamilton-Jacobi equation.
Before continuing, we wish to stress a point that has not been often fully appreciated, namely, that equation (42) is a quantum equation and converting it to equation (45) requires no appeal whatsoever to classical physics. It is true that in the traditional approach to this equation, the BI has made use of the relation $`p=(S_x/x)`$ by appealing to classical canonical theory. However, this is an unnecessary backward step.
It is because equations (22) and (42) are part of the quantum formalism that we were able to derive quantum effects such as the Aharonov-Bohm, Aharonov-Casher and Berry phases from equation (22). In passing it should also be noted that both the $`x`$\- and $`p`$-representations of this equation (i.e., (42) and (43) ) contain a term which we have called the quantum potential. This potential is modified by the presence of the gauge effects as was first shown by Philippidis, Bohm and Kaye . The quantum potential is central to ensuring energy is conserved and, furthermore, it encapsulates quantum non-separability or quantum non-locality . The quantum potential plays a key role in our approach and must be distinguished from Bohmian mechanics as advocated by Dürr et al..
If we now turn to the $`p`$-representation, i.e., equation (43), we can write it in the form
$$\frac{S_p}{t}+\frac{p^2}{2m}+\frac{K}{2}x_r^2\frac{K}{2R_p}\left(\frac{^2R_p}{p^2}\right)=0$$
(46)
by introducing
$$x_r=\frac{\mathrm{}\left[\psi ^{}(p)X\psi (p)\right]}{|\psi (p)|^2}=\left(\frac{S_p}{p}\right)$$
(47)
Here $`x_r`$ is the position “beable”, which now supplements the momentum $`p`$. Again in the classical limit, we have $`S_pS_{cl}`$, $`x_rx_{cl}=(S_{cl}/p)`$ and the last term on the RHS of (46) becomes negligible. It should be noted that in this limit equations (45) and (46) reduce to the same equation giving rise to a unique phase space, which is identical to the classical phase space.
All the above equations are part of standard quantum mechanics. Although we have drawn attention to the significance of equations (45) and (46) to the BI, we have yet to discuss the interpretation in any detail. To do this we first need to find a way to calculate “trajectories”. In the traditional approach to the BI this is done by regarding $`p=(S_x/x)`$ as a “guidance” condition and then using $`\dot{x}=p/m`$ from which one can calculate a set of trajectories. These trajectories are then integrals of the velocity associated with the probability current in the co-ordinate representation. However this is not the general way to do it as can be seen by considering the $`p`$-representation. The analogous expression in this representation is $`x=(S_p/p)`$ and this clearly cannot be regarded as a “guidance” condition. Something is not quite right here. In order to find out what is involved it is necessary to explore the Liouville equation (21) in more detail.
## IV Probability currents
In this section we will focus our attention on the Liouville equation (21). In the $`x`$-representation, this equation gives rise to the conservation of probability equation (23) with the probability current defined by
$$𝐣=\frac{1}{2mi}[\psi ^{}(\psi )(\psi ^{})\psi ].$$
(48)
However, our aim is to find an expression for the current that is not representation specific. To do this we first consider the classical Liouville equation
$$\frac{\rho }{t}+\{\rho ,H\}=0$$
(49)
where $`\{\}`$ is the Poisson bracket. It is easy to verify that this equation can be written in the form
$$\frac{\rho }{t}+\{𝐣_x^c,𝐩\}\{𝐣_p^c,𝐱\}=\mathrm{𝟎}$$
(50)
with
$$𝐣_x^c=\rho _pH\text{and}𝐣_p^c=\rho _xH.$$
(51)
(The Poisson bracket of two vector functions is defined here as $`\{𝐯,𝐰\}=_k\{v_k,w_k\}`$.) If we now follow Dirac’s suggestion by respectively replacing classical variables and Poisson brackets with operators and commutators, we find
$$i\frac{\rho }{t}+[𝐉_X,𝐏][𝐉_P,𝐗]=0$$
(52)
with
$$𝐉_X=_P(\rho H)\text{and}𝐉_P=_X(\rho H)$$
(53)
where the derivatives are on operators. For $`f(\rho ,X,P)=\rho X^kP^n(k1,n1)`$ these are defined as
$$\frac{f}{X}=X^{k1}P^n\rho +X^{k2}P^n\rho X+\mathrm{}+P^n\rho X^{k1}$$
(54)
and
$$\frac{f}{P}=P^{n1}\rho X^k+P^{n2}\rho X^kP+\mathrm{}+\rho X^kP^{n1}.$$
(55)
In the simple case of a free particle of mass $`m`$ we have
$$𝐉_X=\frac{1}{2m}(\rho 𝐏+𝐏\rho )\text{and}𝐉_P=0.$$
(56)
To see how this connects to the conventional results, let us evaluate equation (52) in the $`x`$-representation. Here we find
$$i\frac{𝐱|\rho |𝐱}{t}+𝐱|[𝐉_X,𝐏]|𝐱𝐱|[𝐉_P,𝐗]|𝐱=0.$$
(57)
If $`H=\frac{𝐏^\mathrm{𝟐}}{2m}+V(𝐗)`$ then the first commutator gives $`𝐱|[𝐉_𝐗,𝐏]|𝐱=_x𝐣_x`$ and the second commutator vanishes. Thus, equation (57) becomes
$$\frac{𝒫(𝐱)}{t}+_x𝐣_x=0$$
(58)
which is just equation (23) and
$$𝐣_x=𝐱|_P(\rho \frac{𝐏^\mathrm{𝟐}}{2m})|𝐱$$
(59)
which is gives an expression for the current that is identical to the usual expression given by equation (48). Furthermore, it is unique since it is independent of the form of the potential used in the Hamiltonian.
In the $`p`$-representation we find
$$\frac{𝒫(𝐩)}{t}+_p𝐣_p=0$$
(60)
where
$$𝐣_p=𝐩|_X(\rho V(𝐗))|𝐩$$
(61)
Thus, we can now calculate probability currents in the $`p`$-representation. Unfortunately equation (61) does not give us a model independent expression for the probability current because the specific form of the current depends on the form of $`V(x)`$. On reflection this is not surprising because the rate of change of momentum must depend upon the externally applied potential. We will examine the consequences of these results for the Bohm interpretation in the next section.
## V Re-examination of the Bohm approach
Let us now re-appraise the Bohm approach in the light of the new results presented above.
It has been assumed that it is not possible to construct a BI using any representation other than the $`x`$-representation. This belief arises from an early correspondence between Epstein and Bohm . Epstein suggested that it should be possible to develop an alternative causal interpretation by starting in the momentum representation. Bohm replied agreeing that a new causal interpretation could possibly arise from such a procedure provided the canonical transformation on the particle variables were simultaneously accompanied by a corresponding linear transformation on the wave function. But, he concluded that this did not seem to lead, even in the simplest of cases, to an acceptable causal interpretation. He does not explain why he came to this conclusion but this position has remained the accepted wisdom.
The general results with the harmonic oscillator presented above show that, at least as far as the mathematics is concerned, it does seem possible to develop a causal interpretation in the $`p`$-representation based upon equations (46), (47) and (61). We will illustrate how this can be done using specific examples in section VI, but here we will simply discuss the general principles involved.
We have already pointed out in section III that the so called “guidance” condition is assumed to play a pivotal role in what is known as “Bohmian mechanics” does not generalise to the $`p`$-representation. However what does generalise is a method based on probability currents. Thus, in any $`q`$-representation we use
$$\frac{dq}{dt}=\frac{j_q}{𝒫(q)}$$
(62)
which can be integrated immediately to find a set of trajectories in a general $`(q,t)`$ space<sup>\**</sup><sup>\**</sup>\**Hereafter, the notation is restricted to one degree of freedom. Generalisation to many degrees of freedom is straight forward..
In the $`x`$-representation we have
$$\frac{dx}{dt}=\frac{j_x}{𝒫(x)}$$
(63)
which, when integrated, gives the particle trajectories of the BI. This approach is similar to that used in pragmatic quantum mechanics where the probability current is assumed to describe the flux of particles emerging from, say, a scattering process. Here the flux at a detector is interpreted as the rate of arrival of the scattered particles.
The additional assumption made in the BI is that particles exist with simultaneously well defined positions and momenta and each particle follows one of the one-parameter curves. Such an assumption is clearly excluded in standard quantum mechanics, but this leaves us with the difficulty of understanding how to incorporate the Born probability postulate and its role in the continuity equation (23) except in some abstract sense.
If we do follow the BI in the $`x`$-representation, then the position of the particle is clearly defined and the momentum, $`p_r`$, associated with the particle must be provided through the relation
$$p_r=\frac{\mathrm{}\left[\psi ^{}(x)P\psi (x)\right]}{|\psi (x)|^2}=\left(\frac{S_x}{x}\right)=m\frac{j_x}{𝒫(x)}=m\frac{dx}{dt}$$
(64)
Here the “beable” momentum $`p_r`$ is wholly quantum in origin showing that the BI has its origins entirely within quantum mechanics.
Now in the $`p`$-representation we use
$$\frac{dp}{dt}=\frac{j_p}{𝒫(p)}$$
(65)
to give a set of one-parameter curves in momentum space. In this approach the momentum of the particle has a clear meaning, while the position beable $`x_r`$ is given by equation (47), namely
$$x_r=\frac{\mathrm{}\left[\psi ^{}(p)X\psi (p)\right]}{|\psi (p)|^2}=\left(\frac{S_p}{p}\right)$$
(66)
This means that the derivative in the current $`𝐣_p=𝐩|_X(\rho V(𝐗))|𝐩`$ given by equation (61) must be evaluated at $`x=x_r`$. Thus we again have a specification of the particle with a given momentum at a given “beable” position $`x_r`$.
Thus, the BI based on equations (45) and (46) leads to two distinct phase spaces, one constructed on each representation. Each phase space contains a set of trajectories, one derived from $`j_x`$ and the other from $`j_p`$. Although these phase spaces are actually different, they carry structures that are consistent with the content of the Schrödinger equation. This is in contrast to the classical limit where there is a unique phase space. However we have already noted that equations (45) and (46) reduce to a single equation in the classical limit so the existence of (at least) two phase spaces is a consequence of the quantum formalism.
Now the existence of at least two phase spaces may come as a surprise to those who see the BI as a return to classical or quasi-classical notions. What we have shown here is that the BI enables us to construct what we may call “shadow phase spaces”, a construct that is a direct consequence of the non-commutative nature of the quantum algebra. Giving ontological meaning to the non-commutative algebra implies a very radical departure from the way we think about quantum processes. This was the central theme of Bohm’s work on the implicate order . The work presented in this paper fits directly into this conceptual structure, a point that will be discussed at length elsewhere.
Our present purpose is to clarify the structure of the mathematics lying behind the BI. To this end note that choosing a representation is equivalent to choosing an operator which is to be diagonal. Thus in the phase space described by equation (45) the position eigenvalues are used for the $`x`$ co-ordinates and we then construct the momentum co-ordinate through the condition $`p_r=(S_x/x)`$ to provide the “beable” momentum.
On the other hand, equation (46) describes a phase space constructed using the momentum eigenvalues together with the “beable” position $`x_r`$ defined by $`x_r=(S_p/p)`$. In this way we see exactly how it is possible to construct two different phase spaces, one for each representation.
The fact that we can find a BI in the $`p`$-representation removes the criticism that the BI does not use the full symplectic symmetry of the quantum formalism. But removing this asymmetry might, at first sight, destroy the claim that the Bohm interpretation provides a unique ontological interpretation. This would only be true if we were insisting that the ontology demands a unique phase space. However, as we have already remarked the quantum algebra is non-commutative and a unique phase space is not possible. It was for this reason that Bohm and one of us (BJH) began to explore the possibility of giving ontological significance to the algebra itself. This involves thinking in terms of process rather than particles- or fields-in-interaction and this leads, in turn, to introducing the implicate order mentioned above. This is a very different order from the one assumed by most physicists, which is essentially what we call the Cartesian order.
Once again we contrast our approach with Bohmian mechanics introduced by Dürr et al . Their approach requires the $`x`$-representation to be taken as basic and the guidance relation to be taken as the defining equation of the approach. In view of the results presented here, we see we could have started from the $`p`$-representation. But here the relation $`p=(S/x)`$ cannot play the role of a guidance condition. Hence making the guidance condition as the defining equation in the $`x`$-representation is arbitrary and contrary to what Bohm himself had in mind .
In regard to the lack of $`x`$-$`p`$ symmetry in the traditional approach to the BI, Bohm and Hiley found it necessary to discuss why $`x`$ was the only intrinsic property of the particle, all others depended upon the context. This was certainly felt by one of us (BJH) to be a somewhat arbitrary imposition that did not seem to be a natural consequence of the symplectic invariance of the formalism itself. Had we started with the $`p`$-representation we would have found $`p`$ to be the intrinsic property, while $`x`$ depended upon some context. Thus the restoration of symmetry explains why particular variables become intrinsic and others not.
In the examples we give in this paper, we only consider the two operators $`X`$ and $`P`$. If we regard the change from the $`x`$-representation to the $`p`$-representation as a rotation of $`\pi /2`$ in phase space, we could think about exploring rotations through other angles. Such transformation exist and are known as fractional Fourier transformations which correspond to rotations through any angle $`\alpha `$ in phase space . These allow us to express equations (19) and (20) any arbitrary representation. This generalisation has been investigated and will be reported elsewhere .
All of this shows that the $`x`$ variable is not special as far as the mathematics goes. The real question is why it is necessary to construct different phase spaces in the first place, but before we go into this question we want to present some examples where we can compare in more detail the results obtained from both $`x`$\- and $`p`$-representations.
## VI Specific examples: comparisons of $`x`$\- and <br>$`p`$-representations
### A The free particle described by a Gaussian wave packet
We will start with the simplest case of a particle described by a Gaussian wave packet centred at position $`x=0`$ with mean momentum zero.
The wave packet has the (normalised) Gaussian distribution
$$\varphi (p,t)=\left[\frac{2(\mathrm{\Delta }x)^2}{\pi }\right]^{\frac{1}{4}}\mathrm{exp}\left[p^2(\mathrm{\Delta }x)^2\right]\mathrm{exp}\left[\frac{ip^2}{2m}t\right]$$
(67)
in the momentum representation. In this representation the current $`j_p=0=(dp/dt)`$ so that the trajectories are of constant momentum. Equation (22) gives
$$\frac{S}{t}+\frac{p^2}{2m}=0$$
(68)
which shows that the quantum potential is zero, as is to be expected from the form of the wave function $`\varphi (p,t)`$.
In the $`x`$-representation, the wave packet spreads in the $`x`$-direction, having the wave function
$`\psi (x,t)={\displaystyle \frac{1}{(2\pi (D(t))^{\frac{1}{4}}}}\mathrm{exp}`$ (69)
$`\left[{\displaystyle \frac{x^2}{4D(t)}}+i\left({\displaystyle \frac{x^2t}{8m(\mathrm{\Delta }x)^2D(t)}}{\displaystyle \frac{1}{2}}\mathrm{arctan}\left({\displaystyle \frac{t}{2m(\mathrm{\Delta }x)^2}}\right)\right)\right]`$ (70)
where $`D(t)=(\mathrm{\Delta }x)^2+\left(\frac{t^2}{4m^2(\mathrm{\Delta }x)^2}\right)`$. The corresponding current is
$$j_x=\frac{𝒫(x)}{m}\frac{xt}{4m(\mathrm{\Delta }x)^2D(t)}$$
(71)
and equation (22) yields
$$\frac{S}{t}+\frac{1}{2m}\left(\frac{S}{x}\right)^2+\frac{1}{4mD(t)}\frac{x^2}{8m[D(t)]^2}=0,$$
(72)
where the last two terms constitute the quantum potential.
This result can be easily understood since we are starting with the particle confined in a region $`\mathrm{\Delta }x`$ and, as time progresses, the wave packet spreads out as expected. The current $`j_x=𝒫(x)(dx/dt)`$ and the trajectories calculated from this current fan out in a way that exactly reflects the spread of the wave packet. As the wave packet spreads, the quantum potential reduces eventually to zero. Thus, for a particular trajectory, the energy of the quantum potential is progressively converted to the kinetic energy of the particle, so accelerating it away from its initial position.
### B The quadratic potential
Here we will simply collect the results derived earlier in the paper for ease of comparison.
In the $`x`$-representation, where we write $`\psi (x,t)=R_x\mathrm{exp}[iS_x]`$, the energy equation becomes
$$\frac{S_x}{t}+\frac{1}{2m}\left(\frac{S_x}{x}\right)^2+\frac{K}{2}x^2\frac{1}{2mR_x}\left(\frac{^2R_x}{x^2}\right)=0.$$
(73)
While in the $`p`$-representation, where we now write $`\psi (p,t)=R_p\mathrm{exp}[iS_p]`$, the conservation of energy equation is
$$\frac{S_p}{t}+\frac{p^2}{2m}+\frac{K}{2}\left(\frac{S_p}{p}\right)^2\frac{K}{2R_p}\left(\frac{^2R_p}{p^2}\right)=0$$
(74)
Now we turn to the probability currents and find
$$j_x=\frac{1}{2mi}\left[\psi ^{}(x)\left(\frac{\psi (x)}{x}\right)\left(\frac{\psi ^{}(x)}{x}\right)\psi (x)\right]$$
(75)
$$j_p=\frac{K}{2i}\left[\psi ^{}(p)\left(\frac{\psi (p)}{p}\right)\left(\frac{\psi ^{}(p)}{p}\right)\psi (p)\right],$$
(76)
in which the symmetry of the Hamiltonian is evident. That these currents are in fact different should not be too surprising as they arise in different spaces. Indeed, we can bring this out more clearly by using the respective polar forms of the x- and p-representation wave functions. In the x-representation
$$j_x=\frac{1}{m}R_x^2\left(\frac{S_x}{x}\right)$$
(77)
and so
$$\frac{dx}{dt}=\frac{1}{m}\left(\frac{S_x}{x}\right)=\frac{p_r}{m},$$
(78)
whereas in the p-representation
$$j_p=KR_p^2\left(\frac{S_p}{p}\right)$$
(79)
so that
$$\frac{dp}{dt}=K\left(\frac{S_p}{p}\right)=\left(\frac{V}{x}\right)_{x=x_r}.$$
(80)
Thus, we see that the currents provide the mathematical means of constructing trajectories in the $`x`$-space and $`p`$-space respectively. It is a feature of both the linear potential and the quantum harmonic oscillator that $`\frac{dp}{dt}=\left(\frac{V}{x}\right)_{x=x_r}`$, though this is not generally true.
### C The linear potential
Here the potential is $`V(x)=ax`$. In this case the current operators are
$$J_x=\frac{1}{2m}(\rho P+P\rho )$$
(81)
and
$$J_p=a\rho .$$
(82)
In the $`p`$-representation we find
$$j_p=p|J_p|p=𝒫(p)a=𝒫(p)\frac{dp}{dt}.$$
(83)
This result is identical to that obtained from classical mechanics through the equation
$$\frac{dp}{dt}=\frac{V}{x}=a$$
(84)
and suggests that the p-representation trajectories lie on the corresponding classical manifold. Indeed, equation (22) gives
$$\frac{S_p}{t}+\frac{p^2}{2m}a\frac{S_p}{p}=0.$$
(85)
Now using $`x_r=(S_p/p)`$, we find that the corresponding energy equation is
$$\frac{S_p}{t}+\frac{p^2}{2m}+ax_r=0,$$
(86)
which has the same form as the classical Hamilton-Jacobi equation with $`x_r=x`$. This confirms that, for the p-representation of the linear potential, there is no quantum potential and that the trajectories are indeed classical.
We now compare these results with those for the $`x`$-representation. Here the corresponding Schrödinger equation is
$$\frac{d^2\psi }{dx^2}A^3x\psi =0$$
(87)
where $`A=(2ma)^{\frac{1}{3}}`$. This equation has an Airy function
$$\psi (x)=CAi(Ax)$$
(88)
as a solution, which, being real, implies a zero probability density current
$$j_x=x|J_x|x=\frac{𝒫(x)}{m}\frac{S_x}{x}=𝒫(x)\frac{dx}{dt}=0.$$
(89)
Using this result in equation (22) shows that in the $`x`$-representation the quantum potential is the negative of the classical potential and is not zero as in the $`p`$-representation. This example demonstrates that, while they are consistent with the Schrödinger equation, Bohm trajectories may be representation dependent.
By way of explanation of the latter point, we observe that in the $`p`$-representation the wave function is complex, its incoming and outgoing components being separate on respectively the positive and the negative $`p`$-domains. On the other hand, in the $`x`$-representation, the incoming and outgoing waves combine to produce a real wave function. In particular, the $`x`$-representation solution may be split into incident and reflected components using the relation
$`Ai(Ax)+\mathrm{exp}({\displaystyle \frac{2}{3}}i\pi )Ai(Ax\mathrm{exp}({\displaystyle \frac{2}{3}}i\pi ))`$ (90)
$`+\mathrm{exp}({\displaystyle \frac{2}{3}}i\pi )Ai(Ax\mathrm{exp}({\displaystyle \frac{2}{3}}i\pi ))=0.`$ (91)
Taking the incident and reflected wave function separately, one obtains non-zero probability density currents and a non-zero quantum potential. The resulting trajectories are classical at infinity but are non-classical near the origin, where reflection takes place with an instantaneous change of sign in velocity. This is in contrast to the classical trajectory which turns smoothly at the origin. It is important to note that the trajectories of the incident and reflected waves respectively do not embody the effects of interference. It is this interference, absent in the $`p`$-representation, which produces a stationary trajectory for the combined solution in the $`x`$-representation.
### D The cubic potential
The quantum phase equation (22) in the $`x`$-representation using $`p_r=(S_x/x)`$ gives
$$\frac{S_x}{t}+\frac{p_r^2}{2m}+Ax^3\frac{1}{2mR_x}\frac{^2R_x}{x^2}=0,$$
(92)
while in the $`p`$-representation
$`{\displaystyle \frac{S_p}{t}}+{\displaystyle \frac{p^2}{2m}}+Ax_{r}^{}{}_{}{}^{3}+{\displaystyle \frac{3A}{R_p}}{\displaystyle \frac{^2R_p}{p^2}}\left({\displaystyle \frac{S_p}{p}}\right)`$ (93)
(94)
$`+{\displaystyle \frac{3A}{R_p}}\left({\displaystyle \frac{R_p}{p}}\right)\left({\displaystyle \frac{^2S_p}{p^2}}\right)+A\left({\displaystyle \frac{^3S_p}{p^3}}\right)=0,`$ (95)
where we have used $`x_r=(S_p/p)`$. This clearly gives a far more complicated quantum potential. Nevertheless, the content is still consistent with Schrödinger’s equation. Both equations reduce to the same classical Hamilton-Jacobi equation when the quantum potential terms reduce to zero.
The respective currents are
$$j_x=\frac{1}{2mi}\left[\psi ^{}(x)\left(\frac{\psi (x)}{x}\right)\left(\frac{\psi ^{}(x)}{x}\right)\psi (x)\right]$$
(96)
and
$$j_p=\frac{A}{2i}\left[\psi (p)\frac{^2\psi ^{}(p)}{p^2}+\psi ^{}(p)\frac{^2\psi (p)}{p^2}\frac{\psi (p)}{p}\frac{\psi ^{}(p)}{p}\right]$$
(97)
which gives
$$j_p=R_p^2\left(\frac{V}{x}\right)_{x=x_r}^2+A\left[2R_p\left(\frac{^2R_p}{p^2}\right)\left(\frac{R_p}{p}\right)^2\right]$$
(98)
as opposed to the simple expression for $`j_x`$
$$j_x=\frac{1}{m}R_x^2\left(\frac{S_x}{x}\right).$$
(99)
This clearly shows the limitation of using the condition $`p_r=(S_x/x)`$ as the guidance condition. It should also by now be quite clear that the Bohm trajectories in a particular representation are obtained from the probability current for that particular representation and not from any additional guidance condition.
## VII Conclusions
### A Algebraic formulation of the Schrödinger picture
In this paper we have shown how it is possible to write the content of the Schrödinger equation in algebraic form without reference to either Hilbert space or to any specific representation. The resulting two equations are respectively the Liouville equation, equation (19), and an equation that describes the time development of the phase, equation (20), which we have called the quantum phase equation. Furthermore, we have shown that this equation is gauge invariant and from it we calculated the Aharonov-Bohm, the Aharonov-Casher and the Berry phases in a simple and straight forward way.
We have also shown that it is possible to write the probability currents as algebraic operator forms. This allows us to define probability currents in any arbitrary representation. All of these results follow from the quantum formalism without the need to appeal to any classical formalism.
### B The $`x`$ and $`p`$ representations: the quantum potential and the trajectories of probability current
In sections III and IV, we expressed equations (21) and (22) in the $`x`$-representation (equations (58) and (45) respectively) and showed that they are identical to the two defining equations of the traditional Bohm interpretation . The quantum potential emerges from equation (42), which in turn comes directly from equation (22), showing that it cannot be “dismissed as artificial and obscuring the essential meaning of the Bohm approach” without missing some of the essential novel features of quantum processes.
In particular, observed characteristic quantum phase or gauge effects come directly from equation (22). As Philippidis, Bohm and Kaye have shown many years ago, the presence of the AB effect alters the quantum potential, which in turn accounts for the fringe shifts. Furthermore, it is the presence of the quantum potential that offers an explanation of Einstein-Podolsky-Rosen-type correlations , as well as quantum state teleportation .
In section V, we also showed that we can construct a BI in the $`p`$-representation. Comparing representations shows very clearly that the Bohm trajectories are simply the trajectories associated with the probability currents of the standard theory. The only assumption added to the standard quantum theory in the Bohm-Hiley version of the BI is that particles have simultaneously well defined positions and momenta and actually follow these trajectories. Further, we claim that this position is implicit in pragmatic quantum mechanics in which the probability currents are assumed to be related to particle fluxes.
### C Shadow phase spaces
The central point that emerges from our approach is that we can construct two different phase spaces. In the example of the harmonic oscillator, the $`x`$-phase space is based on equations (44) and (45), while the $`p`$-phase space is built using equations (46) and (47). The explication of different phase spaces was further exemplified in section VI. The reason why we must resort to constructing different phase spaces is not too difficult to see once it is realised that we are dealing with a non-commutative structure<sup>††</sup><sup>††</sup>†† A simple example of this kind of structure will be found in Hiley , and Hiley and Monk ..
For a commutative algebra, the Gel’fand construction allows us to start from the algebra and re-construct the underlying manifold . Here the points, the topology and the metric structure of the manifold are all carried by the algebra. No such construction is possible for a non-commutative algebra. Thus, in our case there is no underlying phase space with points that can be specified by the pair of observables $`(x,p)`$. This is just what the uncertainty principle is telling us. This is the physicist’s way of explanation why there can be no single, unique, underlying continuous phase space.
Any attempt to produce a single phase space, such as is done in the Wigner-Moyal approach, must necessarily contain unacceptable features . In this case, the probability distribution can be negative in certain situations. For these reasons, we must follow what is usually done in non-commutative geometry and construct shadow manifolds.
In this context, equation (42) provides an explanation as to why the energy can be conserved when we attribute to the particle at position $`x`$, the beable momentum $`p_r=S_x/x`$. Since it is a constructed momentum and not a measured momentum, the kinetic energy will not have the value necessary to conserve the total energy. Thus we need another term to “carry” this difference. Since equation (45) is the expression for the conservation of energy, the last term on the RHS of equation (45) is the place to “store” this energy difference. This shows that the quantum potential energy is an internal energy, and clearly does not have an external source.
### D Implications for the Bohm Interpretation of quantum mechanics
Finally, we will briefly comment on the implications of the above analysis on the BI. The traditional BI assumes the $`x`$-representation is special, but the reasons for this were never made clear. It was generally assumed that all physics must take place in an a priori given space-time arena, a point of view that we have called the Cartesian order. Hence the attempt to use the guidance condition as a defining equation for Bohmian mechanics . However our mathematical analysis above shows that this condition is a contingent feature, which is dependent on the asymmetry of the Hamiltonian.
On the other hand, if we take the quantum formalism as primary then we must place our emphasis on the non-commutative structure of the algebra of formalism. If we do this then attempts to focus on a single phase space, which is equivalent to giving primary relevance to space-time, will fail. This in turn calls into question the way we think about quantum processes. Indeed Bohm has already argued that we must abandon the Cartesian order and replace it by a radically new approach to quantum phenomena which he called the implicate order . Here the ontology is provided by the concept of process which is to be described by the non-commutative algebra. This is not a process in space-time, but a process from which space-time is to be abstracted. Abstraction here means to ‘make manifest’ and the order that is made manifest is called the explicate order.
The key point about this view is that there may be more than one explicate order and that these explicate orders cannot be made manifest together at the same time. This can be regarded as a direct consequence of the participatory nature of the quantum process. Thus the implicate order contains an ontological complementarity, which is a necessary consequence of the non-commutative structure. In this picture the BI discussed above is said to contain two explicate orders, one depending on the $`x`$-representation and the other on the $`p`$-representation. These are the shadow phase spaces. Both are equally valid descriptions of the outward appearance of a quantum process within a given context.
Since our classical world is dominated by appearances in space-time, we would expect the most relevant explicate order to be that based on the $`x`$-representation, with the context being provided by the classical world. This is the world in which we place our apparatus and where our measurements take place. But clearly we need to explore these ideas further as a number of questions remain unanswered. We will leave this discussion for another paper.
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# LANCS-THMay 2000 Gravity localization with a domain wall junction in six dimensions
## 1 Introduction
Possibilities of extra dimensions have been studied for a long time, and the existence is well motivated by superstring theories. A traditional method to hide extra dimensions is an idea of compactification in which the extra dimensions are supposed to be extremely small.
In a couple of years, it has been recognized that the extra dimensions may have sub-millimeter size or infinitely large volume if the standard model fields are confined at a three brane. In particular Randall and Sundrum recently proposed an scenario alternative to compactification using exponentially warped metric in a five-dimensional gravity model . Due to the warped metric, the four-dimensional massless graviton are localized on a brane, and the four-dimensional Newton law is approximately realized on the brane at low energy. Supersymmetrization of this scenario has been discussed in Ref. .
In the Randall-Sundrum scenario, cosmological constants are introduced for both the five-dimensional bulk and the four-dimensional branes, and the warped metric is derived as a solution to the Einstein equation . In order to obtain the solution with the four-dimensional Poincaré invariance, however, the cosmological constants have to be specially tuned. Stabilization mechanism of the extra dimension has been discussed by introducing a bulk scalar field .
In order to have a natural explanation of this tuning, we have to discuss the origin of the brane cosmological constant. Instead of pursuing string theories, we consider a the field theoretic approach in this paper. Several works using domain wall solutions in five-dimensional gravity models have been done . In these analyses, a supergravity-motivated scalar potential is introduced , and the domain wall solution of the scalar fields produces an effective cosmological constant on the brane to implement the warped metric in the Randall-Sundrum scenario. This line is a gravity version of a previous idea of living in a domain wall . An analysis on domain walls in arbitrary dimensions is given in Ref. . Domain wall solutions in four-dimensional supergravity models have been studied in Ref. .
However it has recently been pointed out that smooth domain wall solutions interpolating between supersymmetric vacua can not exist in odd-dimensional supergravity models . This naturally leads us to work in a framework of six-dimensional supergravity models . In this case we need a two-dimensional topological solution like a vortex solution . Gravity localization on a string-like defect in six dimensions has been studied in Ref. .
There is another interesting two-dimensional stable solitonic solution, namely, a domain wall junction solution in supersymmetric models . The domain wall junction preserves one-quarter of the underlying supersymmetry , and satisfies a first order BPS equation . With domain wall junctions, general analyses on gravity localization in infinitely large extra dimensions have been given in Refs. .
In this paper we study gravity localization in the context of a six-dimensional gravity model coupled with complex scalar fields. Similar analyses have been done in Refs. . With a supergravity-motivated scalar potential, we derive first order equations which the metric and the scalar fields should satisfy. We calculate the tensions of domain wall and juncion. Finally we study a warped metric and discuss gravity localization on the junction.
## 2 Set-up
We consider a six-dimensional gravity coupled with complex scalar fields $`\varphi ^i`$ ($`i`$ $`=`$ 1, $`\mathrm{}`$, $`N`$). We use coordinates $`x^M`$ $`=`$ $`(x^\mu ,x^m)`$, where $`M`$ $`=`$ $`0,1,2,3,5,6`$, $`\mu `$ $`=`$ $`0,1,2,3`$ and $`m`$ $`=`$ $`5,6`$. $`x^\mu `$ are coordinates for the observable four spacetime dimensions. $`x^m`$ are coordinates for the two extra dimensions, where $`\mathrm{}`$ $`<`$ $`x^m`$ $`<`$ $`\mathrm{}`$. The action is given by
$`S`$ $`=`$ $`{\displaystyle d^6x\sqrt{g}\left[\frac{1}{2\kappa ^2}R+K_{ij^{}}g^{MN}_M\varphi ^i_N\varphi ^jV(\varphi ,\varphi ^{})\right]},`$ (1)
where $`g_{MN}`$ is the metric in a time-like signature convention ($`+`$, $``$, $``$, $``$, $``$, $``$). $`V(\varphi ,\varphi ^{})`$ is a scalar potential. The scalar kinetic term has field dependent coefficients $`K_{ij^{}}(\varphi ,\varphi ^{})`$ which are derived from the Kähler potential as $`K_{ij^{}}`$ $`=`$ $`^2K(\varphi ,\varphi ^{})/\varphi ^i\varphi ^j`$. In the following, we adopt the six-dimensional Planck mass unit $`\kappa ^2`$ $`=`$ 1 unless otherwise stated. We put the following ansatz for the background metric
$`ds^2`$ $`=`$ $`g_{MN}dx^Mdx^N`$ (2)
$`=`$ $`e^{2\sigma (x^5,x^6)}\left[\eta _{\mu \nu }dx^\mu dx^\nu (dx^5)^2(dx^6)^2\right],`$
where $`\eta _{\mu \nu }`$ is a four-dimensional flat metric. This metric guarantees the four-dimensional Poincaré invariance. We look for static scalar configurations which depend on $`x^m`$ only: $`\varphi ^i`$ $`=`$ $`\varphi ^i(x^5,x^6)`$. Then the equations of motion can be written as
$`_m^2\sigma +4(_m\sigma )^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2\sigma }V,`$ (3)
$`4_m_n\sigma +4(_m\sigma )(_n\sigma )`$ $`=`$ $`K_{ij^{}}(_m\varphi ^i_n\varphi ^j+_n\varphi ^i_m\varphi ^j),`$ (4)
$`_m^2\varphi ^i+4(_m\sigma )(_m\varphi ^i)`$
$`+K^{ij^{}}(_kK_{lj^{}})(_m\varphi ^k)(_m\varphi ^l)`$ $`=`$ $`e^{2\sigma }K^{ij^{}}{\displaystyle \frac{V}{\varphi ^j}},`$ (5)
where repeated indices are summed. The first two equations correspond to $`(\mu ,\nu )`$ and $`(m,n)`$ components of the Einstein equations, respectively. The last one is the scalar field equation.
For the scalar potential, we assume the following form motivated by supergravity models :
$`V`$ $`=`$ $`e^K\left[K^{ij^{}}(D_iW)(D_jW)^{}{\displaystyle \frac{5}{2}}|W|^2\right],`$ (6)
where $`D_iW`$ $`=`$ $`W/\varphi ^i`$ $`+`$ $`(K/\varphi ^i)W`$ and $`W`$ is an arbitrary function of $`\varphi ^i`$ which may be interpreted as a superpotential if we can derive this potential from a supergravity model. A five-dimensional gauged supergravity model has a similar potential . Given the potential (6), we can obtain the first order equations which the metric and the scalar fields should satisfy . We will see this in the next section. These first order equations make it easy to solve the equations of motion (3)-(5).
In general the potential (6) is not bounded below because of $`|W|^2`$ term. Furthermore in some higher-dimensional supergravity models, there are no local minima and all the extrema are local maxima or saddle points. However it has been shown that even a vacuum at a maximum is stable under local fluctuations around an AdS background unless the curvature of the potential at the maximum is too negative . In this sense the potential (6) is sensible.
No six-dimensional supergravity models which provide the potential (6) have been constructed, so we can not work in a supergravity context. In this paper we just assume the potential (6) from the beginning.
## 3 Tensions of domain wall and junction
In this section, we derive the first order equations for the metric and the scalar fields. We also obtain the formula for the tensions of domain wall and junction.
With the potential (6), we can write the action as a sum of perfect squares up to total derivative terms . For this purpose, it is convenient to define a complex coordinate variable $`z`$ $`=`$ $`x^5`$ $`+`$ $`ix^6`$ and derivatives $``$ $`=`$ $`(_5i_6)/2`$, $`\overline{}`$ $`=`$ $`(_5+i_6)/2`$. In terms of this variable, the action (1) is written as
$`S`$ $`=`$ $`{\displaystyle }d^6xe^{6\sigma }[20e^{2\sigma }(\overline{}\sigma +2\sigma \overline{}\sigma )2e^{2\sigma }K_{ij^{}}(\varphi ^i\overline{}\varphi ^j+\overline{}\varphi ^i\varphi ^j)`$ (7)
$`e^K\{K^{ij^{}}(D_iW)(D_jW)^{}\frac{5}{2}|W|^2\}].`$
We can make a perfect square, e.g., from a part of the $`\sigma \overline{}\sigma `$ term and the $`|W|^2`$ term. In doing that we can introduce a complex phase $`\theta (x^5,x^6)`$ as follows:
$`\left|\sigma \frac{1}{4}e^\sigma e^{K/2}e^{i\theta }W^{}\right|^2.`$ (8)
Taking into account this phase, the action can be written as a sum of a local contribution and a topological term
$`S`$ $`=`$ $`S_{\mathrm{local}}+S_{\mathrm{topological}}.`$ (9)
The local contribution is given by a sum of perfect squares
$`S_{\mathrm{local}}`$ $`=`$ $`{\displaystyle }d^6xe^{4\sigma }[40|D\sigma |^24K_{ij^{}}(D\varphi ^i)(D\varphi ^j)^{}\{4i(D\theta )(D\sigma )^{}+\mathrm{h}.\mathrm{c}.\}]`$ (10)
$`=`$ $`{\displaystyle d^6xe^{4\sigma }\left[40\left|D\sigma \frac{i}{10}D\theta \right|^24K_{ij^{}}(D\varphi ^i)(D\varphi ^j)^{}\frac{2}{5}\left|D\theta \right|^2\right]},`$
where the symbols $`D\sigma `$, $`D\varphi ^i`$ and $`D\theta `$ are defined by
$`D\sigma `$ $`=`$ $`\sigma \frac{1}{4}e^\sigma e^{K/2}e^{i\theta }W^{},`$ (11)
$`D\varphi ^i`$ $`=`$ $`\varphi ^i+\frac{1}{2}K^{ij^{}}e^\sigma e^{K/2}e^{i\theta }(D_jW)^{},`$ (12)
$`D\theta `$ $`=`$ $`\theta \frac{i}{2}\left(K_i\varphi ^iK_i^{}\varphi ^i\right).`$ (13)
The perfect squares in the integrand imply that the configurations which extremize the action satisfy the first order equations $`D\sigma `$ $`=`$ $`D\varphi ^i`$ $`=`$ $`D\theta `$ $`=`$ 0, namely,
$`\sigma `$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{4}}e^\sigma e^{\kappa ^2K/2}e^{i\theta }W^{},`$ (14)
$`\varphi ^i`$ $`=`$ $`{\displaystyle \frac{1}{2}}K^{ij^{}}e^\sigma e^{\kappa ^2K/2}e^{i\theta }(D_jW)^{},`$ (15)
$`\theta `$ $`=`$ $`{\displaystyle \frac{i}{2\kappa ^2}}\left(K_i\varphi ^iK_i^{}\varphi ^i\right).`$ (16)
Here we have recovered the gravitational coupling explicitly. Note that the solutions to equations (14)-(16) automatically satisfy the equations of motion (3)-(5). In the four-dimensional supergravity model, similar equations are derived from conditions for the existence of unbroken supersymmetry. The equation (15) implies that the scalar fields $`\varphi ^i`$ should depend on both $`z`$ and $`z^{}`$ in order to have a non-trivial configuration.
Because of the positive contribution $`40|D\sigma |^2`$ in eq. (10), the solutions to these first order equations do not minimize the energy in general. However on an AdS background, the vacuum at an extremum is stable under local fluctuations unless the curvature of the potential at the maximum is too negative .
The equation (16) can be written as $`_m\theta `$ $`=`$ $`\mathrm{Im}(K_i_m\varphi ^i)`$. In the four dimensional supergravity model, the quantity $`\mathrm{Im}(K_i_m\varphi ^i)`$ corresponds to an auxiliary vector field in a supergravity multiplet. This field acts as a gauge field for the Kähler-Weyl invariance.
The topological term consists of two kinds of total derivative terms
$`S_{\mathrm{topological}}`$ $`=`$ $`{\displaystyle d^4x\underset{m}{}(Z_m+Y_m)}.`$ (17)
The first contributions $`Z_m`$ are given by
$`Z_5`$ $`=`$ $`{\displaystyle 𝑑x^5𝑑x^6_5\left[e^{4\sigma }(5_5\sigma )e^{5\sigma }e^{K/2}(e^{i\theta }W+e^{i\theta }W^{})\right]},`$
$`Z_6`$ $`=`$ $`{\displaystyle 𝑑x^5𝑑x^6_6\left[e^{4\sigma }(5_6\sigma )+ie^{5\sigma }e^{K/2}(e^{i\theta }We^{i\theta }W^{})\right]}.`$ (18)
These include the domain wall tensions and the metric terms. The second contributions are
$`Y_5`$ $`=`$ $`{\displaystyle 𝑑x^5𝑑x^6_5\left[e^{4\sigma }\left\{_6\theta +\mathrm{Im}(K_i_6\varphi ^i)\right\}\right]},`$
$`Y_6`$ $`=`$ $`{\displaystyle 𝑑x^5𝑑x^6_6\left[e^{4\sigma }\left\{_5\theta +\mathrm{Im}(K_i_5\varphi ^i)\right\}\right]}.`$ (19)
These include the ordinary domain wall junction tentions $`\mathrm{Im}(K_i_m\varphi ^i)`$ and ‘gravitational’ contributions $`_m\theta `$. Note that $`Y_m`$ can have non-vanishing values only when the fields have two-dimensional non-trivial configurations, while $`Z_m`$ can be non-vanishing even when the fields depends on only one of $`x^m`$.
In the absence of gravity, the derivatives of $`\theta `$ do not contribute, and the terms $`\mathrm{Im}(K_i_m\varphi ^i)`$ give the domain wall junction tension. Note that $`_m\theta `$ terms in eq. (19) are proportional to $`\kappa ^2`$ if we write the gravitational coupling explicitly. In the presence of gravity, however, the derivatives of $`\theta `$ contribute to $`Y_m`$. Substituting the first order equation (16) into eq. (19), we find that the two contributions $`_m\theta `$ and $`\mathrm{Im}(K_i_m\varphi ^i)`$ cancel locally in the integrand with each other. Therefore the junction tentions $`Y_m`$ vanish
$`Y_m`$ $`=`$ $`0.`$ (20)
This implies that in the presence of the gravitational degrees of freedom, a constant $`\theta `$ configuration can not extremize the action. The system requires a non-trivial $`\theta `$ to extremize the action in such a way that $`_m\theta `$ cancel the ordinary contributions $`\mathrm{Im}(K_i_m\varphi ^i)`$ to the junction tension.
Let’s explain briefly the reason why $`_m\theta `$ and $`\mathrm{Im}(K_i_m\varphi ^i)`$ appear in the vanishing combination $`D\theta `$. In order to see this, it is essential to observe that the total derivative in the domain wall tension term produce the vanishing combination as follows:
$`_m\left[e^{K/2}e^{i\theta }W\right]`$ $`=`$ $`e^{K/2}e^{i\theta }\left[(D_iW)_m\varphi ^iiW(_m\theta +\mathrm{Im}K_i_m\varphi ^i)\right].`$ (21)
We can see that in the other parts of calculations, $`_m\theta `$ and $`\mathrm{Im}(K_i_m\varphi ^i)`$ always appear in the same combination. If we work in the four-dimensional supergravity model, we may have deeper understanding of this result referring to the Kähler-Weyl invariance.
As for the wall contributions $`Z_m`$, the situation is different. Substituting the first order equation (14) into eq. (18), we see that $`Z_m`$ are written as
$`Z_5`$ $`=`$ $`{\displaystyle 𝑑x^5𝑑x^6_5\left[\frac{1}{2}e^{5\sigma }e^{K/2}\mathrm{Re}(e^{i\theta }W)\right]},`$
$`Z_6`$ $`=`$ $`{\displaystyle 𝑑x^5𝑑x^6_6\left[\frac{1}{2}e^{5\sigma }e^{K/2}\mathrm{Im}(e^{i\theta }W)\right]}.`$ (22)
Unlike the junction tensions, the wall tensions $`Z_m`$ do not vanish in general.
## 4 Warped metric from domain wall junction
In this section, we discuss the solution to the first order equations (14)-(16). In most cases the analytic solution is not available, so in this papar we only discuss rough behavior of the solution. As an example for the function $`W`$, let’s consider a quartic form for a single complex scalar field $`\varphi `$
$`W`$ $`=`$ $`{\displaystyle \frac{1}{4}}\varphi ^4+\varphi .`$ (23)
For the Kähler potential $`K`$, we take the minimal form $`K`$ $`=`$ $`\varphi \varphi ^{}`$. It is known that in the global supersymmetric model such a quartic superpotential $`W`$ allows static domain wall junction solutions . In the global case, the scalar potential has three isolated degenerate minima at $`\varphi `$ $`=`$ 1, $`\omega `$, $`\omega ^2`$, where $`\omega `$ $`=`$ $`e^{2\pi i/3}`$. Note that the potential does not allow static domain wall solutions in the four-dimensional supergravity model . This comes from the reality of the metric. In the domain wall junction case, however, the same discussion cannot be applied since the first order equation (14) for the metric includes the derivative with respect to the complex coordinate $`z`$.
With the function $`W`$ in eq. (23) and the minimal Kähler potential $`K`$, we find three vacua $`\varphi `$ $`=`$ $`k`$, $`k\omega `$, $`k\omega ^2`$ by solving $`D_\varphi W`$ $`=`$ 0. Here $`k`$ $``$ 1.2 is a single solution of a quintic equation $`k^5`$ $`+`$ $`4k^3`$ $`4k^2`$ $`4`$ $`=`$ 0. The potential (6) and the first order equations (14)-(16) in this case are invariant under $`𝐙_3`$ action $`z`$ $``$ $`\omega z`$, $`\varphi `$ $``$ $`\omega \varphi `$, $`\sigma `$ $``$ $`\sigma `$, $`\theta `$ $``$ $`\theta `$. Therefore we expect a $`𝐙_3`$ invariant domain wall junction solution. The solution describes the two-dimensional space separated into three regions by three walls, and these regions are labeled by the vacua $`k`$, $`k\omega `$, $`k\omega ^2`$. Equations (15) and (16) imply that outside the walls and the junction the scalar fields $`\varphi `$ and the phase $`\theta `$ are almost constant. We consider the solution with arg$`\varphi `$ $`=`$ arg$`(z)`$ outside the wall so that the junction is centered at $`z`$ $`=`$ 0. Then the first order equation (14) leads to the asymptotic behavior of the metric at spacial infinity in the extra dimensions
$`e^{2\sigma }`$ $``$ $`{\displaystyle \frac{C^2}{[\mathrm{Re}(\omega ^nz)]^2}},|z|\mathrm{},`$ (24)
where $`n`$ $`=`$ 0, 1, 2 label the vacua $`k`$, $`k\omega `$, $`k\omega ^2`$, respectively. The constant $`C`$ is given by $`C`$ $`=`$ $`2e^{k^2}(kk^4/4)^1`$, where we have chosen $`\theta `$ $`=`$ 0 outside the wall. The above metric (24) describes an AdS background. This asymptotic form is consistent with analyses in Refs. . From this behavior it follows that the extra dimension is infinitely large, since the volume $`V`$ $`=`$ $`𝑑x^5𝑑x^6`$ $`e^{2\sigma }`$ $``$ $`𝑑r/r`$ is divergent where $`r`$ $`=`$ $`|z|`$.
Note that this behavior (24) can be applied only outside the walls. Inside the wall the scalar field is not constant any more, hence we have to solve the coupled equations (14)-(16) to know the profile. However it seems natural to assume that for a fixed $`r`$, the wall essentially has a kink-like profile. Then the evolution equation (15) along the direction perpendicular to the wall is given by $`\varphi `$ $``$ $`r^1`$ $`(D_\varphi W)^{}`$ inside the wall. This is consistent with the metric behavior (24) outside the wall. The metric may have a peak structure in the wall , but the height of the peak scales as $`1/r^2`$. Thus the scaling property $`e^{2\sigma }`$ $``$ $`1/r^2`$ holds even in the walls. Also this assumption means that the wall width grows up like $``$ $`r`$ far from the origin. On the other hand the ‘height’ of the wall remains the same even far from the origin. Consequently, the wall structure almost disappears far from the junction due to the metric suppression.
It is difficult to discuss the solution inside the junction, too. However we can see the behavior of the metric near the origin $`x^m`$ $`=`$ 0 in a weak gravitational coupling approximation. In the $`\kappa ^2`$ $`=`$ 0 limit, the metric $`\sigma `$ and the phase $`\theta `$ are constant, and eq. (15) reduces to a simpler equation. From symmetry consideration the scalar field must vanish at the origin. Then it follows that $`\varphi `$ $``$ $`z`$ from eq. (15). In the first order in $`\kappa ^2`$, the metric equation (14) becomes $`\sigma `$ $`=`$ $`W^{}/4`$ where we have substituted the zeroth order result in the right-hand side. Using this equation, we obtain $`\sigma `$ $``$ $`zz^{}/4`$ near the origin so that
$`e^{2\sigma }`$ $``$ $`1{\displaystyle \frac{zz^{}}{2}},|z|1.`$ (25)
The curvature near the origin is constant $`R`$ $``$ $`10`$ up to the order of $`|z|^2`$.
Let’s estimate the topological charges $`Z_m`$ in this example. Unlike the case of $`Y_m`$, the integrand of $`Z_m`$ in eq. (22) does not vanish locally. In this example, however, the integrand goes to zero at spatial infinity in the extra dimensions faster than $`1/r`$ due to the asymptotic power law suppression of the metric (24). Therefore, using the Stokes’s theorem, the domain wall tensions also vanish
$`Z_m`$ $`=`$ $`0.`$ (26)
Given $`Y_m`$ $`=`$ 0 in eq. (20), we find that all the topological contribution to the action vanish. Namely, the four-dimensional cosmological constant is zero. This is consistent with our ansatz of the flat four-dimensional spacetime in eq. (2).
The discussion here is based on a qualitative consideration. In particular the behavior of the solution inside the walls and the junction has not been obtained. We need to study the first order equations numerically to know the precise behavior of the solution.
We normally need a fine tuning of parameters to obtain a vanishing four-dimensional cosmological constant. This is also the case in our model. The source of the fine tuning in our model lies in choosing the particular form of the scalar potential (6) rather than the detail of the function $`W`$. The same situation has been observed in five-dimensional models with supergravity-motivated scalar potential in Refs. . At present, no symmetries to guarantee the form of the potential (6) have been found. So we cannot solve the cosmological constant problem in our model.
## 5 Gravity localization
In this section we study the fluctuations of the metric defined by
$`ds^2`$ $`=`$ $`g_{MN}^{(0)}dx^Mdx^N+h_{\mu \nu }dx^\mu dx^\nu ,`$ (27)
where the first term in the right-hand side is the background metric in eq. (2). In this analysis, we focus on the transverse traceless modes which satisfy $`h_\mu ^\mu `$ $`=`$ $`^\mu h_{\mu \nu }`$ $`=`$ 0.
In order to obtain the linealized equation for the fluctuation, we expand the action (1) in terms of $`h_{\mu \nu }`$ around the background metric up to the second order of the fluctuations. The second order terms are given by
$`S^{(2)}`$ $`=`$ $`{\displaystyle }d^6x[\frac{1}{8}h^{\mu \nu }\{\mathrm{}_4_m^2+2(_m^2\sigma )+4(_m\sigma )^2\}h_{\mu \nu }`$ (28)
$`+\{_m^2\sigma +\frac{3}{2}(_m\sigma )^2+\frac{1}{4}(K_{ij^{}}_m\varphi ^i_m\varphi ^j+e^{2\sigma }V)\}h^{\mu \nu }h_{\mu \nu }],`$
where upper Lorentz indices are defined by the flat metric $`\eta _{\mu \nu }`$. The terms in the second curly bracket in the right-hand side cancel because of the first order equations (14)-(16). Variation of $`S^{(2)}`$ with respect to $`h_{\mu \nu }`$ gives rise to the linearized equation of motion for the metric fluctuation
$`\left[\mathrm{}_4_m^2+2(_m^2\sigma )+4(_m\sigma )^2\right]h_{\mu \nu }=0.`$ (29)
For the fluctuation with the four-dimensional dependence of a plane wave $`e^{ipx}`$, this equation can be written as a two-dimensional Schrödinger equation $`(_m^2+V_{\mathrm{QM}})`$ $`h_{\mu \nu }`$ $`=`$ $`p^2h_{\mu \nu }`$ with a potential
$`V_{\mathrm{QM}}`$ $`=`$ $`2(_m^2\sigma )+4(_m\sigma )^2.`$ (30)
In the following we again consider the quartic function $`W`$ in eq. (23) and the minimal Kähler potential. Using eq. (24), the asymptotic behavior of this potential far from the junction is given by
$`V_{\mathrm{QM}}`$ $``$ $`{\displaystyle \frac{6}{[\mathrm{Re}(\omega ^nz)]^2}}.`$ (31)
In five-dimensional models, the solutions to the corresponding Schrödinger equation for $`p^2`$ $`>`$ 0 are described by the Bessel functions $`\sqrt{x}J_2(px)`$ and $`\sqrt{x}Y_2(px)`$ . In our case, the equation involves the two variables, and the complete solution is not available. However the potential (31) shows that outside the wall, the Schrödinger equation is essentially one dimensional equation. Therefore the solutions for $`p^2`$ $`>`$ 0 can be described by the Bessel functions $`\sqrt{x}J_{5/2}(px)`$ and $`\sqrt{x}Y_{5/2}(px)`$ outside the wall, where $`x`$ $`=`$ $`\mathrm{Re}(\omega ^nz)`$. These can be written by trigonometric functions only. The similar solutions have been found in Ref. .
In the case of the massless fluctuation $`p^2`$ $`=`$ 0, the solution to eq. (29) is given by
$`h_{\mu \nu }=e^{2\sigma }e^{ipx}\eta _{\mu \nu }.`$ (32)
This means that the massless fluctuation has the same configuration as the background metric in the extra dimensions. From eq. (24) and our assumption on the wall profile, we find that the massless graviton is localized on the junction. This agrees with the result in Ref. . We see that the massless mode (32) is normalizable on our curved background
$`{\displaystyle 𝑑x^5𝑑x^6e^{2\sigma }h^{\mu \nu }h_{\mu \nu }}`$ $`<`$ $`\mathrm{}.`$ (33)
In the transverse traceless components which we have considered, there are no tachyonic modes. This is shown by writing eq. (29) as follows
$`Q_m^{}Q_mh_{\mu \nu }=p^2h_{\mu \nu },`$ (34)
where $`Q_m`$ $``$ $`_m`$ $`+`$ $`2(_m\sigma )`$. In the flat space, $`Q_m^{}`$ $``$ $`_m`$ $`+`$ $`2(_m\sigma )`$ is the adjoint of $`Q_m`$. Similar equation appears in supersymmetric quantum mechanics. The solutions to the Schrödinger equation for $`p^2`$ $`<`$ 0 are described by the modified Bessel functions $`\sqrt{x}I_{5/2}(px)`$ and $`\sqrt{x}K_{5/2}(px)`$. Following the discussion in Ref. , we can see that normalizable modes always satisfy $`p^2`$ $``$ 0 for the transverse traceless components.
## 6 Summary
In summary, we have studied a six-dimensional gravity coupled with complex scalar fields. With a supergravity-motivated scalar potential, the domain wall junction solutions localize a four-dimensional massless graviton. We have shown that unlike the global supersymmetric model, contributions to the junction tension cancel locally with gravitational contributions. The wall tension vanishes due to the metric suppression.
## Acknowledgements
The author would like to thank L. Roszkowski and H.B. Kim for useful conversations. This work was supported in part by PPARC grant PPA/G/S/1998/00646.
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# Functional Approach to Quantum Decoherence and the Classical Final Limit.
## I Introduction.
Following the idea that the interplay of observables and states is the fundamental ingredient of quantum mechanics <sup>*</sup><sup>*</sup>*According to W. Zurek:…”the only sensible subject of consideration aimed at the interpretation of quantum theory… is the relation between the universal state vector and the state memory (records) of somewhat special system - such as observers \- which are, for necessity, perceiving the Universe from within. It is the inability to appreciate the consequences of this rather simple but fundamental observation that has led to such desperate measures as the search of an alternative quantum physics ”. we have developed paper where we have studied the relation of the state vectors $`\rho `$ of a close isolated quantum system (that belong to a convex set of states $`𝒮`$), to the observables $`O`$ within this closed system (that belong to a space of observables $`𝒪`$). We consider that the essence of this relation is the mean value of an observable $`O`$ in a state $`\rho ,`$ which is given by the equation:
$$O_\rho =Tr(\rho O)=(\rho |O)$$
(1)
In fact at the statistical level what we actually measure in an ensemble of identical states, are these kind of averages, since we cannot either measure directly the state $`\rho `$ or measure it with an infinite precision . Moreover, these averages can be considered, as in the r. h. s. of eq. (1), the result of a linear functional $`(\rho |𝒮`$ acting on a vector $`|O)𝒪,`$ and therefore we can say that $`𝒮𝒪^{^{}}`$, being $`𝒪^{^{}}`$ the dual of space $`𝒪`$. While for the usual states (mixed or pure) we can use $`Tr(\rho O)`$, there are generalized states that can be defined as the functional $`(\rho |O)`$ as explained in papers . Many results were obtained using this formalism (see e. g.: , , ).
In this paper we will use the formalism of paper to study the so called ”classical limit problem”, namely the statistical quantum mechanics $``$ classical mechanics limit that appears in some quantum systems when observed using certain spaces of observables $`𝒪`$ . For conceptual reasons we will divide the problem in two different processes (that may or may not happen simultaneously):
(a)- Statistical process: Namely the limit statistical quantum mechanics $``$ statistical classical mechanics, where the phenomenon of decoherence combined with the disappearance of the uncertainty relations in the limit $`\mathrm{}0,`$ originates the classical final stationary state. Almost all the paper will be devoted to this problem . We will see how when $`t\mathrm{}`$ the quantum system reaches a classical final stationary state $`\rho _{}(q,p)`$ where the statistical dynamics is trivial, since $`\rho _{}(q,p)`$ is time independent, but the systems of the ensemble move according to the non-trivial laws of classical dynamics. In general we will have an unlocalized statistical classical state of many identical systems moving in phase space.
(b)- Localization process: It is the evolution statistical classical mechanics $``$ classical mechanics. In some special cases the evolution privileges a single space-time trajectory, in such a way that all trajectories (endowed with a non negligible positive probability) concentrate around itIn some cases this phenomenon does not happen for all the systems but only for a subsystem.. In this case we will have correlations and localization. Then we have the statistical classical state of all the systems practically moving along the same trajectory in such a way that we may consider that we are dealing with a single classical system. We will discuss this process in section IV and appendix B.
The usual technique to solve these problems is coarse-graining. But in our method we will consider not only the coarse-graining average but all possible averages made using the observables of space $`𝒪`$, thus we are generalizing the coarse-graining idea At least the ”coarse graining alla Zurek”.. In fact, among the observables of $`𝒪`$ there are some that, from the density matrix $`\rho ,`$ take into account only some component $`\rho _r`$, the so called relevant part of $`\rho ,`$ and completely neglect the complementary component $`\rho _i,`$ the so called irrelevant part of $`\rho ,`$ i. e. these observables only measure (macroscopic) properties of what it is considered as the ”system” (contained in $`\rho _r)`$ and neglect or average the (microscopic) properties of the ”environment” (contained in $`\rho _i).`$ But we will consider not only this kind of observables but all observables in $`𝒪.`$ Therefore the interplaying of observables and states will take the role of the coarse-graining in this paper (see also the end of section IIA). With this strategy we cannot only obtain all the old results, but also we will find some new ones.
We will use this method to study the process (a) and to prove that certain quantum systems evolve from a statistical quantum state to the statistical classical final stationary state. In the same framework we will study the process (b) obtaining the classical motion of a single system.
The paper is organized as follows:
In section II we will see, using the Riemann-Lebesgue theorem, that transition (a) takes place in close systems endowed with a continuous spectrum and with just one bound state (as in the classical mixing systems). More general cases will be considered in section II C. The main characteristics of the quantum laws are:
1.- The non-boolean nature of the way to find the probability of two exclusive events (this probability is the square modulus of the sum of their amplitudes and not the sum of the probabilities).
2.- The uncertainty relations.
In the evolution from quantum mechanics to classical mechanics the first characteristic disappears (and the boolean way of adding probabilities is established) by the process of decoherence and the uncertainty relations can be neglected in the limit $`\mathrm{}0.`$ Then we can use the laws of classical statistical mechanics.
At this stage four remarks are in order:
i.- Using our language the generalized idea of decoherence can be introduced in the following way: At the quantum level the average (1) reads:
$$O_\rho ^{(q)}=\underset{\omega ,\omega ^{}}{}\rho _{\omega \omega ^{}}O_{\omega ^{}\omega }$$
(2)
where $`\rho _{\omega \omega ^{}}`$ and $`O_{\omega \omega ^{}}`$ are the components in some basis of the operators $`\rho `$ and $`O`$ respectively. Eq. (2) can be considered as the average of some quantities $`O_{\omega \omega ^{}}`$ weighted by some generalized correlations $`\rho _{\omega \omega ^{}}`$ (since the $`\rho _{\omega \omega }`$ are probabilities but the $`\rho _{\omega \omega ^{}}`$, with $`\omega \omega ^{},`$ are quantum correlations). On the other hand, at the classical level we also have some quantities $`O_\omega `$ that correspond to a set $`\{\omega \}`$ of the exhaustive and exclusive alternatives, each one with a (boolean) probability $`p_\omega `$ of measure $`\omega `$ for the observable $`O.`$ The corresponding classical weighted average is:
$$O_\rho ^{(cl)}=\underset{\omega }{}p_\omega O_\omega $$
(3)
where $`_\omega p_\omega =1.`$ The transition from the quantum phase to the classical one is therefore:
$$\underset{\omega ,\omega ^{}}{}\rho _{\omega \omega ^{}}O_{\omega ^{}\omega }\underset{\omega }{}p_\omega O_\omega $$
(4)
at least for some $`O`$ which belong to a preferred sub space of $`𝒪`$ (i.e. to a subspace expanded by a complete set of commuting observables, a CSCO, that we will define below; the eigenbasis of this set will be the so called final pointer basis). If in (4) we take $`\rho _{\omega \omega }=p_\omega `$ and $`O_{\omega \omega }=O_\omega `$, the matrix $`\rho _{\omega \omega ^{}}`$ must become diagonal in the final pointer basis. This is the essence of the transition (a), since the above relation will be valid for all observables of the CSCO and we will have:
$$O_\rho ^{(q)}O_\rho ^{(cl)}$$
(5)
If this transition takes place, boolean logic is established in the statistical classical system, if we perform the measurement with the observables of the preferred CSCO. In the usual parlance we will then say that the density matrices that contain quantum interference terms become diagonal, in such away that these interferences are suppressed. Then the quantum way to find probabilities of exclusive and exhaustive alternatives, i. e.: adding the corresponding amplitudes and computing the norm, becomes the classical boolean way: just adding the probabilities.
ii.- In this paper decoherence is essentially studied in systems with continuous spectrum. The case of the discrete spectrum, and the causes of decoherence in this case, are discussed in section II C.
iii.- In the case of the continuous spectrum the essence of the method is the following: If $`\omega ^+`$ are the eigenvalues of $`H`$ and we call $`\nu =\omega \omega ^{},`$ the $`\rho _{\omega \omega ^{}}`$ of eq. (2) is a function $`\rho (\nu ,\mathrm{}).`$ Then the time limit of its evolution is given by the Riemann-Lebesgue theorem, that prescribes that:
$$\underset{t\mathrm{}}{lim}_a^ae^{i\nu t}\rho (\nu ,\mathrm{})𝑑\nu =0$$
(6)
if $`\rho (\nu ,\mathrm{})`$ is integrable. So all the diagonal terms ($`\nu =0)`$ and all the off-diagonal terms ($`\nu 0)`$ vanish. Therefore this theorem cannot be used as a computation method in the case of continuous spectrum. Nevertheless when we consider the problem within a cube of size $`L`$, we define $`\rho _{\omega \omega ^{}}`$ there, and when we make $`L\mathrm{}`$, it can be shown that a singular structure appears for $`\rho (\nu ,\mathrm{})`$ and the corresponding singular diagonal term remains as it should. The method introduced in paper is precisely designed to rigorously deal with these singular structures. It has yielded good results in papers , , .
iv.- Before the classical stationary state limit is reached usually the system goes through a ”classical phase” where the state can be considered as classical but it is not yet in its final classical stationary state. But our method can only be used when $`t\mathrm{}`$. It only allows to find the ”statistical classical final limit”. So, we essentially study this final stationary state but we believe that our method can be generalized to cover the classical phase before the final stationary state, so we will discuss these matters in section V. Moreover, we believe that the understanding of the final limit will enhance the chances to understand the much more difficult problem of the classical phase, in the clearest and concise way.
In section III we reach to the principal aim of the formalism of transition ”a” which is to create a bridge between quantum and classical mechanics, precisely between quantum mechanics and classical statistical mechanics at equilibrium. We know that the uncertainty relations disappear, when $`\mathrm{}0`$, (more precisely when the characteristic dimension of the system makes $`\mathrm{}`$ a negligible quantity). Then, let us consider a system where the quantum state is defined by a density matrix $`\rho ,`$ and a set of classical trajectories in phase space labelled by some constants $`x,l_1,\mathrm{},l_N`$, $`a_1,\mathrm{},a_N,`$ where $`x`$ corresponds to the energy, $`l_1,\mathrm{},l_N`$ to other dynamical momentum variables, and, $`a_1,\mathrm{},a_N`$ to configuration variables. The aim of the theory is:
1.- To transform the matrix $`\rho `$ into a classical density function in phase space $`\rho (q,p)`$ when $`\mathrm{}0.`$
2.- To decompose $`\rho (q,p)`$ as:
$$\rho (q,p)=\underset{x,l_1,\mathrm{},l_N,a_1,\mathrm{},a_N}{}p_{x,l_1,\mathrm{},l_N,a_1,\mathrm{},a_N}\rho _{x,l_1,\mathrm{},l_N,a_1,\mathrm{},a_N}(q,p)$$
(7)
where $`q`$ and $`p`$ are the position and momentum coordinates and the classical densities $`\rho _{x,l_1,\mathrm{},l_N,a_1,\mathrm{},a_N}(q,p)`$ would correspond to each classical trajectory <sup>§</sup><sup>§</sup>§The dimension of the phase space considered is $`2(N+1)`$. Then there are $`(N+1)`$ momenta and ($`N+1)`$ coordinates. So $`N`$ $`+1`$ is the number of parameters necessary to label the momenta of the classical space-time trajectories, and $`N`$ the number necessary to label the origins of the trajectories. ( in the classical sense that it is peaked in the trajectory and thus it rapidly vanishes when going from the near vicinity of the trajectory to the far zones of the phase space) and $`p_{x,l_1,\mathrm{},l_N,a_1,\mathrm{},a_N}`$ is the probability of each trajectory.
We will obtain (when $`\mathrm{}0`$) these results as follows:
1.-$`\rho (q,p)`$ will be the Wigner function corresponding to the matrix $`\rho .`$
2.-$`\rho _{x,l_1,\mathrm{},l_N,a_1,\mathrm{},a_N}(q,p)`$ will be the Wigner functions of the wave packets going along the classical trajectories labelled by the constant of the motion $`x,l_1,\mathrm{},l_N,`$ and passing by the initial point of coordinates $`a_1,\mathrm{},a_N.`$
We will see that all this happens after a convenient decoherence time and we will obtain the last expansion (cf. eq. (50)) and therefore what we consider the best bridge between classical and quantum statistical concepts (see paper for a very similar conclusion).
We will devote section IV to discuss transition (b), namely the localization process. Eventually in some cases this process takes place and correlations appear and we reach to a single classical state if the localization process is efficient enough. Then we can use the laws of classical mechanics. This phenomenon happens if the dynamic of the system and the initial conditions are such that some canonically conjugated variables correlate (see appendix B). We will see how this fact can be incorporated in our formalism.
We will draw our main conclusions and comments in section V.
Appendix A is devoted to compare our results with those in the literature. In appendix B we deal with correlations and localization. Finally, in appendix C we translate the results into the language of usual decoherence of histories.
## II Decoherence.
### A Decoherence in the energy.
Let us consider an isolated quantum system with $`N+1`$ dynamical variables and a Hamiltonian endowed with a continuous spectrum and just one bounded state. So the discrete part of the spectrum of $`H`$ has only one value $`\omega _0`$ and the continuous spectrum is let say $`0\omega <\mathrm{}`$ (how the discrete spectrum behaves in the continuous limit can be seen in papers , ). Eventually we will give the collective name $`x`$ to both $`\omega _0`$ and $`\omega .`$ Let us assume that it is possible to diagonalize the Hamiltonian $`H`$, together with $`N`$ observables $`O_i`$ ($`i=1,\mathrm{},N)`$. The operators ($`H`$, $`O_1`$,…,$`O_N`$) form a complete set of commuting observables (CSCO). For simplicity we also assume a discrete spectrum for the $`N`$ observables $`O_i`$. Therefore we write
$$H=\omega _0\underset{m}{}|\omega _0,m\omega _0,m|+_0^{\mathrm{}}\omega \underset{m}{}|\omega ,m\omega ,m|d\omega $$
(8)
where $`\omega _0<0`$ is the energy of the ground state, and $`m\{m_1,\mathrm{},m_N\}`$ labels a set of discrete indexes which are the eigenvalues of the observables $`O_1`$,…,$`O_N`$. $`\{|\omega _0,m,|\omega ,m\}`$ is a basis of simultaneous generalized eigenvectors of the CSCO:
$`H|\omega _0,m=\omega _0|\omega _0,m,H|\omega ,m=\omega |\omega ,m,`$
$`O_i|\omega _0,m=m_i|\omega _0,m,O_i|\omega ,m=m_i|\omega ,m.`$
The most general observable that we are going to consider in our model reads:
$`O`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}O(\omega _0)_{mm^{}}|\omega _0,m\omega _0,m^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega O(\omega )_{mm^{}}|\omega ,m\omega ,m^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega O(\omega ,\omega _0)_{mm^{}}|\omega ,m\omega _0,m^{}|+`$ (10)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}O(\omega _0,\omega ^{})_{mm^{}}|\omega _0,m\omega ^{},m^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑\omega 𝑑\omega ^{}O(\omega ,\omega ^{})_{mm^{}}|\omega ,m\omega ^{},m^{}|,`$
where $`O(\omega )_{mm^{}}`$, $`O(\omega ,\omega _0)_{mm^{}}`$, $`O(\omega _0,\omega )_{mm^{}}`$ and $`O(\omega ,\omega ^{})_{mm^{}}`$ are ordinary functions of the real variables $`\omega `$ and $`\omega ^{}`$(these functions must have some mathematical properties in order to develop the theory; these properties are listed in paper ). Namely, the most general observables have a singular component (the second term of the r.h.s. of the last equation) and a regular part (all the other terms). If the singular term would be missing the Hamiltonian (8) would not belong to the space of the chosen observables . We will say that these observables belong to a space $`𝒪`$. This space has the basis $`\{|\omega _0,mm^{})`$, $`|\omega ,mm^{})`$, $`|\omega \omega _0,mm^{})`$, $`|\omega _0\omega ^{},mm^{})`$, $`|\omega \omega ^{},mm^{})\}`$:
$`|\omega _0,mm^{})|\omega _0,m\omega _0,m^{}|,|\omega ,mm^{})|\omega ,m\omega ,m^{}|,|\omega \omega _0,mm^{})|\omega ,m\omega _0,m^{}|,`$
$$|\omega _0\omega ^{},mm^{})|\omega _0,m\omega ^{},m^{}|,|\omega \omega ^{},mm^{})|\omega ,m\omega ^{},m^{}|.$$
(11)
The quantum states $`\rho `$ are measured by the observables just defined, computing the mean values of these observables in the quantum states, i. e. in the usual notation: $`O_\rho =Tr(\rho ^{}O)`$ . These mean values, generalized as in paper , can be considered as linear functionals $`\rho ,`$ mapping the vectors $`O`$ on the real numbers, that we can call $`(\rho |O)`$ . In fact, this is a generalization of the usual mean value definition. Then $`\rho 𝒮𝒪^{^{}},`$ where $`𝒮`$ is a convenient convex set contained in $`𝒪^{^{}}`$, the space of linear functionals over $`𝒪`$ , . The basis of $`𝒪^{}`$ (that can also be considered as the co-basis of $`𝒪)`$ is $`\{(\omega _0,mm^{}|`$, $`(\omega ,mm^{}|`$, $`(\omega \omega _0,mm^{}|`$, $`(\omega _0\omega ^{},mm^{}|`$, $`(\omega \omega ^{},mm^{}|\}`$ defined as functionals by the equations:
$`(\omega _0,mm^{}|\omega _0,nn^{})=\delta _{mn}\delta _{m^{}n^{}},(\omega ,mm^{}|\eta ,nn^{})=\delta (\omega \eta )\delta _{mn}\delta _{m^{}n^{}},(\omega \omega _0,mm^{}|\eta \omega _0,nn^{})=\delta (\omega \eta )\delta _{mn}\delta _{m^{}n^{}},`$
$$(\omega _0\omega ^{},mm^{}|\omega _0\eta ^{},nn^{})=\delta (\omega ^{}\eta ^{})\delta _{mn}\delta _{m^{}n^{}},(\omega \omega ^{},mm^{}|\eta \eta ^{},nn^{})=\delta (\omega \eta )\delta (\omega ^{}\eta ^{})\delta _{mn}\delta _{m^{}n^{}}.$$
(12)
and all other $`(.|.)`$ are zero. In particular we have
$$(\omega _0,mm^{}|O)=O(\omega _0)_{mm^{}}=\omega _0,m|O|\omega _0,m^{}$$
(13)
for any $`O𝒪`$. But $`(\omega ,mm^{}|O)=O(\omega )_{mm^{}}`$ is not equal to $`\omega ,m|O|\omega ,m^{}`$, which is not even defined if $`O`$ is given by eq.(10). Therefore $`(\omega ,mm^{}|`$ can only be considered as a functional, being a typical generalized state. Then, a generic quantum state reads:
$`\rho `$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\overline{\rho (\omega _0)}_{mm^{}}(\omega _0,mm^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega )}_{mm^{}}(\omega ,mm^{}|++{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega ,\omega _0)}_{mm^{}}(\omega \omega _0,mm^{}|+`$ (15)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega _0,\omega ^{})}_{mm^{}}(\omega _0\omega ^{},mm^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega {\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega ,\omega ^{})}_{mm^{}}(\omega \omega ^{},mm^{}|,`$
where $`\overline{\rho (\omega _0)}_{mm}`$ and $`\overline{\rho (\omega )}_{mm}`$ are real and non negative, $`\overline{\rho (\omega ,\omega _0)}_{mm^{}}=\rho (\omega _0,\omega )_{m^{}m}`$ and $`\overline{\rho (\omega ,\omega ^{})}_{mm^{}}=\rho (\omega ^{},\omega )_{m^{}m}`$. Moreover, $`\rho (\omega _0)_{mm^{}}`$ and $`\rho (\omega )_{mm^{}}`$ satisfy the total probability condition
$$(\rho |I)=\underset{m}{}\rho (\omega _0)_{mm}+\underset{m}{}_0^{\mathrm{}}𝑑\omega \rho (\omega )_{mm}=1,$$
(16)
where $`I=_m|\omega _0,m\omega _0,m|+_0^{\mathrm{}}𝑑\omega _m|\omega ,m\omega ,m|`$ is the identity operator in $`𝒪`$. Eq. (16) is the extension to state functionals of the usual condition $`Tr\rho ^{}=1`$, used when $`\rho `$ is a density operator.
The time evolution of the quantum state $`\rho `$ reads:
$`\rho (t)`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\overline{\rho (\omega _0)}_{mm^{}}(\omega _0,mm^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega )}_{mm^{}}(\omega ,mm^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega ,\omega _0)}_{mm^{}}e^{i(\omega \omega _0)t}(\omega \omega _0,mm^{}|+`$ (18)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega _0,\omega ^{})}_{mm^{}}e^{i(\omega _0\omega ^{})t}(\omega _0\omega ^{},mm^{}|+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega {\displaystyle _0^{\mathrm{}}}d\omega ^{}\overline{\rho (\omega ,\omega ^{})}_{mm^{}}e^{i(\omega \omega ^{})t}(\omega \omega ^{},mm^{}|`$
The mean value of an observable $`O`$ in a quantum state $`\rho `$ reads:
$`O_{\rho (t)}`$ $`=`$ $`(\rho (t)|O)=`$ (19)
$`=`$ $`{\displaystyle \underset{mm^{}}{}}\overline{\rho (\omega _0)_{mm^{}}}O(\omega _0)_{mm^{}}+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}d\omega \overline{\rho (\omega ,\omega ^{}})_{mm^{}}O(\omega )_{mm^{}}+`$ (22)
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega \overline{\rho (\omega ,\omega _0)}_{mm^{}}e^{i(\omega \omega _0)t}O(\omega ,\omega _0)_{mm^{}}+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}\overline{\rho (\omega _0,\omega ^{})}_{mm^{}}e^{i(\omega _0\omega ^{})t}O(\omega _0,\omega ^{})_{mm^{}}`$
$`+{\displaystyle \underset{mm^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}\overline{\rho (\omega ,\omega ^{})}_{mm^{}}e^{i(\omega \omega ^{})t}O(\omega ,\omega ^{})_{mm^{}}.`$
Using the Riemann-Lebesgue theorem we obtain the weak limit, for all $`O𝒪`$
$$\underset{t\mathrm{}}{lim}O_{\rho (t)}=O_\rho _{}$$
(23)
where we have introduced the diagonal asymptotic or final stationary state functional
$$\rho _{}=\underset{mm^{}}{}\overline{\rho (\omega _0)}_{mm^{}}(\omega _0,mm^{}|+\underset{mm^{}}{}_0^{\mathrm{}}d\omega \overline{\rho (\omega )}_{mm^{}}(\omega ,mm^{}|$$
(24)
Therefore, in a weak sense we have:
$$W\underset{t\mathrm{}}{lim}\rho (t)=\rho _{}$$
(25)
Thus, any quantum state weakly goes to a linear combination of the energy diagonal states $`(\omega _0,mm^{}|`$ and $`(\omega ,mm^{}|`$ (the energy off-diagonal states $`(\omega \omega _0,mm^{}|`$, $`(\omega _0\omega ^{},mm^{}|`$ and $`(\omega \omega ^{},mm^{}|`$ are not present in $`\rho _{}`$). This is the case if we observe and measure the system evolution with any possible observable of space $`𝒪`$. Then, from the observational (or generalized coarse-graining) point of view, we have decoherence of the energy levels when $`t\mathrm{}`$, even that, from the strong limit (fine-graining) point of view the off-diagonal terms never vanish, they just oscillate, since we cannot directly use the Riemann-Lebesgue theorem in the operator equation (18).
Some observations are in order:
i.- The real existence of the two singular parts of $`O`$ and $`\rho `$ is assured by the physic of the problem. The singular part of the observables is just a necessary generalization of the singular part of the Hamiltonian, which has a singular part $`|\omega )`$ (eq. (8)). The states must also be singular objects since, intuitively, we realize that a continuous by continuous matrix will decohere in a matrix with some kind of singularity in the diagonal. The method is precisely designed to deal with this object.
ii.- From eq. (23) we can again see that what we are doing is just a generalized version of coarse graining, where a projector on the ”relevant” part of the system is defined. The ”relevant” part of the states $`(\rho |`$ is in our case $`(\rho |O)`$ for all $`O𝒪`$, i.e. the ”projection” of $`\rho `$ on the class of observables of the form given in eq. (10). An ”irrelevant” projection would be a $`(\rho |O^{})`$ where $`O^{}𝒪`$.
### B Decoherence in the other ”momentum” dynamical variables.
Having established the decoherence in the energy levels we must consider the decoherence in the other dynamical variables $`O_i`$, of the CSCO where we are working. We will call these variables ”momentum variables”. For the sake of simplicity we will consider, as in the previous section, that the spectra of these dynamical variables are discrete. As the expression of $`\rho _{}`$ given in eq. (24) involves only the time independent components of $`\rho (t)`$, it is impossible that a different decoherence process would take place to eliminate the off-diagonal terms in the remaining $`N`$ dynamical variables. Therefore, the only thing to do is to find if there is a basis where the off-diagonal components of $`\rho (\omega _0)_{mm^{}}`$ and $`\rho (\omega )_{mm^{}}`$ vanish at any time before the final state is reached. This basis in fact exists, it is constant in time, and it will be called the final pointer basis.
Let us consider the following change of basis
$$|\omega _0,r=\underset{m}{}U(\omega _0)_{mr}|\omega _0,m,|\omega ,r=\underset{m}{}U(\omega )_{mr}|\omega ,m,$$
(26)
where $`r`$ and $`m`$ are short notations for $`r\{r_1,\mathrm{},r_N\}`$ and $`m\{m_1,\mathrm{},m_N\}`$, and $`\left[U(x)^1\right]_{mr}=\overline{U(x)}_{rm}`$ ($`x`$ denotes either $`\omega _0<0`$ or $`\omega ^+`$).
The new basis $`\{|\omega _0,r,|\omega ,r\}`$ verifies the generalized orthogonality conditions
$`\omega _0,r|\omega _0,r^{}=\delta _{rr^{}},\omega ,r|\omega ^{},r^{}=\delta (\omega \omega ^{})\delta _{rr^{}},\omega _0,r|\omega ,r^{}=\omega ,r|\omega _0,r^{}=0.`$
It is easy to obtain the components of the states $`\rho 𝒮`$ in the new basis
$`\rho (\omega _0)_{rr^{}}`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\left[U(\omega _0)^1\right]_{rm}\rho (\omega _0)_{mm^{}}\left[U(\omega _0)\right]_{m^{}r^{}},`$
$`\rho (\omega )_{rr^{}}`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\left[U(\omega )^1\right]_{rm}\rho (\omega )_{mm^{}}\left[U(\omega )\right]_{m^{}r^{}},`$
$`\rho (\omega ,\omega ^{})_{rr^{}}`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\left[U(\omega )^1\right]_{rm}\rho (\omega ,\omega ^{})_{mm^{}}\left[U(\omega ^{})\right]_{m^{}r^{}},`$
$`\rho (\omega _0,\omega ^{})_{rr^{}}`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\left[U(\omega _0)^1\right]_{rm}\rho (\omega _0,\omega ^{})_{mm^{}}\left[U(\omega ^{})\right]_{m^{}r^{}},`$
$`\rho (\omega ,\omega _0)_{rr^{}}`$ $`=`$ $`{\displaystyle \underset{mm^{}}{}}\left[U(\omega )^1\right]_{rm}\rho (\omega ,\omega _0)_{mm^{}}\left[U(\omega _0)\right]_{m^{}r^{}},`$
As $`\overline{\rho (\omega _0)}_{mm^{}}=\rho (\omega _0)_{m^{}m}`$ and $`\overline{\rho (\omega )}_{mm^{}}=\rho (\omega )_{m^{}m}`$, it is possible to choose $`U(\omega _0)`$ and $`U(\omega )`$ in such a way that the off-diagonal parts of $`\rho (\omega _0)_{rr^{}}`$ and $`\rho (\omega )_{rr^{}}`$ would vanish, i.e.
$`\rho (\omega _0)_{rr^{}}=\rho _r(\omega _0)\delta _{rr^{}},\rho (\omega )_{rr^{}}=\rho _r(\omega )\delta _{rr^{}}.`$
Therefore, there is a final pointer basis for the observables given by $`\{|\omega _0,rr^{})`$, $`|\omega ,rr^{})`$, $`|\omega \omega _0,rr^{})`$, $`|\omega _0\omega ^{},rr^{})`$, $`|\omega \omega ^{},rr^{})\}`$ and defined as in eq. (11). The corresponding final pointer basis for the states $`\{(\omega _0,rr^{}|`$, $`(\omega ,rr^{}|`$, $`(\omega \omega _0,rr^{}|`$, $`(\omega _0\omega ^{},rr^{}|`$, $`(\omega \omega ^{},rr^{}|\}`$ diagonalizes the time independent part of $`\rho (t)`$ and therefore it diagonalizes the final state $`\rho _{}`$
$$\rho _{}=W\underset{t\mathrm{}}{lim}\rho (t)=\underset{r}{}\rho _r(\omega _0)(\omega _0,rr|+\underset{r}{}_0^{\mathrm{}}d\omega \rho _r(\omega )(\omega ,rr|.$$
(27)
Now we can define the final exact pointer observables
$$P_i=\underset{r}{}P_r^i(\omega _0)|\omega _0,r\omega _0,r|+_0^{\mathrm{}}𝑑\omega \underset{r}{}P_r^i(\omega )|\omega ,r\omega ,r|.$$
(28)
As $`H`$ and $`P_i`$ are diagonal in the basis $`\{|\omega _0,r`$, $`|\omega ,r\}`$, the set $`\{H,P_i,\mathrm{}P_N\}`$ is precisely the complete set of commuting observables (CSCO) related to this basis, where $`\rho _{}`$ is diagonal in the corresponding co-basis for states. For simplicity we define the operators $`P_i`$ such that $`P_r^i(\omega _0)=P_r^i(\omega )=r_i`$, thus
$$P_i|\omega _0,r=r_i|\omega _0,r,P_i|\omega ,r=r_i|\omega ,r.$$
(29)
Therefore $`\{|\omega _0,r`$, $`|\omega ,r\}`$ is the final observers’ pointer basis where there is a perfect decoherence in the corresponding state co-basis. Moreover the generalized states $`(\omega _0,rr|`$ and $`(\omega ,rr|`$ are constants of the motion, and therefore these exact pointer observables have a constant statistical entropy and will be ”at the top of the list” of Zurek’s ”predictability sieve” . The final pointer basis is therefore defined by the dynamics of the model and by the quantum state considered.
Therefore:
i.- Decoherence in the energy is produced by the time evolution when $`t\mathrm{}`$.
ii.- Decoherence in the other dynamical variables can be seen if we choose an adequate basis, namely the final pointer basis.
Essentially we have given a partial answer, for this kind of models, to the fundamental question of Gell-Mann and Hartle (precisely only an answer in the case when $`t\mathrm{})`$: For each $`H`$ and each initial state $`\rho `$ there is only one final pointer basis and therefore only one ”quasi-classical domain or realm”. But of course this unique consistent set depends of the chosen space of observable $`𝒪`$ (see more in appendix III)..
Our main result is eq. (27): When $`t\mathrm{}`$ then $`\rho (t)\rho _{}`$ and in this state the dynamical variables $`H,P_1,\mathrm{},P_N`$ are well defined. Therefore the eventual conjugated variables to these momentum variables (namely: configuration variables, if they exist) are completely undefined.
In fact, calling $`𝕃_i`$ the generator of the displacements along the eventual configuration variable conjugated to $`P_i`$, we have $`(𝕃_i\rho _{}|O)=(\rho _{}|𝕃_i^{}O)=(\rho _{}|[P_i,O])=0`$ for all $`O𝒪`$ as it can be proved by direct computation using eqs. (10),(12), (24), and (28). Then $`𝕃_i\rho _{}=0,`$ and $`\rho _{}`$ is homogeneous in these configuration variables.
### C Decoherence characteristic decaying time, the permanent quantum states case, and the role of the environment.
From the preceding section we may have the feeling that the process of decoherence must be found in all the physical systems. It is not so and there are two reasons:
i.- Characteristic decaying times can be computed using analytic continuation technics, as in paper . E. g. in particular models we can find the characteristic times for the system (e. g. an oscillator) and the field (e. g. the environments or bath) as below eq. (56) of the last quoted paper. If the maximal characteristic time $`\gamma ^1`$ is very large, even if theoretically the decoherence process will always take place, it will be so slow that the system will behave as a quantum one for a very long time. Then there will be no measurable decoherence.
ii.- It may also happen that more than one of the $`\gamma `$ would be zero. Then, Hamiltonian $`H,`$ has more than one bound state, let us say $`n`$ (or even part of its spectrum is discrete). Then the first term of the r. h. s. of eq. (22) must be changed to
$$\underset{ij}{}\rho _{ji}O_{ij}e^{i(\omega _i\omega _j)t}=\underset{i}{}\rho _{ii}O_{ii}+\underset{ij}{}\rho _{ji}O_{ij}e^{i(\omega _i\omega _j)t}$$
(30)
where $`i,j=1,\mathrm{}n,`$ and as the second term of the r. h. s. does not vanish, decoherence does not take place. This is the case of a theoretical atom, not coupled to the electromagnetic field, where the electrons will remain for ever in their exited states, and they will never decay. Then the atom never goes to a decohered state. But if the atom is coupled to an electromagnetic field (that usually it is called the ”environment”, as in appendix II) there will be only one bound state, the second term of the r. h. s. of eq. (30) will be absent, and decoherence will occur. In fact, in many examples the role of the ”environment” is just to introduce a continuous spectrum to be coupled in such a way that only one bound state remains and the decoherence is complete. In other cases fluctuations (or imperfections) of continuous nature take the role of the continuous spectrum and produce the average and make the off diagonal term disappear. This is the case of the spin recombination experiment ( page. 180) that takes place in a single crystal interferometer.
iii.- More generally, using only observables from a subset $`\mathrm{\Omega }𝒪`$ we may only involve some components of the state functional, e. g. those constructed with the eigenvectors of $`H`$ that eventually expand the space $`\mathrm{\Omega }.`$ Then if we only consider the observables of $`\mathrm{\Omega }`$ it may be that the components of the state related with these observables become decohered, because their decoherence times are small, while the other components remain undecohered, because they have a larger decoherence time. Then we will have a system which is partially decohered and partially not decohered, (which in fact is the case of the universe where there are both classical and quantum phenomena).
## III The classical statistical limit.
### A Expansion in sets of classical motions.
In this section we will use the Wigner integrals that introduce an isomorphism between quantum observables $`O`$ and states $`\rho `$ and their classical analogues $`O^W(q,p)`$ and $`\rho ^W(q,p)`$ :
$`O^W(q,p)`$ $`=`$ $`{\displaystyle 𝑑\lambda q\frac{\lambda }{2}|O|q+\frac{\lambda }{2}\mathrm{exp}(\frac{i\lambda p}{\mathrm{}})}`$ (31)
$`\rho ^W(q,p)`$ $`=`$ $`{\displaystyle \frac{1}{\pi \mathrm{}}}{\displaystyle }d\lambda (\rho ||q+\lambda q\lambda |)\mathrm{exp}({\displaystyle \frac{2i\lambda p}{\mathrm{}}}).`$ (32)
It is possible to prove that $`𝑑q𝑑p\rho ^W(q,p)=(\rho |I)=1`$, but $`\rho ^W`$ is not in general non negative. It is also possible to deduce that
$$(\rho ^W|O^W)=𝑑q𝑑p\rho ^W(q,p)O^W(q,p)=(\rho |O),$$
(33)
and therefore to the mean value in the classical Liouville space it corresponds the mean value in the quantum Liouville space. Moreover, calling $`L`$ the classical Liouville operator, and $`𝕃`$ the quantum Liouville-Von Neumann operator, we have
$$L\left[\rho ^W(q,p)\right]=\left[𝕃\rho \right]^W(q,p)+O(\mathrm{}),$$
(34)
where $`L\rho ^W(q,p)=i\{H^W(q,p),\rho ^W(q,p)\}_{PB}`$ and
$$(𝕃\rho |O)=(\rho |[H,O]).$$
(35)
Finally, if $`O=O_1O_2`$, where $`O_1`$ and $`O_2`$ are two quantum observables, we have
$$O^W(q,p)=O_1^W(q,p)O_2^W(q,p)+O(\mathrm{}).$$
(36)
We will prove that the distribution function $`\rho _{}^W(q,p)`$, that corresponds to the state functional $`\rho _{}`$ via the Wigner integral is a non negative function of the classical constants of the motion, in our case $`H^W(q,p)`$, $`P_1^W(q,p)`$,…, $`P_N^W(q,p),`$ obtained from the corresponding quantum operators $`H`$, $`P_1`$,…, $`P_N`$.
From eq. (27) we have:
$$\rho _{}=W\underset{t\mathrm{}}{lim}\rho (t)=\underset{r}{}\rho _r(\omega _0)(\omega _0,rr|+\underset{r}{}_0^{\mathrm{}}d\omega \rho _r(\omega )(\omega ,rr|,$$
(37)
so we must compute:
$$\rho _{\omega r}^W(q,p)\left(\frac{1}{\pi \mathrm{}}\right)^{N+1}(\omega ,rr||q+\lambda q\lambda |)e^{2ip\lambda }d\lambda $$
(38)
We know from section II C (or we can directly prove from eqs. (27-29)) that
$$(\omega _0,rr|H^n)=\omega _0^n,(\omega ,rr|H^n)=\omega ^n,(\omega _0,rr|P_i^n)=r_i^n,(\omega ,rr|P_i^n)=r_i^n,$$
(39)
for $`i=1,\mathrm{},N`$ and $`n=0,1,2,\mathrm{}`$ Using the relation (36) between quantum and classical products of observables and relation (33) between quantum and classical mean values, in the limit $`\mathrm{}0`$ (we will consider that we always take this limit when we refer to classical equations below) we deduce that the characteristic property of the distribution $`\rho _{\omega r}^W(q,p)`$, that corresponds to the state functional $`(\omega ,rr|`$, is:
$$\rho _{\omega r}^W(q,p)[H^W(q,p)]^n𝑑q𝑑p=\omega ^n,\rho _{\omega r}^W(q,p)[P_i^W(q,p)]^n𝑑q𝑑p=r_i^n,$$
(40)
for any natural number $`n.`$ Thus $`\rho _{\omega r}^W(q,p)`$ must be the functional We have omitted the $`O(\mathrm{})`$ of eqs. (34) and (36). If we reintroduce these $`O(\mathrm{})`$ we will see that eqs. (41) and (42) are only valid in the limit $`\mathrm{}0.`$ If $`\mathrm{}`$ is only very small the $`\delta `$ are just functions strongly peaked at the zero value of their variables.:
$$\rho _{\omega r}^W(q,p)=\delta (H^W(q,p)\omega )\delta (P_1^W(q,p)r_1)\mathrm{}\delta (P_N^W(q,p)r_N).$$
(41)
For the distribution $`\rho _{\omega _0r}^W(q,p)`$, that corresponds to the state functional $`(\omega _0,rr|`$, we obtain
$$\rho _{\omega _0r}^W(q,p)=\delta (H^W(q,p)\omega _0)\delta (P_1^W(q,p)r_1)\mathrm{}\delta (P_N^W(q,p)r_N).$$
(42)
Therefore, going back to eq. (37) and since the Wigner relation is linear, we have:
$$\rho _{}^W(q,p)=\underset{r}{}\rho _r(\omega _0)\rho _{\omega _0r}^W(q,p)+\underset{r}{}_0^{\mathrm{}}𝑑\omega \rho _r(\omega )\rho _{\omega r}^W(q,p).$$
(43)
Also we obtain $`\rho _{}^W(q,p)0`$, because $`\rho _r(\omega _0)`$ and $`\rho _r(\omega )`$ are non negative.
Therefore, the classical state $`\rho _{}^W(q,p)`$ is a linear combination of the generalized classical states $`\rho _{xr}^W(q,p)`$ (where $`x`$ is either $`\omega _0`$ or $`\omega `$), having well defined values $`x`$, $`r_1`$,…, $`r_N`$ of the classical observables $`H^W(q,p)`$, $`P_1^W(q,p)`$,…, $`P_N^W(q,p)`$ and the corresponding classical canonically conjugated variables completely undefined since the $`\rho _{xr}^W(q,p)`$ are not functions of these variables. So we reach, in the classical case, to the same conclusion than in the quantum case (see end of subsection II B). But now all the classical canonically conjugated variables $`a_0,a_1,\mathrm{},a_N`$ do exist since they can be found solving the corresponding Poisson brackets differential equations.
As the momenta $`H^W,P_1^W,\mathrm{},P_N^W`$, or any function of these momenta, that we will call generically $`\mathrm{\Pi },`$ are also constant of the motion, then we have $`\frac{d}{dt}\mathrm{\Pi }=H/\alpha =0`$, where $`\alpha `$ is the classically conjugated variable to $`\mathrm{\Pi }.`$ So $`H`$ is just a function of the $`\mathrm{\Pi }`$ and:
$$\frac{d}{dt}\alpha =\frac{H(\mathrm{\Pi })}{\mathrm{\Pi }}=\varpi (\mathrm{\Pi })=const.$$
(44)
so:
$$\alpha _j(t)=\varpi _j(\mathrm{\Pi })t+\alpha _j(0),j=0,1,\mathrm{},N.$$
(45)
Thus (going back to the old coordinates) in the set of classical motions contained in the densities (41) and (42) the momenta $`H,P_1,\mathrm{},P_N`$, are completely defined and the origin of the corresponding motions, that we will respectively call $`a_0(0)`$, $`a_1(0)`$,…and $`a_N(0)`$, are completely undefined, in such a way that the motions represented in the last equation homogeneously fill the surface where $`H^W`$, $`P_1^W`$,…, and $`P_N^W`$, have constant values, which now turns out to be a usual torus of phase space <sup>\**</sup><sup>\**</sup>\**If $`H^W`$, $`P_1^W`$ ,…,$`P_N^W`$ are isolating constants of the motion and the tori are not broken .. This is the interpretation that we give to the densities (41) and (42) which are just functions of the variables $`H^W`$, $`P_1^W`$,…, $`P_N^W,`$ but they are not of the classical conjugated variables $`a_0`$, $`a_1`$,…, $`a_N`$.
Then, eq. (43) can be considered as the expansion of $`\rho _{}^W(q,p)`$ in the sets of classical motions contained in $`\rho _{xr}^W(q,p),`$ each one with a probability $`\rho _r(x)`$ ($`x=\omega _0,\omega `$).
Summing up:
i.- We have shown that the quantum state functional $`\rho (t)`$ evolves to a final diagonal state $`\rho _{}`$.
ii.- This quantum state $`\rho _{}`$ has $`\rho _{}^W(q,p)`$ as its corresponding classical density.
iii.- This classical density can be decomposed in sets of classical motions where $`H^W`$, $`P_1^W`$,…, $`P_N^W`$ remain constant. The origin of these motions: $`a_0(0),a_1(0),\mathrm{},a_N(0)`$ are homogeneously distributed.
iv.- From eqs. (41-43) we obtained that $`\rho _{}^W(q,p)=`$
$`f(H^W(q,p),P_1^W(q,p),\mathrm{},P_N^W(q,p))0`$.
### B Expansion in terms of classical motions.
We can now expand the densities given in eqs. (41-43) in terms of classical motions. In fact, since
$$\underset{i=0}{\overset{N}{}}\delta (a_i(q,p)a_i(t))\underset{i=0}{\overset{N}{}}da_i(0)=1$$
(46)
where $`a_j(t)=\varpi _j(P^W)t+a_j(0)`$, we can write eq. (43) as:
$`\rho _{}^W(q,p)`$ $`=`$ $`{\displaystyle \underset{r}{}\rho _r(\omega _0)\rho _{\omega _0r}^W(q,p)\underset{i=0}{\overset{N}{}}\delta (a_i(q,p)a_i(t))\underset{i=0}{\overset{N}{}}da_i(0)}+`$ (48)
$`{\displaystyle \underset{r}{}_0^{\mathrm{}}𝑑\omega \rho _r(\omega )\rho _{\omega r}^W(q,p)\underset{i=0}{\overset{N}{}}\delta (a_i(q,p)a_i(t))\underset{i=0}{\overset{N}{}}da_i(0)}.`$
We define
$$\rho _{x,r,a(0)}^W(q,p,t)\delta (H^W(q,p)x)\delta (P_1^W(q,p)r_1)\mathrm{}\delta (P_N^W(q,p)r_N)\delta (a_0(q,p)a_0(t))\mathrm{}\delta (a_N(q,p)a_N(t)),$$
(49)
which corresponds to the classical distribution of a motion with momenta $`x`$, $`r_1`$,…, $`r_N`$ and initial conditions $`a_0(0)`$,…, $`a_N(0)`$, and therefore to a single classical motion. So we can write eq. (48) as
$$\rho _{}^W(q,p)=\underset{r}{}\rho _r(\omega _0)\rho _{\omega _0,r,a(0)}^W(q,p,t)\underset{i=0}{\overset{N}{}}da_i(0)+\underset{r}{}_0^{\mathrm{}}𝑑\omega \rho _r(\omega )\rho _{\omega ,r,a(0)}^W(q,p,t)\underset{i=0}{\overset{N}{}}da_i(0).$$
(50)
So we have proved eq. (7) as we promised in the introduction.
The densities $`\rho _{x,r,a(0)}^W(q,p,t)`$ represent a point in phase space with momenta $`H^W=x`$, $`P_1^W=r_1`$,…, $`P_N^W=r_N`$ and coordinates $`a_j(t)=\varpi _j(P^W)t+a_j(0)`$, i.e. they represent single classical trajectories.
Then we have obtained the final classical limit. When $`t\mathrm{}`$ the quantum state functional $`\rho `$ becomes a diagonal state $`\rho _{}.`$ The corresponding classical distribution $`\rho _{}^W(q,p)`$ can be expanded as a linear combination of density functions $`\rho _{\omega _0,r,a(0)}^W(q,p,t)`$ and $`\rho _{\omega ,r,a(0)}^W(q,p,t)`$, representing classical trajectories, each one weighted by their corresponding probabilities $`\rho _r(\omega _0)`$ and $`\rho _r(\omega )`$. As the limit when $`t\mathrm{}`$ of our quantum model we have obtained a statistical classical mechanical model, , and the classical statistical realm is obtained.
## IV Correlations and Localization.
From many examples (e. g. ) we know that eventually correlations and the localization appear when $`t\mathrm{},`$ at least in some variables and in some quantum systems. E. g., in appendix B we give an example obtained using our method, where we can see that correlations appear in variables $`Q`$ and $`P,`$ when $`t\mathrm{}`$ (see eq. (B23)). As this state with correlations is a final state let as call it $`\rho _{}`$ and let us see how it can be incorporated in our formalism.
As $`\rho _{}`$ is a final stationary state it can be decomposed as in eq. (27). From eq. (12) we have:
$$(\omega _0,rr|\omega _{0,}r^{}r^{})=\delta _{rr^{}},(\omega ,rr|\omega _,r^{}r^{})=\delta _{rr^{}},(\omega _0,rr|\omega _,r^{}r^{})=0$$
(51)
Thus from eq. (51) we have:
$$(\rho _{}|\omega _{0,}rr)=\rho _r(\omega _0),(\rho _{}|\omega _,rr)=\rho _r(\omega )$$
(52)
So, given $`\rho _{},`$ endowed with correlations and computed by any method (including ours, see appendix B) we can find the corresponding initial conditions $`\rho _r(\omega _0),\rho _r(\omega )`$ that yield, when $`t\mathrm{}`$, to this correlated state<sup>††</sup><sup>††</sup>††The remaining initial conditions: $`\rho _{rr^{}}(\omega _0),`$ $`\rho _{rr^{}}(\omega ),`$ $`\rho _{rr^{}}(\omega ,\omega ^{}),`$ $`\rho _{rr^{}}(\omega _0,\omega ^{}),`$ $`\rho _{rr^{}}(\omega ,\omega _0)`$ are irrelevant since the corresponding terms disappear when $`t\mathrm{}.`$. In general, all final decohered stationary states (but not any quantum state) can be decomposed in this way, in particular our correlated state.
We can repeat all these formulae in the classical perspective of section III using the relation (33) between quantum and classical symbols, computing the initial conditions $`\rho _r(\omega _0),\rho _r(\omega )`$ using classical formulas:
$$(\rho _{}^W(q,p)|\rho _{\omega _0r}^W(q,p))=\rho _r(\omega _0),(\rho _{}^W(q,p)|\rho _{\omega r}^W(q,p))=\rho _r(\omega )$$
(53)
In this way the correlation and localization phenomena can be incorporated in our formalism. But it is difficult to use coordinates $`x,r,a`$ to directly obtain the final state $`\rho _{}`$ since this state looks quite unfamiliar in these coordinates but it turns out to be the minimal uncertainty wave packet if we study the problem in the usual coordinates $`q,`$ $`p,`$ as we prove in the appendix B via an example<sup>‡‡</sup><sup>‡‡</sup>‡‡In fact, it is not possible to formulate a general theory in $`q,`$ $`p`$ coordinates because, in order to make the computations, we must know the relation of these coordinates with the energy and other momenta, and this is only defined in specific models.. Furthermore the correlation phenomenon only appears if the potential and the initial conditions are such that all the trajectories with non negligible probability are concentrated by the dynamics eventually yielding a ”maximally localized” or ”minimal uncertainty” wave packet (as in the example of appendix B). It is difficult to see this fact in the abstract unfamiliar frame of the coordinates $`x,r,a`$, because the potentials are hidden by the diagonalization even if the initial conditions are obviously present (i. e,: in the choice $`\rho _r(\omega _0),`$ $`\rho _r(\omega )).`$ Anyhow the phenomenon is there, since we obtain localization when $`t\mathrm{}.`$ In this way we can consider the localized wave packet like a single classical system and the limit classical statistical mechanics $``$ classical mechanics is obtained because the motion of the wave packet satisfies the classical equations (45) as in all the trajectories. Now the processes ”a” and ”b” are explained and the limit statistical quantum mechanics$``$classical mechanics is completed. The classical realm is appeared.
## V Comments and conclusions.
Some observation are in order:
### A Sketch of the classical limit.
Using the final pointer basis obtained in section IIB, the time dependant Wigner function (namely the diagonalized version of eq. (18)) is
$`\rho ^W(q,p,t)`$ $`=`$ $`{\displaystyle \underset{r}{}}\overline{\rho _r(\omega _0)}\rho _{\omega _0r}^W(q,p)+{\displaystyle \underset{r}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega \overline{\rho _r(\omega )}\rho _{\omega r}^W(q,p)+{\displaystyle \underset{rr^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega \overline{\rho (\omega ,\omega _0)}_{rr^{}}e^{i(\omega \omega _0)t}\rho _{\omega \omega _0rr^{}}^W(q,p)+`$ (55)
$`{\displaystyle \underset{rr^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}\overline{\rho (\omega _0,\omega ^{})}_{rr^{}}e^{i(\omega _0\omega ^{})t}\rho _{\omega _0\omega rr^{}}^W(q,p)+{\displaystyle \underset{rr^{}}{}}{\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle _0^{\mathrm{}}}𝑑\omega ^{}\overline{\rho (\omega ,\omega ^{})}_{rr^{}}e^{i(\omega \omega ^{})t}\rho _{\omega \omega ^{}rr^{}}^W(q,p)`$
$`=`$ $`\rho _{}^W(q,p)+\mathrm{\Delta }\rho (q,p,t),`$ (56)
where the coefficients $`\rho _r(\omega _0)`$ and $`\rho _r(\omega )`$ are the probabilities of each ”classical” final history, and $`\mathrm{\Delta }\rho (q,p,t)`$ corresponds to ”quantum” non-decohered histories<sup>\**</sup><sup>\**</sup>\**We do not write these histories in detail since they will disappear in a moment.. It is clear that when $`t\mathrm{}`$ (really after a decoherence time $`\gamma ^1)`$ the terms corresponding to these histories vanish according to Riemann-Lebesgue theorem (see paper ).
Now we know:
i.- That when $`\mathrm{}0,`$ $`\rho ^W(q,p,t)`$ satisfies the classical Liouville equation, namely the laws of classical mechanics.
ii.- That $`\rho _{}^W(q,p)0`$, but that the second term $`\mathrm{\Delta }\rho (q,p,t)`$ is not positive definite, so $`\rho ^W(q,p,t)=\rho _{}^W(q,p)+\mathrm{\Delta }\rho (q,p,t)`$ is not positive definite and therefore cannot be considered as a classical density. Nevertheless $`\mathrm{\Delta }\rho `$ vanishes when $`t\mathrm{}`$, so that $`\rho ^W(q,p,t)`$ is ”almost” positive definite. Therefore with an adequate ”coarse-graining” <sup>\*†</sup><sup>\*†</sup>\*†Or, in our language, observed by an observer space smaller than $`𝒪`$., the averaged $`\rho ^W(q,p,t)`$ may be positive definite, for $`t\gamma ^1`$, and also would satisfy the classical Liouville equation, in its way towards classical equilibrium. This $`\rho ^W(q,p,t)`$ would be a ”classical limit” before equilibrium. We will follow this line of research elsewhere. For the moment it is clear that the non-final pointer basis has for limit the final pointer basis when $`t\mathrm{}.`$ This fact may help us to find both the non final pointer basis and the classical limit.
### B Local vs. global equilibrium.
In the last subsection we have considered the case where classicality is reached before classical equilibrium. In this subsection we will see this process in a different way. We can decompose the global system in a set of local systems. In these subsystems, classical local equilibrium can be reached after a time $`\gamma ^1`$ with positive definite local equilibrium densities (and, among other things we will be able to define a classical local equilibrium entropy ). Since the classical local equilibrium densities of each subsystem are positive definite the classical density of the whole system will be positive definite. But the whole system is out of equilibrium and its evolution is defined by the classical Liouville equation. Now, as the system is already classical, since all its parts are classical, we can study its evolution towards equilibrium with a classical relaxation time $`\tau ,`$ that we are supposing $`\tau >\gamma ^1,`$ with the usual classical methods (and the global entropy of classical phenomenological thermodynamics will become maximal).
Of course, if the interaction is such that $`\tau <\gamma ^1`$ this process is not possible and we will directly reach to the final equilibrium without a previous stage of local equilibrium. But the case $`\tau >\gamma ^1`$ is the usual one and it takes place if the global-long-range interactions are unimportant with respect to the local-short-range interactions. Therefore the system can be decomposed into a set of many weakly interacting subsystems that can be considered as quasi-isolated. The local interactions transform these quantum subsystems in classical subsystems in equilibrium, and the system can be described as a non homogeneous distribution of local subsystem with classical momenta $`x,`$ $`l_1,\mathrm{},l_N.`$ These momenta, and the subsystems density are different in each of them, producing a state of classical non-equilibrium, that will reach equilibrium, due to the global-long-range forces at a time $`\tau >\gamma ^1(`$see ).
### C Final Conclusion.
Using the interplay of observables and states, that we have considered as functionals over the space of observables, we have found an exact final pointer basis and an intrinsically consistent set of final histories. So, given a Hamiltonian $`H`$ and a state $`\rho `$ we have found the exact final pointer basis $`\{|x,r_1,\mathrm{},r_N\}`$ and we have shown that $`\rho _{}^W`$, the Wigner function of $`\rho _{}`$, can be expanded in Wigner functions corresponding to the co-basis $`(x,r_1,\mathrm{},r_N|`$. To obtain this results or similar ones almost all the authors use coarse graining methods based in projectors and try to obtain a limit. So they essentially use the weak limit of eq. (23), namely:
$`\underset{t\mathrm{}}{lim}(\rho (t)|O)=(\rho _{}|O),O𝒪`$
But, at least in the classical case, we know that this weak limit exists if and only if the system is mixing . And the system is mixing if it has a continuous spectrum ( papers , , ) and the present paper can be considered as an extension of the theorem that says that the mixing evolutions have a weak limit towards equilibrium but now formulated in the quantum case. Thus the only way to deal with the problem (at least in the limit $`t\mathrm{})`$ in an exact way is to use a method , as ours, specially adapted to deal with the singularities inherent to that continuous spectrum. If not we are condemned to only do approximate calculations.
Nevertheless approximated methods are important and, in some cases, unavoidable to obtain the non-final pointer basis, but they can be better understood if they are compared with exact methods. We will continue our research following this line.
ACKNOWLEDGMENTS.
We are very grateful for the hospitality of Jonathan Halliwell, and the Imperial College (London), where the research of this subject was began by one of us (M.C.). This work was partially supported by grants CI1-CT94-0004 and PSS$`{}_{}{}^{}0992`$ of the European Community, PID 3183/93 of CONICET, EX053 of the Buenos Aires University, and also grants from Fundación Antorchas and OLAM Foundation.
## A Comparison with the literature.
In this appendix we would like to compare our method with those that can be found in the literature, where the models are studied using the variables $`Q`$ and $`P.`$
Let us first see what is the shape of the diagonal states $`(\omega _0|,`$ $`(\omega |`$ or $`\rho _{}`$ in $`Q`$ and $`P.`$ In order to determinate the diagonal states, like $`\rho _{},`$ in the configuration and momentum basis (and also to find the correlations in these variables) we are forced to go to a particular model were the relation among $`H,Q,`$and $`P`$ is defined. We consider a coupled system of an oscillator and a bath such that the Hamiltonian reads :
$$H=\frac{1}{2}\mathrm{\Omega }(p^2+q^2)+\frac{1}{2}\omega (p_\omega ^2+q_\omega ^2)𝑑\omega +\lambda V(\omega )(qq_\omega +pp_\omega )𝑑\omega $$
(A1)
where the first term corresponds to the oscillator (with bound eigenstates $`|\omega _i`$ and ground state $`|\omega _0`$), the second term to a field: the ”bath” (with eigenstates $`|\omega _1,\omega _2,\mathrm{},\omega _n`$), and the third is an interaction term (with a $`pq`$ symmetry) that leaves just one set of possible final states
$$\rho =(\omega _0|=\underset{n}{}d\omega _id\omega _i^{}\rho _{0\omega _1,\omega _2,\mathrm{},\omega _n\omega _1^{},\omega _2^{},\mathrm{},\omega _n^{}}(\omega _{0,}\omega _1,\omega _2,\mathrm{},\omega _n\omega _1,\omega _2^{},\mathrm{},\omega _n^{}|$$
(A2)
namely, the oscillator in the ground state and the bath in any state (see ). The $`(\omega _{0,}\omega _1,\omega _2,\mathrm{},\omega _n\omega _1^{},\omega _2^{},\mathrm{},\omega _n^{}|`$ are generated by the dressed operators that we will define in eq. (B7).
Then let us try to get an idea of the form of $`\rho _{}`$ via a heuristic reasoning based on the symmetry of the Hamiltonian (A1). If we consider a quantum state $`\rho `$ and the position operator $`Q`$ (that symbolizes either the operator $`q`$ or any of the operators $`q_\omega )`$ and the momentum operator $`P`$ (that symbolizes either $`p`$ or $`p_\omega )`$ in the usual case we will have:
$`(\mathrm{\Delta }Q)^2=Tr(Q^2\rho )[Tr(Q\rho )]^2=Q^2Q^2`$
$$(\mathrm{\Delta }P)^2=Tr(P^2\rho )[Tr(P\rho )]^2=P^2P^2$$
(A3)
When $`\rho `$ is a functional, we can generalize these equation as:
$`(\mathrm{\Delta }Q)^2=(\rho |Q^2)[(\rho |Q)]^2=Q^2Q^2`$
$$(\mathrm{\Delta }P)^2=(\rho |P^2)[(\rho |P)]^2=P^2P^2$$
(A4)
and in general $`(\mathrm{\Delta }Q)^2(\mathrm{\Delta }P)^2.`$ But if $`\rho `$ is the diagonal state $`\rho _{}`$ of eq. (24) (or the states $`(\omega _0|,`$ $`(\omega |)`$ we will have:
$$(\rho |Q)=(\omega _0|Q),\text{ }(\rho |Q^2)=(\omega _0|Q^2),\text{ }(\rho |P)=(\omega _0|P),\text{ }(\rho |P^2)=(\omega _0|P^2)$$
(A5)
So, as $`H`$ has a $`qp`$symmetry everything is symmetric under the transformation $`pq`$ (or $`qi\frac{}{q},i\frac{}{q}q)`$ and therefore $`(\mathrm{\Delta }Q)^2=(\mathrm{\Delta }P)^2.`$ This would not be the case if $`\rho `$ would not be diagonal in the basis where $`H`$ is diagonal, since $`Q`$ and $`P`$ are not diagonal in this basis, e. g. $`\rho `$ could commute with $`P`$ but not with $`Q`$ showing, in this case, a clear asymmetry $`PQ.`$ Thus our diagonal states are states such that $`\mathrm{\Delta }Q=\mathrm{\Delta }P`$. Namely: for the ground state we have $`\mathrm{\Delta }p=\mathrm{\Delta }q,`$ really a well known fact. If now we introduce in (A1) a small asymmetric interaction $`\lambda ^{}W`$ $`(\lambda ^{}1)`$ we will have $`\mathrm{\Delta }Q\mathrm{\Delta }P`$. On the contrary if the interaction is $`\lambda ^{}W(q,q_\omega )`$ ($`\lambda ^{}1)`$ the Hamiltonian $`H`$ can be neglected and the diagonal states will be position eigenvalues (so our results coincide with those of refs. and , see a detailed example below).
## B An example of correlations and localization.
There are systems, e. g. the one of appendix A, with variables $`Q`$ and $`P`$ and a bath, where the interaction is such that $`Q`$ and $`P`$ become correlated. Namely the evolution makes both $`\mathrm{\Delta }Q`$ and $`\mathrm{\Delta }P`$ bounded, and a wave packet appears that eventually becomes a minimal uncertainty wave packet, when $`t\mathrm{},`$ then maximal localization appears in the usual way as promised in section IV. As an example, let us now find the correlations between $`Q`$ and $`P`$ in the model of Hamiltonian (A1) using our method as explained above and in the paper . Let us first write the Hamiltonian (A1) using creation and annihilation operators.
$$H=\mathrm{\Omega }b^{}b+𝑑𝐤\omega _ka_𝐤^{}a_𝐤+𝑑𝐤V_k(a_𝐤^{}b+b^{}a_𝐤),\omega _k=k,k=\left|𝐤\right|$$
(B1)
The coordinate $`q`$ and the momentum $`p`$ of the oscillator can be expressed as a function of the $`b^{}`$ and $`b`$ as:
$$q=\left(\frac{\mathrm{}}{2m\mathrm{\Omega }}\right)^{\frac{1}{2}}(b^{}+b);p=i\left(\frac{m\mathrm{}\mathrm{\Omega }}{2}\right)^{\frac{1}{2}}(b^{}b).$$
(B2)
We can adimensionalize the last equation defining $`Q`$ and $`P`$ such that:
$$q=\left(\frac{\mathrm{}}{m\mathrm{\Omega }}\right)^{\frac{1}{2}}Q,p=\left(m\mathrm{}\mathrm{\Omega }\right)^{\frac{1}{2}}P$$
(B3)
Then:
$$Q=\frac{1}{\sqrt{2}}(b^{}+b),P=i\frac{1}{\sqrt{2}}(b^{}b)$$
(B4)
$$b=\frac{1}{\sqrt{2}}(Q+iP),b^{}=\frac{1}{\sqrt{2}}(QiP)$$
(B5)
In the Heisenberg representation the operator $`Q`$ evolves as:
$$Q(t)=\frac{1}{\sqrt{2}}(b^{}(t)+b(t))=\frac{1}{\sqrt{2}}𝑑𝐤V_k\left(\frac{1}{\eta _{}(k)}A_𝐤^{}e^{i\omega _kt}+\frac{1}{\eta _+(k)}A_𝐤e^{i\omega _kt}\right),$$
(B6)
where:
$$A_𝐤^{}=a_𝐤^{}+\frac{V_k}{\eta _+(k)}\left(b^{}+\frac{d𝐤^{}V_k^{}a_𝐤^{}^{}}{\omega _k\omega _k^{}+i0}\right),A_𝐤=a_𝐤+\frac{V_k}{\eta _{}(k)}\left(b+\frac{d𝐤^{}V_k^{}a_𝐤^{}}{\omega _k\omega _k^{}+i0}\right).$$
(B7)
The functions $`\eta _{}(k)`$ and $`\eta _+(k)`$ and all the details of the calculations can be found in paper .
Let us consider the initial conditions
$$Q_{t=0}=Q_0,P_{t=0}=P_0$$
(B8)
for the oscillator, and also
$$a_𝐤^{}_{t=0}=a_𝐤_{t=0}=0,$$
(B9)
which corresponds to the field being initially in its ground state. Therefore:
$$b_{t=0}=\frac{1}{\sqrt{2}}(Q_0+iP_0),b^{}_{t=0}=\frac{1}{\sqrt{2}}(Q_0iP_0)$$
(B10)
and for the time evolution of the mean value of the coordinate and momentum of the oscillator we obtain
$$Q_t=\frac{1}{\sqrt{2}}𝑑𝐤\frac{V_k^2}{\eta _{}(k)\eta _+(k)}\left(e^{i\omega _kt}b^{}_{t=0}+e^{i\omega _kt}b_{t=0}\right),$$
(B11)
$$P_t=\frac{i}{\sqrt{2}}𝑑𝐤\frac{V_k^2}{\eta _{}(k)\eta _+(k)}\left(e^{i\omega _kt}b^{}_{t=0}e^{i\omega _kt}b_{t=0}\right).$$
(B12)
The oscillating time dependent factors inside the integrals produce the vanishing of both $`Q_t`$ and $`P_t`$ for very long times. We can study the poles of the analytic extension of the factor $`\frac{V_k^2}{\eta _{}(k)\eta _+(k)}`$, as in paper and prove that the trajectory of $`Q_t`$ and $`P_t`$ in the phase space of the oscillator is a spiral ending at $`Q=P=0`$.
Now we would like to compute $`\mathrm{\Delta }q`$ and $`\mathrm{\Delta }p`$ as a function of time. In addition to equations (B8) and (B9) let us assume the following initial conditions for the oscillator
$$bb^{}_{t=0}=\beta ,bb_{t=0}=\alpha ,b^{}b_{t=0}=1\beta ,b^{}b^{}_{t=0}=\alpha ^{},$$
(B13)
being $`\alpha `$ and $`\beta `$ some arbitrary constants. If the field is in its ground state, we also have
$$a_𝐤a_𝐤^{}^{}_{t=0}=\delta ^3(𝐤𝐤^{}),a_𝐤^{}a_𝐤^{}_{t=0}=a_𝐤^{}a_𝐤^{}^{}_{t=0}=a_𝐤a_𝐤^{}_{t=0}=0.$$
(B14)
All other initial mean values of products of pairs of creation or annihilation operators are zero. This means that we have taken the oscillator in an arbitrary state and the field in the ground state as initial conditions.
Therefore we have
$`A_𝐤^{}A_𝐤^{}^{}_{t=0}`$ $`=`$ $`{\displaystyle \frac{V_kV_k^{}}{\eta _+(k)\eta _+(k^{})}}\alpha ^{},`$ (B15)
$`A_𝐤^{}A_𝐤^{}_{t=0}`$ $`=`$ $`{\displaystyle \frac{V_kV_k^{}}{\eta _+(k)\eta _{}(k^{})}}(\beta 1),`$ (B16)
$`A_𝐤A_𝐤^{}_{t=0}`$ $`=`$ $`{\displaystyle \frac{V_kV_k^{}}{\eta _{}(k)\eta _{}(k^{})}}\alpha ,`$ (B17)
$`A_𝐤A_𝐤^{}^{}_{t=0}`$ $`=`$ $`\delta ^3(𝐤𝐤^{})+{\displaystyle \frac{V_k^{}V_k}{\eta _+(k^{})(\omega _k^{}\omega _k+i0)}}+{\displaystyle \frac{V_kV_k^{}}{\eta _{}(k)(\omega _k\omega _k^{}i0)}}+`$ (B19)
$`+{\displaystyle \frac{V_kV_k^{}}{\eta _{}(k)\eta _+(k^{})}}\beta +{\displaystyle \frac{V_kV_k^{}}{\eta _{}(k)\eta _+(k^{})}}{\displaystyle \frac{d𝐤^{\prime \prime }V_{k^{\prime \prime }}^2}{(\omega _k\omega _{k^{\prime \prime }}i0)(\omega _k^{}\omega _{k^{\prime \prime }}+i0)}}.`$
The time evolution of the mean value of $`Q(t)^2`$ is given by
$`Q(t)^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d𝐤V_k({\displaystyle \frac{1}{\eta _{}(k)}}A_𝐤^{}e^{i\omega _kt}+{\displaystyle \frac{1}{\eta _+(k)}}A_𝐤e^{i\omega _kt})\times `$ (B21)
$`\times {\displaystyle }d𝐤^{}V_k^{}({\displaystyle \frac{1}{\eta _{}(k^{})}}A_𝐤^{}^{}e^{i\omega _k^{}t}+{\displaystyle \frac{1}{\eta _+(k^{})}}A_𝐤^{}e^{i\omega _k^{}t})`$
Then, replacing equations (B17) in equation (B21) and always using the Riemann-Lebesgue theorem we have:
$$\underset{t\mathrm{}}{lim}Q(t)^2=\frac{1}{2}𝑑𝐤\frac{V_k^2}{\eta _{}(k)\eta _+(k)}=\frac{1}{2},$$
(B22)
and therefore
$`\underset{t\mathrm{}}{lim}\left[\mathrm{\Delta }Q(t)\right]^2=\underset{t\mathrm{}}{lim}(Q(t)Q(t))^2={\displaystyle \frac{1}{2}}.`$
Making an analogous calculation for $`P`$, we have:
$$\underset{t\mathrm{}}{lim}\mathrm{\Delta }Q=\underset{t\mathrm{}}{lim}\mathrm{\Delta }P=\frac{1}{\sqrt{2}}$$
(B23)
so the wave packet around the spiral trajectory evolves to a minimal uncertainty symmetrical wave packet, showing the localization process in the usual way. This proves the presence of correlations in our model. Reestablishing the units, when $`t\mathrm{},`$ and introducing the velocity $`v`$ we have:
$$\mathrm{\Delta }q=\left(\frac{\mathrm{}^2}{2m\mathrm{\Omega }}\right)^{\frac{1}{2}},\mathrm{\Delta }v=\left(\frac{\mathrm{}^2\mathrm{\Omega }}{2m}\right)^{\frac{1}{2}}$$
(B24)
This fact shows that the wave packet is more peaked for big a $`m`$ than for small a $`m.`$ Then, in some models big mass particles can be considered as classical while other remain quantum. Moreover if $`\mathrm{}0`$ the uncertainties disappear. Also the classical limit of the oscillator has a spiral motion in phase space. In fact, if using eq. (32) we compute the Wigner function corresponding to the matrix density, we will find the motion of this classical density that will be centered in the spiral trajectory and having, when $`t\mathrm{},`$ a symmetrical circle of diameter ($`\frac{1}{2}\mathrm{})^{\frac{1}{2}}`$ as support <sup>\*‡</sup><sup>\*‡</sup>\*‡The cosmological models of papers are other examples that we will further develop elsewhere..
## C Decoherence of histories.
From the section VA we can conclude that our notion of history of the system is essentially contained in the state $`\rho (t).`$ This history ends in the final equilibrium state $`\rho _{}`$. In this appendix we will study the relation of this notion with the usual histories formalism and compare the results. The computation will turn out to be very simple for two reasons:
i.- As in all the paper we will work only in the limit $`t\mathrm{}`$, where the existence of an exact final pointer basis will make all the computations quite trivial.
ii.- Also we will only consider one space of observables $`𝒪`$ and therefore just one set of final consistent histories . The case of many sets will be considered elsewhere.
Nevertheless we think that the results are of some interest since:
i.- For times $`t\gamma ^1`$ all the exact result obtained in the limit $`t\mathrm{}`$ can be considered as good approximations.
ii.-The existence of a space of observables $`𝒪,`$ where we can use the Riemann-Lebesgue theorem, perhaps can be considered as a selection principle to choose the physically relevant consistent set .
So let us begin giving the main definitions. Let us consider a time depending basis of $`:`$ $`\{|\alpha (t)\},`$ and the projectors:
$$P_\alpha (t)=|\alpha (t)\alpha (t)|$$
(C1)
such that they represent exhaustive and exclusive alternatives:
$$\underset{\alpha }{}P_\alpha =1,\text{ }P_\alpha P_\beta =\delta _{\alpha \beta }P_\alpha $$
(C2)
We will call (fine-grained) history $`\stackrel{}{\alpha }`$ to a string of time depending projectors We can consider a more general case were the exclusive and exhausting set of histories is different at every time $`t_i`$ and therefore the projectors are $`P_{\alpha _i}^i(t_i).`$ But this is not the usual case.:
$$C_\stackrel{}{\alpha }=P_{\alpha _1}(t_1)\mathrm{}P_{\alpha _n}(t_n),\text{ }t_1<\mathrm{}<t_n$$
(C3)
For a state $`\rho `$ we will call decoherence matrix to:
$$M(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})=C_\stackrel{}{\alpha }^{}\rho C_\stackrel{}{\alpha }^{}=P_{\alpha _n}(t_n)\mathrm{}P_{\alpha _1}(t_1)\rho P_{\alpha _1^{}}(t_1^{})\mathrm{}P_{\alpha _n^{}}(t_n^{})$$
(C4)
We introduce this matrix because we consider it as the natural generalization of the usual density matrix to the case where single projectors are changed by histories.
We will call decoherence functional to:
$$D(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})=TrM(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})$$
(C5)
that would be the generalization of the trace of an ordinary matrix If some of the $`\alpha `$ are continuous indices, for them we must use the generalization of the trace introduced in paper .
We will call candidate probability for the history $`\stackrel{}{\alpha }`$ to:
$$p(\stackrel{}{\alpha })=TrM(\stackrel{}{\alpha },\stackrel{}{\alpha })$$
(C6)
that would be the generalization of the usual probability. It is only a ”candidate probability” because, at this stage, it does not satisfy the axioms of the usual boolean probability theory.
If
$$ReD(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})=0$$
(C7)
for $`\stackrel{}{\alpha }\stackrel{}{\alpha }^{}`$ we will say that the set of histories is consistent or weakly decoherent, in this case it is proved that the set can be in principle submitted to the ordinary boolean logic , and the candidate probability can be considered as the probability of each history.
If
$$D(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})=0$$
(C8)
for $`\stackrel{}{\alpha }\stackrel{}{\alpha }^{}`$ we will say that the set has medium decoherence. Theorems about records can be proved if the set of histories has this type of decoherence .
If
$$M(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})=0$$
(C9)
for $`\stackrel{}{\alpha }\stackrel{}{\alpha }^{}`$ we will say that the set is intrinsically consistent or that it has matrix decoherence. Of course matrix decoherence implies medium decoherence, and medium decoherence implies weak decoherence.
Let us now compare all these concepts with our formalism.
We choose:
$$|\alpha (t_1)=|x,r_1,\mathrm{},r_N;t_1$$
(C10)
where we have used the shorthand notation introduced above. The set of operators $`P_{\alpha (t)}=|\alpha (t)\alpha (t)|=|x,r_1,\mathrm{},r_N;tx,r_1,\mathrm{},r_N;t|`$ will be our ”final pool of operators” if we use the language of paper . The evolution of these operators will be:)
$`P_\alpha (t)=e^{iH(tt_1)}P_\alpha (t_1)e^{iH(tt_1)}=e^{ix(tt_1)}P_\alpha (t_1)e^{ix(tt_1)}=`$
$$P_\alpha (t_1)=P_\alpha =|\alpha (0)\alpha (0)|$$
(C11)
i. e.: these operators are constant. Then the projectors are time constant and:
$$C_\stackrel{}{\alpha }=P_\alpha $$
(C12)
and these histories can be labeled with the ordinary $`\alpha `$ instead of the $`\stackrel{}{\alpha }`$ with the arrow.
In more detail let us first study our ”pool” of projectors to compute eq. (C4) in our formalism and when $`t\mathrm{},`$
$$P_{\alpha (t)}=P_\alpha =|x,r_1,\mathrm{},r_Nx,r_1,\mathrm{},r_N|=|x,r_1,\mathrm{},r_N)$$
(C13)
1.-$`r_1,\mathrm{},r_N`$ are discrete indices and the final stationary state $`\rho _{}`$ is diagonal in these indices, so this part of the problem is trivial.
2.-$`x`$ symbolizes $`(\omega _0,\omega )`$ where only $`\omega `$ is continuous, so the treatment of $`\omega _0`$ is also trivial.
The problem is only $`\omega `$ so, for simplicity, let us only consider this index. The projector reads:
$$P_\omega =|\omega \omega |=|\omega )$$
(C14)
So let us compute:
$$P_\omega \rho _{}P_\omega ^{}=|\omega \omega |\rho _{}|\omega ^{}\omega ^{}|$$
(C15)
but first we must find the meaning of this symbol. In the discrete case we have:
$$|ab|\rho |cd|=|aTr(\rho |cb|)d|$$
(C16)
that can be generalized to the continuous case as:
$$|ab|\rho |cd|=|a(\rho ||cb|)d|$$
(C17)
Thus:
$$P_\omega \rho _{}P_\omega ^{}=||\omega \omega |[\rho _{\omega \mathrm{"}}(\omega \mathrm{"}|d\omega \mathrm{"}]|\omega ^{}\omega ^{}|=|\omega [\rho _{\omega \mathrm{"}}(\omega \mathrm{"}|\omega ^{},\omega )d\omega \mathrm{"}]\omega ^{}|$$
(C18)
So, from eqs. (12) we have:
1.- If $`\omega \omega ^{}`$ it is $`P_\omega \rho _{}P_\omega ^{}=0.`$
2.- If $`\omega =\omega ^{}`$ it is:
$$P_\omega \rho _{}P_\omega ^{}=|\omega \left[\rho _{\omega \mathrm{"}}(\omega \mathrm{"}|\omega )𝑑\omega \mathrm{"}\right]\omega |=|\omega \left[\rho _{\omega \mathrm{"}}\delta (\omega \mathrm{"}\omega )𝑑\omega \mathrm{"}\right]\omega |=\rho _\omega |\omega \omega |$$
(C19)
So with a symbolic obvious notation (that we will use from now on) we can say that:
$$P_\omega \rho _{}P_\omega ^{}=|\omega \rho _\omega \delta _{\omega \omega ^{}}\omega ^{}|$$
(C20)
If now we repeat the reasoning including all the trivial discrete indices we will obtain the same result since $`\rho _{}`$ is diagonal in these indices. Then, when $`t\mathrm{}`$ we have that
$$M(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})\delta _{\alpha \alpha ^{}}\rho _\alpha |\alpha \alpha |$$
(C21)
and therefore we have final matrix decoherence in a time long enough. Then we have found the final ”statistical classical domain or realm” of Gell-Mann and Hartle. In this way final classical behavior emerges from quantum behavior and transition 2a” appears in the histories formalism. Essentially we have used the weak limit of eq. (25) and the fact that it is the only possible limit we can use, since $`\rho `$ is a functional over the space $`𝒪.`$ But the choice of eq. (C10) has an extra bonus: it decomposes the density matrix just in the way that was announced in the introduction.
From the matrix decoherence we have medium decoherence and weak decoherence, so we have proved that any quantum system, fulfilling the conditions required in section II, has a set of final intrinsically consistent histories, the essential conditions being the continuous spectrum and the existence of just one ground state. Classically these histories will be the $`\rho _{xr}^W(q,p)`$ of eq.(56). This exact final decoherence has being obtained using the basis {$`|xr\}`$, other near bases obviously yield final approximate decoherence. Also basis {$`|xr\}`$ will give approximate decoherence in a time long enough.
But we must observe that in all cases where $`P_\alpha (t)=P_\alpha =const.`$ (even if eq. (C10) is not satisfy) we can immediately prove medium decoherence with no reference to matrix decoherence. In fact, if $`P_\beta =|\beta \beta |=const.`$ we have:
$$D(\stackrel{}{\beta },\stackrel{}{\beta }^{})=D(\beta ,\beta ^{})=Tr(|\beta \beta |\rho |\beta ^{}\beta ^{}|)=\beta |\beta ^{}\beta |\rho |\beta ^{}=\delta _{\beta \beta ^{}}p(\beta )$$
(C22)
These would be the case with the $`P_\alpha `$ of this section and also for any constant $`P_\beta `$. This result seems very trivial but it is not. The essential property of projector (C13) is that it is time constant, but our formalism contains other time-constant projectors. If we go back to section IIA we find:
$$P_{\beta (t)}=P_\beta =|x,m_1,\mathrm{},m_Nx,m_1,\mathrm{},m_N|=|x,m_1,\mathrm{},m_N)$$
(C23)
namely the projectors related with the basis $`\{x,m\}`$ before the diagonalization (26) that yields the basis $`\{x,r\}.`$ The $`P_\beta `$ are also time constants and yield medium decoherence (only the $`P_a`$ yield matrix decoherence). The main fact is that in order to reach the classical statistical mechanics of section III we must use the basis $`\{x,r\}`$ that diagonalize $`\rho _{}`$ in all indices (see eqs. (37) to (43)). Thus, since our demonstration is based in the matrix decoherence in the basis $`\{x,r\},`$ these objects are essential for us. Only after this demonstration we can speak of classical constants of the motion and classical trajectories because only then we can pass from the quantum formulae to the classical ones.
Then the last result can be translated as follows:
1.- There is final matrix decoherence between any pair of different sets of constants ($`x,r)`$ i. e. between any pair of sets of classical trajectories in the phase space. This set of sets of trajectories is intrinsically consistent (see eq. (27)).
2.- But, of course, any set of functions of the $`\mathrm{"}r\mathrm{"}`$, such as the $`\mathrm{"}m\mathrm{"}`$, will define equally well the set of classical trajectories. But the $`\mathrm{"}m\mathrm{"}`$ do not provide a basis with good defined probabilities, as the $`\mathrm{"}r\mathrm{"}`$ does, since in the basis $`\mathrm{"}m\mathrm{"}`$ the $`\rho _{}`$ is not diagonal (see eq. (24)). In this case the set of histories is consistent but not intrinsically consistent.
So our point of view is that, even if all sets endowed of medium decoherence can be considered as consistent sets, there is only one with physical importance, the one with matrix decoherence, the only one which is an ”intrinsically consistent set”. This idea may help to find the selection principle searched in papers .
Finally, if the potential and the initial conditions are such to privilege a history (as in appendix B) the locations process ”b” will take place and we will have a unique classical object with a unique history $`P_\alpha .`$ Then we would find the final ”classical domain or realm” of Gell-Mann and Hartle.
We will end this section showing how several requirements necessary for a efficient histories decoherence are satisfied by our formalism:
### 1 Griffiths-Omnès condition
The Griffiths-Omnès condition for consistency , is automatically satisfied since:
$$ReTr[|\alpha \alpha |\rho (1|\alpha \alpha |)|\alpha \alpha |]=0$$
(C24)
### 2 Permanence of the past.
If we take our projectors from the pool of the projectors $`|\alpha \alpha |`$ the condition of permanence of the past is trivially satisfied , since a chain with a certain number of $`|\alpha \alpha |`$ can only be continued repeating this projector. This is the most important property required in papers .
### 3 Insensitivity.
While quantum states are modified by the measurement processes, classical states are not sensible to these measurements. This property of classical states is called insensitivity . The projector $`P_{\alpha _i}`$ =$`|\alpha _i\alpha _i|`$ can be considered as a measurement operator, so if $`\rho _{before}`$ is the state before the measurement and $`\rho _{after\text{ }}`$is the state after the measurement, we will have:
$$\rho _{after}=\underset{i}{}P_{\alpha _i}\rho _{before}P_{\alpha _i}=\underset{i}{}|\alpha _i\alpha _i|\rho _{before}|\alpha _i\alpha _i|=\underset{i}{}P_{\alpha _i}\rho _{before}P_{\alpha _i}$$
(C25)
where $`p_i`$ is the probability to measure $`\alpha _i`$ . Now if, after the decoherence process, $`\rho _{before}`$ is a diagonal matrix, precisely $`\rho _{}`$ i.e.
$$\rho _{before}=\underset{i}{}p_i|\alpha _i\alpha _i|$$
(C26)
and we only measure the observers in the CSCO $`\{H,P_1,\mathrm{}P_n\},`$ so the $`P_{\alpha _i}`$ are just the $`P_\alpha =|\alpha \alpha |,`$ we have:
$$\rho _{after}=\underset{i}{}|\alpha _i\alpha _i|(\underset{j}{}p_j|\beta _j\beta _j|)|\alpha _i\alpha _i|=\underset{i}{}p_i|\alpha _i\alpha _i|=\rho _{before}$$
(C27)
So, in fact, the matrix $`\rho _{}`$ is insensitive to the measurement of the CSCO $`\{H,P_1,\mathrm{}P_n\}`$ $`(`$ and also the CSCO $`\{H,O_1,\mathrm{},O_N\}`$ where the operators $`O`$ are related with the constants $`m)`$. This is the maximum insensitivity we can get.
### 4 Strong decoherence and records.
If for any history $`\stackrel{}{\alpha }`$ there is a projector $`R_\alpha `$ such that $`\{R_\alpha \}`$ is not necessarily a complete set of projectors in $``$, in the sense that $`\{R_\alpha |\psi \}`$ is not necessarily a basis of $``$, and for any state $`\rho `$ it is:
$$C_\stackrel{}{\alpha }\rho =R_\alpha \rho ,\text{ }R_\alpha R_\beta =\delta _{\alpha \beta }P_\alpha $$
(C28)
we will say the we have strong decoherence (, eq. (2.4)). As the $`R_\alpha `$ are timeless entities and as $`C_\stackrel{}{\alpha }R_\alpha ,`$ $`R_\alpha `$ can be considered as the $`record`$ of the history $`\stackrel{}{\alpha }.`$ But $`R_\alpha `$ can also be considered as the record of not one but several decohered histories, associated by unitary transformations . So really $`R_\alpha `$ is the record of an equivalent class of histories.
It is clear that if these records exist we have medium decoherence. In fact:
$$D(\stackrel{}{\alpha },\stackrel{}{\alpha }^{})=Tr(C_\stackrel{}{\alpha }^{}\rho C_\stackrel{}{\alpha }^{})=Tr(R_\alpha \rho R_\alpha ^{})=Tr(\rho R_\alpha ^{}R_\alpha )=\delta _{\alpha \alpha ^{}}p(\stackrel{}{\alpha })$$
(C29)
so strong decoherence implies medium decoherence.
In our case these final $`R_\alpha `$ exist and they are:
$$R_\alpha =|\alpha \alpha |=|xrxr|=P_\alpha $$
(C30)
Thus the numbers $`x,r_1,\mathrm{},r_N`$ can be considered as the record of the corresponding final history.
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# The action of the Frobenius map on rank 2 vector bundles in characteristic 2
## 1 Introduction
Let $`X`$ be a smooth algebraic curve of genus $`g`$ over a field $`k`$ of characteristic $`p>0`$. The behaviour of semi-stable bundles with respect to the absolute Frobenius $`F_a`$ remains mysterious if $`g2`$. Let us briefly explain why this question should be of interest. Start with a continuous representation $`\rho `$ of the algebraic fundamental group in $`\mathrm{GL}_r(\overline{k})`$, where $`\overline{k}`$ is the algebraic closure of $`k`$. Let $`E_\rho `$ be the corresponding rank $`r`$ bundle over $`X`$. Then all the bundles $`F_a^{(n)}E_\rho `$, for $`n>0`$, where $`F_a^{(n)}`$ denotes the $`n`$-fold composite $`F_a\mathrm{}F_a`$, are semi-stable. Conversely, assuming that $`k`$ is finite, let $`E`$ be a semi-stable rank $`r`$ bundle defined over $`\overline{k}`$. Because the set of isomorphism classes of semi-stable bundles of degree $`0`$ over $`X_k^{}`$, where $`k^{}`$ is any finite extension of $`k`$, is finite, one observes (see \[LS\]) that some twist of $`E`$ comes from a representation as above. Therefore, if one is interested in unramified continuous representations of the Galois group over $`k`$ of a global field $`k(X)`$ in characteristic $`p`$, it is natural to look at Frobenius semi-stable bundles, that is those whose pull-backs by $`F_a^{(n)}`$ are all semi-stable. This condition is stable by tensor product (\[Mi\] section 5), which is not usually the case for ordinary semi-stability in positive characteristic.
Assume that $`k`$ is arbitrary of characteristic $`p`$. Let us emphasize that the locus of Frobenius semi-stable bundles of degree $`0`$ is a countable intersection of open subsets of the coarse moduli scheme of semi-stable vector bundles of degree $`0`$ over $`X`$. In particular, it is not clear at all if the Frobenius semi-stable points are dense in general. This could a priori depend on the arithmetic of the base field $`k`$. We would like to study the dynamics of the Frobenius, namely the sequence $`nF_a^{(n)}E`$ for a given vector bundle $`E`$ over $`X`$.
If $`k`$ is a discrete valuation field of characteristic $`p`$ and if $`X`$ is a Mumford curve, the situation is well understood \[F\], \[G\]: among other results, it is shown for arbitrary genus and characteristic that there exist semi-stable bundles which are destabilized by the Frobenius $`F_a`$, that $`EF_a^{}E`$ induces a surjective (rational) map on the moduli space of semi-stable rank $`r`$ vector bundles of degree $`0`$, and that the set of bundles coming from continuous representations of the algebraic fundamental group is dense. Another case, which is studied in the literature, are elliptic curves (see e.g. \[O\],\[S\]). For example, it can easily be shown that over an elliptic curve semi-stability is preserved under pull-back by Frobenius and that a stable bundle of rank $`r`$ and degree $`d`$ is Frobenius stable if and only if $`pd`$ and $`r`$ are coprime. But in general, not much seems to be known.
In this paper we study the action of the Frobenius map $`F_a`$ in a very particular case: $`X`$ is an ordinary curve of genus $`2`$ defined over an algebraically closed field $`k`$ of characteristic $`2`$. In that case the coarse moduli space $`\mathrm{M}_X`$ of semi-stable rank $`2`$ vector bundles of trivial determinant is known to be the projective space $`^3`$ and we show (Proposition 6.1) that the Frobenius map identifies with a rational map $`^3^3`$ given by quadratic polynomials, which are explicitly computed in terms of generalized theta constants of the curve $`X`$. This result allows us to deduce that, if the determinant is fixed, the Frobenius map is not surjective and that the set of Frobenius semi-stable bundles is Zariski dense in $`^3`$ (Proposition 8.1). On the other hand, the Frobenius map becomes surjective if the determinant is not fixed (Proposition 6.4).
As a by-product of our computations, we also get the explicit equation of Kummer’s quartic surface (section 4) and a description of the Frobenius map acting on the moduli space of rank $`2`$ vector bundles of fixed odd degree determinant (section 7).
For other aspects of this problem see the recent article \[JX\].
We would like to thank M.S. Narasimhan and C.S. Seshadri who explained to the authors why the classical identification $`\mathrm{M}_X=^3`$ of \[NR\] remains valid in characteristic $`2`$. We especially thank M. Raynaud for having introduced us to these questions and for his encouragements.
## 2 Review of Theta groups and Theta divisors in characteristic $`p`$
Let $`k`$ be a perfect field of characteristic $`p>0`$. By schemes we implicitly mean $`k`$-scheme.
### 2.1 Relative Frobenius
Let $`A`$ be an ordinary abelian variety of dimension $`g`$. We shall denote by $`F`$ the relative Frobenius morphism
$$F:AA_1$$
which is a purely inseparable $`k`$-isogeny of degree $`p^g`$. Its kernel $`\widehat{G}`$ is a subgroup scheme of the group scheme $`A[p]`$ of $`p`$-torsion points of $`A`$ and $`F`$ identifies to the quotient morphism $`AA/\widehat{G}`$. Because
$$\mathrm{ker}F=\widehat{G}A[p]=\mathrm{ker}[p],$$
the diagram
$$AA/\widehat{G}A/A[p]$$
becomes
$$[p]:A\stackrel{F}{}A_1\stackrel{V}{}A.$$
###### 2.1 Remark.
In the case where $`A`$ is a Jacobian, the map $`V`$, called Verschiebung, is simply the pull-back by the relative Frobenius $`F:XX_1`$.
Let $`G`$ be the reduced part of $`A[p]`$. Because $`A`$ is ordinary, $`G`$ intersects $`\widehat{G}`$ at the origin and one gets a canonical decomposition
$$A[p]=G\times \widehat{G}$$
(2.1)
and the relative Frobenius induces an isomorphism
$$F:\mathrm{ker}V=A[p]/\widehat{G}\stackrel{}{}G.$$
### 2.2 Theta group scheme
We will use the basic results (and notations) of Theta groups associated to line bundles on abelian varieties as in \[Mu2\]. We assume that $`A_1`$ is principally polarized and that $`k`$ is big enough in order to have an isomorphism $`G(/p)^g`$ (however, we do not choose such an isomorphism). Let $`\mathrm{\Theta }_1`$ be a (not necessarily symmetric) Theta divisor representing the polarization. We denote by $`M_1`$ (resp. $`L_1`$) the line bundle $`𝒪(\mathrm{\Theta }_1)`$ (resp. $`𝒪(p\mathrm{\Theta }_1)`$). We consider the Theta group $`G(L_1)`$ associated to the line bundle $`L_1`$ (see \[Mu2\], page 221). This group scheme is a central extension
$$0𝔾_mG(L_1)K(L_1)0$$
and $`K(L_1)=A_1[p]`$. We denote by $`e^{L_1}:K(L_1)\times K(L_1)𝔾_m`$ the skew-symmetric form on $`K(L_1)`$ given by the commutator taken in the Theta group $`G(L_1)`$.
###### 2.2 Lemma.
There exists a unique ample line bundle $`M`$ on $`A`$ such that $`V^{}M=L_1.`$
###### Proof.
Because $`K(L_1)`$ is self-dual and contains $`\widehat{G}`$, it is the whole $`A[p]`$. One has to show (theorem 2, page 231 of \[Mu2\]) that the restriction of $`G(L_1)`$ to $`G=\mathrm{ker}V`$ is split (the uniqueness is obvious). Since $`G`$ is reduced, it is enough to define the splitting at the level of $`k`$-points. Let $`gG(k)`$. Since $`\mu _p(k)=\{1\}`$ and $`k^{}`$ is divisible, there exists a unique $`\stackrel{~}{g}G(L_1)`$ over $`g`$ satisfying $`\stackrel{~}{g}^p=1`$. Since $`\mu _p(k)=\{1\}`$ again, the restriction of $`e^{L_1}`$ to $`G(k)\times G(k)`$ is trivial and therefore $`G`$ is isotropic for $`e^{L_1}`$. This implies that any two elements $`\stackrel{~}{g},\stackrel{~}{h}`$ commute and the map $`g\stackrel{~}{g}`$ is a morphism of group schemes. ∎
The bundle $`M`$ defines a principal polarization on $`A`$ and the $`G`$-invariant morphism
$$V^{}M^p=V^{}[(V_{}L_1)^G]^pL_1^p=M_1^{p^2}=V^{}F^{}M_1$$
defines an isomorphism
$$M^p\stackrel{}{}F^{}M_1.$$
(2.2)
We denote by $`G(L)`$ the Theta group scheme associated to the line bundle $`L:=F^{}M_1=M^p`$.
###### 2.3 Lemma.
The restriction of $`G(L)`$ to both $`G`$ and $`\widehat{G}`$ is canonically split.
###### Proof.
Because $`L`$ is the pull-back of $`M_1`$ by $`F`$, the restriction of $`G(L)`$ to both $`\widehat{G}`$ splits (see Theorem 2, page 231 of \[Mu2\] again). Observe however that the splitting is determined not only by $`L`$ but by $`M_1`$. For the splitting over $`G`$, proceed as in lemma 2.2. ∎
Therefore, the decomposition $`A[p]=K(L)=\widehat{G}\times G`$ of (2.1) is symplectic in the sense that both $`G`$ and $`\widehat{G}`$ are isotropic for $`e^L`$. In particular, $`e^L`$ identifies $`\widehat{G}`$ and the Cartier dual of $`G`$.
###### 2.4 Remark.
One can interpret more geometrically the action of $`\widehat{G}`$ as follows. By \[K\], the line bundle $`L`$ comes with a canonical $`p`$-integrable connection $``$. The Lie algebra $`𝔤`$ of $`\widehat{G}`$ is the whole tangent space $`𝐓_0A`$. The link between the infinitesimal action of $`𝔤`$ on $`L`$ and $``$ is simply given by
$$v.l=_vl$$
where $`l`$ is a local section of $`L`$, $`v`$ a point of $`𝔤`$ and $`_vl`$ is the $``$-derivative of $`l`$ with respect to the invariant vector field defined by $`v`$.
Let $`\theta `$ be a non-zero global section of $`M`$. By the remark above, since the section $`\theta ^p`$ is annihilated by $``$, it defines a $`\widehat{G}`$-invariant section of $`F^{}\mathrm{\Theta }_1`$, namely a non-zero section $`\theta _1`$ of $`M_1`$ such that
$$\theta ^p=F^{}\theta _1.$$
(2.3)
The main theorem on the group scheme $`G(L)`$ and its representations is the following structure theorem due to Mumford (for the characteristic $`p`$ version we need, we refer to \[Sek\]):
###### 2.5 Theorem.
The space of global sections $`H^0(A,L)`$ is the unique irreducible representation of weight $`1`$ (i.e. $`𝔾_m`$ acts by its natural character) of the Theta group scheme $`G(L)`$.
### 2.3 Canonical Theta basis of $`|p\mathrm{\Theta }|`$ and $`|p\mathrm{\Theta }_1|`$
In \[Mu1\] Mumford constructs canonical bases for any linear system $`H^0(A,L)`$ where $`L`$ is a line bundle of separable type. Because of Theorem 2.5, we can adapt his construction to line bundles $`L`$ not of separable type. In our situation, namely $`L=𝒪_A(p\mathrm{\Theta })`$, where $`\mathrm{\Theta }`$ is not necessarily symmetric, we get the following
###### 2.6 Lemma.
(i) There exists a basis $`\{X_g\}_{gG}`$ of $`H^0(A,L)`$, unique up to a multiplicative scalar, which satisfies the following relations
$$a.X_g=X_{g+a}\alpha .X_g=e^L(\alpha ,g)X_ga,gG,\alpha \widehat{G}$$
(2.4)
(ii) For every $`gG`$, there exists a unique section $`Y_gH^0(A_1,L_1)`$ such that $`X_g^p=F^{}Y_g`$. The family $`\{Y_g\}_{gG}`$ is a basis of $`H^0(A_1,L_1)`$.
###### Proof.
Let us construct geometrically the basis $`\{X_g\}_{gG}`$. We define $`X_g`$ by
$$X_g=g.F^{}\theta _1=g.\theta ^p.$$
By construction, one has
$$a.X_g=X_{g+a}a,gG.$$
Because $`X_0`$ comes from $`A_1=A/\widehat{G}`$, it is $`\widehat{G}`$ invariant. The relations (2.4) follow. Because $`H^0(A,L)`$ is irreducible, $`\{X_g\}_{gG}`$ is a basis. If we have another such basis $`X_g^{}`$, the endomorphism given by $`X_gX_g^{}`$ is $`G(L)`$-equivariant and therefore a scalar by Schur’s lemma, which proves (i).
The sections $`X_g^p`$ live in $`H^0(A,L^p)=H^0(A,F^{}L_1)^{\widehat{G}}`$ by (i) and (ii) follows. ∎
We will identify $`H^0(A,L)`$ and $`H^0(A_1,L_1)`$ with their duals (more precisely $`H^0(A,L)`$ and $`H^0(A,L)`$) using these bases.
###### 2.7 Remark.
More generally, we can construct Theta bases for any power $`L=𝒪(p^l\mathrm{\Theta })`$ with $`l1`$. We denote by $`G_l`$ the reduced part of $`\mathrm{ker}[p^l]`$. Note that we have $`G_l=(/p^l)^g`$, $`G_1=G`$ and, for any $`l1`$, we have an exact sequence $`0G_1G_{l+1}\stackrel{p}{}G_l0`$. As above, we prove that there exist canonical bases $`\{X_g^{(l)}\}_{gG_l}`$ of $`H^0(A,L)`$ and $`\{Y_g^{(l)}\}_{gG_l}`$ of $`H^0(A_1,L_1)`$ satisfying relations (2.4).
### 2.4 Addition formula in characteristic $`2`$
From now we assume that $`M`$ is symmetric and that $`p=2`$. We will explain how to obtain an addition formula in this context. The method essentially goes as in the proof of Lemma 1.2 of \[Sek\], which is a generalization of Mumford’s arguments.
Consider the homomorphism
$$\xi :\{\begin{array}{ccc}A\times A& & A\times A\\ (x,y)& & (xy,x+y)\end{array}$$
(2.5)
(in the notation of loc. cit., $`a=b=1`$). By the See-Saw theorem, we have an isomorphism
$$\xi ^{}(MM)M^2M^2$$
(2.6)
hence we get an injection
$$\xi ^{}:H^0(A,M)H^0(A,M)H^0(A,M^2)H^0(A,M^2)$$
Let $`\theta `$ be the unique (up to a scalar) section of $`H^0(A,M)`$ and consider the Theta basis $`\{X_g\}_{gG}`$ (2.4) of $`H^0(A,L)=H^0(A,M^2)`$. We have an exact sequence
$$1𝔾_mG(M^2M^2)A[2]\times A[2]1.$$
The kernel $`K=\mathrm{ker}\xi =A[2]`$ sits diagonally in $`A[2]\times A[2]`$ and the isomorphism (2.6) determines a lift $`K^{}G(M^2M^2)`$ of $`K`$.
The main point is that $`K=A[2]`$ has a Göpel system in classical terminology of theta functions, namely
$$KG\times \widehat{G}$$
where both $`G`$ and $`\widehat{G}`$ are isotropic. In loc. cit., this existence follows from the condition $`pa+b`$, which is the only reason to put this arithmetic condition, which of course is not fulfilled here. Via the projection map
$$G(M^2)\times G(M^2)G(M^2M^2)$$
the lift $`K^{}`$ induces lifts of $`G`$ and $`\widehat{G}`$ into $`G(M^2)`$. Because the $`\widehat{G}`$-invariant part of $`H^0(A,M^2)`$ is generated by $`X_0`$, Sekiguchi’s result gives
###### 2.8 Lemma (Sekiguchi).
Normalizing $`\xi `$ suitably, one has the formula
$$\xi ^{}(\theta \theta )=\underset{hG}{}X_hX_h$$
(2.7)
In other words, we have
$$\theta (xy)\theta (x+y)=\underset{gG}{}X_g(x)X_g(y),x,yA.$$
Let us define the Kummer morphism $`\mathrm{Kum}_A:`$ $`A|2\mathrm{\Theta }|`$, $`yT_y^{}\mathrm{\Theta }+T_y^{}\mathrm{\Theta }`$, where $`T_y`$ denotes translation by $`y`$. Then we can write
$$\mathrm{div}(\underset{g}{}X_g(y)X_g)=T_y^{}\mathrm{\Theta }+T_y^{}\mathrm{\Theta }=\mathrm{Kum}_A(y),yA.$$
Using the relation $`F^{}\mathrm{\Theta }_1=2\mathrm{\Theta }`$, one gets the analogous relation
$$\mathrm{div}(\underset{g}{}Y_g(y)Y_g)=T_y^{}\mathrm{\Theta }_1+T_y^{}\mathrm{\Theta }_1=\mathrm{Kum}_{A_1}(y),yA_1.$$
The element (2.7) induces a linear isomorphism $`H^0(A,M^2)^{}\stackrel{}{}H^0(A,M^2)`$, which allows us to identify both spaces.
###### 2.9 Corollary.
With the identification above, the complete linear system $`\phi _L`$ (resp. $`\phi _{L_1}`$) is the Kummer morphism $`\mathrm{Kum}_A`$ (resp. $`\mathrm{Kum}_{A_1}`$).
### 2.5 The Theta divisor $`\mathrm{\Theta }`$ in characteristic $`2`$
From now we assume that $`p=2`$ and that $`X`$ is an ordinary curve of genus $`g`$, whose Jacobian is denoted by $`J`$. Let $`B`$ be the theta-characteristic of $`X_1`$ defined by the exact sequence (\[R\] section 4)
$$0𝒪_{X_1}F_{}𝒪_XB0$$
(2.8)
and we denote by $`\mathrm{\Theta }_1J_1`$ the symmetric Theta divisor determined by $`B`$, i.e.,
$$\mathrm{supp}\mathrm{\Theta }_1=\{NJ_1:h^0(X_1,BN)>0\}.$$
(2.9)
We denote by $`\mathrm{\Theta }`$ the Theta divisor on $`J`$ obtained from $`\mathrm{\Theta }_1`$ (Lemma 2.2). Note that we have $`\mathrm{\Theta }_1=\iota ^{}\mathrm{\Theta }`$, where $`\iota :J_1J`$ is the $`k`$-semi-linear isomorphism. Let $`R`$ be the ring of dual numbers $`k[ϵ]`$ with $`ϵ^2=0`$. We recall that the Lie algebra $`𝔤=\widehat{G}(R)`$ identifies with the tangent space $`𝐓_0J`$. For any tangent vector $`v𝔤`$ and $`gG`$ we still denote by $`e^L(v,g)k`$ the derivative of
$$\widehat{G}𝔾_m,\overline{g}e^L(\overline{g},g).$$
Writing $`\overline{g}=1+ϵv`$, with $`v𝔤`$, we obtain $`v𝔤`$,$`g,hG`$, $`e^L(v,g+h)=e^L(v,g)+e^L(v,h)`$. We recall some well-known facts about $`\mathrm{\Theta }`$.
###### 2.10 Lemma.
* The Theta divisor $`\mathrm{\Theta }_1`$ (and $`\mathrm{\Theta }`$) passes through any non-zero $`gG`$.
* A non-zero $`g`$ is a smooth point of $`\mathrm{\Theta }_1`$ (and $`\mathrm{\Theta }`$) if and only if $`h^0(X_1,Bg)=1`$.
* Assuming (ii), the tangent space at $`g`$ $`𝐓_g\mathrm{\Theta }_1𝐓_gJ_1`$ is defined by the linear form $`e^L(.,g)`$ on $`𝔤`$. Alternatively, identifying the tangent space $`𝐓_gJ`$ with $`H^1(X,𝒪_X)`$, the tangent space $`𝐓_g\mathrm{\Theta }`$ is the image of the injective $`k`$-linear map induced by the relative Frobenius,i.e.,
$$𝐓_g\mathrm{\Theta }=\mathrm{im}F:H^1(X_1,g)H^1(X,𝒪_X).$$
###### Proof.
Part (i) follows immediately from (2.8) and (2.9). Part (ii) is a special case of Riemann’s singularity theorem (see e.g. \[Ke1\]). The differential at the origin of the separable isogeny $`V:J_1J`$ is an isomorphism $`dV:𝐓_0J_1\stackrel{}{}𝐓_0J`$, which identifies with the Hasse-Witt map $`F:H^1(X_1,𝒪_{X_1})H^1(X,𝒪_X)`$. Given $`g0`$, it will be enough to compute the tangent space to the divisor $`V^{}(T_g^{}\mathrm{\Theta })`$ at the origin. Let $`\{Y_g\}_{gG}`$ be the canonical Theta basis of $`H^0(J_1,𝒪(\mathrm{\Theta }_1+T_g^{}\mathrm{\Theta }_1))`$ (Apply Lemma 2.6(ii) to $`L_1=𝒪(\mathrm{\Theta }_1+T_g^{}\mathrm{\Theta }_1)=𝒪(2T_h^{}\mathrm{\Theta }_1)`$ with $`2h=g`$). Then, by the isogeny formula (2.11), we have (up to a scalar) $`V^{}(T_g^{}\mathrm{\Theta })=_{gG}Y_g`$. Let $`\varphi _v:\mathrm{Spec}(R)J_1`$ be a tangent vector to $`J_1`$ at the origin. Then we compute, using $`v.Y_g=e^L(v,g)Y_g`$, $`v𝔤`$
$$\varphi _v^{}(\underset{gG}{}Y_g)=ϵe^L(v,g)\underset{hG/g}{}Y_h(0)$$
Form this we deduce that $`\mathrm{\Theta }_1`$ is singular at $`g`$ if and only if $`_{hG/g}Y_h(0)=0`$ and, assuming $`g`$ smooth, the equation of $`𝐓_g\mathrm{\Theta }_1`$. The second description of $`𝐓_g\mathrm{\Theta }`$ is a consequence of \[Ke1\]. ∎
###### 2.11 Question.
Are there other principally polarized abelian varieties $`(A,\mathrm{\Theta })`$ than Jacobians which have property (i) of Lemma 2.10?
###### 2.12 Definition.
We say that $`X`$ has no vanishing theta-null if $`X`$ is ordinary and if all theta characteristics $`\kappa `$ different from $`B`$ satisfy $`h^0(X,\kappa )=1`$, or equivalently (Lemma 2.10(ii)) all non-zero $`2`$-torsion points are smooth points of $`\mathrm{\Theta }`$.
In the next section we will see (Proposition 3.1(1)) that a generic curve has no vanishing theta-null.
### 2.6 Isogeny formulae
Given an isogeny $`f:XY`$ and a line bundle $`L`$ on $`Y`$, the isogeny formula gives the linear map $`f^{}:H^0(Y,L)H^0(X,f^{}L)`$ in terms of the canonical Theta bases. Although originally proved for separable isogenies and line bundles of separable type, we can extend the isogeny formula to more general line bundles. We assume $`p=2`$ and we present (without proof) the three cases needed in this paper.
1. The separable isogeny $`V:J_1J`$ with kernel $`G`$. Let $`\{X_g\}_{gG}`$ be the Theta basis of $`H^0(J,2\mathrm{\Theta })`$ and $`\{Y_u^{(2)}\}_{uG_2}`$ be the Theta basis of $`H^0(J_1,4\mathrm{\Theta }_1)`$ (Remark 2.7). Then we have
$$gGV^{}X_g=\underset{\genfrac{}{}{0pt}{}{uG_2}{2u=g}}{}Y_u^{(2)}.$$
(2.10)
2. Isogeny $`V`$ as in 1, with $`V^{}:H^0(J,T_g^{}\mathrm{\Theta })H^0(J_1,𝒪(\mathrm{\Theta }_1+T_g^{}\mathrm{\Theta }_1))`$, for $`gG`$. Then
$$V^{}X_0=\underset{gG}{}Y_g.$$
(2.11)
3. The inseparable isogeny $`\xi :J_1\times J_1J_1\times J_1`$ defined in (2.5) with kernel $`J_1[2]=G\times \widehat{G}`$. Let $`A`$ be the quotient $`J_1\times J_1/\widehat{G}`$. Then $`\xi `$ factorizes through a separable isogeny (with kernel $`G`$) $`\overline{\xi }:AJ_1\times J_1`$ and we identify $`H^0(A,\overline{\xi }^{}(2\mathrm{\Theta }_12\mathrm{\Theta }_1))`$ with the $`\widehat{G}`$-invariant subspace of $`H^0(J_1\times J_1,4\mathrm{\Theta }_14\mathrm{\Theta }_1)`$. A canonical basis of the latter space is given by the tensors $`\{Y_u^{(2)}Y_v^{(2)}\}`$ with $`u,vG_2`$ such that $`u+vG_1G_2`$. The isogeny formula, applied to $`\overline{\xi }`$, gives
$$g,hG\xi ^{}(Y_hY_{h+g})=\underset{\genfrac{}{}{0pt}{}{u,vG_2}{\xi (u,v)=(h,h+g)}}{}Y_u^{(2)}Y_v^{(2)}.$$
(2.12)
## 3 Extending Frobenius to $`|2\mathrm{\Theta }|`$ in characteristic $`2`$
We consider a principally polarized Jacobian $`(J,\mathrm{\Theta })`$ of an ordinary curve $`X`$, with $`\mathrm{\Theta }`$ the symmetric Theta divisor defined by $`B`$ (section 2.5), and the morphism $`\phi `$ induced by the linear system $`|2\mathrm{\Theta }|`$ on $`J`$ (resp. $`\phi _1`$ on $`J_1`$), which we identify with the Kummer morphism $`\mathrm{Kum}_J`$ (Corollary 2.9).
### 3.1 Factorization
###### 3.1 Proposition.
With the notation as above, we have
1. There exists a non-empty open set of the moduli space of genus $`g`$ curves parametrizing curves $`X`$ with no vanishing theta-null.
2. Suppose $`X`$ has no vanishing theta-null. Then the isogeny $`V`$ can be “extended” to a rational map $`\stackrel{~}{V}`$ such that $`\phi V=\stackrel{~}{V}\phi _1`$,i.e., the diagram
$$\begin{array}{ccc}J_1& \stackrel{V}{}& J\\ \phi _1& & \phi & & \\ |2\mathrm{\Theta }_1|^{}& \stackrel{\stackrel{~}{V}}{}& |2\mathrm{\Theta }|^{}\end{array}$$
(3.1)
is commutative.
3. In terms of the canonical Theta bases of $`|2\mathrm{\Theta }|`$ and $`|2\mathrm{\Theta }_1|`$ the equations of $`\stackrel{~}{V}`$ are given by $`2^g`$ quadrics
$$\stackrel{~}{V}:|2\mathrm{\Theta }_1|^{}|2\mathrm{\Theta }|^{}x:=(x_0:\mathrm{}:x_g:\mathrm{})(\lambda _0P_0(x):\mathrm{}:\lambda _gP_g(x):\mathrm{})$$
where
* the constants $`\{\lambda _g\}_{gG}`$, $`\lambda _gk`$, satisfy $`gG,g0`$
$$\lambda _g=0g\text{is a singular point of}\mathrm{\Theta }$$
In particular, if $`X`$ has no vanishing theta-null, the $`\lambda _g`$’s are all non-zero.
* the polynomials $`P_g`$ are given by
$$P_g(x)=\underset{hG/g}{}x_{g+h}x_h.$$
(3.2)
###### Proof.
In order to show commutativity of the diagram (3.1) it suffices to show that the image of the injection $`V^{}:H^0(J,2\mathrm{\Theta })H^0(J_1,4\mathrm{\Theta }_1)`$ is contained in the image of the multiplication map
$$\mathrm{Sym}^2H^0(J_1,2\mathrm{\Theta }_1)H^0(J_1,4\mathrm{\Theta }_1).$$
(3.3)
Let $`\{X_g\}_{gG}`$ and $`\{Y_g\}_{gG}`$ be the canonical Theta bases of $`H^0(J,2\mathrm{\Theta })`$ and $`H^0(J_1,2\mathrm{\Theta }_1)`$ and consider the pull-back of equality (2.7), written for $`A=J_1`$, by the morphism $`\psi _g:J_1J_1\times J_1`$, $`\psi _g(x)=(x+g,x)`$,
$$\psi _g^{}\xi ^{}(\theta _1\theta _1)=\theta _1(g)V^{}X_g,\psi _g^{}(\underset{hG}{}Y_hY_h)=\underset{hG}{}Y_{h+g}Y_h.$$
(3.4)
If $`g=0`$, we get
$$\underset{hG}{}Y_h^2=\theta _1(0)V^{}X_0$$
(3.5)
with $`\theta _1(0)0`$, and define $`\lambda _0=\theta _1(0)`$. If $`g0`$, we see that both members of (3.4) are zero. In order to get a meaningful statement we restrict to a first infinitesimal neighbourhood of $`\psi _g(J_1)J_1\times J_1`$. The notation is as in section 2.5. We pull-back the morphisms $`\xi `$ and $`\psi _g`$ to $`\mathrm{Spec}(R)`$ replacing $`g`$ by $`(1+ϵv).g`$ and keeping the same notation for the object over $`k`$ and its pull-back to $`R`$. Pulling back equality (2.7) by $`\psi _g`$, for $`g0`$, we get the following two elements in $`H^0(J_1,4\mathrm{\Theta }_1)R`$
$$\psi _g^{}\xi ^{}\theta _1\theta _1=ϵe^L(v,g)\lambda _gV^{}X_g$$
where $`\lambda _g`$ is a scalar which vanishes if and only if $`\mathrm{\Theta }_1`$ is singular at the point $`g`$ (Lemma 2.10), and
$`\psi _g^{}({\displaystyle \underset{hG}{}}Y_hY_h)`$ $`=`$ $`{\displaystyle \underset{hG}{}}[(1+ϵv).Y_{g+h}]Y_h`$
$`=`$ $`{\displaystyle \underset{hG}{}}Y_{g+h}Y_h+ϵ{\displaystyle \underset{hG}{}}(v.Y_{g+h})Y_h`$
$`=`$ $`ϵ{\displaystyle \underset{hG}{}}e^L(v,g+h)Y_{g+h}Y_h`$
$`=`$ $`ϵe^L(v,g){\displaystyle \underset{hG/g}{}}Y_{g+h}Y_h`$
where $`h`$ runs over a set of representatives of $`G/g`$. Since these elements are equal $`v𝔤`$ and there exists $`v𝔤`$ such that $`e^L(v,g)0`$, we obtain (up to a multiplicative non-zero scalar)
$$\underset{hG/g}{}Y_{g+h}Y_h=\lambda _gV^{}X_g.$$
(3.6)
In order to complete the proof it suffices to show that for a general curve $`X`$ the $`\lambda _g`$’s are non-zero, for all $`g0`$. Assume that the contrary holds. Then the equality in $`\mathrm{Sym}^2H^0(J_1,2\mathrm{\Theta }_1)`$ (which we leave as an exercise)
$$\left(\underset{hG_g}{}Y_h\right)\left(\underset{hG_g}{}Y_h\right)=\underset{hG_g}{}P_h$$
where $`G_g`$ is any index $`2`$ subgroup of $`G`$ not containing $`g`$ and the $`P_h`$ are the quadrics (3.2), shows that either $`_{hG_g}Y_h`$ or $`_{hG_g}Y_h`$ is zero, since the RHS is mapped to zero in $`H^0(J_1,4\mathrm{\Theta }_1)`$. But this is impossible, since $`\{Y_h\}`$ is a basis. Hence for any curve $`X`$ there exists $`g0`$ such that $`\lambda _g0`$ and we conclude by a monodromy argument found in \[E\].
###### 3.2 Remark.
We notice that $`\stackrel{~}{V}`$ is uniquely defined only up to a degree $`2`$ equation of the image $`\phi _1(J_1)|2\mathrm{\Theta }_1|^{}`$. We will show uniqueness of $`\stackrel{~}{V}`$ (Proposition 4.1) for an ordinary genus $`2`$ curve. We expect that there are no quadrics containing $`\phi _1(J_1)`$ for a curve of genus $`g2`$ with no vanishing theta-null.
Equalities (3.5) and (3.6) can be used to define the vector $`(\lambda _g)_{gG}`$ up to a scalar. In the next proposition we give a more direct definition in terms of theta-constants. Let $`\{Y_u^{(2)}\}_{uG_2}`$ be the Theta basis of $`H^0(J_1,4\mathrm{\Theta }_1)`$ and we denote by $`Y_u^{(2)}(0)k`$ the value of $`Y_u^{(2)}`$ at the origin (after having chosen an isomorphism $`𝒪(4\mathrm{\Theta }_1)_0\stackrel{}{}k`$).
###### 3.3 Proposition.
With the notation as above, we have
$$gG\lambda _g=\underset{uS_g}{}Y_u^{(2)}(0)$$
(3.7)
where $`S_g`$ is a set of representatives of $`\{uG_2:2u=g\}/g`$.
###### Proof.
Let $`i`$ be the inclusion in the first factor $`i:J_1J_1\times J_1`$. A standard computation modelled on \[Ke2\] Proposition 6, which involves (2.12), gives the equality in $`\mathrm{Sym}^2H^0(J_1,4\mathrm{\Theta }_1)`$
$$i^{}\xi ^{}\left(\underset{hG/g}{}Y_hY_{h+g}\right)=\left(\underset{uS_g}{}Y_u^{(2)}(0)\right)\left(\underset{\genfrac{}{}{0pt}{}{uG_2}{2u=g}}{}Y_u^{(2)}\right)$$
We observe that the composite $`\xi i`$ is the diagonal map, hence the LHS can be rewritten as $`Y_hY_{h+g}`$. Applying the isogeny formula (2.10), we get $`gG`$
$$\underset{hG/g}{}Y_hY_{h+g}=\left(\underset{uS_g}{}Y_u^{(2)}(0)\right)V^{}X_g.$$
Comparing with (3.5) and (3.6) we can conclude. Finally we observe that the expression (3.7) is well-defined since $`u`$ such that $`2u=g`$ we have $`Y_u^{(2)}(0)=Y_u^{(2)}(0)=Y_{u+g}^{(2)}(0)`$ ($`Y_0^{(2)}`$ is symmetric). ∎
## 4 Kummer’s quartic surface in characteristic $`2`$
As an application of Proposition 3.1 we shall deduce the equation of the image $`\phi (J)|2\mathrm{\Theta }|^{}`$ for a Jacobian of an ordinary genus $`2`$ curve $`X`$. We observe that by Clifford’s theorem $`h^0(X_1,Bg)1`$, $`gG`$, hence $`X`$ has no vanishing theta-null and we can apply Proposition 3.1(2). Somewhat surprisingly, invariance under the Theta group $`G(L)`$ and the “Frobenius” map $`\stackrel{~}{V}`$ are sufficient to determine the equation. We fix an isomorphism $`G(/2)^2`$.
###### 4.1 Proposition.
In the canonical Theta basis $`\{X_g\}_{gG}`$ of $`H^0(J,2\mathrm{\Theta })`$ the equation of the Kummer surface is (up to a non-zero scalar)
$$\lambda _{10}^2(x_{00}^2x_{10}^2+x_{01}^2x_{11}^2)+\lambda _{01}^2(x_{00}^2x_{01}^2+x_{10}^2x_{11}^2)+\lambda _{11}^2(x_{00}^2x_{11}^2+x_{01}^2x_{10}^2)+\frac{\lambda _{10}\lambda _{01}\lambda _{11}}{\lambda _{00}}x_{00}x_{10}x_{01}x_{11}$$
(4.1)
In particular, $`\phi (J)`$ is not contained in a quadric and therefore the map $`\stackrel{~}{V}`$ is uniquely defined.
###### Proof.
First we observe that the image $`\phi (J)|2\mathrm{\Theta }|`$ is invariant under the Theta group $`G(L)`$. Hence the equation of $`\phi (J)`$ has to be $`G(L)`$-invariant. Since $`(2\mathrm{\Theta })^2=8`$ and $`\mathrm{deg}\phi =2`$ ($`\phi `$ is separable), we have $`\mathrm{deg}\phi (J)=4`$. The unique $`G(L)`$-invariant quadric is $`x_{00}^2+x_{01}^2+x_{10}^2+x_{11}^2=(x_{00}+x_{01}+x_{10}+x_{11})^2`$. Since $`\phi (J)`$ is non-degenerate, we see that it is not contained in a quadric. A straightforward computation shows that a basis of $`G(L)`$-invariant quartics is given by the $`5`$ polynomials
$$S:=x_{00}^4+x_{01}^4+x_{10}^4+x_{11}^4R:=x_{00}x_{01}x_{10}x_{11}$$
$$Q_{10}:=x_{00}^2x_{10}^2+x_{01}^2x_{11}^2Q_{01}:=x_{00}^2x_{01}^2+x_{10}^2x_{11}^2Q_{11}:=x_{00}^2x_{11}^2+x_{01}^2x_{10}^2$$
Let us denote by
$$F=\alpha _{10}Q_{10}+\alpha _{01}Q_{01}+\alpha _{11}Q_{11}+\beta R+\gamma S$$
the equation of $`\phi (J)`$, where $`\alpha _{10},\alpha _{01},\alpha _{11},\beta ,\gamma k`$ need to be determined. First $`\phi (e)=(1:0:0:0)`$, where $`eJ`$ is the origin, implies that $`\gamma =0`$. Since the image is reduced, $`\beta 0`$. Next, the equation of $`\phi _1(J_1)|2\mathrm{\Theta }_1|`$ is given in a Theta basis by
$$F_1=\alpha _{10}^2Q_{10}+\alpha _{01}^2Q_{01}+\alpha _{11}^2Q_{11}+\beta ^2R$$
Since $`\phi _1(J_1)`$ is not contained in a quadric, the rational map $`\stackrel{~}{V}`$ (3.1) is uniquely defined. Since $`\stackrel{~}{V}(\phi _1(J_1))=\phi (J)`$, there exists a $`G(L)`$-invariant quartic $`A`$, such that (up to a non-zero scalar)
$$\stackrel{~}{V}^{}(F)=F_1A$$
(4.2)
We shall use the Theta coordinates $`\{x_g\}_{gG}`$ on both spaces. In order to determine the equation of $`A`$, we restrict equality (4.2) to the hyperplane $`H:_{gG}x_g=0`$ and write the expressions as degree $`8`$ polynomials in $`x_{01},x_{10},x_{11}`$. A straightforward computation, which we omit, leads to
$`\stackrel{~}{V}^{}(F)_{|H}=\lambda _{01}^2\lambda _{10}^2\lambda _{11}^2`$ $`(x_{01}+x_{11})^2(x_{01}+x_{10})^2(x_{11}+x_{10})^2`$
$`\left[x_{01}^2({\displaystyle \frac{\alpha _{10}}{\lambda _{10}^2}}+{\displaystyle \frac{\alpha _{11}}{\lambda _{11}^2}})+x_{10}^2({\displaystyle \frac{\alpha _{01}}{\lambda _{01}^2}}+{\displaystyle \frac{\alpha _{11}}{\lambda _{11}^2}})+x_{11}^2({\displaystyle \frac{\alpha _{10}}{\lambda _{10}^2}}+{\displaystyle \frac{\alpha _{01}}{\lambda _{01}^2}})\right]`$
Suppose that $`A_{|H}0`$, i.e. $`\stackrel{~}{V}^{}(F)_{|H}0`$. Then at least one of the factors $`x_{01}+x_{10},x_{01}+x_{11},x_{11}+x_{10}`$ has to divide $`F_{1}^{}{}_{|H}{}^{}`$. Again elementary computation shows that this can only happen when $`\beta =0`$, which is impossible. Hence $`A`$ is a multiple of $`H`$ and, since $`A`$ is also $`G(L)`$-invariant, we get $`A=H^4=_{gG}x_g^4`$. Moreover $`\stackrel{~}{V}^{}(F)_{|H}=0`$ implies that there exists a $`\mu k^{}`$ such that $`\alpha _g=\mu \lambda _g^2`$ for all $`gG^{}`$.
It remains to determine the constant $`\beta `$. We compute the degree $`8`$ polynomial $`\stackrel{~}{V}^{}(F)`$ (again we omit details)
$`\mu \lambda _{00}^2H^4(\lambda _{10}^4Q_{10}+\lambda _{01}^4Q_{01}+\lambda _{11}^4Q_{11})`$
$`+\mu \lambda _{01}^2\lambda _{10}^2\lambda _{11}^2H^2(x_{00}^2x_{01}^2x_{10}^2+x_{00}^2x_{01}^2x_{11}^2+x_{00}^2x_{10}^2x_{11}^2+x_{01}^2x_{10}^2x_{11}^2)`$
$`+\beta \lambda _{00}\lambda _{01}\lambda _{10}\lambda _{11}H^2(RH^2+x_{00}^2x_{01}^2x_{10}^2+x_{00}^2x_{01}^2x_{11}^2+x_{00}^2x_{10}^2x_{11}^2+x_{01}^2x_{10}^2x_{11}^2)`$
This expression being equal to $`F_1H^4`$, we get the equation mentioned in the proposition. ∎
###### 4.2 Remark.
If the characteristic is different from $`2`$, Kummer’s quartic surface is a much studied object (see e.g.\[H\]). Among many other results, let us just mention that the coefficients of the quartic equation (in a Theta basis) satisfy a cubic equation and are themselves polynomials of degree $`12`$ in the theta-constants. Since we could not find any treatment of the characteristic $`2`$ case in the literature, we decided to include it in our paper.
## 5 The moduli space $`\mathrm{M}_X`$ of rank $`2`$ vector bundles
We assume $`p=2`$. Let $`\mathrm{M}_X`$ denote the moduli space of rank $`2`$ semi-stable vector bundles over $`X`$ with trivial determinant. As was observed in \[MR\], the theta divisor
$$\stackrel{~}{\mathrm{\Theta }}=\{[E]\mathrm{M}_X|h^0(EB)0\}$$
is Cartier (and not only $``$-Cartier) and is ample by GIT. We denote by $`_0`$ the line bundle $`𝒪_{\mathrm{M}_X}(\stackrel{~}{\mathrm{\Theta }})`$. By \[R\] there exists a regular morphism
$$D:\mathrm{M}_XH^0(J,2\mathrm{\Theta })=|2\mathrm{\Theta }|$$
which maps the class of the semi-stable bundle $`E`$ to the divisor $`D(E)`$
$$\mathrm{supp}D(E)=\{LJ|H^0(EBL)0\}.$$
As in the complex case \[NR\] one has
###### 5.1 Proposition (V. Balaji).
If $`g=2`$, the morphism $`D:\mathrm{M}_X^3=|2\mathrm{\Theta }|`$ is an isomorphism.
###### Proof.
(sketch) We proceed in two steps. First we consider a flat family of curves $`𝒳T=\mathrm{Spec}(A)`$, where $`A`$ is a discrete valuation ring such that its residue field at the closed point $`0`$ is $`k`$, its field of fraction $`K`$ is of characteristic zero, and $`𝒳_0=X`$. We consider the moduli scheme $`T`$ of semi-stable rank $`2`$ vector bundles of trivial determinant over the family $`𝒳T`$. Let $`_0`$ be the fibre of $`T`$ over $`0`$. Then, by GIT over arbitrary base \[Ses\], we have a canonical bijective morphism $`i:\mathrm{M}_X_0`$. Moreover on the open set of stable points $`i`$ is an isomorphism since the action of the projective group on the Quot scheme is free. We conclude that $`i:\mathrm{M}_X_0`$ is an isomorphism by Zariski’s Main Theorem.
Secondly, we extend the morphism $`D`$ \[NR\] to the family $`T`$,i.e., we construct a morphism over the base $`T`$
$$𝒟:_T^3$$
such that $`𝒟_0=D`$. In order to show the proposition, it will be enough (again by Zariski’s Main Theorem) to show that $`𝒟`$ is birational and bijective. We consider the fibre $`𝒟_\xi :_\xi _\xi ^3`$ over the generic point $`\xi T`$. Working over an algebraic closure of $`K`$ (of characteristic $`0`$), we see \[NR\] that $`D_\xi `$ is an isomorphism. Hence $`𝒟`$ is birational. It remains to show that $`𝒟_0=D`$ is bijective. Surjectivity is obvious. Let $`H`$ be the hyperplane $`|2\mathrm{\Theta }|`$ of divisors passing through $`𝒪`$. The inverse image of $`H`$ is $`\mathrm{\Theta }`$ showing $`D^{}𝒪(1)_0`$. It follows that $`D`$ is finite because $`D^{}𝒪(1)`$ is ample. Let $`=D^{}𝒪_{_T^3}(1)`$ be the relatively ample line bundle over $`T`$. Consider the canonical inclusion of $`𝒪_{_T^3}`$ in $`𝒟_{}𝒪_{}`$ with cokernel $`𝒬`$.
$$0𝒪_{_T^3}𝒟_{}𝒪_{}𝒬0$$
We twist by $`𝒪_{_T^3}(n)`$. If $`n`$ is large enough, we have $`tT`$ and $`i>0`$, $`h^i(_t,_t^n)=0`$. Hence, since $`_T^3T`$ is flat, we see that
$$h^0(_t^3,𝒪_{_t^3}(n))h^0(_t,_t^n)=\chi (_t^3,𝒪_{_t^3}(n))\chi (_t,_t^n)$$
is constant. If $`t=\xi `$, this number is zero \[NR\]. Hence $`h^0(_0^3,𝒬_0(n))=0`$ and $`Q_0=0`$. So $`𝒪_{_0^3}=D_{}𝒪_{\mathrm{M}_X}`$ and $`D`$ is injective.
###### 5.2 Remark.
As we were told by C.S. Seshadri, the first part of the proof is completely worked out in the PhD thesis of Venkata Balaji (in preparation). A direct proof of this isomorphism along the lines of the original paper \[NR\] was obtained by M.S. Narasimhan.
## 6 Frobenius action on $`\mathrm{M}_X`$ for an ordinary genus $`2`$ curve
The goal of this section is to describe the Frobenius map (more precisely, its separable part, the Verschiebung)
$$V:\mathrm{M}_{X_1}\mathrm{M}_X,EF^{}E.$$
We consider the morphism $`\psi :J\mathrm{M}_X`$, $`L[LL^1]`$ which, when composed with $`D`$, equals the Kummer morphism $`\mathrm{Kum}_J`$. Because of Proposition 5.1 and Corollary 2.9 we can identify $`\psi `$ and $`\phi `$. Since $`H^0(\mathrm{M}_X,_0)=H^0(J,2\mathrm{\Theta })`$ and since $`\stackrel{~}{V}`$ is uniquely defined (Proposition 4.1), we can identify (via $`D`$) the Verschiebung $`V:\mathrm{M}_{X_1}\mathrm{M}_X`$ with the rational map $`\stackrel{~}{V}:|2\mathrm{\Theta }_1|^{}|2\mathrm{\Theta }|^{}`$ given by the equations (3.2). We gather our results in the following proposition.
###### 6.1 Proposition.
Let $`X`$ be an ordinary genus $`2`$ curve.
1. The semi-stable boundary of $`\mathrm{M}_X`$ (resp. $`\mathrm{M}_{X_1}`$) is isomorphic to Kummer’s quartic surface $`\phi (J)`$ (resp. $`\phi _1(J_1)`$) whose equation is given in (4.1). In particular, $`V`$ maps $`\phi _1(J_1)`$ onto $`\phi (J)`$.
2. There exists a unique stable bundle $`E_{BAD}\mathrm{M}_X`$, which is destabilized by the Frobenius map, i.e. $`F^{}E_{BAD}`$ is not semi-stable. We have $`E_{BAD}=F_{}B^1`$ and its projective coordinates are $`(1:1:1:1)`$.
3. The set of bundles $`\{[E]\mathrm{M}_{X_1}|\mathrm{Hom}(E,F_{}B)0\}`$ is the hyperplane $`H_1:_{gG}x_g=0`$. In particular, $`E_{BAD}H_1`$. The restriction of $`V`$ to $`H_1`$ contracts $`H_1`$ to the conic $`\phi (J)H`$, where $`H`$ is the hyperplane $`x_0=0`$ in $`|2\mathrm{\Theta }|`$. In particular, any stable bundle $`EH_1`$ is mapped into the semi-stable boundary of $`\mathrm{M}_X`$.
4. The fiber of $`V`$ over a point $`[E]\mathrm{M}_X`$ is
* a non-degenerate $`G`$-orbit of a point $`[E_1]\mathrm{M}_{X_1}`$ ($`4`$ distinct points), if $`[E]H`$
* empty, if $`[E]H(H\phi (J))`$
* a projective line passing through $`E_{BAD}`$, if $`[E]H\phi (J)`$
In particular, $`V`$ is not surjective and the separable degree of $`V`$ is $`4`$.
###### Proof.
1. This follows immediately from \[NR\] and Proposition 4.1.
2. The base locus of $`\stackrel{~}{V}`$ is given by the intersection $`_{gG}\{P_g=0\}`$, which turns out (after some elementary computations) to be a unique point with projective coordinates $`(1:1:1:1)`$. In terms of vector bundles this point, denoted $`E_{BAD}`$, corresponds to the direct image $`F_{}B^1`$. Indeed, $`F_{}B^1`$ is stable: a nonzero map $`LF_{}B^1`$ is equivalent, by adjunction, to a non-zero map $`F^{}LB^1`$, hence $`2\mathrm{d}\mathrm{e}\mathrm{g}L1`$. Moreover, since we have a canonical nonzero map $`F^{}F_{}B^1B^1`$, this bundle is destabilized. Uniqueness can also be proved without the use of the equations: consider any stable bundle $`E\mathrm{M}_X`$, which is destabilized by $`V`$. By \[LS\] Corollary 2.6 there exists a theta-characteristic $`A`$ which appears as a quotient $`F^{}EA^1`$. By adjunction this quotient map induces a map $`EF_{}A^1`$. Since the two bundles are stable with the same slope, we deduce that they have to be isomorphic. The determinant of $`F_{}A^1`$ being trivial, we get $`A=B`$ and we are done.
3. First we observe that the hyperplane $`H_1`$ is mapped into the hyperplane $`H`$. After fixing an isomorphism $`G(/2)^2`$, we straightforwardly compute that a point in the image $`V(H_1)`$ satisfies the equations
$$\lambda _{10}x_{01}x_{11}+\lambda _{01}x_{10}x_{11}+\lambda _{11}x_{01}x_{10}=0$$
(6.1)
which is precisely (after squaring) the equation of Kummer’s quartic surface (4.1) restricted to the hyperplane $`H`$. Moreover $`V^{}(H)=H_1^2`$ and by adjunction $`\mathrm{Hom}(E,F_{}B)=\mathrm{Hom}(F^{}E,B)`$ which proves the first assertion.
4. Given a point $`[E]\mathrm{M}_X`$ with projective coordinates $`(a_{00}:a_{01}:a_{10}:a_{11})`$, we have to solve the system of quadratic equations
$$P_g=\frac{a_g}{\lambda _g}gG(/2)^2$$
(6.2)
where the quadrics $`P_g`$’s are defined in (3.2). We write $`b_g=\frac{a_g}{\lambda _g}`$. Adding any two $`P_g`$’s with $`g0`$, we find
$$P_{01}+P_{10}=x_{00}x_{01}+x_{10}x_{11}+x_{00}x_{10}+x_{01}x_{11}=(x_{00}+x_{11})(x_{01}+x_{10})=b_{01}+b_{10}$$
(6.3)
and similarly
$$(x_{00}+x_{01})(x_{11}+x_{10})=b_{11}+b_{10}(x_{00}+x_{10})(x_{01}+x_{11})=b_{01}+b_{11}$$
(6.4)
and
$$x_{00}+x_{01}+x_{10}+x_{11}=c$$
(6.5)
with $`c^2=b_{00}`$. We let $`\alpha =x_{00}+x_{11},\beta =x_{01}+x_{10},\gamma =x_{00}+x_{01},\delta =x_{11}+x_{10}`$. Then the equations (6.3), (6.4), (6.5) imply that $`\alpha ,\beta `$ (resp. $`\gamma ,\delta `$) are the roots of the polynomial
$$t^2+ct+(b_{01}+b_{10})\text{resp.}t^2+ct+(b_{11}+b_{10})$$
(6.6)
Substituing $`x_{01},x_{11},x_{01}`$ with $`\gamma +x_{00},\alpha +x_{00}`$ and $`\beta +\gamma +x_{00}`$ respectively in the equation $`P_{01}=b_{01}`$, we find
$$cx_{00}=b_{10}+\gamma \alpha $$
(6.7)
Assuming $`c0`$,i.e. $`[E]H`$, we see that the fiber $`V^1([E])`$ consists of the point $`[E_1]\mathrm{M}_{X_1}`$ with projective coordinates $`(x_{00}:x_{01}:x_{10}:x_{11})`$ given by
$$x_{00}=\frac{1}{c}(b_{10}+\gamma \alpha )x_{01}=\frac{1}{c}(b_{10}+\gamma \beta )x_{10}=\frac{1}{c}(b_{10}+\beta \delta )x_{11}=\frac{1}{c}(b_{10}+\delta \alpha )$$
plus the other $`3`$ points obtained by switching $`\alpha `$ and $`\beta `$ as well as $`\gamma `$ and $`\delta `$. Since $`c0`$, one easily sees that these $`4`$ points are distinct.
Assume $`c=0`$, i.e. $`[E]H`$. We get $`\alpha =\beta ,\gamma =\delta ,\alpha ^2=b_{01}+b_{10},\gamma ^2=b_{11}+b_{10}`$. Hence $`(b_{10}+\gamma \alpha )^2=b_{10}b_{11}+b_{01}b_{10}+b_{10}b_{11}`$. Assume that $`b_{10}b_{11}+b_{01}b_{10}+b_{10}b_{11}0`$,i.e. $`[E]H(H\phi (J))`$. Because of (6.7) the system (6.2) does not have a solution. If $`b_{10}b_{11}+b_{01}b_{10}+b_{10}b_{11}=0`$, i.e. $`[E]H\phi (J)`$, it is easy to check that $`V^1([E])`$ is the $`^1`$, intersection of the two hyperplanes
$$H_1:x_{00}+x_{10}+x_{01}+x_{11}=0H_{\alpha ,\gamma }:(\alpha +\gamma )x_{00}+\alpha x_{01}+\gamma x_{11}=0$$
###### 6.2 Remark.
Let $`L_1JX_1`$ and $`M=F^{}L_1JX`$. We assume that $`M^2𝒪_X`$. Then obviously $`F^{}[L_1L_1^1]=[MM^1]`$ and, as one easily checks, if $`h^0(MB)=0`$, the three isomorphism classes contained in $`[L_1L_1^1]`$, i.e. the decomposable bundle $`L_1L_1^1`$, the two non-split extensions of $`L_1`$ by $`L_1^1`$ and of $`L_1^1`$ by $`L_1`$ (which are interchanged by the hyperelliptic involution $`i`$), are mapped to the corresponding isomorphism class in $`[MM^1]`$. On the other hand, if $`h^0(MB)>0`$, all three isomorphism classes are mapped to $`MM^1`$. Moreover, in that case, by Proposition 6.1 (4), there exist stable bundles $`E`$ such that $`[F^{}E]=[MM^1]`$. Since $`i`$ commutes with $`F`$ and since any stable bundle on $`X_1`$ is $`i`$-invariant, we have $`F^{}EMM^1`$. This shows that the non-split extension of $`M`$ by $`M^1`$ is not of the form $`F^{}E`$ for some semi-stable bundle $`E`$.
###### 6.3 Remark.
It would be interesting to have a coordinate-free proof of the following fact, which follows from Proposition 6.1(3): If $`E_1`$ is stable over $`X_1`$ with $`\mathrm{Hom}(E_1,F_{}B)0`$, then $`F^{}E_1`$ is unstable.
Let $`\mathrm{N}_X`$ (resp. $`\mathrm{N}_{X_1}`$) denote the moduli space of semi-stable rank $`2`$ vector bundles of degree $`0`$ over $`X`$ (resp. $`X_1`$). As a corollary of Proposition 6.1 we have
###### 6.4 Proposition.
For any ordinary curve $`X`$, the rational map $`V:\mathrm{N}_{X_1}\mathrm{N}_X`$ given by $`EF^{}E`$ is surjective.
###### Proof.
It will be enough to show that any point $`[E]\mathrm{M}_X\mathrm{N}_X`$ is in the image, since, twisting by a degree zero line bundle, we can always assume the determinant to be trivial. For any non-zero $`gG`$, we choose a $`hG_2`$ such that $`2h=g`$ and we denote by $`\mathrm{M}_{X_1}(g)`$ the moduli of rank $`2`$ vector bundles with fixed determinant equal to $`g`$. Then we have a commutative diagram
$$\begin{array}{ccc}\mathrm{M}_{X_1}(g)& \stackrel{V}{}& \mathrm{M}_X\\ T_h& & T_g& & \\ \mathrm{M}_{X_1}& \stackrel{V}{}& \mathrm{M}_X\end{array}$$
Now since the vertical maps are isomorphisms, the image of the first horizontal map contains the complement to the hyperplane $`x_g0`$. Since the hyperplanes $`\{x_g=0\}_{gG}`$ have no common point in $`\mathrm{M}_X`$ we get surjectivity. ∎
## 7 The action on the odd degree moduli space $`\mathrm{M}_{X_1}(\mathrm{\Delta })`$
In this section we briefly discuss the action of the Frobenius map on the moduli space $`\mathrm{M}_{X_1}(\mathrm{\Delta })`$ of semi-stable rank $`2`$ vector bundles over $`X_1`$ with fixed determinant equal to $`\mathrm{\Delta }`$, with $`\mathrm{deg}\mathrm{\Delta }=1`$. We will study the rational map
$$V:\mathrm{M}_{X_1}(\mathrm{\Delta })\mathrm{M}_X,EF^{}E\mathrm{\Delta }^1$$
Note that we use the same letter ($`\mathrm{\Delta },B,\mathrm{}`$) for the line bundles over $`X_1`$ and $`X`$, which correspond under the $`k`$-semi-linear isomorphism $`\iota :X_1X`$.
The next proposition holds for any curve of genus $`g2`$.
###### 7.1 Proposition.
The image of $`V`$ is contained in the Theta divisor $`\mathrm{\Theta }\mathrm{M}_X`$, i.e.
$$E\mathrm{M}_{X_1}(\mathrm{\Delta })h^0(X,F^{}E\mathrm{\Delta }^1B)>0$$
In particular, $`V`$ is not dominant.
###### Proof.
We write $`E`$ as an extension of line bundles
$$0LE\mathrm{\Delta }L^10(ϵ)$$
for some line bundle $`L`$ with $`\mathrm{deg}L0`$ and $`ϵH^1(X_1,L^2\mathrm{\Delta }^1)`$. The extension class $`F^{}ϵ`$ of the exact sequence, gotten by pull-back under $`F`$, i.e.
$$0L^2F^{}E\mathrm{\Delta }^2L^20(F^{}ϵ)$$
(7.1)
is obtained from $`ϵ`$ via the linear map
$$H^1(X_1,L^2\mathrm{\Delta }^1)\stackrel{F^{}}{}H^1(X,L^4\mathrm{\Delta }^2)=H^1(X_1,L^2\mathrm{\Delta }^1F_{}𝒪_X)$$
(7.2)
The last map coincides with the induced map on cohomology of the canonical exact sequence
$$0L^2\mathrm{\Delta }^1L^2\mathrm{\Delta }^1F_{}𝒪_XL^2\mathrm{\Delta }^1B0$$
(7.3)
In order to prove that $`h^0(X,F^{}E\mathrm{\Delta }^1B)>0`$, we tensorize (7.1) with $`\mathrm{\Delta }^1B`$. If $`h^0(X,L^2\mathrm{\Delta }^1B)>0`$, we are done. Therefore we assume $`h^0(X,L^2\mathrm{\Delta }^1B)=0`$. It follows from (7.2) and (7.3) that
$$0H^1(X_1,L^2\mathrm{\Delta }^1)\stackrel{F^{}}{}H^1(X,L^4\mathrm{\Delta }^2)H^1(X_1,L^2\mathrm{\Delta }^1B)0$$
(7.4)
On the other hand, we see that $`h^0(X,F^{}E\mathrm{\Delta }^1B)>0`$ if and only if the symmetric coboundary map
$$H^0(X,\mathrm{\Delta }L^2B)\stackrel{F^{}ϵ}{}H^1(X,\mathrm{\Delta }^1L^2B)=H^0(X,\mathrm{\Delta }L^2B)^{}$$
is degenerate. We write $`V:=H^0(X,\mathrm{\Delta }L^2B)`$. It is well-known that the linear map
$$H^1(X,L^4\mathrm{\Delta }^2)=H^0(X,\mathrm{\Omega }_XL^4\mathrm{\Delta })^{}\stackrel{m^{}}{}V^{}V^{},\delta \delta $$
(7.5)
is the dual of the multiplication map of global sections $`\mathrm{Sym}^2V\stackrel{m}{}H^0(X,\mathrm{\Omega }_XL^4\mathrm{\Delta })`$. Let us denote by $`𝒟_2(V^{})`$ the space of divided powers of $`V^{}V^{}`$,i.e. the subspace of tensors invariant under the involution $`\varphi \psi \psi \varphi `$. We have $`𝒟_2(V^{})=(\mathrm{Sym}^2V)^{}`$. We denote by $`F(V)`$ the subspace of $`\mathrm{Sym}^2V`$ generated by the squares $`v^2`$ with $`vV`$. Dually, the kernel of the surjection $`K:=\mathrm{ker}(𝒟_2(V^{})F(V)^{})`$ coincides with the space of alternating maps $`VV^{}`$. The main point, which will be used later, is that, since, by Riemann-Roch, $`\mathrm{dim}V=\mathrm{dim}H^1(X,\mathrm{\Delta }^1L^2B)=12\mathrm{d}\mathrm{e}\mathrm{g}L`$ is odd, any map in $`K`$ is degenerate.
By Serre duality we have $`H^1(X_1,L^2\mathrm{\Delta }^1B)=H^0(X_1,\mathrm{\Delta }L^2B)^{}`$. Put $`V_1:=H^0(X_1,\mathrm{\Delta }L^2B)`$. Then we have a $`k`$-linear isomorphism $`\phi :V_1\stackrel{}{}F(V)`$ defined as follows: we have $`V_1=\iota ^{}V=V_kk`$ and we put $`\phi (vt)=tv^2F(V)`$. We also note that $`F(V)^{}=F(V^{})`$. Again by Serre duality, we observe that the linear maps (7.4) and (7.5) coincide, i.e. we have a commutative diagram
$$\begin{array}{ccccccccc}H^1(X_1,L^2\mathrm{\Delta }^1)& \stackrel{F^{}}{}& H^1(X,L^4\mathrm{\Delta }^2)& & & & H^1(X_1,L^2\mathrm{\Delta }^1B)& & 0\\ & & & & & & & & \\ & & H^0(X,\mathrm{\Omega }_XL^4\mathrm{\Delta })^{}& \stackrel{m^{}}{}& 𝒟_2(V^{})& & F(V)^{}& & \end{array}$$
where the vertical maps are $`k`$-linear isomorphisms. Now we can conclude as follows: by commutativity, any extension class $`F^{}ϵ`$, with $`ϵH^1(X_1,L^2\mathrm{\Delta }^1)`$, is mapped by $`m^{}`$ into $`K`$. Hence the corresponding coboundary map is degenerate. ∎
We come back to an ordinary curve $`X`$ of genus $`2`$. Let us denote by $`\mathrm{Gr}\mathrm{\Lambda }^2H^0(J_1,2\mathrm{\Theta }_1):=^5`$ the Grassmannian of projective lines in $`H^0(J_1,2\mathrm{\Theta }_1)=^3=\mathrm{M}_{X_1}`$. Following \[B\] section 3.4 we consider, for a general point $`qX`$, the morphism
$$\mathrm{M}_{X_1}(𝒪(q))\stackrel{𝕃}{}\mathrm{Gr},E𝕃(E)=\{E^{}\mathrm{M}_{X_1}:E^{}E\}.$$
We choose an isomorphism $`G(/2)^2`$ and consider the Plücker coordinates on $`\mathrm{\Lambda }^2H^0(J_1,2\mathrm{\Theta }_1)`$
$`z_1=x_{00}x_{01},z_2=x_{00}x_{10},z_3=x_{00}x_{11},`$
$`z_4=x_{10}x_{11},z_5=x_{01}x_{11},z_6=x_{01}x_{10}.`$
The equation of the Grassmannian $`\mathrm{Gr}`$ is $`z_1z_4+z_2z_5+z_3z_6=0`$. and the $`G`$-invariant subspace $`Z\mathrm{\Lambda }^2H^0(J_1,2\mathrm{\Theta }_1)`$ is given by the $`3`$ linear equations
$$Z:z_1+z_4=z_2+z_5=z_3+z_6=0$$
###### 7.2 Proposition.
For a general point $`q`$, we have a commutative diagram
$$\begin{array}{ccc}\mathrm{M}_{X_1}(𝒪(q))& \stackrel{V}{}& \mathrm{M}_X\\ 𝕃& & & & \\ \mathrm{\Lambda }^2H^0(J_1,2\mathrm{\Theta }_1)& \stackrel{\stackrel{~}{V}}{}& H^0(J,2\mathrm{\Theta })=^3\end{array}$$
where $`\stackrel{~}{V}`$ is the projection with center $`Z=^2`$. In terms of canonical coordinates on both spaces, we have
$$\stackrel{~}{V}^{}(x_{00})=0,\stackrel{~}{V}^{}(x_{10})=z_2+z_5,\stackrel{~}{V}^{}(x_{11})=z_3+z_6,\stackrel{~}{V}^{}(x_{01})=z_1+z_4.$$
There exists a unique bundle in $`\mathrm{M}_{X_1}(𝒪(q))`$ which is destabilized by $`F`$, namely $`F_{}(B^1(q))`$.
###### Proof.
Since the proof is in the same spirit as the proof of Proposition 6.1, we just give a sketch. Let $`D_q`$ be the vector field on $`J_1`$ associated to $`q`$. We also denote by $`D_q`$ the endomorphism of $`H^0(J_1,4\mathrm{\Theta }_1)`$ obtained via the canonical $`p`$-integrable connection $``$ (Remark 2.4). We observe that we have a commutative diagram
$$\begin{array}{ccc}\mathrm{Sym}^2H^0(J_1,2\mathrm{\Theta }_1)& & \mathrm{\Lambda }^2H^0(J_1,2\mathrm{\Theta }_1)\\ m& & W_{D_q}& & \\ H^0(J_1,4\mathrm{\Theta }_1)& \stackrel{D_q}{}& H^0(J_1,4\mathrm{\Theta }_1)\end{array}$$
(7.6)
The first horizontal map is the canonical projection and $`W_{D_q}`$ is the Wahl map associated to $`D_q`$ (\[B\] section A.10). We consider the map $`J_1\mathrm{M}_{X_1}(𝒪(q))`$ defined in \[B\] section 3 and compose with $`𝕃`$. For $`q`$ general, the composite is non-degenerate and the induced (injective) map on global sections coincides with $`W_{D_q}`$ \[B\]. Using (7.6) we can now deduce the equations of $`\stackrel{~}{V}`$. The last assertion can be proved as in Proposition 6.1 (2). ∎
## 8 Frobenius dynamics
Let $`F_a`$ be the absolute Frobenius map of $`X`$. We write $`F_a^{(n)}`$ for the $`n`$-fold composite $`F_a\mathrm{}F_a`$. We will study the set of Frobenius semi-stable bundles, i.e.
$$\mathrm{\Omega }^{Frob}:=\{[E]\mathrm{M}_X|F_a^{(n)}E\text{semi-stable}n1\}$$
and the set of bundles coming from representations of the algebraic fundamental group of $`X`$,i.e. (see \[LS\] Satz 1.4)
$$\mathrm{\Omega }^{Rep}:=\{[E]\mathrm{M}_X|n>0F_a^{(n)}E\stackrel{}{}E\}$$
We obviously have $`\mathrm{\Omega }^{Rep}\mathrm{\Omega }^{Frob}`$ and we show
###### 8.1 Proposition.
The set $`\mathrm{\Omega }^{Frob}`$ is Zariski dense in $`\mathrm{M}_X=^3`$.
###### Proof.
Since $`F_a=\iota F`$, the action of $`F_a`$ on $`\mathrm{M}_X=^3`$ factorizes as
$$F_a^{}:\mathrm{M}_X\stackrel{i^{}}{}\mathrm{M}_{X_1}\stackrel{V}{}\mathrm{M}_X$$
Since $`i^{}(x_{00}:x_{01}:x_{10}:x_{11})=(x_{00}^2:x_{01}^2:x_{10}^2:x_{11}^2)`$, we see that in terms of the canonical Theta coordinates
$$F_a^{}(x)=(\mathrm{},\lambda _gP_g^2(x),\mathrm{}).$$
Let $`k_0`$ be the subfield of $`k`$ generated by the constants $`(\lambda _g)_{gG}`$ and $`k_0^{(n)}`$, for $`n1`$, be the finite field extension of $`k_0`$ generated by the coordinates $`(x_g)_{gG}`$ such that $`F_a^{(n)}(x)=(1:1:1:1)`$ with $`x=(x_g)_{gG}`$. We obviously have a tower of extensions
$$k_0=k_0^{(1)}k_0^{(2)}\mathrm{}k_0^{(n1)}k_0^{(n)}\mathrm{}$$
and, by the computations carried out in the proof of Proposition 6.1(4), we see that $`\mathrm{deg}[k_0^{(n)}:k_0^{(n1)}]`$ is a power of $`2`$. Hence, by induction, any element $`x^3\mathrm{\Omega }^{Frob}`$ has coodinates $`(x_g)_{gG}`$, which lie in an extension of $`k_0`$ of degree $`2^m`$ for some $`m`$. But elements of odd degree over $`k_0`$ are evidently dense in $`^3`$. ∎
###### 8.2 Question.
Is $`\mathrm{\Omega }^{Rep}`$ Zariski dense?
## 9 list of questions
1. For higher genus curves we no longer have a simple description of $`\mathrm{M}_X`$ as for $`g=2`$ (Proposition 5.1). But we can ask whether the diagram
$$\begin{array}{ccc}\mathrm{M}_{X_1}& \stackrel{V}{}& \mathrm{M}_X\\ D& & D& & \\ |2\mathrm{\Theta }_1|& \stackrel{\stackrel{~}{V}}{}& |2\mathrm{\Theta }|\end{array}$$
(9.1)
where $`\stackrel{~}{V}`$ is defined as in Proposition 3.1, is commutative for $`X`$ general. Note that the “restriction” of $`V`$ to $`J_1`$ commutes. Do we have $`\mathrm{dim}H^0(\mathrm{M}_X,^2)=2^{g1}(2^g+1)`$? Indeed, we could check that the latter equality implies commutativity of (9.1).
2. A straightforward computation shows that, for any genus $`g`$, the map $`\stackrel{~}{V}`$ defined by the $`2^g`$ quadrics (3.2) surjects on the complement of the hyperplane $`H:x_0=0`$. A priori this is not sufficient to deduce that the map $`V:\mathrm{M}_{X_1}\mathrm{M}_X`$ surjects on the complement of the divisor $`\stackrel{~}{\mathrm{\Theta }}=D^{}(H)\mathrm{M}_X`$ as for $`g=2`$. We optimistically conjecture
###### 9.1 Conjecture.
For any semi-stable bundle $`E\mathrm{M}_X`$ satisfying $`h^0(X,EB)=0`$, there exists a semi-stable bundle $`E_1\mathrm{M}_{X_1}`$ such that $`F^{}E_1=E`$.
As in the proof of Proposition 6.4, this conjecture implies that $`V:\mathrm{N}_{X_1}\mathrm{N}_X`$ is surjective.
3. What happens for non-ordinary curves?
We plan to return to these questions in a future work.
Yves Laszlo
Université Paris-Sud
Mathématiques Bâtiment 425
91405 Orsay Cedex France
e-mail: Yves.Laszlo@math.u-psud.fr
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# A cluster algorithm for Potts models with fixed spin densities
## I Introduction
Before 1987, the Potts model was almost exclusively simulated by means of the Metropolis algorithm , in which single-spin updates are proposed and either accepted or rejected depending on the change in energy. This algorithm works quite satisfactorily, except close to the critical point. At the critical temperature, the correlation times increase with system size as $`L^z`$, with a critical dynamic exponent equal to or slightly above two: for the two-states Potts model (Ising model), the critical dynamical exponent is reported to be $`z=2.167\pm 0.001`$ in two, and $`z=2.02\pm 0.02`$ in three dimensions . The introduction of cluster algorithms has greatly advanced the accuracy with which critical properties of the Potts model and many other models in statistical physics can be studied. The first widely used cluster algorithm was introduced by Swendsen and Wang ; we will describe their algorithm in section II A. The dynamic exponent of the Swendsen-Wang algorithm in the two- and three-dimensional Ising model is reported to be $`z=0.25\pm 0.01`$ and $`z=0.54\pm 0.02`$, respectively: cluster algorithms are able to significantly reduce critical slowing down.
To study multi-component lattice gases in the co-existence regime, for instance to study interfaces or equilibrium crystal shapes, one has to fix the number of particles in the lattice gas for each component, i.e., the spin density for each state. One typically resorts to spin-exchange dynamics, with the unfortunate consequence of a critical slowing down at least as severe as experienced with the Metropolis algorithm applied to the regular Potts model. The usual cluster algorithms do not conserve the spin densities. For the conserved-order-parameter Ising model, Heringa and Blöte recently introduced a cluster algorithm, which is in spirit related to the Wolff algorithm . It is reported to have hardly any critical slowing down, with a dynamical exponent of $`z=0.21`$. This algorithm has not been generalized to Potts models with more than two states.
In this paper, we present a modification of the Swendsen-Wang algorithm, to conserve the spin densities. In the coming section, we describe their algorithm, and introduce our modified density-conserving cluster algorithm. In the next section, we present measurements of the critical dynamic exponent for our algorithm, and show its efficiency. The paper is concluded with a summary and conclusions, and a discussion of future work.
## II Cluster algorithms
### A The Swendsen-Wang algorithm
The Swendsen-Wang algorithm is designed to simulate the Potts model, defined by the Hamiltonian
$$H=J\underset{i,j}{}\delta (\sigma _i,\sigma _j),$$
(1)
in which $`J`$ is the coupling constant, $`\delta `$ denotes the Kronecker delta function, and the summation runs over all pairs of nearest-neighbor sites, each having a spin with value ($`\sigma =1\mathrm{}Q`$). We use the usual symbols $`N`$ for the number of lattice sites, $`L`$ for the lateral dimension of the lattice with periodic boundary conditions, and $`\rho _i=(1/N)_k\delta (\sigma _i,k)`$ for the density of spins with value $`i`$.
In this algorithm, the entire lattice is divided into clusters of aligned spins, to each of which a random new value is assigned. In detail, one step of the algorithm proceeds as follows:
1. Visit all nearest-neighbor pairs of lattice sites; do nothing if the two spins are not aligned, but if they are, activate the bond between those two sites with a probability $`P_c=1\mathrm{exp}(\beta J)`$, where $`\beta `$ is the inverse temperature.
2. Group lattice sites that are connected by such activated bonds into clusters.
3. Select a random new spin value for each cluster, and assign this spin value to each of the sites constituting the cluster.
Steps 1, 2 and 3 are to be repeated many times, to obtain a set of sample configurations.
The proof of correctness for our density-conserving cluster algorithm is based on that for the Swendsen-Wang algorithm, which is presented in the remainder of this section. First we show detailed balance, next we discuss ergodicity.
Suppose we denote the spin configuration before and after the move by $`C_a`$ and $`C_b`$, respectively, with total energies $`E_a`$ and $`E_b`$, and the intermediate cluster configuration $`C_m`$ (also known as Fortuin-Kasteleyn representation ); furthermore, we write the probability to move from a configuration $`X`$ to configuration $`Y`$ as $`T(XY)`$. Then, the probability to move from a spin configuration $`C_a`$ to a cluster configuration $`C_m`$ is a product with factors $`P_c`$ over all nearest-neighbor pairs of spins that are connected, times a product with factors $`1P_c`$ over all aligned nearest-neighbor pairs of spins that are disconnected:
$$T(C_aC_m)=\underset{\begin{array}{ccc}i,j& & \\ \sigma _i^{(a)}=\sigma _j^{(a)}& & \\ i\text{}j\text{ conn.}& & \end{array}}{}\left(P_c\right)\underset{\begin{array}{ccc}i,j& & \\ \sigma _i^{(a)}=\sigma _j^{(a)}& & \\ i\text{}j\text{ disconn.}& & \end{array}}{}\left(1P_c\right)$$
(2)
and a similar expression for $`T(C_bC_m)`$. Since spins that are connected, are necessarily aligned both before and after the move, the first product on the right hand side is equal in $`T(C_aC_m)`$ and $`T(C_bC_m)`$. All factors in the second product on the right hand side dealing with pairs of spins that are aligned both before and after the move are also equal in $`T(C_aC_m)`$ and $`T(C_bC_m)`$. That leaves in the ratio of the transition rates only the factors dealing with disconnected pairs of spins that are aligned either in configuration $`C_a`$, or in configuration $`C_b`$, but not both. The ratio of the transition rates $`T(C_aC_m)`$ and $`T(C_bC_m)`$ therefore reduces to
$$\frac{T(C_aC_m)}{T(C_bC_m)}=\underset{\begin{array}{ccc}i,j& & \\ \sigma _i^{(a)}=\sigma _j^{(a)}& & \\ \sigma _i^{(b)}\sigma _j^{(b)}& & \end{array}}{}\left(1P_c\right)/\underset{\begin{array}{ccc}i,j& & \\ \sigma _i^{(a)}\sigma _j^{(a)}& & \\ \sigma _i^{(b)}=\sigma _j^{(b)}& & \end{array}}{}\left(1P_c\right).$$
(3)
Using that $`\mathrm{log}(1P_c)=\beta J`$, in combination with some rewriting, we obtain for the logarithm of this ratio
$`\mathrm{log}(T(C_aC_m))\mathrm{log}(T(C_bC_m))`$ (4)
$`=`$ $`\beta J{\displaystyle \underset{i,j}{}}\left[\delta (\sigma _i^{(a)},\sigma _j^{(a)})\delta (\sigma _i^{(b)},\sigma _j^{(b)})\right].`$ (5)
As can easily been seen from the Hamiltonian eq. (1), this is equal to $`\beta (E_aE_b)`$. Since $`T(C_mC_a)=T(C_mC_b)=2^n`$ where $`n`$ is the number of clusters in $`C_m`$, detailed balance follows:
$`{\displaystyle \frac{T(C_bC_a)}{T(C_aC_b)}}`$ $`=`$ $`{\displaystyle \frac{T(C_bC_m)T(C_mC_a)}{T(C_aC_m)T(C_mC_b)}}`$ (6)
$`=`$ $`\mathrm{exp}\left(\beta (E_aE_b)\right).`$ (7)
In addition to obeying detailed balance, the algorithm is ergodic, since there is a finite probability that in a given move all clusters will contain one site only, to which any value can be assigned. Since this algorithm is ergodic and satisfies detailed balance, it is guaranteed that eventually these sample configurations are drawn from the Boltzmann distribution for the regular Potts model. The densities $`\rho _i`$ are not conserved in the Swendsen-Wang algorithm.
### B Density-conserving cluster algorithm
The topic of this paper is to present a modification to this algorithm, that ensures the conservation of the densities. This modification is made in step 3, in which the new spin values are assigned: rather than assigning random spin values to each cluster, we redistribute spin values over the clusters while conserving the spin densities. As for the original Swendsen-Wang algorithm, the general idea is a two-step approach, $`C_aC_mC_b`$, where all the energetics required for obtaining detailed balance are incorporated in the construction of the clusters, and detailed balance is achieved by conservation of the property $`T(C_mC_a)=T(C_mC_b)`$.
The first step towards such an algorithm is to devise an elementary move. The move we are looking for, is identifying one set of aligned clusters with spin value $`q_1`$ and another such set with spin value $`q_2q_1`$ with exactly the same area (number of sites), and then exchanging the spin values $`q_1`$ and $`q_2`$.
How do we identify such sets? First of all, for each spin value $`i=1\mathrm{}Q`$ we group all clusters with spin value $`i`$ into the set $`S_i`$. Next, within each set we list these clusters in a random order, and keep track of the cumulative area. Every time that in two sets the same value for the cumulative area occurs, we have found an exchange point. If the spin values are exchanged in all clusters up to the exchange point, while the original spin values in all other clusters are conserved, the spin values of two sets of clusters are exchanged without violation of the spin density conservation.
Unless extra measurements are taken, an algorithm based on these elementary moves will not obey detailed balance: the probability of occurrence for an exchange point is not necessarily equal before and after the cluster exchange. We denote the total number of clusters with spin-value $`q`$ before the exchange takes place as $`n_q`$. Suppose that the exchange takes place between clusters with spin $`1`$ and $`2`$, and that the number of clusters with spin $`1`$, $`2`$ that are to be exchanged is $`a_1`$ and $`a_2`$, respectively, while the number of clusters with spin $`1`$, $`2`$ that are not to be exchanged is $`n_1a_1`$ and $`n_2a_2`$, respectively. The likelihood that there is an exchange point exactly between these sets of clusters is then equal to
$$T()=\left[\left(\genfrac{}{}{0pt}{}{n_1}{a_1}\right)\left(\genfrac{}{}{0pt}{}{n_2}{a_2}\right)\right]^1$$
(8)
while after the exchange, this probability becomes
$`T()`$ $`=`$ $`\left[\left({\displaystyle \genfrac{}{}{0pt}{}{n_1^{}}{a_2}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n_2^{}}{a_1}}\right)\right]^1`$ (9)
$`=`$ $`\left[\left({\displaystyle \genfrac{}{}{0pt}{}{a_2+n_1a_1}{a_2}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{a_1+n_2a_2}{a_1}}\right)\right]^1`$ (10)
To restore detailed balance, it suffices to introduce a Metropolis acceptance ratio:
$$P_a=\mathrm{min}[1,\frac{n_1^{}!n_2^{}!}{n_1!n_2!}].$$
(11)
Once this acceptance probability is included, the elementary move can be used for a correct algorithm, since for two configurations $`X`$ and $`Y`$, we now restored the property $`T(C_mX)=T(C_mY)`$.
In an actual implementation, the total procedure is to make for each spin value a cumulative list of clusters, where the clusters are placed in a random order. Next, all exchange points are identified, the corresponding exchanges are accepted with the probability as given in eq. (11). It can be verified that also for the concatenation of exchange points, the product over all exchange points of the ratio of forward and backward acceptance probabilities, as given in eq. (11), equals
$$\underset{q=1}{\overset{Q}{}}\left(n_q^{}!\right)/\underset{q=1}{\overset{Q}{}}\left(n_q!\right),$$
(12)
which exactly cancels the ratio of the number of ways in which the clusters can be sorted, i.e. the ratio of selection probabilities in forward and backward direction. Consequently $`T(C_mC_a)=T(C_mC_b)`$.
The density-conserving algorithm is ergodic for the same reason that the Swendsen-Wang algorithm is ergodic: there is a finite probability that all clusters contain one site only, and then each of these can obtain any spin value (under the constraint on the densities).
Having shown that the basic steps of our algorithm are correct we will now summarize the procedure in the form of a step-wise algorithm. Steps number 1 and 2 of the Swendsen-Wang algorithm remain unchanged. Step 3 becomes:
1. For each state $`q`$, list all clusters with this spin-value in list $`S_q`$, in a random order.
2. Order the lists with respect to the total area of their not-yet-assigned clusters. Use a random order for lists with equal such areas. If the first two lists (those with the largest and next-largest areas) are equal, exchange their colors with the probability as given by eq. (11). Select one cluster from the first list and assign to it a new color. Update the ordering and repeat this step until spin values are assigned to all clusters.
Note that the computational effort required for step 3 scales with the total number of clusters $`n=_qn_q`$, whereas step 2 scales with the number of spins $`N`$ in the system. Since $`nN`$, step 3 is repeated $`N/(2n)`$ times for each time step 2 is performed, and we still have an implementation in which the computational effort per sweep scales linearly with the number of sites; this greatly decreases the auto-correlation time. It actually also reduces the dynamical critical exponent $`z`$, since the ratio $`N/n`$ varies with the system size.
## III Computational properties
In order to compare the efficiency of the density-conserving cluster algorithm presented above with that of non-local spin-exchange (Kawasaki) dynamics, we have computed the energy autocorrelation times in the Ising model at critical temperature and equal spin densities, for several system sizes. Figures 2 and 3 show the correlation times as a function of the linear system size $`L`$ of the two- and three-dimensional Ising model, respectively, both at their critical point. For all data points the correlation time $`\tau `$ was obtained from a least-squares fit of the form $`e^{t/\tau }`$, to the energy autocorrelation function. For the spin-exchange algorithm, these fits were done in the region where the autocorrelation drops from $`e^1`$ to $`e^2`$; for the cluster algorithm, we fitted in a broader region (from 1 to $`e^3`$), in order to have enough points to fit. All runs where started from a random configuration, which was thermalized over a time varying from 6000 MCS to 60000 MCS. In order to generate enough statistics, the total length of the runs was set to ten times the thermalization time. The statistical errors were determined by repeating each run 10 to 50 times.
As expected, we find that the cluster algorithm suffers significantly less from critical slowing down and clearly outperforms spin-exchange dynamics at physically interesting lattice sizes in both two and three dimensions. Since one move in our cluster algorithm takes an amount of CPU-time comparable to what is required for one sweep in the non-local spin exchange, our cluster algorithm outperforms non-local spin exchange by one or two orders of magnitude, depending on the system size.
For the non-local spin-exchange algorithm, we find a critical dynamic exponent of $`z=2.0`$ in both two and three dimensions. This is in good agreement with the exponents of the three-dimensional non-conserving Metropolis algorithm ($`z=2.02\pm 0.02`$) but not with the critical exponent for the two-dimensional non-conserving Metropolis algorithm ($`z=2.167\pm 0.001`$). Perhaps this is an indication that the conservation of the order parameter affects the critical dynamical exponent, but our statistics are not conclusive.
For the critical dynamic exponent for our new density-conserving cluster algorithm, we find values of $`z=0.38\pm 0.01`$ in two, and $`z=0.66\pm 0.02`$ in three dimensions. These values are both slightly larger than non-conserved Swendsen-Wang values ($`z=0.25\pm 0.01`$ and $`z=0.54\pm 0.02`$ for two and three dimensions respectively).
## IV Summary and future work
We have presented a density-conserving cluster algorithm for the Potts model. This algorithm is only moderately sensitive to critical slowing down: its dynamic critical exponent is found to be $`z=0.38\pm 0.01`$ and $`z=0.66\pm 0.02`$ for the two- and three-dimensional two-states Potts model, respectively. It outperforms the traditional algorithm, non-local spin exchange, by one or two orders of magnitude.
In future research, we will use this algorithm to study wetting properties, where the wetting takes place at a curved interface between two co-existing phases; such non-flat interfaces arise for instance between a droplet and a surrounding fluid. Other future applications will include the study of line tension between three co-existing phases, and equilibrium shapes in multi-component mixtures.
## Acknowledgments
Useful discussion with Henk van Beijeren and Matthieu Ernst is gratefully acknowledged.
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# Decays of the 𝐵_𝑐 Meson in a Relativistic Quark–Meson Model
## 1 Introduction
A primary goal in the study of semileptonic decays of heavy mesons is to extract the values of the CKM matrix elements. The great virtue of semileptonic decays is that the effects of the strong interaction can be separated from the effects of the weak interaction into a set of Lorentz invariant form factors . Thus the theoretical problem associated with analysing semileptonic decays is essentially that of calculating the form factors.
The focus of this work is the decay of the $`B_c`$ meson (for a review of the properties of this system see ). This system is unique among mesons made up of heavy (charm or bottom) quarks, it is the only one which is stable with respect to strong and electromagnetic interactions. Therefore, the $`B_c`$ system is the only heavy meson for which form factors (albiet transition form factors rather than elastic) can be measured. These form factors then provide a unique probe of the dynamics of heavy quark systems.
There are many approaches to the calculation of decay form factors, for example, lattice QCD , QCD sum rules , and phenomological modelling . In this work a particular model with an effective quark-meson coupling is adopted. There are many models of this type . The one used here has its genesis in the QCD version of the Nambu-Jona-Lasinio model extended to heavy quarks and is most closely related to the model used recently by Ivanov and Santorelli in a their study of pseudoscalar meson decays .
The advantage of this approach is that it is fully relativistic and very versatile. Quarks and mesons for all masses are treated within the same framework. For light quarks the model has the features of spontaneous chiral symmetry breaking and in the single heavy quark limit the form factor constraints of heavy quark symmetry are obtained.
Our work differs from Ivanov and Santorelli in the choice of the quark-meson vertex function and in the way that parameters are fixed. A number of heavy mesons decays not calculated in Ref. are treated here. The main new results are the extension of the model to include doubly heavy mesons, the calculation of $`B_c`$ semileptonic decays and the electromagnetic vector to pseudoscalar transitions.
This paper is organized as follows: the next section introduces the model, discusses the general method of calculation, and fits the models free parameters. Sect. 3 presents the calculation of the form factors and decay rates for the semileptonic decays of a wide varity of heavy-light pseudoscalar mesons. These calculations are compared with both measured results, and other theoretical approaches. Sect. 4 presents the same set of calculations for the eight primary semileptonic decays of the $`B_c`$ meson. The predictions are compared with other theoretical work, in order to highlight the differences that exist between various approaches. Sect. 5 briefly discusses the electromagnetic decays $`VP+\gamma `$ for a number of vector mesons, including the $`B_c^{}`$. Sect. 6 gives conclusions and directions for future work.
## 2 Quark–Meson Coupling
The particular quark–meson coupling used in this work is based on an effective Lagrangian which models the interaction between mesons and quarks with a non-local interaction vertex . The interaction Lagrangian has the form
$`L_{\mathrm{int}}(x)`$ $`=`$ $`g_MM(x){\displaystyle }dx_1dx_2\delta (x\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{m_1x_1+m_1x_2}{m_1+m_2}}\right))\times `$ (1)
$`f[(x_1x_2)^2]\overline{q}_1(x_1)\mathrm{\Gamma }_Mq_2(x_2),`$
where $`\mathrm{\Gamma }_M`$ is the Dirac matrix appropriate to the meson field M, $`f[(x_1x_2)^2]`$ is a non-local vertex function, which simulates the finite size of the meson, and $`q_1`$ and $`q_2`$ are the quark fields. A condition imposed on the vertex function is that it should render all loop diagrams UV finite. The coupling constant $`g_M`$ is determined by the compositeness condition, which is the requirement that the renormalization constant of the meson fields be zero, i.e.
$$Z_M=1\frac{g_M^2}{2}\frac{d\mathrm{\Pi }_M(p^2)}{dp^2}|_{p^2=M_M^2}=0.$$
(2)
Here $`\mathrm{\Pi }_M(p^2)`$ is the self energy of the meson field, given by
$`\mathrm{\Pi }_M(p^2)`$ $`=`$ $`2N_c{\displaystyle \frac{d^4k}{(2\pi )^4i}f^2(Q^2)}`$ (3)
$`\mathrm{tr}\left\{\mathrm{\Gamma }_M{\displaystyle \frac{1}{m_1(\mathit{}+\mathit{})}}\mathrm{\Gamma }_M{\displaystyle \frac{1}{m_2\mathit{}}}\right\},`$
where $`m_1`$ and $`m_2`$ are the masses of the quarks in the loop and $`Q`$ is a relative momentum chosen to be $`Q=k+\alpha p`$ with $`\alpha =\frac{m_2}{m_1+m_2}`$.
The constituent quark masses in (3) are free parameters. As well, the vertex function will contain a free parameter which reflects the size of the meson. These parameters will be different for the different mesons.
The use of free constituent quark propagators in expressions like (3) can lead to a problem which reflects the lack of quark confinement in the model. If the meson mass $`M_M`$ is greater than the sum of its constituent quark masses loop integrals will develop imaginary parts. This indicates a non-zero amplitude for the creation of a free quark-antiquark pair. There have been some various attempts to obviate this problem within quark-meson effective theories . Here we adopt the approach of Ref. and use free propagators. The constituent quark masses are then fit to allow for the inclusion of as many mesons as possible.
In order to carry out calculations a choice must be made for the vertex function $`f(q^2)`$. The function that was used in this analysis was the dipole
$$f(Q^2)=\frac{\mathrm{\Lambda }^4}{[\mathrm{\Lambda }^2Q^2]^2}.$$
This choice was made for two reasons; first the form of the dipole vertex function is the same as a propagator, allowing standard Feynman parameter techniques to be used in evaluating loop integrals. Second, the vector decay constant $`f_V`$ would diverge if only a monopole vertex function was used. Since one of the primary criteria for the vertex functions is that they should render all diagrams UV finite, a function with UV fall-off as least as fast as a dipole is needed. The parameter $`\mathrm{\Lambda }`$ characterizes the finite size of the meson, and will be different for different mesons. To account for this the various values of $`\mathrm{\Lambda }`$ will be distinguished by subscripts which reflect either the meson type or the quark content, *e.g.* $`\mathrm{\Lambda }_{B_c}`$ and $`\mathrm{\Lambda }_{bc}`$ will be used interchangeably. Further, in expressions involving the vertex form factor, the same comvention will be used. Note that the calculations of Ref. used a Gaussian vertex function so that the parameters used there can not be compared directly with ours.
The parameters of the model were fit to the leptonic decay constants, $`f_P`$ and $`f_V`$. These quantities are defined by
$`0|i\gamma ^\mu \gamma ^5|P`$ $`=`$ $`if_Pp^\mu ,`$ (4)
$`0|i\gamma ^\mu |V,ϵ`$ $`=`$ $`M_V^2f_Vϵ^\mu `$ (5)
where $`M_V`$ is the vector meson mass. The pseudoscalar decay constant is given by the one-loop expression
$`f_Pp^\mu `$ $`=`$ $`N_c{\displaystyle \frac{d^4k}{(2\pi )^4i}g_Pf(Q^2)}`$ (6)
$`{\displaystyle \frac{\mathrm{tr}\left\{\gamma ^\mu \gamma ^5[m_1+\mathit{}+\mathit{}]\gamma ^5[m_2+\mathit{}]\right\}}{\left[m_1^2(k+p)^2\right]\left[m_2^2k^2\right]}}.`$
Here $`m_1`$ and $`m_2`$ refer to the masses of the quarks in the loop, this convention will be used throughout this paper. Using the dipole vertex function, combining the denominators using Feynman parameters, and performing the integration over k yields
$$f_P=g_P\frac{3\mathrm{\Lambda }_P^4}{4\pi ^2}D\stackrel{}{x}\frac{x_1[m_2(1\sigma )+m_1\sigma ]}{\mathrm{\Delta }^2},$$
(7)
with
$`\sigma `$ $`=`$ $`\alpha x_1+x_2,`$
$`\eta `$ $`=`$ $`\alpha ^2x_1+x_2,`$
$`\mathrm{\Delta }`$ $`=`$ $`\mathrm{\Lambda }_P^2x_1+m_1^2x_2+m_2^2x_3+(\sigma ^2\eta )M_P^2,`$
$`\mathrm{}`$ $`=`$ $`k+\sigma p,`$
$`{\displaystyle D\stackrel{}{x}}`$ $`=`$ $`{\displaystyle _0^1}\left({\displaystyle \underset{i=1}{\overset{3}{}}}dx_i\right)\delta \left({\displaystyle \underset{i=1}{\overset{3}{}}}x_i1\right),`$
where $`M_P`$ is the mass of the pseudoscalar meson. Likewise the expression for the vector decay constant is
$$f_V=g_V\frac{3\mathrm{\Lambda }_V^4}{4\pi ^2M_V^2}D\stackrel{}{x}x_1\frac{m_1m_2+\mathrm{\Delta }+\sigma (1\sigma )M_V^2}{\mathrm{\Delta }^2},$$
(8)
where the same defintions have been used, with the obvious change of $`M_P`$ to $`M_V`$ and $`\mathrm{\Lambda }_P`$ to $`\mathrm{\Lambda }_V`$ in the expression for $`\mathrm{\Delta }`$.
To compute the coupling constants $`g_P`$ and $`g_V`$, the self energies and their derivatives must be computed. Then the compositeness condition (2) can be used to find the couplings. The self energy for a pseudoscalar meson is given by
$$\mathrm{\Pi }_P(p^2)=\frac{3\mathrm{\Lambda }_P^8}{2\pi ^2}D\stackrel{}{x}x_1^3\frac{m_1m_2+p^2\sigma (1\sigma )+\frac{2}{3}\overline{\mathrm{\Delta }}}{\overline{\mathrm{\Delta }}^4},$$
(9)
where $`\overline{\mathrm{\Delta }}=\mathrm{\Lambda }_P^2x_1+m_1^2x_2+m_2^2x_3+(\sigma ^2\eta )p^2`$ and all the other quantities are the same as the ones defined above. The self energy for a pseudoscalar meson is given by the tensor $`\mathrm{\Pi }_V^{\mu \nu }`$ which can be expressed as
$$\mathrm{\Pi }_V^{\mu \nu }(p)=\mathrm{\Pi }_V(p^2)g_{\mu \nu }+\overline{\mathrm{\Pi }}_V(p^2)\frac{p_\mu p_\nu }{p^2}.$$
(10)
Unfortunately $`\mathrm{\Pi }_V\overline{\mathrm{\Pi }}_V`$, so this does not have the proper structure for a vector propagator. This problem was solved (following ) by simply dropping the $`\overline{\mathrm{\Pi }}_V`$ term, which would cancel out of any calculation of a physical process at one-loop order (since $`ϵp=0`$). The relevant part of the vector meson self energy is given by
$$\mathrm{\Pi }_V(p^2)=\frac{3\mathrm{\Lambda }_V^8}{2\pi ^2}D\stackrel{}{x}x_1^3\frac{m_1m_2+\frac{1}{3}\overline{\mathrm{\Delta }}+\sigma (1\sigma )p^2}{\overline{\mathrm{\Delta }}^4},$$
(11)
where all the quantities appearing have been defined previously.
The free parameters of the model are fit to the six values of $`f_P`$ and the measured values $`f_{J/\psi }=0.1309`$ and $`f_\mathrm{{\rm Y}}=0.075012`$ . These data, which are displayed in Table 1, fix eight free parameters. In order to reduce the number of free parameters to match the available data the value of the strange quark mass was fixed at 450 MeV and the vertex parameter for a vertex containing only u and d quarks $`\mathrm{\Lambda }_\pi `$ was taken (following ) to be 1 GeV. In addition the following further simplifying assumptions were made
$$\mathrm{\Lambda }_{us}=\mathrm{\Lambda }_{ds}=\mathrm{\Lambda }_{ss}=\mathrm{\Lambda }_K,$$
$$\mathrm{\Lambda }_{uc}=\mathrm{\Lambda }_{dc}=\mathrm{\Lambda }_{sc}=\mathrm{\Lambda }_D,$$
$$\mathrm{\Lambda }_{ub}=\mathrm{\Lambda }_{db}=\mathrm{\Lambda }_{sb}=\mathrm{\Lambda }_B.$$
This leaves the following parameters to be fit, $`m_q`$, $`m_c`$, $`m_b`$, $`\mathrm{\Lambda }_K`$, $`\mathrm{\Lambda }_D`$, $`\mathrm{\Lambda }_B`$, $`\mathrm{\Lambda }_{cc}`$, $`\mathrm{\Lambda }_{bb}`$, and $`\mathrm{\Lambda }_{bc}`$. The parameter $`\mathrm{\Lambda }_{bc}`$ could only be fit to a value for $`f_{B_c}`$ which is not in the values listed in Table 1, hence it is retained as a free parameter, leaving eight to be fit. The fit to the remaining eight parameters is given by (all values in MeV)
$`m_{u,d}`$ $`=`$ $`245,`$
$`m_c`$ $`=`$ $`1800,`$
$`m_b`$ $`=`$ $`5100,`$
$`\mathrm{\Lambda }_K`$ $`=`$ $`1225,`$
$`\mathrm{\Lambda }_D`$ $`=`$ $`1350,`$
$`\mathrm{\Lambda }_B`$ $`=`$ $`1500,`$
$`\mathrm{\Lambda }_{cc}`$ $`=`$ $`1420,`$
$`\mathrm{\Lambda }_{bb}`$ $`=`$ $`2900.`$
The values for the self energies, coupling constants, and leptonic decay constants arising from these parameters are displayed in Tables 2 and 3.
In order to fix $`\mathrm{\Lambda }_{bc}`$ a value of $`f_{B_c}`$ must be given. There is no experimental value for this quantity and theoretical estimates tend to fall in the range $`400\mathrm{MeV}f_{B_c}500\mathrm{MeV}`$ (see, for example, ). A further complication is that the mass of $`M_{B_c}`$ is also not yet known very well. The current measurement is $`M_{B_c}^{CDF}=6.4\pm 0.39(stat)\pm 0.13(sys)\frac{\mathrm{GeV}}{c^2}`$ , which comes from the few confirmed $`B_c`$ events at the Tevatron. Theoretical results tend to lie within this range, so following the potential model prediction of the mass of the $`B_c`$ was chosen to be 6.25 GeV.
One general argument guides the selection of $`\mathrm{\Lambda }_{bc}`$, it should lie between $`\mathrm{\Lambda }_B`$ and $`\mathrm{\Lambda }_{bb}`$. With this in mind, and using the value for $`M_{B_c}`$ above, a number of values of $`\mathrm{\Lambda }_{bc}`$ were tried, spanning the possible range. Fig. 1 shows the value of $`f_{B_c}`$ as a function of $`\mathrm{\Lambda }_{bc}`$. The value selected selected for use in this work was $`\mathrm{\Lambda }_{bc}=2.3`$ GeV, which gives $`f_{B_c}=450`$ MeV, a value in the middle of the range of the theoretical predictions.
## 3 Semileptonic Decays of $`K`$, $`D`$, and $`B`$ Mesons
The model used in this work is phenomenological but having fixed its parameters, the results for semileptonic decays are predictions. Before proceeding to decays of $`B_c`$ it is important to test the model against experimental results where they are available. Therefore several semileptonic decays of $`K`$, $`D`$, and $`B`$ meson are calculated. The formalism for these calculations, presented in this section, extends directly also to the calculation of $`B_c`$ decay.
Some of the decays considered here have already been treated by Ivanov and Santorelli . However, that work does not demonstrate the full applicablity of the approach. Apart from decays to light vector mesons, the model is capable of treating virtually any semileptonic decay (with the restriction that a value for the meson mass must be supplied as input).
The amplitude $`A`$ for a semileptonic decay is given by,
$$A=\frac{G_F}{\sqrt{2}}V_{QQ^{}}L_\mu H^\mu .$$
(12)
Here $`G_F`$ is the Fermi constant, $`V_{QQ^{}}`$ is the relevant CKM matrix element, $`L_\mu `$ is the lepton current
$$L_\mu =\overline{u}_\nu _{\mathrm{}}\gamma _\mu (1\gamma ^5)v_{\mathrm{}},$$
and $`H^\mu `$ is the hadron current
$$H^\mu =k,ϵ|(V^\mu A^\mu )|P,$$
(13)
where $`P`$ is the momentum of the parent meson, $`k`$ is the momentum of the daughter meson, and $`ϵ`$ is the polarization, if the daughter meson is a vector. The two currents in (13) are the vector $`V^\mu `$ and axial $`A^\mu `$. If the final state is a pseudoscalar the hadron current can be decomposed as follows,
$`k|A^\mu |P`$ $`=`$ $`0,`$
$`k|V^\mu |P`$ $`=`$ $`f_+(q^2)(P+k)^\mu +f_{}(q^2)(Pk)^\mu ,`$
where $`f_+(q^2)`$ and $`f_{}(q^2)`$ are Lorentz invariant form factors. Likewise, if the final state is a vector meson,
$`k,ϵ|A^\mu |P`$ $`=`$ $`f(q^2)ϵ^\mu +a_+(q^2)(ϵ^{}P)(P+k)^\mu +`$
$`a_{}(q^2)(ϵ^{}P)(Pk)^\mu ,`$
$`k,ϵ|V^\mu |P`$ $`=`$ $`ig(q^2)ϵ^{\mu \nu \rho \sigma }ϵ_\nu ^{}(P+k)_\rho (Pk)_\sigma ,`$
where the form factors are $`g`$, $`f`$, $`a_+`$, and $`a_{}`$. In each of these expressions $`q=(Pk)`$ is the momentum transfer.
For a decay to a pseudoscalar meson (with mass denoted by $`M_P^{}`$) the differential decay rate can be reduced to
$$\frac{d\mathrm{\Gamma }}{dq}=\frac{G_F^2|V_{QQ^{}}|^2M_P^2K^3}{24\pi ^3}|f_+(q^2)|^2.$$
(14)
where,
$$K=\frac{M_P}{2}\sqrt{\left[1\frac{M_P^{}^2}{M_P^2}y\right]^24\frac{M_P^{}^2}{M_P^2}y}.$$
(15)
The lepton spectrum is given by,
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}`$ $`=`$ $`{\displaystyle \frac{G_F^2|V_{QQ^{}}|^2M_P^5}{16\pi ^3}}(12x)`$ (16)
$`{\displaystyle _0^{y_{max}(x)}}\left([y_{max}(x)y]|f_+(q^2)|^2\right)𝑑y,`$
where $`y_{max}(x)=\frac{4x(x_{max}x)}{12x}`$ with $`x_{max}=\frac{M_P^2M_P^{}^2}{2M_P^2}`$. If the final state is a vector meson (with mass $`M_V`$)the corresponding differential decay rate is,
$$\frac{d\mathrm{\Gamma }}{dy}=\frac{G_F^2|V_{QQ^{}}|^2KM_P^2y}{96\pi ^3}\left(|\overline{H}_+|^2+|\overline{H}_{}|^2+|\overline{H}_0|^2\right),$$
(17)
where
$`\overline{H}_\pm `$ $`=`$ $`f(q^2)2M_PKg(q^2),`$
$`\overline{H}_0`$ $`=`$ $`\left[{\displaystyle \frac{M_P}{2M_V\sqrt{y}}}\right]\left[\left(1{\displaystyle \frac{M_V^2}{M_P^2}}y\right)f(q^2)+4K^2a_+(q^2)\right],`$
and the final mass $`M_V`$ should be subsititued for $`M_P^{}`$ in (15). The expression for the lepton spectrum is given by
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}`$ $`=`$ $`{\displaystyle \frac{G_F^2|V_{QQ^{}}|^2M_P^5}{32\pi ^3}}{\displaystyle _0^{y_{max}(x)}}\{{\displaystyle \frac{\alpha (y)}{M_P^2}}y+`$ (18)
$`2(12x)[y_{max}(x)y]\beta _{++}(y)+`$
$`\gamma (y)y[2x_{max}4x+y]\},`$
where the following definitions were made
$`\alpha (q^2)`$ $`=`$ $`|f(q^2)|^2+\lambda |g(q^2)|^2,`$
$`\gamma (q^2)`$ $`=`$ $`2f(q^2)g(q^2),`$
$`\beta _{++}(q^2)`$ $`=`$ $`{\displaystyle \frac{1}{4M_V}}\{|f(q^2)|^2+\lambda |a_+(q^2)|^24M_V^2q^2|g(q^2)|^2`$
$`+2(M_P^2M_V^2q^2)f(q^2)a_+(q^2)\},`$
$`\lambda (q^2)`$ $`=`$ $`(M_P^2M_V^2q^2)^24M_V^2q^2,`$
$`x_{max}`$ $`=`$ $`{\displaystyle \frac{M_P^2M_V^2}{2M_P^2}}`$
Note that all of these expressions assume that lepton mass $`m_{\mathrm{}}`$ is zero.
Our model can be used to calculate all of these form factors. In all the following expressions the inital meson is composed of quarks with masses $`m_1`$ and $`m_2`$ and the final meson is composed of quarks with masses $`m_3`$ and $`m_2`$ (*i.e.* $`m_1`$ is the mass of the quark which decays to a new quark with mass $`m_3`$, and $`m_2`$ is the mass of the spectator).
The form factors for decay to a pseudoscalar meson are
$`f_+(q^2)`$ $`=`$ $`g_Pg_P^{}{\displaystyle \frac{9\mathrm{\Lambda }_P^4\mathrm{\Lambda }_P^{}^4}{\pi ^2}}`$ (19)
$`{\displaystyle D\stackrel{}{x}x_1x_2\frac{\chi _+\frac{3}{4}\overline{\mathrm{\Delta }}(\sigma _1+\sigma _2)+\overline{\mathrm{\Delta }}}{\overline{\mathrm{\Delta }}^5}},`$
$`f_{}(q^2)`$ $`=`$ $`g_Pg_P^{}{\displaystyle \frac{9\mathrm{\Lambda }_P^4\mathrm{\Lambda }_P^{}^4}{\pi ^2}}`$ (20)
$`{\displaystyle D\stackrel{}{x}x_1x_2\frac{\chi _{}\frac{3}{4}\overline{\mathrm{\Delta }}(\sigma _1\sigma _2)}{\overline{\mathrm{\Delta }}^5}}.`$
The following definitions were made to simplify the expressions
$`{\displaystyle D\stackrel{}{x}}`$ $`=`$ $`{\displaystyle _0^1}\left({\displaystyle \underset{i=1}{\overset{5}{}}}dx_i\right)\delta \left({\displaystyle \underset{i=1}{\overset{5}{}}}x_i1\right),`$
$`\mathrm{\Delta }`$ $`=`$ $`\mathrm{\Lambda }_P^{}^2x_1+\mathrm{\Lambda }_P^2x_2+m_3^2x_3+m_1^2x_4+m_2x_5,`$
$`\mu _{ij}`$ $`=`$ $`{\displaystyle \frac{m_i}{m_i+m_j}},`$
$`\sigma _{1,(2)}`$ $`=`$ $`x_{4,(3)}+\mu _{12,(23)}x_{2,(1)},`$
$`\eta _{1,(2)}`$ $`=`$ $`x_{4,(3)}+\mu _{12,(23)}^2x_{2,(1)},`$
$`\overline{\mathrm{\Delta }}`$ $`=`$ $`\mathrm{\Delta }+(\sigma _1^2\eta _1+\sigma _1\sigma _2)M_P^2`$
$`+(\sigma _2\eta _2+\sigma _1\sigma _2)M_P^{}^2\sigma _1\sigma _2q^2,`$
$`\kappa `$ $`=`$ $`m_3(m_2m_1)+m_1m_2+{\displaystyle \frac{1}{2}}(M_P^2+M_P^{}^2q^2),`$
$`ϵ`$ $`=`$ $`(\sigma _1+\sigma _2)\sigma _1M_P^2+(\sigma _1+\sigma _2)\sigma _2M_P^{}^2\sigma _1\sigma _2q^2,`$
$`\zeta _1`$ $`=`$ $`m_1m_2\left[\sigma _1\left(\sigma _11\right)+\sigma _2\left(\sigma _1{\displaystyle \frac{1}{2}}\right)\right]M_P^2`$
$`\sigma _2\left(\sigma _1+\sigma _2{\displaystyle \frac{1}{2}}\right)M_P^{}^2+\sigma _2\left(\sigma _1{\displaystyle \frac{1}{2}}\right)q^2,`$
$`\zeta _2`$ $`=`$ $`m_2m_3\sigma _1\left(\sigma _1+\sigma _2{\displaystyle \frac{1}{2}}\right)M_P^2`$
$`\left[\sigma _2\left(\sigma _21\right)+\sigma _1\left(\sigma _2{\displaystyle \frac{1}{2}}\right)\right]M_P^{}^2`$
$`+\sigma _1\left(\sigma _2{\displaystyle \frac{1}{2}}\right)q^2,`$
$`\chi _\pm `$ $`=`$ $`(ϵ\kappa )(\sigma _1\pm \sigma _2)\pm \zeta _1+\zeta _2.`$
These definitions (in addition to $`\alpha `$ and $`\mu `$) will be used throughout the rest of this paper, with the obvious substitution of $`M_V`$ and $`\mathrm{\Lambda }_V`$ for $`M_P^{}`$ and $`\mathrm{\Lambda }_P^{}`$ when the final state is a vector meson.
The form factors for decays to vector mesons are given by
$`g(q^2)`$ $`=`$ $`g_{M_P}g_{M_V}{\displaystyle \frac{9\mathrm{\Lambda }_P^4\mathrm{\Lambda }_V^4}{\pi ^2}}{\displaystyle D\stackrel{}{x}x_1x_2}`$ (21)
$`{\displaystyle \frac{\sigma _2(m_2m_3)+\sigma _1(m_2m_1)m_2}{\overline{\mathrm{\Delta }}^5}},`$
$`f(q^2)`$ $`=`$ $`g_{M_P}g_{M_V}{\displaystyle \frac{18\mathrm{\Lambda }_P^4\mathrm{\Lambda }_V^4}{\pi ^2}}{\displaystyle }D\stackrel{}{x}x_1x_2\times `$ (22)
$`{\displaystyle \frac{1}{\overline{\mathrm{\Delta }}^5}}[m_1m_2m_3{\displaystyle \frac{1}{4}}(m_2m_12m_3)\overline{\mathrm{\Delta }}`$
$`+[\xi _1+\xi _3]M_P^2+[\xi _2+\xi _3]M_V^2\xi _3q^2],`$
$`a_\pm (q^2)`$ $`=`$ $`g_{M_P}g_{M_V}{\displaystyle \frac{18\mathrm{\Lambda }_P^4\mathrm{\Lambda }_V^4}{\pi ^2}}{\displaystyle D\stackrel{}{x}x_1x_2}`$ (23)
$`{\displaystyle \frac{\frac{1}{2}(\beta _1\pm \beta _2\sigma _2m_3)}{\overline{\mathrm{\Delta }}^5}}.`$
The following further definitions have been made,
$`\xi _1`$ $`=`$ $`\sigma _1(m_3m_2)(1\sigma _1)m_1\sigma _1^2,`$
$`\xi _2`$ $`=`$ $`\sigma _2(m_1m_2)(1\sigma _2)m_3\sigma _2^2,`$
$`\xi _3`$ $`=`$ $`m_2\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}(\sigma _1+\sigma _2)+\sigma _1\sigma _2\right]+m_1\sigma _1\left({\displaystyle \frac{1}{2}}\sigma _2\right)`$
$`+m_3\sigma _2\left({\displaystyle \frac{1}{2}}\sigma _1\right),`$
$`\beta _1`$ $`=`$ $`2\sigma _1[m_1\sigma _1+m_2(1\sigma _1)],`$
$`\beta _2`$ $`=`$ $`m_2(\sigma _1+\sigma _212\sigma _1\sigma _2)m_1\sigma _1(12\sigma _2).`$
Excluding the $`B_c`$ decays a total of sixteen pseudoscalar to pseudoscalar decays were considered. Due to the difficulty with confinement the corresponding number of pseudoscalar to vector decays that could be treated was only four. Table 4 shows the predictions for the decay rates and branching ratios for all of the decays considered. The values of the CKM matrix elements, and the necessary lifetimes were taken from .
Many of the decay rates treated in this section have been measured, hence most of the predictions can be compared to observed quantities. Table 4 shows the predicted and measured results for the branching ratios. The experimental results are taken from and the errors in the predictions represent the uncertainties in the CKM matrix elements. Overall, the agreement with experiment is reasonable which increases the level of confidence in the areas where direct comparison with experiment is not possible.
Table 5 shows values of $`f_+(0)`$ as computed in this work and in various other theoretical approaches. The other approaches are widely varied: uses the ISGW model, uses the WBS model, gives results from a bag model, uses a Dyson-Schwinger equation approach and gives lattice QCD results. Of particular interest is the work of Ivanov *et al.* which uses the quark confinement model. This quark–meson model is based on similar considerations to the model used in this work so its predictions should be close to ours.
As well, the decay $`BD+\mathrm{}^++\nu _{\mathrm{}}`$ can be treated in a model independent way using the HQET. Ivanov *et al.* have shown in several papers that quark–meson models of the type used here give the correct tree level HQET relations in the infinite mass limit. Nevertheless a direct check with finite quark mass is useful. The HQET gives the prediction
$$f_+^{HQET}(q_{max}^2)=\frac{m_{B^0}+m_D^{}}{2\sqrt{m_{B^0}m_D^{}}}=1.138,$$
which compares well with our value $`f_+(q_{max}^2)=1.133.`$
The most important comparison that can be made is with Ref.. This paper uses a different vertex function to treat B and D decays. This serves as a check on the dependence of the model on the choice of vertex function. Apart from the case $`B\pi +\mathrm{}+\nu `$ agreement with is very good. In addition presents the values of $`f_+(q^2)`$ over the full range of $`q^2`$. Overall agreement is good between the two calculations, Fig. 2 illustrates the agreement in the case $`D^0K^{}+\mathrm{}^++\nu _{\mathrm{}}`$. Fig. 3 shows the case $`B^0\pi ^{}+\mathrm{}^++\nu _{\mathrm{}}`$, for which the agreement is better over the whole range than indicated in Table 5.
Due to lack of confinement very few pseudoscalar to vector decays can be calculated. Of the few decays treated in this work only the decay $`BD^{}+\mathrm{}+\nu _{\mathrm{}}`$ has been studied extensively. Table 6 compares our predictions with some other calculations. Overall the agreement is reasonable.
## 4 Semileptonic Decays of the $`B_c`$ Meson
The methods of the previous section can be directly applied to the semileptonic $`B_c`$ decays. Using the procdeure outline above, decay rates, lepton spectra, and branching ratios can be computed. In this work the lifetime of the $`B_c`$ was taken to be $`0.5`$ ps, which agrees with the CDF value of $`\tau _{B_c}^{CDF}=0.46_{0.16}^{+0.18}\pm 0.03`$ ps . Table 7 shows $`f_+(0)`$, $`f_+(q_{max}^2)`$, the total decay rate $`\mathrm{\Gamma }`$ and the branching ratio for the four pseudoscalar decays. For the decays to vector mesons, values of the form factors at $`q^2=0`$ as well as total decay rates and branching ratios are displayed in Tables 8 and 9.
There are a number of other calculations of the semileptonic decays of $`B_c`$. A comparison of some results for the dominant decay modes is given in Table 10. In contrast to the situation in Sect. 3 where our quark-meson model predictions, for the most part, agreed with other models and the various other models agreed with each other, there are substantial differences between calculations of $`B_c`$ decays. The clearest examples of this are the predictions for the decays to the $`B_s^{}`$ and $`J/\psi `$. These two decays are expected to be the most important semileptonic decay channels However there is disagreement not only over the values of the branching ratios but also as to which decay will be favoured. For example, the quark-meson model used in this work predicts the decay to the $`J/\psi `$ to be slightly favoured over the decay to the $`B_s^{}`$ while the heavy quark approach used in and predicts the decay to $`B_s^{}`$ to dominate. This divergence of predictions may be expected; the heavy–heavy quark content of the $`B_c`$ poses a challenge for models. Light–quark mesons may be constrained by chiral symmetry and heavy–light mesons by heavy quark symmetry. On the other hand the physics of heavy–heavy systems is less constrained by symmetries so extending models into this domain provides a severe test.
## 5 Electromagnetic Decays $`VP+\gamma `$
In addition to semileptonic decays the electromagnetic decays of vector mesons can be treated within our effective quark-meson coupling model. Since the amplitude involves the matrix element $`V|V^\mu |P`$ it is clear this process will be related to the form factor $`g(q^2)`$. The the amplitude for this process is
$$A=2eϵ^{\mu \nu \alpha \beta }\stackrel{~}{ϵ}_\mu ϵ_\nu p_{V\alpha }p_{P\beta }[Q_1g_1(0)+Q_2g_2(0)],$$
(24)
where $`Q_{1,(2)}`$ is the charge of $`q_{1,(2)}`$, and $`p_{P,V}`$ are the momenta of the pseudoscalar and vector mesons. The functions $`g_i(0)`$ are the form factors given by (21), with the appropriate masses inserted, and with $`q^2=0`$. The appropriate masses in these functions are given by the interchange of $`M_P`$ and $`M_V`$ and the subscript which denotes which of the quark lines the gauge field is coupled to (*i.e.* for $`g_1(0)`$ the appropriate expression sets $`m_3`$=$`m_1`$). Defining $`g_{VP\gamma }=2[Q_1g_1+Q_2g_2]`$, and summing over initial and final polarizations gives
$$|A|^2=\frac{2\alpha \pi }{3}M_V^4\left[1\frac{M_P^2}{M_V^2}\right]^2g_{VP\gamma }^2,$$
(25)
where $`\alpha =\frac{1}{137}`$ is the fine structure constant. Standard techniques yield the total rate
$$\mathrm{\Gamma }_{VP\gamma }=\frac{\alpha }{24}M_V^3g_{VP\gamma }^2\left(1\frac{M_P^2}{M_V^2}\right)^3.$$
(26)
Electromagetic decays have been the subject of several theoretical studies. As well the decay $`J/\psi \eta _c+\gamma `$ has been measured. Table 11 shows our predictions for $`g_{VP\gamma }`$ along with the single experimental result and the predictions of some other models. In and two different heavy quark approaches were used. The quark confinement model , which has some similarity to the quark-meson model used in this work, gives results which are quite close to ours.
There are measured branching ratios for the $`D^{}`$ decays, however no lifetime measurement has been made. Therefore our predictions (which do not include the lifetime) cannot be compared directly with experiment. In order to obtain branching ratios a theoretical estimate of the lifetime must be used. The quark confinement model is the ideal choice, since its predictions are closest to our work. Using the results from and our predictions for the total rates (obtained from (26)) the following branching ratios are obtained:
$`BR\left[(D^{})^0D^0+\gamma \right]`$ $`=`$ $`33.0\%,`$
$`BR\left[(D^{})^+D^++\gamma \right]`$ $`=`$ $`1.43\%.`$
These compare well with the experimental values
$`BR_{expt.}\left[(D^{})^0D^0+\gamma \right]`$ $`=`$ $`38.1\%,`$
$`BR_{expt.}\left[(D^{})^+D^++\gamma \right]`$ $`=`$ $`1.1\%.`$
In order to treat the electromagnetic decay $`B_c^{}B_c+\gamma `$ the mass of the $`B_c^{}`$ meson must be specified. Theoretical estimates indicate that the mass difference should be small; $`M_{B_c^{}}M_{B_c}<100`$ MeV. To examine the effect of a small change in the $`B_c^{}`$ mass, the self energy and coupling constant were calculated over a range of masses. These results were used to calculate $`g_{B_c^{}B_c\gamma }`$ and are displayed in Table 12. The decay rate is shown in Fig. 4.
The radiative decay of $`B_c^{}`$ has not been studied extensively. A QCD sum rule approach , using $`M_{B_c^{}}=6.6`$ GeV and $`M_{B_c}=6.3`$ GeV, gives the result $`g_{B_c^{}B_c\gamma }^{SR}=0.270\pm 0.095GeV^1`$. Using these masses our prediction is $`g_{B_c^{}B_c\gamma }=0.2196GeV^1`$. The two values are in agreement.
## 6 Conclusion
A Lagrangian which models mesons in terms of an effective non-local quark-meson interaction vertex was extended in this paper to describe mesons such as $`B_c^{}`$ composed of two heavy quarks. The model has the advantage of treating all quarks (heavy and light) on the same footing, thereby permitting a unified investigation of heavy $``$ heavy, and heavy $``$ light quark decays. The model does not provide a complete dynamical description of quark interactions, meson masses can not be calculated and must be introduced as input parameters. Quark masses and the parameters associated with the quark-meson vertex were determined by fitting the pseudoscalar and vector meson decay constants $`f_P`$ and $`f_V`$. Due to lack of confinement in this approach some light vector mesons had to be excluded from the analysis.
To both test the model, and demonstrate its versatility, a large number of semileptonic decays of $`K`$, $`D`$, and $`B`$ mesons were analysed. Agreement with measured results and other theoretical approaches was good.
The main focus of this work was on the analysis of the semileptonic decays of the doubly-heavy $`B_c`$ meson. The form factors characterizing the strong interactions of the $`B_c`$ system were computed over the entire available range of momentum transfer. Using these results, decay rates and branching ratios were computed for all eight decay going to mesonic ground states.
A comparison with some other approaches highlighted the significant differences among various model predictions concerning the $`B_c`$. On very general grounds one would expect the decays to $`J/\psi `$ and $`(B_s^{})^0`$ states to be the most important. However, there is no agreement from models which of the channels dominates and absolute rates differ by a factor of 3 to 4. This is in contrast to the situation in K, D and B meson decays where much smaller differences between different models are found.
As a further illustration of the versatility of the model, the electromagnetic decays $`VP+\gamma `$ were investigated. A reasonable description of $`J/\psi `$ and $`D^{}`$ radiative decay was found and the rate for $`B_c^{}B_c+\gamma `$ was calculated for a range of $`B_c^{}`$ masses.
There is room for further analysis within the model considered in this work. Hadronic decays, such as $`B_c^+J/\psi +\pi ^+`$ can also be treated. A detailed analysis of all of these decays, combined with the results for the leptonic and semileptonic decays, could be used to make a prediction for the lifetime $`\tau _{B_c}`$. A further area that needs work is the difficulty with confinement. The natural solution to this problem appears to be provided by the quark confinement model . Recently Ivanov *et al.* have proposed a modification of the quark confinement model which may aid its application to heavy mesons .
*Acknowledgements.* We would like to thank M.A. Ivanov and P. Santorelli for a helpful communication. This work was supported in part by the Natural Sciences and Engineering Research Council of Canada.
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# Linear repetitivity, I. Uniform subadditive ergodic theorems and applications
## 1. Introduction
In a recent paper, Lagarias and Pleasants studied linearly and densely repetitive tilings . It was shown that these structures are diffractive and they proposed to consider linearly repetitive tilings as models of “perfectly ordered quasicrystals.”
In fact, several special classes of linearly repetitive tilings have attracted much attention. One such class is given by tilings arising from primitive substitutions. They have been studied in several contexts , including random Schrödinger operators and lattice gas models. Both the study of lattice gas models and the study of random Schrödinger operators require a uniform subadditive ergodic theorem. The appropriate theorem has been established in . In the one-dimensional case there is another important class of examples of linearly repetitive structures, namely, Sturmian dynamical systems whose rotation number has bounded continued fraction expansion. Again, this class allows for a uniform subadditive ergodic theorem. This has been shown by one of the authors (cf. for applications). These results immediately raise the following question:
* Does linear repetitivity imply a uniform subadditive ergodic theorem?
This question is answered in the affirmative by Theorem 1 in Section 3 of this paper (cf. as well). This theorem generalizes the theorem of . Moreover, combined with the known linear repetitivity of tilings generated by primitive substitution , it gives a conceptual proof for the subadditive ergodic theorem of . Of course, this theorem also implies an additive ergodic theorem. However, this additive ergodic theorem is not as effective as the corresponding theorem of , as it does not contain an error estimate (cf. Section 3).
We defer discussion of the methods used in the proofs of our results to the corresponding sections. However, we would like to emphasize the following perspective in our considerations: Our point of view is a purely local one. Thus, the key object of our studies is neither a tiling nor a species of tilings but rather certain sets of pattern classes. The appropriate sets are defined in Definition 2.1 and termed admissible. The advantage of this point of view is twofold. Firstly, in this approach, the uniformity of results is built in as the local structure is uniform for all tilings in the species. Secondly, the role of asymptotic translation invariance appearing in the subadditive ergodic theorems is clarified (cf. Section 4).
The article is organized as follows. In Section 2 we review basic facts on tilings and fix some notation. Section 3 contains a rather general form of a subadditive ergodic theorem. This is the main result of this paper. It gives an affirmative answer to Question (Q). In Section 4 we specialize the main theorem to various situations. This recovers several known (sub)additive ergodic theorems. Finally, in Section 5, we sketch applications of the foregoing results in the study of random operators associated to tilings.
## 2. Preliminaries
The aim of this section is to introduce certain notions and to fix some notation.
Consider a set consisting of subsets of $`^d`$ which are homeomorphic to the closed unit ball in $`^d`$ and pairwise disjoint up to their boundaries. Such a set of sets will be called a pattern if it is finite. It will be called a tiling (of $`^d`$) if the union of its elements equals the whole space. Its elements will be called tiles. For certain applications, it is useful to consider decorated tiles and patterns. A pattern with decorations from a set $`\mathrm{\Gamma }`$ is a set $`M`$ of pairs $`m=(a_m,c_m)`$ with $`a_m^d`$ and $`c_m\mathrm{\Gamma }`$ such that $`\{a_m:mM\}`$ is a pattern. One should think of a pair $`(a_m,c_m)`$ as a tile colored or decorated by $`c_m`$. The following definitions apply to both patterns and decorated patterns. However, to avoid tedious repetitions, they are phrased in terms of patterns only. If a pattern $`M`$ is contained in a pattern or tiling $`N`$, we write $`MN`$ and say that $`M`$ is a subpattern of $`N`$. Similarly, if a tile $`t`$ belongs to a pattern or tiling $`N`$, we write $`tN`$. For a pattern $`M`$, we define the underlying set $`s(M)`$ by
$$s(M)=_{mM}m^d.$$
The inner radius $`r_{\mathrm{in}}(M)`$ of a pattern $`M`$ is defined by
$$r_{\mathrm{in}}(M)=\mathrm{max}\{r:x^d,K(x,r)s(M)\},$$
and the outer radius $`r_{\mathrm{out}}(M)`$ of a pattern $`M`$ is defined by
$$r_{\mathrm{out}}(M)=\mathrm{min}\{r:x^d,K(x,r)s(M)\},$$
where $`K(x,r)`$ denotes the closed ball around $`x`$ with radius $`r`$. The existence of the minimum and maximum in question follows by compactness of $`s(M)`$. For a pattern $`M`$ and a closed set $`B`$ homeomorphic to the unit ball with $`s(M)B`$, define the restriction $`MB`$ of $`M`$ to $`B`$ by
(1)
$$MB=\{mB:mM,m\mathrm{int}(B)\mathrm{}\}.$$
In the applications we have in mind, $`B`$ will be either a box (cf. Section 3) or a closed ball.
There exists a natural equivalence relation on the set of patterns. Two patterns are equivalent if and only if they agree up to translation. The class of a pattern will also be called a pattern class or an abstract pattern. Similarly, an abstract tile is the class of a tile up to translation. The relations “$``$” and “$``$” (resp., the functions $`r_{\mathrm{in}}`$ and $`r_{\mathrm{out}}`$) give rise to relations (resp., functions) on abstract patterns in the obvious way. The induced relations (resp., functions) will be denoted by the same symbols. Similarly, concepts such as connectedness of patterns, disjointness, or distance of tiles in patterns, etc. can easily be carried over to abstract patterns. This will tacitly be done in the sequel, whenever necessary. Moreover, we will sometimes omit the word abstract in abstract patterns if no confusion can arise.
Our point of view is a purely local one. Thus, the following definition introduces the main object of our studies.
###### Definition 2.1.
A set $`𝒫`$ of abstract patterns in $`^d`$ is called admissible if it satisfies the following conditions.
* $`P𝒫`$, $`QP`$ implies $`Q𝒫`$.
* There exist $`0<r_{\mathrm{min}},r_{\mathrm{max}}<\mathrm{}`$ with $`r_{\mathrm{min}}r_{\mathrm{in}}(a)r_{\mathrm{out}}(a)r_{\mathrm{max}}`$ for all abstract tiles $`a𝒫`$.
* Let $`P𝒫`$ with representative $`\dot{P}`$ with $`0\dot{P}`$ and $`r>0`$ be given. Then, there exists a $`Q𝒫`$ with representative $`\dot{Q}`$ with $`K(0,r)s(\dot{Q})`$ and $`\dot{P}\dot{Q}`$.
In the sequel we will be exclusively concerned with admissible sets $`𝒫`$. For each admissible set, there is a natural set of tilings associated with it. Conversely, to a tiling of $`^d`$, one can associate a set of abstract patterns. This is the content of the next definition.
###### Definition 2.2.
(a) Let $`T`$ be a tiling of $`^d`$. The set $`𝒫(T)`$ of abstract patterns associated to $`T`$ is defined to be the set of classes of subpatterns of $`T`$.
(b) Let $`𝒫`$ be an admissible set of abstract patterns. A tiling $`T`$ is said to be associated to $`𝒫`$ if $`𝒫(T)𝒫`$.
(c) Let $`𝒫`$ be an admissible set of abstract patterns. The set of all tilings $`T`$ associated to $`𝒫`$ with the topology induced by the metric
$$\mathrm{d}(T,S)=inf\{ϵ:TB(0,\frac{1}{ϵ})=(S+t)B(0,\frac{1}{ϵ}),t^d,tϵ\}$$
is a topological space denoted by $`\mathrm{\Omega }(𝒫)`$ (cf. ).
This article is centered around the notion of linear repetitivity. This notion has been studied in for Delone sets in $`^d`$. In our context it is given in the following definition.
###### Definition 2.3.
An admissible set $`𝒫`$ of patterns is called linearly repetitive if there exists a constant $`c_{\mathrm{LR}}>0`$ such that every $`P𝒫`$ with $`r_{\mathrm{out}}(P)1`$ is contained in every $`Q𝒫`$ with $`r_{\mathrm{in}}(Q)c_{\mathrm{LR}}r_{\mathrm{out}}(P)`$.
An important property of the tiling space associated to a linearly repetitive $`𝒫`$ is the following:
###### Proposition 2.4.
If the admissible $`𝒫`$ is linearly repetitive, then $`\mathrm{\Omega }(𝒫)`$ is compact.
Proof. This follows by rather standard arguments once it is realized that linear repetitivity implies finiteness of the number of pattern classes with a prescribed maximal outer radius. For the reader’s convenience, we include a proof in Appendix A. $`\mathrm{}`$
Let us finish this section by discussing the role of Delone sets and the Voronoi construction in our context. Recall that a subset $`D`$ of $`^d`$ is called a Delone set if there exist positive constants $`r_0`$ and $`r_1`$ such that each ball in $`^d`$ of radius at least $`r_1`$ contains a point of $`D`$ and each ball of radius at most $`r_0`$ does not contain more than one point of $`D`$. The Voronoi construction assigns to each $`x`$ in a given Delone set $`D`$ the set $`V(x)=\{y^d:\mathrm{dist}(x,y)\mathrm{dist}(z,y),zD\}`$, where $`\mathrm{dist}(,)`$ denotes Euclidean distance. Then, $`V(D)=\{V(x):xD\}`$ is a tiling of $`^d`$ by convex polytopes (cf. ). Proposition 5.2 of says that $`𝒫(V(D))`$ is admissible for a Delone set $`D`$. Thus, Delone sets give rise to admissible sets. This motivates the following definition.
###### Definition 2.5.
A Delone set $`D`$ is called linearly repetitive if $`𝒫(V(D))`$ is linearly repetitive.
###### Remark 1.
Using the following proposition, it is not hard to show that this definition of linear repetitivity for Delone sets agrees with the definition of .
###### Proposition 2.6.
Let $`D`$ be a Delone set. Then for each $`xD`$, the tile $`V(x)`$ is determined by the points of $`D`$ lying inside a ball of radius $`2r_1`$ around $`x`$.
Proof. This is just Corollary 5.1 in $`\mathrm{}`$
## 3. The Main Theorem
This section is devoted to a proof of a rather general uniform subadditive ergodic theorem. The proof is similar to that of , which in turn uses ideas of (cf. for further details). The formulation relies on patterns on boxes. Thus, we will start this section with a discussion of boxes.
A box $`B`$ in $`^d`$ is a subset of the form $`B=\{(x_1,\mathrm{},x_d):a_jx_jb_j,j=1,\mathrm{},d\}`$, where $`a_j<b_j`$ for each $`j`$. The length of the $`j`$-th side is denoted by $`l_j`$, that is, $`l_j=b_ja_j`$. The volume and the surface area of a box $`B`$ are denoted by $`|B|`$ and $`\sigma (B)`$, respectively. Moreover, let the width $`\omega (B)`$ of a box $`B`$ be defined by $`\omega (B)=\mathrm{min}\{l_j:j=1,\mathrm{},d\}`$. For $`r^+`$, an $`r`$-box is a box whose sidelengths satisfy
$$rl_j2r,j=1,\mathrm{},d.$$
The set of all boxes (resp., $`r`$-boxes) is denoted by $`(^d)`$ (resp., $`(r)`$). A box-pattern (resp., $`r`$-pattern) is a pattern $`M`$, where $`s(M)`$ is a box (resp., $`r`$-box). For a box $`B`$ and a pattern (or tiling) $`M`$ with $`s(M)B`$, the box-pattern derived from $`M`$ by restricting to $`B`$ denoted by $`MB`$ has been defined in (1).
Now, let an admissible set of abstract patterns $`𝒫`$ be given. The set $`𝒫_\mathrm{b}`$ of abstract box-patterns derived from $`P`$ consists of all abstract patterns $`Q`$ which have representatives $`\dot{Q}`$ of the form $`\dot{Q}=\dot{P}B`$, where $`B`$ is a box and $`\dot{P}`$ is a representative of $`P𝒫`$. If $`B`$ is an $`r`$-box, the abstract pattern $`Q`$ is called an abstract $`r`$-pattern. The set of all abstract $`r`$-patterns derived from $`𝒫`$ is denoted by $`𝒫(r)`$. Moreover, let $`𝒫(\mathrm{})`$ be defined by
$$𝒫(\mathrm{})=\underset{r>0}{}𝒫(r).$$
The functions $`l_j`$, $`||`$, $`\sigma `$, and $`\omega `$ induce functions on $`𝒫_\mathrm{b}`$ in the obvious way, which will be denoted by the same symbols. The inclusion relation $``$ on the set of boxes induces a relation on $`𝒫_\mathrm{b}`$, again denoted by $``$. That is, the relation $`PQ`$ for $`P,Q𝒫_\mathrm{b}`$ holds if and only if there exist boxes $`B_P,B_Q`$ and representatives $`\dot{P},\dot{Q}`$ of $`P`$ and $`Q`$, respectively, with $`s(\dot{P})=B_P`$, $`s(\dot{Q})=B_Q`$ and $`\dot{P}=\dot{Q}B_P`$. Similarly, the equation
$$P=\underset{j=1}{\overset{n}{}}P_j$$
for $`P,P_j𝒫_\mathrm{b}`$, $`j=1,\mathrm{},n`$ is defined to hold if and only if there exist representatives $`\dot{P}`$ of $`P`$ and $`\dot{P}_j`$ of $`P_j`$, $`j=1,\mathrm{},n`$, with
$$s(\dot{P})=\underset{j=1}{\overset{n}{}}s(\dot{P}_j).$$
Here, the equation $`B=_{j=1}^nB_j`$ for boxes $`B,B_j`$, $`j=1,\mathrm{},n`$ is defined to hold if and only if the $`B_j`$ are pairwise disjoint up to their boundaries and their union is $`B`$. Equations of the form $`P=_{j=1}^nP_j`$ (resp., $`B=_{j=1}^nB_j`$) are called decompositions or partitions of patterns (resp., boxes). The notion of linear repetitivity appropriate to box-patterns is contained in part (ii) of the next proposition.
###### Proposition 3.1.
Let $`𝒫`$ be admissible and $`𝒫_\mathrm{b}`$ be as above. Then the following are equivalent:
(i) $`𝒫`$ is linearly repetitive, that is, there exists a constant $`c_{\mathrm{LR}}`$ such that every $`Q𝒫`$ with $`r_{\mathrm{out}}(Q)1`$ is contained in every $`P𝒫`$ with $`r_{\mathrm{in}}(P)c_{\mathrm{LR}}r_{\mathrm{out}}(Q)`$.
(ii) There exists a constant $`c_{\mathrm{LR},\mathrm{b}}`$ such that every $`P𝒫(r)`$ with $`r1`$ is contained in every $`Q𝒫(c_{\mathrm{LR},\mathrm{b}}r)`$.
Proof. This is straightforward. $`\mathrm{}`$
We can now introduce the class of subadditive functions.
###### Definition 3.2.
Let $`𝒫`$ be admissible.
(a) A function $`F:𝒫_\mathrm{b}`$ is called subadditive if there exist nonnegative constants $`d_\mathrm{F}`$ and $`r_\mathrm{F}`$ and a nonincreasing function $`c_\mathrm{F}:[r_\mathrm{F},\mathrm{})`$ with $`lim_r\mathrm{}c_\mathrm{F}(r)=0`$ such that
* $`F(P)_{j=1}^nF(P_j)+_{j=1}^nc_\mathrm{F}(\omega (P_j))|P_j|`$ for $`P=_{j=1}^nP_j`$ with $`\omega (P_j)r_\mathrm{F}`$,
* $`|F(P)|d_\mathrm{F}|P|`$.
(b) A function $`F:𝒫_\mathrm{b}\{\pm \mathrm{}\}`$ is called additive if both $`F`$ and $`F`$ are subadditive.
We will be interested in means of subadditive functions. Our main result states the existence of a certain limit of means of this kind. These means are introduced in the next definition.
###### Definition 3.3.
Let $`F`$ be subadditive on $`𝒫`$. For $`rr_F`$ the means $`F^+(r)`$ and $`F^{}(r)`$ are defined by
$$F^+(r)=sup\{\frac{F(P)}{|P|}:P𝒫(r)\},F^{}(r)=inf\{\frac{F(P)}{|P|}:P𝒫(r)\}.$$
The following proposition is well known. In the context of subadditive functions on Delone sets, it was proved in . For the convenience of the reader, we include a sketch of the proof.
###### Proposition 3.4.
Let $`F`$ be a subadditive function. Then, the following equation holds,
$$\underset{r\mathrm{}}{lim}F^+(r)=\underset{rr_\mathrm{F}}{inf}\{F^+(r)+c_\mathrm{F}(r)\}.$$
Proof. Denote the infimum by $`\overline{F}`$. We show (i) $`\overline{F}lim\; inf_r\mathrm{}F^+(r)`$ and (ii) $`lim\; sup_r\mathrm{}F^+(r)\overline{F}`$.
(i) This is clear by
$$\underset{r\mathrm{}}{lim\; inf}F^+(r)=\underset{r\mathrm{}}{lim\; inf}\left(F^+(r)+c_\mathrm{F}(r)\right)\underset{rr_\mathrm{F}}{inf}\{F^+(r)+c_\mathrm{F}(r)\}.$$
(ii) Fix an arbitrary $`r_0r_\mathrm{F}`$. Now, every $`P𝒫(r)`$ with $`rr_0`$ arbitrary can be written as a sum of patterns in $`𝒫(r_0)`$. The subadditivity condition together with a short calculation then implies
(2)
$$\frac{F(P)}{|P|}F^+(r_0)+c_\mathrm{F}(r_0).$$
As $`P𝒫(r)`$ was arbitrary, equation (2) implies
$$F^+(r)F^+(r_0)+c_\mathrm{F}(r_0)$$
for all $`rr_0`$. This proves (ii) and finishes the proof of the proposition. $`\mathrm{}`$
We can now prove the main theorem of this section.
###### Theorem 1.
Let $`𝒫`$ be admissible and linearly repetitive and let $`F`$ be subadditive on $`𝒫_\mathrm{b}`$. Then the limits $`lim_r\mathrm{}F^+(r)`$ and $`lim_r\mathrm{}F^{}(r)`$ exist and are equal. In particular, the equation
$$\underset{r\mathrm{}}{lim}F^+(r)=\underset{|P|\mathrm{},P𝒫(\mathrm{})}{lim}\frac{F(P)}{|P|}$$
is valid.
Proof. This is proved by contraposition. So, assume $`lim\; inf_r\mathrm{}F^{}(r)<lim\; sup_r\mathrm{}F^+(r)`$. Thus, by Proposition 3.4 there exist a $`\delta >0`$, a sequence $`n(k)`$ with $`n(k)\mathrm{}`$ for $`k\mathrm{}`$, and $`Q_k𝒫(n(k))`$ with
(3)
$$\frac{F(Q_k)}{|Q_k|}F^+(n(k))\delta .$$
W.l.o.g. we can assume $`n(k)r_\mathrm{F}`$. Choose an arbitrary $`k`$ and consider some arbitrary $`P𝒫(3c_{\mathrm{LR},\mathrm{b}}n(k))`$. Here $`c_{\mathrm{LR},\mathrm{b}}`$ is as defined in Proposition 3.1. Let $`\dot{P}`$ be an arbitrary representative of $`P`$ with underlying box $`B=s(\dot{P})`$. By partitioning each side of $`B`$ into three parts of equal length, the box $`B`$ can be decomposed into $`3^d`$ congruent smaller boxes, all belonging to $`(c_{\mathrm{LR},\mathrm{b}}n(k))`$. There is only one of these smaller boxes which does not intersect the boundary of $`B`$. Call it $`B_{\mathrm{int}}`$. The decomposition of $`B`$ into smaller boxes induces a decomposition of $`\dot{P}`$ into $`(c_{\mathrm{LR},\mathrm{b}}n(k))`$-patterns. Denote by $`\dot{P}_{\mathrm{int}}`$ the pattern with $`s(\dot{P}_{\mathrm{int}})=B_{\mathrm{int}}`$. By linear repetitivity, $`\dot{P}_{\mathrm{int}}`$ contains a representative $`\dot{Q}_k`$ of $`Q_k`$. As the distance of $`B_{\mathrm{int}}`$ to the boundary of $`B`$ is bigger than or equal to $`c_{\mathrm{LR},\mathrm{b}}n(k)`$, the same is true for the distance of $`s(\dot{Q}_k)`$ to the boundary of $`B`$. Thus, $`B`$ can be written as
$$B=\underset{j=0}{\overset{n}{}}B_j$$
with suitable $`B_j(n(k))`$, $`j=1,\mathrm{},n`$, and $`B_0=s(\dot{Q}_k)`$. This induces a decomposition of $`P`$ of the form $`P=_{j=0}^nP_j`$ with $`P_j𝒫(n(k))`$, $`j=1,\mathrm{},n`$, and $`P_0=Q_k`$. By subadditivity of $`F`$ and $`(\text{3})`$ this implies
$`{\displaystyle \frac{F(P)}{|P|}}`$ $``$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{F(P_j)}{|P_j|}}{\displaystyle \frac{|P_j|}{|P|}}+{\displaystyle \frac{F(Q_k)}{|Q_k|}}{\displaystyle \frac{|Q_k|}{|P|}}+{\displaystyle \underset{j=0}{\overset{n}{}}}c_\mathrm{F}(\omega (P_j)){\displaystyle \frac{|P_j|}{|P|}}`$
$``$ $`{\displaystyle \underset{j=0}{\overset{n}{}}}F^+(n(k)){\displaystyle \frac{|P_j|}{|P|}}\delta {\displaystyle \frac{|Q_k|}{|P|}}+c_\mathrm{F}(n(k))`$
$``$ $`F^+(n(k))\delta {\displaystyle \frac{|Q_k|}{|P|}}+c_\mathrm{F}(n(k)).`$
Here we used the bound $`\frac{F(P_j)}{|P_j|}F^+(n(k))`$, valid for arbitrary $`P_j𝒫(n(k))`$. Since $`Q_k`$ belongs to $`𝒫(n(k))`$ and $`P`$ belongs to $`𝒫(3c_{\mathrm{LR},\mathrm{b}}n(k))`$, we can estimate
$$\frac{|Q_k|}{|P|}\frac{(n(k))^d}{(23c_{\mathrm{LR},\mathrm{b}}n(k))^d}=\frac{1}{(6c_{\mathrm{LR},\mathrm{b}})^d}.$$
Putting all this together, we arrive at
$$\frac{F(P)}{|P|}F^+(n(k))\frac{1}{(6c_{\mathrm{LR},\mathrm{b}})^d}\delta +c_\mathrm{F}(n(k)).$$
Since $`P𝒫(3c_{\mathrm{LR},\mathrm{b}}n(k))`$ was arbitrary, this implies
$$F^+(3c_{\mathrm{LR},\mathrm{b}}n(k))F^+(n(k))\frac{1}{(6c_{\mathrm{LR},\mathrm{b}})^d}\delta +c_\mathrm{F}(n(k)).$$
As this holds for arbitrary $`k`$, we can now take the limit on both sides using Proposition 3.4 and obtain $`\overline{F}\overline{F}\delta \frac{1}{(6c_{\mathrm{LR},\mathrm{b}})^d}`$, a contradiction. This finishes the proof. $`\mathrm{}`$
As a corollary we get an additive ergodic theorem. This is our version of Theorem 4.1 of . Note, however, that we are not able to estimate the convergence rate (cf. Remark 2 below).
###### Corollary 3.5.
Let $`𝒫`$ be linearly repetitive and let $`F`$ be an additive function on $`𝒫_\mathrm{b}`$. Then the following equation holds,
$$\underset{r\mathrm{}}{lim}F^+(r)=\underset{\omega (P)\mathrm{},P𝒫_\mathrm{b}}{lim}\frac{F(P)}{|P|}.$$
Proof. The decomposition technique of the proof of Proposition 3.4 applied to the subadditive function $`F`$ gives
(4)
$$\underset{\omega (P)\mathrm{},P𝒫_\mathrm{b}}{lim\; sup}\frac{F(P)}{|P|}\underset{r\mathrm{}}{lim\; inf}F^+(r).$$
Since $`F`$ is subadditive as well, this equation immediately implies
(5)
$$\underset{\omega (P)\mathrm{},P𝒫_\mathrm{b}}{lim\; sup}\frac{F(P)}{|P|}\underset{r\mathrm{}}{lim\; inf}(F)^+(r).$$
Multiplying by $`(1)`$ and using $`(F)^+(r)=F^{}(r)`$, we get
(6)
$$\underset{\omega (P)\mathrm{},P𝒫_\mathrm{b}}{lim\; inf}\frac{F(P)}{|P|}\underset{r\mathrm{}}{lim\; sup}F^{}(r).$$
By $`(\text{4})`$, $`(\text{6})`$, and the foregoing theorem, the corollary follows. $`\mathrm{}`$
###### Remark 2.
It is not possible to derive an estimate on the rate of convergence in the subadditive theorem. This can be seen from the following example. Let $`f:^+[0,\mathrm{}]`$ be an arbitrary monotonically decreasing function. Let $`𝒫_\mathrm{b}`$ be an arbitrary set of box-patterns derived from an admissible $`𝒫`$. Define $`F:𝒫_\mathrm{b}`$ by $`F(P)=|P|f(|P|)`$. As $`f`$ is decreasing, the function $`F`$ is subadditive. Moreover, we have $`\frac{F(P)}{|P|}=f(|P|)`$. Since $`f`$ was an arbitrary decreasing function, this shows that the rate of convergence in the subadditive ergodic theorem cannot be estimated.
###### Remark 3.
It appears that uniform subadditive ergodic theorems are special features of linearly repetitive structures. To support this, in the appendix we will exhibit examples of strictly ergodic structures for which a uniform subadditive ergodic theorem does not hold. The examples will be given by Sturmian subshifts whose rotation number has rapidly increasing continued fraction coefficients.
## 4. Specializing the Main Theorem
In this section we derive various corollaries from the subadditive ergodic theorem. First, we consider (sub)additive functions on boxes on a conrete tiling or Delone set. We then use our methods to give a direct proof of the (known) unique ergodicity of dynamical systems arising from linearly repetitive tilings. Finally, we discuss how the theorems of fit into our context.
Let us first introduce the appropriate notion of subadditivity and translation invariance.
###### Definition 4.1.
Let $`w`$ be a function on the set $`B(^d)`$ of all boxes in $`^d`$.
(a) The function $`w`$ is called subadditive if there exists a constant $`r_w`$ and a nonincreasing function $`c_w:[r_w,\mathrm{})`$ with $`lim_r\mathrm{}c_w(r)=0`$ such that
* $`w(B)_{j=1}^nw(B_j)+_{j=1}^nc_w(\omega (B_j))|B_j|`$ for $`B=_{j=1}^nB_j`$, with $`\omega (B_j)r_w`$,
* $`|w(B)|d_w|B|`$.
(b) Let $`T`$ be a Delone set or a tiling. The function $`w`$ is called asymptotically $`T`$-invariant if there exists a constant $`r_w`$ and a nonincreasing function $`e_w:[r_w,\mathrm{})`$ with $`lim_r\mathrm{}e_w(r)=0`$ such that
* $`|w(B)w(B+t)|e_w(\omega (B))|B|`$ if $`(BT)+t=(B+t)T`$ and $`\omega (B)r_w`$.
We can now easily derive subadditive theorems for functions on Delone sets or tilings.
###### Corollary 4.2.
Let $`T`$ be a linearly repetitive tiling in $`^d`$. Let $`w`$ be a subadditive, asymptotically $`T`$-invariant function. Then for every sequence $`B_n`$ with $`B_n(r_n)`$ and $`r_n\mathrm{}`$, the limit
$$\underset{n\mathrm{}}{lim}\frac{w(B_n)}{|B_n|}$$
exists and is independent of the sequence.
Proof. The strategy of the proof is simple. We will construct a subadditive function $`F=F_w`$ on $`𝒫_\mathrm{b}(T)`$ and show that the limit in question equals the limit of $`F^+(r)`$, whose existence is guaranteed by Theorem 1.
Define $`F`$ on $`𝒫_\mathrm{b}(T)`$ by
$$F(P)=sup\{w(s(\dot{P})):\dot{P}=Ts(\dot{P}),\dot{P}\text{ representative of }P\}.$$
By properties (i) and (ii) of $`w`$, the function $`F`$ is subadditive on $`𝒫_\mathrm{b}(T)`$. By construction of $`F`$, we have
$$\frac{w(B)}{|B|}F^+(r)$$
for $`B(r)`$ with $`rr_w`$ arbitrary. This immediately implies
(7)
$$\underset{n\mathrm{}}{lim\; sup}\frac{w(B_n)}{|B_n|}\underset{r\mathrm{}}{lim\; sup}F^+(r).$$
Moreover, by property (iii) of $`w`$, we have
$$\frac{w(B)}{|B|}+e_w(\omega (B))F^{}(r)$$
for $`B(r)`$. This yields
(8)
$$\underset{n\mathrm{}}{lim\; inf}\frac{w(B_n)}{|B_n|}\underset{r\mathrm{}}{lim\; inf}F^{}(r).$$
By $`(\text{7})`$, $`(\text{8})`$, and Theorem 1, the statement of the theorem follows. $`\mathrm{}`$
###### Corollary 4.3.
Let $`D`$ be a linearly repetitive Delone set in $`^d`$. Let $`w`$ be a subadditive, asymptotically $`D`$-invariant function. Then for every sequence $`B_n`$ with $`B_nB(r_n)`$ and $`r_n\mathrm{}`$, the limit
$$\underset{n\mathrm{}}{lim}\frac{w(B_n)}{|B_n|}$$
exists and is independent of the sequence.
Proof. This follows from the foregoing corollary applied to the a colored version of the Voronoi construction $`V(D)`$ (cf. Section 2). Here, each tile in $`V(D)`$ is colored by the unique point of $`D`$ in its interior. To emphasize the coloring, we denote the colored tiling by $`V(D,C)`$. The coloring implies that the function $`w`$ is asymptotically $`V(D,C)`$-invariant. Thus, the result follows from the foregoing corollary. $`\mathrm{}`$
Let us now discuss two classes of examples of the above theorems. They are given by tilings arising from primitive substitutions and tilings arising from Sturmian dynamical systems whose rotation number has bounded continued fraction expansion.
We start by considering primitive substitutions. They give rise to linearly repetitive tilings . Thus, we immediately get the following result.
###### Corollary 4.4.
Let $`S`$ be a primitive substitution and let $`T`$ be a tiling associated to $`𝒫(S)`$ with vertex set $`E`$. Let $`w`$ be an asymptotically $`E`$-invariant subadditive function on boxes in $`^d`$. Then, the limit $`lim_n\mathrm{}\frac{w(B_n)}{|B_n|}`$ exists for every sequence $`B_n`$ of boxes with $`B_nB(r_n)`$ and $`r_n\mathrm{}`$, and it is independent of the sequence.
This is essentially the subadditive ergodic theorem of . The theorem of is slightly more general in that the sequences $`(B_n)`$ considered there are only required to be cube-like van Hove sequences. On the other hand, the notion of subadditivity used there is more restrictive than the notion used here. There, $`w`$ is required to satisfy a subadditivity condition on unions of quite general disjoint (up to their boundary) sets with the constant $`c_w`$ being zero. In fact, under these assumptions, one should be able to extend our theorem to hold for arbitrary cube-like van Hove sequences. However, our theorem is good enough to cover the desired applications.
The other example is given by certain Sturmian dynamical systems; see Appendix B for some background. As shown in , a Sturmian dynamical system is linearly repetitive if and only if its rotation number has bounded continued fraction expansion. Thus, we immediately obtain the following corollary of Theorem 1 which generalizes Theorem 2 of (cf. ) as well).
###### Corollary 4.5.
Let an irrational $`\alpha (0,1)`$ with bounded continued fraction expansion be given. Let $`𝒲(\alpha )`$ be the set of pattern classes of the Sturmian dynamical system with rotation number $`\alpha `$ (cf. for details). Then for every subadditive function $`F`$ on $`𝒲(\alpha )`$, the limit $`lim_n\mathrm{}\frac{F(w_n)}{|w_n|}`$ exists for every sequence $`(w_n)`$ with $`|w_n|`$ going to infinity. Moreover, the limit is independent of the sequence.
Of course, one could use Corollary 3.5 instead of Theorem 1 to obtain an additive ergodic theorem. However, this kind of result falls clearly short of the additive theorem of , as it does not allow one to estimate the rate of convergence. This has been discussed in Remark 2 in Section 3.
We close this section by sketching a direct derivation of the unique ergodicity of dynamical systems associated to linearly repetitive tilings. In fact, the result uses only the compactness of the underlying space and an additive ergodic theorem. Thus, the proof applies verbatim to more general systems. The “inner box” technique given below applies to several contexts (cf. for further discussion). It will be used in the next section as well.
###### Corollary 4.6.
Let $`𝒫`$ be linearly repetitive. Then the tiling dynamical system $`(\mathrm{\Omega }(𝒫),^d)`$ is uniquely ergodic. Here, $`^d`$ acts on $`\mathrm{\Omega }(𝒫)`$ in the canonical way via translation.
Proof. We have to show that for any continuous $`f`$ on $`\mathrm{\Omega }(𝒫)`$, the limits $`\frac{1}{B}_Bf(Tt)𝑑t`$ converge uniformly in $`T`$ for $`\omega (B)`$ going to infinity. The strategy is similar to the proof of Corollary 4.2 above. We will associate to $`f`$ additive functions $`F_{\mathrm{sup}}`$ and $`F_{\mathrm{inf}}`$ on $`𝒫`$. They are defined as follows:
$$F_{\mathrm{sup}}(P)=sup\{_{s(\dot{P})}f(Tt)𝑑t:Ts(\dot{P})=\dot{P}\},$$
$$F_{\mathrm{inf}}(P)=inf\{_{s(\dot{P})}f(Tt)𝑑t:Ts(\dot{P})=\dot{P}\},$$
where $`\dot{P}`$ is an arbitrary representative of $`P`$ (it is not hard to check that these definitions are independent of the actual choice of $`\dot{P}`$). Apparently,
(9)
$$F_{\mathrm{sup}}\left(\underset{j=1}{\overset{n}{}}P_j\right)\underset{j=1}{\overset{n}{}}F_{\mathrm{sup}}(P_j),F_{\mathrm{inf}}\left(\underset{j=1}{\overset{n}{}}P_j\right)\underset{j=1}{\overset{n}{}}F_{\mathrm{inf}}(P_j).$$
Moreover, the following is valid,
(10)
$$|F_{\mathrm{sup}}(P)F_{\mathrm{inf}}(P)|o(\omega (P)),$$
where the little $`o`$ function only depends on the continuity properties of $`f`$. To prove (10) we use an “inner box” argument. Recall that $`f`$ is continuous and thus uniformly continuous since $`\mathrm{\Omega }(𝒫)`$ is compact by Proposition 2.4. Therefore, for each $`ϵ>0`$, there exists $`R`$ such that $`|f(T)f(S)|ϵ`$ whenever $`TB(0,R)=SB(0,R)`$. This implies that for all $`ts(\dot{P})`$ with $`\mathrm{dist}(t,(s(\dot{P})^c)R`$, the difference of the integrands $`|f(Tt)f(St)|`$ is smaller than $`ϵ`$. For large enough $`\omega (P)`$, the set of those $`t`$ agrees with the size of $`P`$ up to a boundary term. This proves (10). By (9) and (10), the functions $`F_{\mathrm{sup}}`$ and $`F_{\mathrm{inf}}`$ are additive. Thus, the additive ergodic theorem implies the existence of the limits $`lim_{\omega (P)\mathrm{}}\frac{F_{\mathrm{sup}}(P)}{|P|}`$ and $`lim_{\omega (P)\mathrm{}}\frac{F_{\mathrm{inf}}(P)}{|P|}`$. By (10), the limits are equal and the corollary follows. $`\mathrm{}`$
## 5. Applications
In this section we consider applications to random operators associated to tilings. In this context, there are two important quantities whose existence is established by a subadditivity argument, namely, the Lyapunov exponent in the one-dimensional case and the integrated density of states in arbitrary dimensions.
The existence of the integrated density of states for Schrödinger-type operators associated to primitive substitutions is thoroughly discussed in . The discussion given there relies on abstract operator theory together with a subadditive ergodic theorem. Thus, it gives essentially the existence of the integrated density of states for Schrödinger-type operators associated to arbitrary linearly repetitive structures. In fact, the argument of can be improved and strengthened in several respects . In particular, it turns out that the existence proof can actually be reduced to an additive ergodic theorem. This is interesting due to the existence of an error estimate in the additive ergodic theorem. This might have useful applications.
Let us be more precise. For a tiling or pattern $`M`$, the space $`l^2(M)`$ is defined to be the space of all square summable sequences indexed by the elements of $`M`$. Let $`A`$ be a selfadjoint operator on a linearly repetitive tiling $`T`$ with matrix elements $`A(x,y)`$ for $`x,yT`$. (Here, the tiling $`T`$ is called linearly repetitive if $`𝒫(T)`$ is linearly repetitive.) We will assume that $`A`$ satisfies the following finite range (FR) and invariance (I) properties: There exists some $`R0`$ with
* $`A(x,y)`$ vanishes for $`\mathrm{dist}(x,y)R`$.
* The value of $`A(x,y)`$ is completely determined by the pattern class $`[\{tT:\mathrm{dist}(t,\{x,y\})R\}]`$.
In fact, the invariance condition implies that the operator $`A`$ can be defined on every tiling $`T`$ of the species $`\mathrm{\Omega }(T)`$. To emphasize this, we will sometimes write $`A(T)`$ for the manifestation of $`A`$ on $`l^2(T)`$.
For a box $`B`$ in $`^d`$, the restriction $`A(T)|_B`$ of $`A`$ to $`B`$ is the operator on $`l^2(BT)`$ with matrix elements
$$A(T)|_B(\stackrel{~}{x},\stackrel{~}{y})=A(x,y),\text{for }\stackrel{~}{x}=xT\text{ and }\stackrel{~}{y}=yT\text{ }.$$
For a box $`B`$ in $`^d`$ and $`\lambda `$, define the function $`k_\lambda ^T(B)`$ by
$$k_\lambda ^T(B)=\frac{1}{|B|}\mathrm{\#}\{\lambda _n:\lambda _n\lambda ,\lambda _n\text{eigenvalues of}A(T)|_B\},$$
where the number of elements of a finite set $`S`$ is denoted by $`\mathrm{\#}S`$. Then the following holds.
###### Theorem 2.
The limit $`lim_{\omega (B)\mathrm{}}k_\lambda ^T(B)`$ exists and is independent of $`T`$. In fact, the convergence (in $`\omega (B)`$) is uniform in $`T`$.
Proof (sketch). By Corollary 4.2 it is enough to show that the map $`Bk_\lambda ^T(B)`$ is translation-invariant and additive. But this follows from the finite range condition together with the invariance condition. Details can be found in . $`\mathrm{}`$
Theorem 2 generalizes the corresponding theorem of , where Penrose tilings are considered. Moreover, it only relies on an additive ergodic theorem, whereas uses a subadditive theorem.
Let us now turn to the study of the Lyapunov exponent. The sketch below follows the detailed discussion of the Sturmian case in . An admissible set $`𝒫`$ of abstract patterns in one dimension over a finite set of tiles can easily be identified with a set $`𝒲`$ consisting of finite words over a finite alphabet $`A`$. The study of one-dimensional Schrödinger operators associated to $`𝒲`$ can be based on the study of the so-called transfer matrices. For each $`E`$, the transfer matrix $`M(E)`$ gives a map $`M(E):𝒲\mathrm{SL}(2,)`$, defined by $`M(E)(w)=T(E,w_n)\times \mathrm{}\times T(E,w_1)`$ for $`w=w_1\mathrm{}w_n`$, where for $`a`$ and $`E`$, the matrix $`T(E,a)`$ is defined by
(11)
$$T(E,a)=\left(\begin{array}{cc}Ea& 1\\ 1& 0\end{array}\right).$$
This map is antimultiplicative if the operation on $`𝒲`$ is standard concatenation of words. Since the standard norm $``$ on $`\mathrm{SL}(2,)`$ is submultiplicative, the function
$$F:𝒲,F(w)=\mathrm{ln}M(w)$$
is subadditive. Thus, the results of Section 3 give the following theorem.
###### Theorem 3.
Let $`𝒲`$ and $`F`$ be as above. If $`𝒲`$ is linearly repetitive, then the limit $`lim_n\mathrm{}\frac{F(w_n)}{|w_n|}`$ exists for each sequence $`(w_n)`$ in $`𝒲`$ with $`|w_n|`$ going to infinity and the limit does not depend on the sequence.
The limit in the theorem is called the Lyapunov exponent. It plays an important role in the study of one-dimensional Schrödinger operators. The theorem applies in particular to systems arising from primitive substitutions and to Sturmian dynamical systems whose rotation number has bounded continued fraction expansion. This is due to the fact that these systems are linearly repetitve, as discussed in Section 4 (cf. and as well). Thus, the theorem generalizes the corresponding theorems of and .
Let us close this section by pointing out that there is a theory of lattice gas models for tilings arising from primitive substitutions . This theory is built upon the subadditive theorem of . Thus, it is very likely that considerable portions of it can be carried over to gas models on linearly repetitive tilings.
Acknowledgments. D. D. was supported by the German Academic Exchange Service through Hochschulsonderprogramm III (Postdoktoranden) and D. L. received financial support from Studienstiftung des Deutschen Volkes (Doktorandenstipendium), both of which are gratefully acknowledged.
## Appendix A Compactness of Linearly Repetitive Tiling Spaces
In this section we sketch a proof of Proposition 2.4. It consists of two steps, namely, establishing a finiteness condition and performing a standard diagonalization procedure; compare .
Let $`𝒫`$ be a linearly repetitive admissible set of abstract patterns. We want to show that $`\mathrm{\Omega }(𝒫)`$ is compact. Observe first that for every $`r0`$, there are only finitely many pattern classes $`P𝒫`$ with $`r_{\mathrm{out}}(P)r`$. To see this, consider any pattern class $`Q`$ such that $`K(0,c_{\mathrm{LR}}r)s(\dot{Q})`$ for some representative $`\dot{Q}`$ of $`Q`$. Delete from $`\dot{Q}`$ all the tiles which have empty intersection with $`K(0,c_{\mathrm{LR}}r)`$ and call the resulting pattern $`\dot{Q}^{}`$ and its pattern class $`Q^{}`$. It is clear that $`\dot{Q}^{}`$ has finite volume and that $`Q^{}`$ contains every abstract pattern $`P𝒫`$ with $`r_{\mathrm{out}}(P)r`$. This proves the assertion.
Let us now consider a sequence $`(T_n)_n`$ in $`\mathrm{\Omega }(𝒫)`$. It suffices to prove that $`(T_n)_n`$ has a convergent subsequence. To find this subsequence, we will inductively define sequences $`(T_n^m)_n`$, $`m`$, such that $`(T_n^1)_n`$ is a subsequence of $`(T_n)_n`$ and for $`m_1m_2`$, $`(T_n^{m_1})_n`$ is a subsequence of $`(T_n^{m_2})_n`$. This will be done in a way such that $`(T_n^n)_n`$ converges. Choose any monotonically increasing sequence $`r_m\mathrm{}`$. Essentially, we will force the sequence $`(T_n^m)_n`$ to converge on $`K(0,r_m)`$. It is then obvious from the definition of $`d(,)`$ that the diagonal sequence $`(T_n^n)_n`$ will be $`d`$-Cauchy with obvious limit tiling.
To define the refinement $`(T_n^m)_n`$ of $`(T_n^{m1})_n`$ (think of $`(T_n)_n`$ as $`(T_n^0)_n`$), we will proceed in two steps. First, consider the pattern $`\dot{P}_n`$ of tiles in $`T_n^{m1}`$ having nonempty intersection with $`K(0,r_m)`$. The patterns $`\dot{P}_n`$ have outer radius bounded by $`r_m+2r_{\mathrm{max}}`$ (with $`r_{\mathrm{max}}`$ from Definition 2.1) and hence, by the above observation, their abstract pattern classes $`P_n`$ belong to a finite set. Hence, one of them, say $`P`$, occurs infinitely often. Delete all the tilings from the sequence $`(T_n^{m1})_n`$ which have $`P_nP`$. By the Selection Theorem , the remaining sequence has a subsequence such that the corresponding sets $`\dot{P}_{n_k}`$ converge with respect to standard Hausdorff metric. Call this sequence $`(T_n^m)_n`$. By the above remarks, it is easy to see that $`(T_n^n)_n`$ is a convergent subsequence of $`(T_n)_n`$.
## Appendix B Strictly Ergodic Subshifts for Which the Uniform Subadditive Ergodic Theorem Fails
In this section we present one-dimensional examples which show that our main result fails if we only require strict ergodicity rather than linear repetitivity. We will consider a standard symbolic form of Sturmian tilings of the real line, that is, we study two-sided sequences over the alphabet $`A=\{0,1\}`$; see for background on Sturmian sequences.
Let us first recall some standard notation. Given a finite alphabet $`A`$, we denote by $`A^{}`$ the set of finite words over $`A`$ and by $`A^{}`$ (resp., $`A^{}`$) the set of one-sided (resp., two-sided) sequences over $`A`$, both called infinite words. Given a finite or infinite word $`w`$, we denote by $`\mathrm{Sub}(w)`$ the set of all finite subwords of $`w`$. Finally, given two finite words $`v,w`$, $`\mathrm{\#}_v(w)`$ denotes the number of occurrences of $`v`$ in $`w`$.
Fix some irrational $`\alpha (0,1)`$ and define the words $`s_n`$ over the alphabet $`A`$ by $`s_1=1`$, $`s_0=0`$, $`s_1=s_0^{a_11}s_1`$, and $`s_n=s_{n1}^{a_n}s_{n2}`$, $`n2`$, where the $`a_n`$ are the coefficients in the continued fraction expansion of $`\alpha `$. By definition, for $`n2`$, $`s_{n1}`$ is a prefix of $`s_n`$. Therefore, the following (“right”-) limit exists in an obvious sense, $`c_\alpha =lim_n\mathrm{}s_nA^{}`$.
Define the associated set of pattern classes $`𝒲(\alpha )A^{}`$ by $`𝒲(\alpha )=\mathrm{Sub}(c_\alpha )`$. The associated symbolic dynamical system $`(\mathrm{\Omega }(\alpha ),T)`$ is then given by $`\mathrm{\Omega }(\alpha )=\{xA^{}:\mathrm{Sub}(x)𝒲(\alpha )\}`$ and $`(Tx)_n=x_{n+1}`$. $`(\mathrm{\Omega }(\alpha ),T)`$ is strictly ergodic for every irrational $`\alpha `$. It is linearly repetitive if and only if the sequence $`(a_n)_n`$ is bounded.
We will prove the following theorem.
###### Theorem 4.
There exist $`\alpha (0,1)`$ irrational and a subadditive function $`F`$ on $`𝒲(\alpha )`$ with the following property: There exist sequences $`(w_n^k)_n`$ in $`𝒲(\alpha )`$, $`k=1,2`$, such that $`|w_n^k|\mathrm{}`$ as $`n\mathrm{}`$, $`k=1,2`$, and
$$\underset{n\mathrm{}}{lim\; sup}\frac{F(w_n^1)}{|w_n^1|}<\underset{n\mathrm{}}{lim\; inf}\frac{F(w_n^2)}{|w_n^2|}.$$
In particular, the limit $`lim_{|w|\mathrm{}}\frac{F(w)}{|w|}`$ does not exist, that is, the uniform subadditive ergodic theorem does not hold for $`𝒲(\alpha )`$.
The following properties of the words $`s_n`$ are well known and will be useful in the proof of Theorem 4.
###### Proposition B.1.
(i) For all $`n2`$, the word $`s_n`$ is a prefix of the word $`s_{n1}s_n`$.
(ii) For every $`n`$, there is no nontrivial occurrence of $`s_n`$ in $`s_ns_n`$, that is, $`s_ns_n=w_1s_nw_2`$ implies $`w_1=\epsilon `$ or $`w_2=\epsilon `$.
We are now in position to give the
Proof of Theorem 4. Define the function $`G`$ on $`𝒲(\alpha )`$ by
$$G(w)=\underset{n=1}{\overset{\mathrm{}}{}}\mathrm{\#}_{s_{n1}s_n}(w)(|s_{n1}|+|s_n|).$$
It is clear that all but finitely many of the terms are zero. Moreover, it is obvious that $`G`$ is superadditive. Thus, by Theorem 2 of , $`\overline{G}=lim_n\mathrm{}\frac{G(s_n)}{|s_n|}`$ exists, but it is possibly infinite. Observe that $`G(s_n)`$ only depends on the numbers $`a_1,\mathrm{},a_n`$. Hence, by (using Proposition B.1)
$$\frac{G(s_{n+1})}{|s_{n+1}|}\frac{a_{n+1}G(s_n)+G(s_{n1})+a_{n+1}_{i=1}^{n1}\left(|s_{i1}|+|s_i|\right)}{a_{n+1}|s_n|},$$
we see that we can force $`\overline{G}`$ to be finite if $`a_n\mathrm{}`$ sufficiently fast. But then we have
$$\frac{G(s_{n1}s_n)}{|s_{n1}s_n|}=\frac{G(s_{n1})}{|s_{n1}s_n|}+\frac{G(s_n)}{|s_{n1}s_n|}+\frac{|s_{n1}s_n|}{|s_{n1}s_n|}\frac{G(s_n)}{|s_n|\left(1+\frac{|s_{n1}|}{|s_n|}\right)}+1,$$
that is, $`\frac{G(s_{n1}s_n)}{|s_{n1}s_n|}`$ does not converge to $`\overline{G}`$. We can therefore conclude by setting $`F=G`$, $`w_n^1=s_{n1}s_n`$, and $`w_n^2=s_n`$. $`\mathrm{}`$
###### Remark 4.
The proof of Theorem 4 actually provides an uncountable set of numbers $`\alpha `$ such that the uniform subadditive ergodic theorem fails for $`𝒲(\alpha )`$. This set, however, has Lebesgue measure zero. By a more sophisticated argument (see ), one may prove this result for all $`\alpha `$’s obeying $`_{n=1}^{\mathrm{}}\frac{1}{a_na_{n+1}}<\mathrm{}`$. Since this set still has Lebesgue measure zero (cf. methods in ), it may be interesting to establish results for the Lebesgue-generic set of $`\alpha `$’s with intermediate $`(a_n)`$-behavior.
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# Spin-orbit and exchange interaction in surface quantum wells on gapless semimagnetic semiconductor HgMnTe
## 1 Introduction
There are two aspects that cause the peculiar features of, and the interest in, the two-dimensional (2D) electron gas in narrow-gap diluted magnetic (semimagnetic) semiconductors (DMS). One stems from the $`s,pd`$ exchange interaction between band electrons and localized magnetic moments . This interaction results in changed spin splitting of the band states, which can be varied by such external factors as a magnetic field and temperature. The other is due to the peculiarities inherent to Kane semiconductors with small gap, leading to such relativisticlike effects as non-parabolicity, kinetic confinement (motional binding ), spin-orbit (SO) splitting and resonant interband mixing by surface electric field . What is important in the context of the specificity of 2D electronic systems involving DMS, is that both exchange and SO interaction cause the rearrangement of the spin structure of Landau levels (LL’s). Although historically the first studies of 2D electron gas in DMS were performed for metal-insulator-semiconductor (MIS) structures based on HgMnTe, the available experimental results are mostly for the grain boundaries in HgMnTe and HgCdMnTe with positive Kane gap $`E_g>100`$ meV at Mn content typically $`x=0.02`$ (for the higher $`x`$, the exchange interaction exhibited itself poorly, that was attributed to antiferromagnetic interaction betwen Mn ions) . This is due to low electron mobility in previously investigated MIS structures. At the same time, the inversion layers in MIS structures are of particular experimental interest because of the possibility of controlling the depth of surface quantum well by gate voltage, and because of the relative ease and accuracy of the description of surface potential (in the case of the bicrystals, the additional poorly verified assumptions have to be used to describe a self-consistent potential near grain boundaries ). An important point is that these results can be compared with the data for MIS structures based on narrow-gap HgCdTe , being non-magnetic analogue of narrow-gap DMS.
As regards their theoretical description, the subband calculations were carried out only for DMS with direct but not inverted bands and without allowance for spinorlike effects . However, the SO splitting in asymmetrical quantum wells at zero magnetic field (this phenomenon in itself is of much current interest ) leads to the rearrangement of subband magnetic levels. In narrow-gap semiconductors, the perturbation of magnetic spectrum is so drastic that SO interaction cannot be neglected in the theoretical treatment. It must be stressed that SO splitting, as we shall see in Sec. III, exceeds by far exchange interaction contribution so it cannot be considered as a correction to the exchange interaction. It is clear also that a treatment based on the quasiclassical quantization in a magnetic field of subband spectrum (calculated at $`B=0)`$ is unsuitable for the description of exchange interaction effects. A more rigorous theoretical consideration of the LL’s structure is required.
In this paper, the peculiarities of 2D electron gas due to exchange and SO interaction are studied in inversion layers on Hg<sub>1-x</sub>Mn<sub>x</sub>Te. with a small Mn content. At $`x<0.08`$ HgMnTe has inverted bands (i.e. becomes semimetal) and traditional galvanometric methods cannot be used because of the shunting of surface conductance by the bulk. We employed the magnetocapacitance spectroscopy method, which is applicable to semiconductors with any sign of the Kane gap. The parameters of the samples and the experimental data relating to the capacitance oscillations versus gate voltage and magnetic field and their temperature evolution are presented in Sec. II. In Sec. III, we present the theoretical model. The treatment of LL’s in 2D subbands is based on the further development of the concept we offered previously for the description of subband spectrum at $`B=0`$. The density of states (DOS) in a magnetic field is described neglecting the mixing between LL’s and assuming a Gaussian shape of each level. In this section an analytical expression for oscillations of the differential capacitance of space charge region in the low-temperature range is also obtained using WKB approach. In Sec. IV, the results of the computer modeling of capacitance oscillations are presented. The results of a comparison of the experimental data and theoretical calculations for different temperatures and parameters of exchange interaction are discussed. The parameters of the broadening of LL’s are determined from a fitting of the amplitudes of calculated oscillations to their experimental values. The dominant mechanisms of the scattering responsible for broadening of LL’s are discussed.
## 2 Samples and experimental results
### 2.1 Samples and experimental methods
In this work the inversion layers in MIS structures fabricated from $`p`$-Hg<sub>1-x</sub>Mn<sub>x</sub>Te single crystals were investigated. No impurities were introduced intentionally and no post growth annealing was performed. Deviations from stoichiometry, and thus the type of the conductivity, were controlled by mercury partial pressure during the growth process. The Hall-effect measurements were performed at variable temperatures and magnetic field strengths. After removing the Hall and tunnel contacts (see below), the substrates were mechanically polished and etched in a 0.5 bromine-methanol solution. Several methods such as anodic oxide formation, silicon oxide and Al<sub>2</sub>O<sub>3</sub> deposition, and the Langmuir-Blodgett film technique have been used for forming an insulating film in MIS structures. The gate electrodes of the typical area $``$ $`5\times 10^4`$ cm<sup>2</sup> were formed evaporating Pb. The differential capacitance $`C`$ and derivative $`dC/dV_g`$ on gate voltage $`V_g`$ of the capacitors were measured in the dark, typically at 1 MHz and a test signal amplitude of 5 mV.
The capacitance magnetooscillations due to the magnetic quantization of 2D electron gas were observed in all the above HgMnTe MIS structures. It was shown that the general shapes of the oscillations at the same carrier surface density and Mn content are similar. In the following we present the results for the structures with a $`80`$ nm thick anodic oxide film, grown in a solution of 0.1M KOH in 90% ethylene glycol / 10% H<sub>2</sub>O at $`0.1`$ mA$``$cm<sup>-2</sup>. There are several reasons for such a choice: (i) the amplitudes of oscillations in these structures are the highest owing to the large value of insulator capacitance (this is caused by the large value of dielectric constant of the anodic oxide), (ii) the highest surface carrier densities are achieved at low gate voltages $`V_g=1015`$ V, and (iii) the dielectric constant of oxide is close to that of a semiconductor so the contribution of image forces in surface potential can be neglected in the calculations.
We investigated samples with different Mn content ($`x=0.024,0.040,0.060`$ and $`0.1`$). Kane gap $`E_g`$ and Kane effective mass $`m_b`$ (and therefore $`x`$) were determined independently by the tunnel spectroscopy method for a comparison of band parameters in the bulk with those in the vicinity of the surface. The discrepancy is within the accuracy of the analysis ($`\mathrm{\Delta }x0.002÷0.003`$). Because the tunnel contacts and studied MIS capacitors were produced using identical technology and differ only by thickness of insulator (Langmuir-Blodgett film or an oxide) this agreement testifies that the surface layers are chemically close to the bulk. The similarity of the results for the structures with different insulators (with different fabrication methods) supports this conclusion. The fact that in the small surface concentration range the measured cyclotron masses in 2D subbands extrapolate to the bulk value $`m_b`$ is direct evidence of an absence of noticeable decomposition in 2D layer during the structure fabrication process.
In the following we shall restrict our consideration to the results for HgMnTe with $`x0.04`$ ($`Eg=100\pm 5`$ meV). The amplitudes of the oscillations for other samples are much less even at 4.2 K and rapidly decrease with increasing temperature. In the case of $`x0.024`$ this is caused by the small cyclotron energy due to a large Kane gap. In the cases $`x0.06`$ and $`x0.1`$ it is due to the large doping level of available materials. As a result, we could not obtain reasonably accurate information about the oscillation’s temperature evolution, in which the specificity of DMS is manifested. As to the measurements at $`T=4.2`$ K, the subband parameters extracted from oscillations for these samples are similar to those for HgCdTe with the same band parameters and agree well with the theory. On the other hand, the samples with $`x=0.04`$ are best suited to the purpose of this first study aimed at investigating the peculiarities of 2D electron gas in DMS with inverted bands, in which (i) the effects of SO and exchange interaction are expected to be more clearly pronounced and (ii) the results can be compared with those for well studied surface layers on gapless HgCdTe with $`E_g(50÷100)`$ meV . For small gap $`\left|E_g\right|<100`$ meV the parameters of 2D subbands depend only weakly on $`E_g`$ (except in the case of small subband occupancies). By contrast, the subband parameters are more sensitive to the doping level. For this reason, we present the results for two samples with $`N_AN_D=1.2\times 10^{16}`$ cm<sup>-3</sup> (sample $`S1`$) and $`N_AN_D=1.5\times 10^{17}`$ $`cm^3`$ (sample $`S2`$). The Hall mobility of holes is of the order of 2500 cm<sup>2</sup>/Vs at 4.2 K. The experimental data are presented below for capacitors with the gate area $`S=7.2\times 10^4`$ cm<sup>2</sup> and insulator capacitance $`C_{ox}=136.5`$ pF for sample $`S1`$, and with $`S=7.7\times 10^4`$ and $`C_{ox}=155.1`$ pF for sample $`S2`$.
### 2.2 Capacitance measurements in perpendicular magnetic fields
Fig. 1 shows the capacitance-voltage characteristics at $`T=4.2`$K in magnetic field $`B=4.5`$ T perpendicular to the 2D layer for the sample $`S2`$. $`C(V_g)`$ characteristics are those of typical low-frequency behavior. This is due to the absence of a gap between conduction and valence (heavy hole) band in gapless semiconductors. As a result, 2D electrons in inversions layers are in equilibrium with the ac ripple and contribute predominantly to the measured capacitance under inversion band bending. The low-frequency conditions with respect to the minority carriers are satisfied in all the investigated frequency range 30 kHz$`÷`$5 MHz. The wide hysteresis loop and the dependence of $`C(V_g)`$ characteristics on the rate of voltage sweep are observed. The capacitance changes in time because of flat-band voltage shift $`\mathrm{\Delta }V_{fb}`$, which is close to logarithmic in time. The time constant is of the order of a few minutes and is almost independent of the temperature. The hysteresis effects point to charge tunnel exchange between the semiconductor and the slow traps in insulator . A history dependence and instability are manifested in all the investigated HgMnTe-based MIS structures. Such behavior is contrary to that of HgCdTe and HgTe-based structures with the same insulators.
The voltage dependence of the charge density $`eN_s(V_g)`$ induced in inversion layer is sublinear. This is well demonstrated by the non-equidistance of quantum oscillations of the capacitance $`C(V_g)`$ (see Fig. 1). The tunneling of electrons from the 2D layer into oxide causes a saturation of $`N_s(V_g)`$ dependence at $`V_gV_{fb}`$ $`(10÷15)`$ V. As a result, the $`N_s`$ range accessible for investigations is limited by the value $`(3÷4)\times 10^{12}`$ cm<sup>-2</sup> (in HgCdTe, values of $`N_s`$ up to 10<sup>13</sup> cm<sup>-2</sup> can be obtained). Although the hysteresis effects hamper the measurements, the discussed physical results are not affected by the instability of band bending. Such instability is caused by the transient processes but not by degradation. In order to assure the stability of band bending during the measurement of $`C(B)`$ oscillations, the sample was held at given capacitance (or voltage) for 5-15 minutes. The identity of $`C(B)`$ plots registered at increasing and decreasing magnetic field (i.e. at different times) was examined for each $`C(B)`$ curve. When the temperature (or angle) dependencies of $`C(B)`$ oscillations were measured, the long term stability was checked by the repetitive measurement of initial (for given measurement cycle) $`C(B)`$ plots (see Fig. 5). Although the $`C(V_g)`$ and $`dC/dV_g(V_g)`$ characteristics are history dependent, they are completely repeatable, if the voltage range, rate and direction of the sweep are the same (see Fig. 6).
As can be seen in Fig. 1, the magnitudes of capacitance in the $`C(V_g)`$ oscillation extrema, corresponding to the same LL’s number, are the same for the curves with different $`V_{fb}`$. $`C(B)`$ oscillations (and consequently subband occupancy and surface potential) measured at the same magnitude of capacitance in a zero magnetic field $`C(0)`$ are also identical no matter what the voltage (the value of the latter for any given $`C(0)`$ is determined by the flat-band voltage, which is a history- and time-dependent). When the dc gate voltage (or flat-band voltage at the same $`V_g`$) is changed, the filling of interface states is also changed but does not respond to the ac ripple, i.e., the interface states do not contribute to the capacitance. This takes place for all frequencies and temperatures and testifies that the high-frequency conditions with respect to interface states are satisfied . Thus there is “one-to-one correspondence” between $`C(0)`$, band bending and surface density of 2D electrons $`N_s=`$ $`N_i`$ ($`i`$ is the 2D subband number). Owing to this we can reproduce $`C(B)`$ curves (characterized by the $`C(0)`$ values) at all times. The subband parameters are presented below as functions of $`N_s`$. Contrary to dependencies vs. $`V_g`$, these dependencies are not affected by the hysteresis effects or any specific parameters of MOS capacitors and are common to the given HgMnTe sample. It may be noted that the hysteresis has some positive points also. We have the possibility to investigate 2D electron gas in the same surface quantum well on the same sample but with a different interface charge. Particularly, it is important for the investigation of scattering mechanisms.
Typical $`C(B)`$ oscillations for both samples at almost equal $`N_s`$ are presented in Fig. 2 together with their $`1/B`$ Fourier transforms. As in the case of gapless HgCdTe , the individual spin components have not been observed in the oscillations at any $`N_s`$ even for lowest LL’s. On the other hand, the oscillation beats and the Fourier spectra demonstrate distinctly the presence of two frequencies connected with the SO splitting of each 2D subband. The additional structure in Fig. 2 near the Fourier lines corresponding to $`i=0`$ subband for sample $`S2`$ results from mixed harmonics. This structure is suppressed if the range of small magnetic fields (in which the oscillations relevant to the higher subbands are dominant) is excluded from Fourier analysis. For a sample $`S1`$, the intensities of $`i=1`$ and $`i=2`$ lines are small, and the above features are practically indistinguishable in the scale of Fig. 2. The surface densities in the spin-split subbands $`N_i^+`$ and $`N_i^{}`$ determined from Fourier transforms are plotted in Fig. 3 as functions of $`N_s`$. The carrier distribution among 2D subbands is different for the two samples. The concentrations $`N_s`$, corresponding to the “starts” of excited subbands, increase with increasing doping level and agree well with the theoretical calculations, in which the bulk values of $`N_AN_D`$ are used. This fact also testifies that the disruption of stoichiometry in surface layers because of probable migration of atoms is insignificant. A discrepancy with theory is detectable only in the relative differences (splitting) of occupancies $`\mathrm{\Delta }N_i/N_i=(N_i^{}N_i^+)/(N_i^{}+N_i^+)`$ in the small $`N_s`$ range (see Fig. 4). Similar disagreement takes place for inversion layers on HgCdTe also. The possible reasons for such behavior are discussed in Ref. .
The intensities of Fourier lines for the high-energy branch $`I_i^+`$ and low-energy branch $`I_i^{}`$ are different (see Fig. 2 and Fig. 7). At low $`N_i`$, the intensities $`I_i^+`$ and $`I_i^{}`$ are close for $`i=0`$ subband, and $`I_i^{}>I_i^+`$ for excited subbands. The ratio $`I_i^{}/I_i^+`$ as a function of $`N_s`$ for $`i=0`$ is shown in Fig. 4. The scattering of the points is due mainly to the fact that $`N_i^\pm `$ values and especially the ratios $`I_i^{}/I_i^+`$ depend strongly (and non-monotonically) on a magnetic field range used in Fourier transform (for instance, it is clear that Fourier analysis will reveal no splitting for the range between beat nodes). Nevertheless, a decrease in $`I_i^{}/I_i^+`$ with increasing $`N_s`$ is clearly visible. It is valid for excited subbands also.
### 2.3 Measurements in tilted magnetic fields
Although there is no doubt that we are dealing with a 2D system (the existence of magnetooscillation effect in the capacitance and the observation of magnetooscillations versus gate voltage in itself testify to it), the fact experiments in tilted magnetic field were also performed. Some results for the sample $`S2`$ are presented in Fig. 5. The magnetic field positions of the oscillation extrema and the fundament fields in the Fourier spectra (to a smaller extent) vary only roughly as cosine of the angle $`\mathrm{\Theta }`$ between $`𝐁`$ and normal to the 2D layer. The clearly distinguishable deviations from such behavior are observed. Namely, the experimental angle dependencies are stronger.
There are several reasons for this deviation from classical cosine dependence, because a number of physical factors are ignored in the simplified model . Firstly, in the strictest sense, such behavior, even in the case of parabolic dispersion, is valid only for an ideal 2D system. A condition to be satisfied for cosine dependence is $`<r>/<z>1`$, where $`<r>`$ and $`<z>`$ are the mean sizes of wave function in the 2D plane and in the confinement direction. In the case of narrow-gap semiconductors, the width of surface quantum well is relatively large and such a strong requirement may be not fulfilled (note also that $`<z>`$ is energy dependent in this case). In strong magnetic field and at small surface concentration, the cyclotron radius and the width of 2D layer may be comparably-sized (especially, for excited subbands) and diamagnetic shift must weaken the angle dependence. This is contrary to the experimental behavior. Secondly, the cosine relation is obtained for spinless particles. This is not the case in a real system. Thirdly, the SO interaction is neglected in this simple consideration. Undoubtedly, the spinlike effects can reflect on a spectrum in a tilted magnetic field and modify the angle dependence. Lastly, the exchange interaction can also make an additional contribution to the deviation from simple angle dependence. This assumption has experimental support. For comparison we investigated the HgCdTe-based samples in a tilted magnetic field. Data for HgCdTe with $`E_g=95`$ meV and $`N_AN_D=2\times 10^{17}`$ cm<sup>-3</sup> are given in Fig. 5. Under the same conditions they also manifest a deviation from cosine behavior. However, the deviation is weaker and opposite in sign to the case of gapless HgMnTe. At the same time, the samples based on HgCdTe with $`E_g>0`$ show a deviation of the same sign as in HgMnTe, but smaller in magnitude. Contrary to HgCdTe samples, changes in the structure of oscillations are observed in HgMnTe inversion layers. Namely, the beat nodes in oscillations $`C(B_{})`$ ($`B_{}=B\mathrm{cos}\mathrm{\Theta }`$) are shifted to the lower LL’s numbers with increasing $`\mathrm{\Theta }`$ (i.e., with increasing total magnetic field $`B`$) (see Fig. 5).
These experimental observations testify that the behavior in tilted magnetic fields is markedly affected by both SO interaction (which essentially depends on $`E_g`$ sign, see Sec. III. and Fig. 8) and exchange interaction. For narrow-gap semiconductors, the theoretical analysis requires a consideration of spin from the outset. Strong SO and exchange interaction and resonant effects lead to serious complication of the theoretical description even for perpendicular orientation (see Sec. III.). The calculations in tilted magnetic fields are troublesome even for the simplest parabolic Hamiltonian with a $`k`$-linear Rashba term. At present we cannot make A reasonable theoretical analysis of effects in tilted fields we will restrict our consideration to the case of perpendicular orientation.
### 2.4 Temperature effects
The structure of oscillations and the subband parameters extracted from oscillations are identical to those in HgCdTe. No features due to exchange interaction are manifested. Because the exchange effects are determined by a magnetization and can be varied by the temperature, the investigation of temperature evolution of oscillations is of prime interest. The results for $`dC/dV_g(V_g)`$ and $`C(B)`$ oscillations are shown in Fig. 6 and Fig. 7 respectively. As it may be seen, no pronounced changes in the position of either $`dC/dV_g(V_g)`$ or $`C(B)`$ oscillations are observed. The shift of beat nodes to the high gate voltages and to low magnetic fields (to the greater LL’s numbers) with increasing temperature (and hence with decreasing magnetization) is the sole temperature effect, besides the usual diminution of oscillation amplitudes. (Notice, the direction of shift with the increasing magnetization is similar to the one observed with increasing total $`B`$ (at the same $`B_{})`$ in the tilted magnetic field experiments). This shift must be attributed to the features inherent in semimagnetic semiconductors because in HgCdTe-based structures neither the positions of the oscillations nor those of the beat nodes change with the temperature.
The analysis based on the Fourier transform of oscillations for different temperatures cannot yield any information about exchange interaction. On the other hand, such data cannot be obtained from spin splitting either because, as noted above, the separate spin components are not observed in the oscillations at any temperatures. Thus we must settle the question by capacitance magnetooscillations modeling.
## 3 Theoretical Consideration
### 3.1 Landau levels structure
The direct numerical solution of an initial matrix equation (especially in the case of non-zero magnetic field) faces obstacles of a fundamental nature, such as singularity problems and the ambiguity of boundary conditions on the interface and in the bulk, especially in the case of gapless semiconductors . In this work, we employ a concept based on the reduction of the matrix equation to a Schrodinger-like equation with an effective potential . At $`B=0`$, the results are quite close to those obtained by numeric calculations . The line of attack of Ref. seems to be a reasonable compromise between accuracy and the ease of the calculation for a non-zero magnetic field. This simplicity of method is of considerable advantage for the purposes of oscillations modeling.
Under homogeneous magnetic field $`𝐁(0,0,B)`$ parallel to the direction of the confinement (surface potential $`V=V(z))`$ a motion in the 2D plane can be quantized, using the mean field approximation for exchange interaction. It can be shown that in the framework of a six-band Kane model with an allowance for a magnetic field and exchange interaction, the following matrix equation determines the subband LL’s energy $`E_n(B)`$ (we do not write out the usual expressions, describing the behavior of envelopes in perpendicular to magnetic field 2D plane)
$$\left[\begin{array}{cccccc}E_{}+\alpha & \frac{E_B\sqrt{3(n1)}}{2}& \frac{E_B\sqrt{n}}{2}& 0& 0& s_b\mathrm{}\widehat{k}_z\\ \frac{E_B\sqrt{3(n1)}}{2}& E_++3\beta & 0& 0& 0& 0\\ \frac{E_B\sqrt{n}}{2}& 0& E_+\beta & s_b\mathrm{}\widehat{k}_z& 0& 0\\ 0& 0& s_b\mathrm{}\widehat{k}_z& E_{}\alpha & \frac{E_B\sqrt{3(n+1)}}{2}& \frac{E_B\sqrt{n}}{2}\\ 0& 0& 0& \frac{E_B\sqrt{3(n+1)}}{2}& E_+3\beta & 0\\ s_b\mathrm{}\widehat{k}_z& 0& 0& \frac{E_B\sqrt{n}}{2}& 0& E_++\beta \end{array}\right]\left(\begin{array}{c}f_1^{n1}(z)\hfill \\ f_3^{n2}(z)\hfill \\ f_5^n(z)\hfill \\ f_2^n(z)\hfill \\ f_4^{n1}(z)\hfill \\ f_6^{n1}(z)\hfill \end{array}\right)=0,$$
(1)
where $`E_\pm =E_nV\pm E_g/2,`$ $`s_b=\sqrt{\left|E_g\right|/2m_b}`$ is Kane velocity, $`n`$ is LL number. The envelopes $`f_{1,4}`$ correspond to $`\mathrm{\Gamma }_6`$ symmetry band, $`f_{3,6}`$ and $`f_{2,5}`$ to heavy and light branches of $`\mathrm{\Gamma }_8`$ band. “Magnetic energy” $`E_B=\sqrt{2m_bs_b^2\mathrm{}\omega _b}=\sqrt{2}s_b\mathrm{}/\lambda `$ ($`\mathrm{}\omega _b=\mathrm{}eB/m_bc`$ is cyclotron energy, $`\lambda =\sqrt{c\mathrm{}/eB}`$ is magnetic length) practically does not depend on band parameters because $`s_b`$ is practically the same for all Kane semiconductors . We denote $`\alpha =\frac{1}{2}xN\alpha ^{}S_z`$ and $`\beta =\frac{1}{6}xN\beta ^{}S_z`$, where $`x`$ is the MnTe mole fraction, $`N`$ is the number of unit cells per unit volume, $`\alpha ^{}`$ and $`\beta ^{}`$ are the exchange integrals for $`\mathrm{\Gamma }_6`$ and $`\mathrm{\Gamma }_8`$ bands respectively. The thermodynamically average $`S_z`$ of the $`z`$ component of a localized spin $`S`$ (for Mn<sup>2+</sup> ions $`S=\frac{5}{2}`$) can be described via normalized Brillouin function $`B_S\left(x\right)`$:
$$S_z=S(1x)^{18}B_S\left(\frac{2\mu _BB}{k_B(T+T_N)}\right),$$
(2)
where $`T_N`$ is effective temperature, arising from antiferromagnetic interaction between Mn<sup>2+</sup> ions . This factor defines the magnetic field and temperature dependency of exchange effects.
Resolving the systems (1) with respect to the components $`f_5`$ and $`f_6`$ we obtain the following set of two “Schrodinger-like” equations (for the envelopes $`\phi _5^n=f_5^n/\sqrt{H_n^+}`$ and $`\phi _6^{n1}=f_6^{n1}/\sqrt{H_n^{}})`$ for the description of the magnetic spectrum of 2D electrons in surface layers on DMS with inverted bands (called Kane “$`p`$\- electrons” as in Ref. )
$$\left|\begin{array}{cc}\frac{\mathrm{}^2\widehat{k}_z^2}{2m_b}E_{eff}+U^+& iU_{so}^+C_g^+s_b\mathrm{}\widehat{k}_z\\ iU_{so}^{}+C_g^{}s_b\mathrm{}\widehat{k}_z& \frac{\mathrm{}^2\widehat{k}_z^2}{2m_b}E_{eff}+U^{}\end{array}\right|\left(\begin{array}{c}\phi _5^n\hfill \\ \phi _6^{n1}\hfill \end{array}\right)=0$$
(3)
with effective energy
$$E_{eff}=(E^2m_b^2s_b^4)/2m_bs_b^2,$$
and effective potential
$$U^\pm =U_0+U_B^\pm +U_{exc}^\pm +U_R^\pm ),$$
in which we single out the following parts: -the spin independent “Klein-Gordon“ term
$$U_0=(V^22EV)/2m_bs_b^2,$$
and spin-like terms: “magnetic potential”
$$U_B^\pm =E_B^2[g^2nR^\pm +3(n\pm 1)(E_+\pm \beta )/(E_+\pm 3\beta )]/2m_bs_b^2,$$
”exchange potential”
$$U_{ex}^\pm =[\alpha \beta \pm (\alpha E_++\beta E_{})]/2m_bs_b^2,$$
and “resonant” term describing “spin-interband“ interaction, arising from the mixing of $`\mathrm{\Gamma }_6`$ and $`\mathrm{\Gamma }_8`$ bands by an electric field
$`U_R^\pm `$ $`=`$ $`{\displaystyle \frac{s_b^2\mathrm{}^2}{2m_bs_b^2}}\left[{\displaystyle \frac{3}{4}}\left({\displaystyle \frac{1+L_n^\pm }{H_n^\pm }}\right)^2+{\displaystyle \frac{L_n^\pm }{H_n^\pm (E_+\pm 3\beta )}}\right]\left({\displaystyle \frac{dV}{dz}}\right)^2`$
$`+{\displaystyle \frac{s_b^2\mathrm{}^2}{4m_bs_b^2}}{\displaystyle \frac{(1+L_n^\pm )}{H_n^\pm }}{\displaystyle \frac{d^2V}{dz^2}},`$
where the following designations are used:
$$L_n^\pm =3E_B^2(n\pm 1)/4(E_+\pm 3\beta )^2,H_n^\pm =E_{}\pm \alpha L_n^\pm (E_+\pm 3\beta ),$$
$$R_n^\pm =H_n^\pm /H_n^{},g=1$$
The spin-orbit terms and the coefficients at linear in $`\widehat{k}_z`$ terms are determined by
$$U_{so}^\pm =C_gs_b\mathrm{}\sqrt{R_n^{}}\left[\frac{1+L_n^\pm }{H_n^\pm }+\frac{1+L_n^{}}{2H_n^{}}(R_n^\pm 1)\right]\frac{dV}{dz},$$
$$C_g^\pm =C_g\sqrt{R_n^{}}(R_n^\pm 1),C_g=g\frac{E_B\sqrt{n}}{4m_bs_b^2}.$$
The second term in an expression for $`U_B^\pm `$ and the dimensionless parameter $`L_n^\pm `$ arise from the interaction with heavy hole branch. It must be stressed that the exchange interaction causes not only the appearance of an exchange term in the effective potential, but also a modification of the terms, describing the ”resonant” and SO interaction.
The LL’s of 2D electrons in Kane semiconductors with $`E_g>0`$ (“$`s`$-electrons”) are described by the same set (3) (for the envelopes $`\phi _2^n=f_2^n/\sqrt{H_n^+}`$ and $`\phi _1^{n1}=f_1^{n1}/\sqrt{H_n^{}}`$) but with $`g=+1`$, $`L_n^\pm =0`$, $`H_n^\pm =E_+\pm \beta `$. It can be shown that the equations for Dirac like electrons in a magnetic field are the same as for Kane “s- electrons” but with $`g=+2`$. It should be emphasized that unlike the $`B=0`$ case , the set (3) cannot be separated and reduced to independent equations for individual spin components because of the SO interaction $`U_{so}^\pm `$ (for “$`p`$\- electrons”, because of linear in $`\widehat{k}_z`$ terms also).
From this point on, we shall restrict our consideration to the semiclassical approximation just as in quantization of spectrum described by Eqs. (3) (it is clear that the use of the semiclassical approximation immediately in (1) results in the loss of spinorlike effects), so also in calculation of the surface potential $`V(z)`$. A semiclassically self-consistent potential in such treatment is calculated in the frame of quasirelativistic modification of the Thomas-Fermi method. The validity of such an approach in narrow-gap semiconductors was argued and demonstrated by a comparison with numerical self-consistent calculations in many papers (see also Ref. and references therein). Substituting as usual $`\phi _i^m=C_i^m\mathrm{exp}(ik_z(z)𝑑z)`$ and neglecting the proportional to $`i(zdk_z/dz+k_z)^2`$ terms (higher-order terms in the expansion of the action in powers of $`\mathrm{})`$ we obtain from (3) the quasiclassical expression for “spin-split” $`z`$components of wave vector
$$k_z^\pm =\frac{\sqrt{2m_bs_b^2}}{s_b\mathrm{}}\{K[K^2(E_{eff}U^+)(E_{eff}U^{})+U_{so}^+U_{so}^{}]^{\frac{1}{2}}\}^{\frac{1}{2}}$$
(4)
with
$$K=E_{eff}(U^++U^{})/2m_bs_b^2C_g^2(R_n^+1)(R_n^{}1).$$
Together with the Bohr-Sommerfeld quantization rule
$$_{V(z=0)}^{V(k_z=0)}k_z(E,V)\left(\frac{dV}{dz}\right)^1𝑑V=\pi (i+\frac{3}{4})$$
(5)
they define the magnetic levels $`E_n^\pm (i,B)`$ in surface quantum well $`V(z)`$.
For the Dirac-like electrons, the derived equations may be considered as a generalization of the result of Ref. that allows for magnetic quantization and spin-like effects. In the case of Kane “$`s`$ electrons“ and without exchange interaction, the Eqs. (3) and (4) coincide with the corresponding expressions in Ref. . As might be expected, at $`\alpha =0,\beta =0,n\mathrm{},E_B\sqrt{n\pm 1}`$ and $`E_B\sqrt{n}s_b\mathrm{}k_s`$ ($`k_s`$ is 2D wave vector) the Eq. (5) with $`k_z`$ from Eq. (4) is reduced to the corresponding equation in Ref. , describing subband dispersions $`E_i^\pm (k_s)`$ at $`B=0`$. Lastly, the Eq. (4) at $`V(z)=const`$ determines the Landau subbands $`E_n^\pm (B,k_z)`$ in the bulk of Kane (at $`\alpha =0`$ and $`\beta =0`$ non-magnetic) semiconductors.
It must be noted that the proportional to $`C_g^2`$ term in expression for $`K`$ (arising from linear in $`\widehat{k}_z`$ terms in (3)) has little effect in comparison with the spin-orbit term in Eq. (4). However, ignoring this term in calculations at $`U_{so}^\pm =0`$ can introduce large error up to the inversion of the order of spin sublevels. In pseudo-ultrarelativistic limit $`E_g=0`$ and without exchange interaction, the spectra in magnetic fields for both “$`s`$ ” and “$`p`$-electrons” are scale invariant with respect to the surface band bending $`\mu _s`$ as well as in the case of $`B=0`$ . The results for ground subband are plotted in Fig. 8 in dimensionless (normalized to $`\mu _s`$) coordinates $`E/\mu _sB/B_s=(E_B/\mu _s)^2`$ $`(B_s=c\mu _s^2/2es_b^2\mathrm{})`$. It is seen that the SO interaction leads to so drastic a reconstruction of 2D spectrum in magnetic fields, that the description of spin splitting by such a non-relativistic parameter as $`g`$-factor loses in essence its physical meaning. The SO splitting far exceeds a contribution due to exchange interaction (in the scale of Fig. 8, the changes induced by exchange interaction are practically indistinguishable for reasonable exchange coupling parameters). This is true for narrow-gap DMS with $`E_g>0`$ also. In connection with this, the results of the analysis of 2D systems in asymmetrical quantum wells in these materials are to be revised, because they ignore the SO interaction.
### 3.2 Density of states and capacitance in magnetic field
For 2D systems with the multisubband spectrum, the differential capacitance of space charge region with 2D electron gas is determined by the sum $`C_{sc}=C_{i\sigma }`$ of “partial” subband capacitances
$$C_{i\sigma }=e^2\frac{dN_{i\sigma }}{d\mu _s}=\frac{e^2}{\pi s^2\mathrm{}^2}\frac{d\mu _{Fi}}{d\mu _s}\frac{d}{d\mu _{Fi}}D_{i\sigma }(E)f(E\mu _{Fi})𝑑E,$$
(6)
where $`N_{i\sigma }`$ is the surface concentration in spin branch $`\sigma =\pm `$ of $`i`$-th subband, $`f(E\mu _{Fi})`$ is the Fermi-Dirac distribution function and $`\mu _{Fi}`$ is the subband Fermi energy. The magnetooscillations of capacitance are dictated by peculiarities in the DOS $`D_{i\sigma }(E)`$ in a magnetic field, i.e. by the energy position and broadening of LL’s. The rigorous quantitative analysis of disorder broadened LL’s is a complicated problem even in one-band approximation . We shall perform the treatment by neglecting the mixing between LL’s and assuming a Gaussian shape of each level
$$D_{i\sigma }=D_i^\pm (E)=\frac{eB}{c\mathrm{}\sqrt{2\pi ^3}}\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{\mathrm{\Gamma }_i^\pm }\mathrm{exp}\left[2\left(\frac{EE_{ni}^\pm }{\mathrm{\Gamma }_i^\pm }\right)^2\right]$$
(7)
with broadening parameters $`\mathrm{\Gamma }_i^\pm `$, independent of LL’s number, but dependent on $`B`$ and (for discussed non-parabolic system) on energy. The analytical description of magnetooscillation phenomena, if possible, is preferable because it is more transparent, and easy to interpret in contrast to direct numerical methods. In addition, the peculiar features of oscillation phenomena in 2D systems with quasirelativistic spectra as compared with those of parabolic spectra can thus be clarified. For this purpose we approximate the spin-split subband dispersions by analytical relativistic-like expression with the subsubband masses $`m_{0i}^\pm `$ and velocities $`s_i^\pm `$. Such approximation holds for practically all the energy range of interest. The replacement in Eq.(7) of numerical solutions of Eq. (5) with $`k_z`$ from Eq. (4) by the corresponding quasiclassical spectrum in a magnetic field (for simplicity we shall drop the spin subband indices $`\sigma =\pm `$ in the following)
$$E_{ni}=\sqrt{E_{Bi}^2\left(n_i+\delta _i\right)+m_{0i}^2s_i^4}m_{0i}s_i^2$$
(8)
($`E_{Bi}=\sqrt{2}s_i\mathrm{}/\lambda `$) does not noticeably reflect on the shape of the DOS. As to the LL’s energy positions (and the positions of oscillations on $`B`$), their values are to be calculated on the basis of Eqs. (4) and (5). In the Born approximation for short-range scattering, the parameter $`\mathrm{\Gamma }_i`$ is related to the classical scattering time in zero magnetic field $`\tau _i`$ by $`\mathrm{\Gamma }_i^2=\sqrt{2/\pi }\mathrm{}^2\omega _{ci}(E)/\tau _i`$ ($`\omega _{ci}(E)=eB/cm_{ci}(E)`$, $`m_{ci}(E)=m_{0i}+E/s_i^2`$) . Inserting the above expressions for $`E_{ni}`$ and $`\mathrm{\Gamma }_i`$ in (7) and using a Poisson summation formula, we arrive (under the assumptions $`\mathrm{\Gamma }_i<<E_{Bi}`$ and $`\mathrm{\Gamma }_i<<E+m_{0i}s_i^2`$, which are justified even for the states nearby the subband bottom) at a “harmonic” representation for DOS (convenient for the description of oscillation effects)
$$D_i(E,B)=D_i(E,0)[1+2\underset{j=1}{\overset{\mathrm{}}{}}(1)^j\mathrm{exp}\left(\frac{j^2\pi }{\omega _{ci}(E)\tau _i}\right)\mathrm{cos}(2\pi jn_i(E))],$$
(9)
where $`D_i(E,0)=m_{di}(E)/2\pi \mathrm{}^2=(E+m_{0i}s_i^2)/2\pi s_i^2\mathrm{}^2`$ is DOS at $`B=0`$. Notice, that in the case of a two-dimension system, as is easy to show, the effective mass of DOS $`m_{di}`$ coincides with the cyclotron effective mass $`m_{ci}`$ for any dispersion law. The LL “number” $`n_i(E)`$ in Eq. (9) is regarded as an arbitrary quantity (not necessarily an integer) and is determined at a given energy by the solutions of Eqs. (4) and (5). One can readily see that in “non-relativistic” limit ($`s_i\mathrm{})`$, Eq. (9) is reduced to the Ando formula if the terms with $`j>1`$ are neglected. It has been shown that for all parameters and regimes of practical interest, the DOS given by (7) and (9) are practically equivalent. This is true in the case of both the sinusoidal shape of DOS (when only ground harmonic predominates in sum (9)) and the nonsinusoidal shape of $`D(E)`$ dependence. The latter case takes place at weakly broadened LL’s, when the individual spin components can be resolved in total DOS.
Substituting Eq. (9) in Eq. (6) and neglecting the integrals of odd functions (this is justified at $`kT<<\mu _{Fi}`$ ) we obtain an expression for magnetocapacitance
$`{\displaystyle \frac{C_i(B)}{C_i(0)}}`$ $`=`$ $`12{\displaystyle \underset{j}{}}\left(1\right)^j\mathrm{exp}\left({\displaystyle \frac{j^2\pi }{\omega _{ci}\tau _i}}\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{cos}(jb_iy)}{2\mathrm{cosh}^2(y/2)}}`$
$`\times `$ $`\{\mathrm{cos}(2\pi jn_i)[\mathrm{cos}(2\pi jc_{Bi}^2y^2)(jc_{\tau i}+yc_{Ti}\mathrm{tan}(jb_iy))\mathrm{sin}(2\pi jc_{Bi}^2y^2)]`$
$``$ $`\mathrm{sin}(2\pi jn_i)[\mathrm{sin}(2\pi jc_{Bi}^2y^2)+(jc_{\tau i}+yc_{Ti}\mathrm{tan}(jb_iy))\mathrm{cos}(2\pi jc_{Bi}^2y^2)]\}dy,`$
where
$$y=(E\mu _{Fi})/kT,\omega _{ci}=\omega _{ci}(\mu _{Fi}),b_i=2\pi kT/\mathrm{}\omega _{ci},$$
$$c_{Ti}=K_{mi}\frac{kT}{\mu _{Fi}},c_{\tau i}=K_{mi}\frac{\mathrm{}}{2\tau _i\mu _{Fi}},c_{Bi}=\frac{kT}{E_{Bi}}=\frac{kT\lambda }{\sqrt{2}s_i\mathrm{}},$$
$$K_{mi}=\left[1+\frac{d(m_{0i}s_i^2)}{d\mu _{Fi}}\right]\left[1+\frac{m_{0i}s_i^2}{\mu _{Fi}}+\frac{d(m_{0i}s_i^2)}{d\mu _{Fi}}\right]^1,$$
and
$$C_i(0)=e^2\left[D_i(\mu _{Fi},0)+\frac{\mu _{Fi}}{2\pi s_i^2\mathrm{}^2}\frac{d(m_{0i}s_i^2)}{d\mu _{Fi}}\right]\frac{d\mu _{Fi}}{d\mu _s}.$$
Note that the identity between the capacitance in a zero magnetic field $`C_i(0)`$ and the quantity $`e^2D(\mu _F,0)d\mu _F/d\mu _s`$ breaks down in 2D systems with relativistic-like spectra . As compared to the semiconductors with standard bands an expression (3.2) contains new parameters of the theory $`c_T`$, $`c_\tau `$ and $`c_B`$. The first and the second are due to the energy dependence of $`m_{di}`$ in an expression for DOS at $`B=0`$ and of $`m_{ci}`$ in exponential (Dingle) factor in Eq. (9). This energy dependence is the result of two effects. The first is the non-parabolicity of subband dispersions. The second is due to the change of the parameters of subsubband dispersions (basically $`m_{0i}^\pm (\mu _s)`$ ) under the modulation of surface potential. The parameter $`c_B`$ is determined in fact only by the ratio $`T/\sqrt{eB}`$ and is independent of subband parameters, because $`s_i^\pm `$ differs only slightly from the ”universal” value of $`s_b`$. (The arising of a $`c_B`$ parameter for weak relativistic Fermi gas $`\mu _F<ms^2`$ of Dirac electrons is reported in Ref. ).
For strong enough magnetic field, the faze $`2\pi jc_B^2y^2`$ is close to zero in the range of small $`y`$ (this range gives a dominant contribution to the integral in Eq. (3.2)) and Eq. (3.2) can be evaluated as
$$\frac{C_i(B)}{C_i(0)}12\underset{j}{}\left(1\right)^j\frac{j\pi b_i}{\mathrm{sinh}(j\pi b_i)}\mathrm{exp}\left(\frac{j^2\pi }{\omega _{ci}\tau _i}\right)\times \left[\mathrm{cos}(2\pi jn_i)jc_{\tau i}\mathrm{sin}(2\pi jn_i)\right]$$
(11)
We point out that in the two-dimension case, the argument of the Dingle exponent is quadratic in harmonic number $`j`$. In ”non-relativistic” limit $`s_i\mathrm{}`$, the values $`c_{Ti}=c_{\tau i}=c_{Bi}=0`$ and Eq. (11) reduces to the Ando expression for oscillations in 2D systems with parabolic spectra . Thus the parabolic approach is appropriate in the case of strong enough magnetic fields $`E_{Bi}>(5÷10)kT`$ and not-too-broadened LL’s $`(\mathrm{}/2\tau _i\mu _{Fi})`$.
## 4 Results of modeling and discussion
The capacitance of the MOS structures $`C(B)=C_{ox}C_{sc}(B)/(C_{ox}+C_{sc}(B))`$ was calculated using the value of oxide capacitance $`C_{ox}`$, determined from the capacitance in a strong accumulation regime. The change of charge in the depletion layer with $`\mu _s`$ in the inversion band bending range (in narrow-gap semiconductors such changes can be quite significant) is taken into account in the calculation. As a rule, the theoretical capacitances $`C(0)`$ are in good agreement with the ones measured at the same surface density. In a modeling, the surface potential and subband Fermi energies are supposed to be constant when a magnetic field is changing. The alternative model is based on the assumption that the surface density is fixed. However, both models give indistinguishable results at large enough LL’s broadening (this is manifested by the cosine form of experimental oscillations) . The temperature dependencies of band parameters and bulk Fermi energy are accounted for in the calculations.
Although the general shapes of simulated and measured oscillations $`C(B)`$ are well matched, the exact magnetic field positions of the peaks and beat nodes are somewhat different. This is not surprising, because a number of physical factors are ignored or cannot be exactly taken into account in a theory (the contribution of remote bands, interface contribution to the SO interaction (see below), the deviation of real surface potential and Landau level shape from those calculated, the superposition of oscillations from different subbands). At the same time the positions of oscillations and especially beat nodes are very sensitive to each of these factors. The adjustable phase factor (a correction $`\mathrm{\Delta }n_i`$ to LL’s “number” $`n_i`$ in the expressions (9), (3.2) and (11)) was introduced for reasons of convenience for a comparison of the temperature evolution in the measured and calculated oscillations. Its magnitude was chosen to fit the high-field node position of the beat pattern at T=4.2 K. Any physically meaningful results discussed are not affected by the choice of this factor.
The oscillations calculated with this correction and their Fourier transforms are plotted in Fig. 7. The agreement is quite good with respect to the structure of oscillations as well as the amplitudes. However, a distinguishable difference in the “number” of oscillations between beat nodes for measured and calculated plots is observed. These results, as well as the similar data on $`dC/dV_g(V_g)`$ oscillations (see Fig. 6), testify to the small (but distinguishable) underestimation of SO splitting by the theory. This conclusion is valid for HgCdTe also. Note that a treatment based on the analysis of Fourier spectra does not give a clearly detectable discrepancy between experiment and theory (excluding the small $`N_s`$ range (see Fig. 4)). This inconsistency with theory can be caused by the interface contribution to the SO interaction , which cannot be treated in the framework of effective mass method. In accordance with the experiment, the individual spin components are not exhibited in simulated $`C(B)`$ or $`dC/dV_g(V_g)`$ oscillations even for the lowest LL’s at any reasonable broadening parameters and magnetic fields of experimental interest. The results of modeling based on Eq. (3.2) do not differ from the results obtained in the approximation (11) practically for all $`N_s`$, $`T`$ and $`B`$ at which the oscillations are observed experimentally. As a rule, the contribution from the sine term in Eq. (11) is small.
### 4.1 Exchange interaction effects
A modeling shows that the magnetic field positions of the oscillations beyond the neighborhood of beat nodes are unaffected by the exchange effects even at the lowest temperatures. As a result these positions are temperature independent, as occurs experimentally. This is true for any available values of the exchange constants $`N\beta ^{}`$ and $`N\alpha ^{}`$ (literature data vary markedly, see Refs. and and references in these works; note that in different papers the notations for exchange constants differ in sign ). As noted in Sec. II, an exchange interaction is very weakly manifested in the studied system, showing itself as only a slight temperature shift of beat nodes. Because the oscillation amplitudes in the neighborhood of nodes are small even at $`T=4.2`$ K and they decrease drastically with temperature, the narrow range of $`T<10÷15`$ K is accessible to the quantitative analysis. Thus the results are not critically sensitive to a choice of $`N\beta ^{}`$ and $`N\alpha ^{}`$. Secondly, the rate of shift depends on the product of the exchange parameters $`N\beta ^{}`$ and $`N\alpha ^{}`$ and the magnetization $`S_z`$. So the variations in $`N\beta ^{}`$ and $`N\alpha ^{}`$can be cancelled out by the variation in $`T_N`$, which is used as an adjustable parameter (see below). It must be stressed that the terms containing a parameter $`\beta `$ play the dominant role in Eqs. (1), (3) and (4) for $`p`$-electrons. Under the conditions of experimental interest, the temperature shift of beat nodes is also slightly sensitive to the variations $`N\alpha ^{}`$ in a wide range even if a small $`N\beta ^{}`$ is chosen.
Although we performed the calculations for a different net of exchange parameters, the results discussed in this Section correspond to $`N\beta ^{}=1.5`$ eV and $`N\alpha ^{}=0.4`$ eV, unless otherwise specified. The results are only slightly sensitive to variances of $`N\beta ^{}`$ in the range $`1.35÷1.65`$ eV and do not differ at all for $`N\alpha ^{}`$ varied through $`(0.35÷0.50)`$ eV range. These values are close to those obtained in Refs. and for narrow-gap and gapless HgMnTe with small $`\left|E_g\right|`$ by the tunnel spectroscopy method. We suppose that these data (similar values for gapless HgMnTe have been obtained in many works (see Refs. and and references in these works) are more suitable for the purposes of this work, because in studied surface quantum wells the typical electron energies are of the order of or even considerably more than $`\left|E_g\right|`$. In tunnel experiments, the LL’s energy positions of “$`p`$\- electrons” as a function of magnetic field are measured at energies up to 150 meV. At the same time, the states with the energy near band bottom are tested by the traditional methods.
It should be noted in connection with this that a decrease of $`|N\alpha ^{}|`$ for $`s`$-electrons in in wide-gap CdMnTe-CdMgMnTe quantum well with increasing energy is reported in a recent paper Ref. . The effect is attributed to the admixing of $`\mathrm{\Gamma }8`$ band states to $`\mathrm{\Gamma }6`$ band at finite $`𝐤`$-vectors which leads to switching-on of a kinetic exchange for electrons of $`\mathrm{\Gamma }6`$ band with the $`d`$ electrons of Mn ions. Note that in narrow-gap semiconductors, the interband mixing, described by the Kane Hamiltonian (1), results in a strong (and energy dependent) contribution of the $`N\beta ^{}`$ containing terms to a spectrum of $`\mathrm{\Gamma }6`$ band. This is true without allowing for the energy dependence of parameter $`N\alpha ^{}`$. As for electrons of $`\mathrm{\Gamma }8`$ band (“$`p`$\- electrons“), the value of the exchange parameter ($`N\beta ^{}`$) is from the outset governed mainly by kinetic exchange (at any $`𝐤`$-vector). In this case, an increase of $`𝐤`$-vector cannot play a critical role. At present, we do not have evidence of energy dependence of parameter $`N\beta ^{}`$. The absence of an essential change in the value of $`N\beta ^{}`$ is noted in Refs. and . Together with the weak sensitivity of the observed effects to a choice of $`N\beta ^{}`$ and $`N\alpha ^{}`$, this suggests that any possible energy dependence of exchange parameters cannot markedly reflect on the results.
Once the exchange parameters are chosen, two parameters can be obtained when the modeling fits experimental data: the effective temperature $`T_N`$, which describes the temperature shift of beat nodes, and the Dingle temperature $`T_D=\mathrm{}/k_B2\pi \tau `$ (used by us as the characteristic of the scattering instead of collision time $`\tau `$), which determines the oscillation amplitudes.
Although a decrease of magnetization with increasing temperature results in the slight energy shift of spin sublevels, the position of resulting oscillations on a magnetic field is almost unchanged (see Fig. 9). The rate of shift with temperature depends on the $`N_i`$ and node number, as is shown in Fig. 10. At the same time, the value of $`T_N`$, extracted from the fit of the temperature evolution of oscillations, is almost the same for different nodes and $`N_i`$, which counts in favor of model used. The results of the simulation are not critically dependent on the exact value of $`T_N`$ chosen. However, the “best fitting” value $`T_N=10\pm 1.5`$ K must be a fairly good estimation. Unfortunately, as far as we know, the low-temperature data on $`T_N`$ value for bulk HgMnTe with $`x=0.04`$ are absent. Most of the literature data are obtained either for high temperatures or for samples with Mn content $`x0.025`$. However the value $`T_N=10`$ K does not contradict other published data. If the sample-independence of spin-spin interaction is postulated, $`T_N`$ is nearly proportional to $`x(1x)^{18}`$ (see Ref. ). Using the low temperature data from Ref. for a sample with $`x=0.01`$ ($`T_N=2.9`$ K at $`T=2`$ K ) we can estimate the value of $`T_N`$ for samples with $`x=0.04`$ as $`T_N8`$ K. This is somewhat less than the measured value, but $`T_N`$ can also be temperature dependent. For example, for the same sample with $`x=0.01`$, $`T_N`$ is equal to 7 K in a high temperature range . It must be noted that the above estimations are based on assumptions which can be violated (including a phenomenological expression itself (2)) for $`x>0.02`$ and low temperature.
The SO interaction not only contributes predominantly to the spin splitting but also suppresses the splitting due to exchange interaction. As an example, the SO splitting, corresponding to the first beat node in Fig. 7, is 24.9 meV. If we take exchange interaction into account, the splitting increases by only 3.7 meV even at T=4.2 K. At the same time, the exchange splitting, calculated without allowance for SO interaction, is 6.4 meV. That is why the exchange effects show themselves only as a weak change in the structure of oscillations near the beat nodes, where the oscillations from different spin branches quench each other.
Let us now return to the dependence of observed exchange effects (the temperature shift of beat nodes) on the value of exchange parameters. Only the shift of beat nodes to low LL’s numbers is observed at low temperatures with decreasing $`N\beta ^{}`$ (the shift is slightly sensitive to the variations $`N\alpha ^{}`$ in $`(0.25÷0.5)`$ eV range). As a result, the rate of temperature shift of nodes decreases and becomes less than the one calculated at $`N\beta ^{}=1.5`$ eV. However, at $`N\beta ^{}>0.75`$ eV, this decrease can be cancelled out by a decrease in $`T_N`$. For $`N\beta ^{}=1.0`$ eV and $`N\alpha ^{}=0.4`$ eV the shifts coincide with those found for $`N\beta ^{}=1.5`$ eV and $`N\alpha ^{}=0.4`$ eV if the value $`T_N=4`$ K is chosen. The oscillations in both cases are practically the same at all $`B`$ (including the ranges nearby the beat nodes) and $`T`$. Although the data do not allow a clear choice between the two, the value $`T_N=4`$ K seems to be too small for $`x=0.04`$.
At the same time, the experimental results cannot be described at $`N\beta ^{}<0.7`$ eV. The measured rate of shift is nearly twice as large as that calculated at $`N\beta ^{}=0.6`$ eV and $`N\alpha ^{}=0.4`$ eV (the values given in Ref. ) even if $`T_N=0`$ is chosen. Although the exchange effects in studied systems with a strong interband mixing are suppressed by SO splitting, such a discrepancy is beyond the limits of experimental error. It is easy to verify that the experimental data (energy position of LL’s and its temperature shift) presented in Refs. and for bulk HgMnTe with small $`|E_g|`$ also cannot be explained at $`N\beta ^{}<1,0÷1.2`$ eV even for $`T_N=0`$. As already noted, the value of $`N\beta ^{}`$ reported in works on gapless HgMnTe falls typically within $`0.9÷1.6`$ eV.
### 4.2 Dingle temperature and scattering
At calculations we suppose that $`T_D`$ for both spin-orbit branches is the same, as it occurs for light and heavy holes in the bulk of a semiconductor. This assumption has supporting experimental evidence. When three or more beat nodes are observed in the oscillations, the “partial” oscillations relating to different spin branches can be extracted from experimental $`C(B)`$ traces, using Fourier filtration and inverse Fourier transform. The $`T_D`$ values determined from fitting turn out to be close for both branches within the accuracy of the analysis. At the same time, the amplitudes corresponding to these branches can differ considerably (up to several times). Such a difference is not surprising. It is clear that both the DOS and cyclotron energy are different for two branches of a spectrum having significantly different dispersions. As a result, the “partial” capacitance oscillations for these branches differ not only by the period (which leads to the beat of oscillations) but also by the amplitudes, even if the relaxation times are equal. Although DOS at $`B=0`$ in a low-energy branch is higher, the corresponding amplitudes can be less, because the lower cyclotron energy in this branch leads to a lesser amplitude factor in Eqs. (3.2) and (11).
The relation between amplitudes depends on the subband occupancies (via the effective mass), magnetic field, temperature and broadening parameters. For different subbands, $`B`$ and $`N_s`$, the ratio of amplitudes $`A_c^{}/A_c^+`$ can be below as well as above unity. However, the value of $`A_c^{}/A_c^+`$ decreases rapidly with increasing $`N_i`$. Such behavior correlates well with decreasing ratio of Fourier line intensities $`I_i^{}/I_i^+`$ experimentally observed (see Fig. 4). Thus the difference in the amplitudes for different spin components of oscillations mentioned in Refs. and is to be expected for 2D systems with strong SO interaction (without invoking spin-dependent scattering).
The Dingle temperature $`T_D`$, determined from the fitting, slightly increases (from $`8÷9`$ to $`13÷15`$ K for ground subband) with the increasing $`N_s`$. Such behavior is inherent in surface roughness scattering . The $`T_D(N_s)`$ dependencies are essentially sublinear. This testifies that the efficiency of scattering is suppressed with the increasing of the Fermi wave vector. This is possible if the correlation length $`\mathrm{\Lambda }`$ is large enough. The best agreement between experimental and calculated values of $`T_D`$ is achieved at $`\mathrm{\Lambda }(110÷120)`$ Åand at the average interface displacement $`\mathrm{\Delta }(20÷25)`$ Å. Note that the screening effects contribute significantly to the scattering, because in surface layers on narrow-gap semiconductors, the Fermi wavelength turns out to be of the same order of magnitude as 2D Thomas-Fermi screening length and $`z`$-size of wave function. The above values of $`\mathrm{\Lambda }`$ and $`\mathrm{\Delta }`$ are many times larger than tuose found for silicon and substantially exceed the corresponding values in III-V semiconductors also . This suggests much more disorder in the interface between ternary compounds and their oxides. Using the $`T_D`$ values, the electron mobility can be estimated as $`0.8\times 10^4`$ cm<sup>2</sup>/Vs in $`i=0`$ subband and $`1.5\times 10^4`$ cm<sup>2</sup>/Vs in $`i=1`$ subband for sample $`S1`$ at $`N_s10^{12}`$ cm<sup>-2</sup> that is close to a value $`1\times 10^4`$ cm<sup>2</sup>/Vs measured for grain boundaries in $`p`$-HgMnTe with $`x=0.1`$ .
Somewhat larger values of $`T_D`$ are detected at small surface densities $`N_s<5\times 10^{11}`$ cm<sup>-2</sup>. The Coulomb scattering from the chargers in oxide cannot cause this, because the theoretical calculations give the values of relaxation times, which are larger by at least two orders of magnitude. This conclusion has direct experimental evidence. It can be seen in Fig. 1, that the charge localized in oxide differs by a factor of several times for different sweep cycles. If the Coulomb scattering were important, the amplitudes of oscillations corresponding to different cycles (different $`V_{fb}`$) but to the same $`N_s`$ (to the same LL’s number at fixed magnetic field) would be different. Nevertheless, the oscillation amplitudes are practically the same. A possible cause for the increase of LL’s broadening at small $`N_i`$ is intersubband scattering .
Acknowledgements. This work was supported in part by the project Esprit N28890 NTCONGS EC (Euro Community) and by the Grant from the Education Committee of Russian Federation.
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# Classification of Incidence Scrolls(I)Mathematics Subject Classifications: 14J26 (14H25, 14H45)
## 1 Ruled Surfaces
A ruled surface is a surface $`X`$ together with a surjective morphism $`\pi :XC`$ to a (connected nonsingular) curve $`C`$ such that the fibre $`X_y`$ is isomorphic to $`𝐏^1`$ for every point $`yC`$, and such that $`\pi `$ admits a section. There exists a locally free sheaf $``$ of rank 2 on $`C`$ such that $`X𝐏()`$ over $`C`$. Conversely, every such $`𝐏()`$ is a ruled surface over $`C`$. Let $`𝒪_X(1)`$ be the invertible sheaf $`𝒪_{𝐏()}(1)`$. Then there is an one-to-one correspondence between sections $`\sigma :CX`$ and surjections $`0`$, where $``$ is an invertible sheaf on $`C`$, given by $`\sigma ^{}(𝒪_X(1))`$. Moreover, if $`D`$ is any section of $`X`$, corresponding to a surjection $`0`$, and if $`=(\text{d})`$ for some divisor d on $`C`$, then $`deg(\text{d})=C_oD`$, and $`D`$ and $`C_o+(\text{d}\text{e})f`$ are linearly equivalent, written $`DC_o+(\text{d}\text{e})f`$.
It is possible to write $`X𝐏()`$ where $``$ is a locally free sheaf on $`C`$ with the property that $`H^0()0`$ but for all invertible sheaves $``$ on $`C`$ with $`deg<0`$, we have $`H^0()=0`$. In this case we say $``$ is normalized. Let e be the divisor on $`C`$ corresponding to the invertible $`^2`$ and let $`e=deg\text{e}`$. Furthermore there is a section $`\sigma _0:CX`$ with image $`C_o`$ such that $`(C_o)𝒪_X(1)`$. Fix such a section $`C_o`$ of $`X`$. If b is any divisor on $`C`$, then we denote the divisor $`\pi ^{}\text{b}`$ on $`X`$ by $`\text{b}f`$, by abuse of notation. Thus any element of Pic($`X`$) can be written $`aC_o+\text{b}f`$ with $`a`$ and $`\text{b}`$ Pic($`C`$).
Let $`X`$ be a ruled surface over the curve $`C`$ of genus $`g`$, determined by a normalized locally free sheaf $``$. If $``$ is decomposable, then $`𝒪_C`$ for some $``$ with $`deg0`$. All values of $`e0`$ are possible. If $``$ is indecomposable, then $`ge2g2`$.
A scroll is a ruled surface embedded in $`𝐏^N`$ in such a way that the fibres $`f`$ have degree 1. If we take a very ample divisor on $`X`$, $`DaC_o+\text{b}f`$, then the embedding $`\mathrm{\Phi }_D:XR_g^d𝐏^N=𝐏(H^0(𝒪_X(D)))`$ determines a scroll when $`Df=a=1`$. Throughout this work, we will restrict our attention to very ample divisors so that $`\mathrm{\Phi }_D(X)`$ be a scroll.
## 2 Definition of Incidence Scroll
###### Definition 2.1
A scroll $`R𝐏^n`$ is said to be an incidence scroll if $`R`$ is generated by the lines which meet a certain set $``$ of linear spaces in $`𝐏^n`$. Such a set is called a base of the incidence scroll.
Since a scroll $`R_g^d𝐏^n`$ is represented by a curve $`C_g^dG(1,n)`$, there is another definition of incidence scroll which is equivalent to the previous one.
###### Definition 2.2
A scroll $`R𝐏^n`$ is an incidence scroll if the correspondent curve in $`G(1,n)`$ is an intersection of special Schubert varieties $`\mathrm{\Omega }(𝐏^r,𝐏^n)`$, $`0r<n1`$.
Such a base will be denoted by:
$$=\{𝐏^{n_1},𝐏^{n_2},\mathrm{},𝐏^{n_r}\}$$
or we will write it simply $``$ when no confusion can arise, where $`n_1n_2\mathrm{}n_r`$, i.e., $`C=_{i=1}^r(\mathrm{\Omega }(𝐏^{n_i},𝐏^n))G(1,n)`$. If $`n_1=\mathrm{}=n_i,n_{i+1}=\mathrm{}=n_j,\mathrm{}\text{and}n_{j+1+k}=\mathrm{}=n_r`$, then we write it $`=\{i𝐏^{n_1},(ji)𝐏^{n_j},\mathrm{},(rjk)𝐏^{n_r}\}`$, for short.
We can certainly assume that $`n3`$. This involves no loss of generality because the only scroll in $`𝐏^2`$ is $`𝐏^2`$. The plane can be obtained as incidence scroll in $`𝐏^n`$ if we take $`=\{𝐏^0,(n2)𝐏^{n2}\}`$. If we want to see it as nondegenerate surface, we must work in $`𝐏^2`$ with $`=\{𝐏^0\}`$ (i.e., $`𝐏^1=\mathrm{\Omega }(𝐏^0,𝐏^2)`$).
We are interested in incidence scrolls where the base is formed by linear spaces in general position. By general position we will mean that $`(𝐏^{n_1},\mathrm{},𝐏^{n_r})𝒲=G(n_1,n)\times \mathrm{}\times G(n_r,n)`$ is contained in a nonempty open subset of $`𝒲`$. Therefore, unless otherwise stated, we will work with general base spaces. For simplicity of notation, we abbreviate it to base in general position.
###### Theorem 2.3
The intersection $`_{j=1}^m\mathrm{\Omega }(𝐏^{n_j},𝐏^n)G(l,n)`$ of the special Schubert varieties associated to linear subspaces $`𝐏^{n_j},j=1,\mathrm{},m`$, of dimension $`n_j`$ such that $`(l+1)(nl)_{j=1}^m(nn_j1)>0`$, is connected. Moreover, the intersection is irreducible of dimension $`(l+1)(nl)_{j=1}^m(nn_j1)`$ for a general choice of the subspaces $`𝐏^{n_j}`$.
Proof. See sols , Theorem 1.1.
With the previous theorem, we can prove immediately an important consequence which provides a characterization of incidence scrolls.
###### Proposition 2.4
The intersection $`C=_{i=1}^r\mathrm{\Omega }(𝐏^{n_i},𝐏^n)`$ of special Schubert varieties associated to linear spaces $`𝐏^{n_i},i=1,\mathrm{},r`$ in general position is a conex irreducible curve of $`G(1,n)`$ if and only if it verifies the following equality
$$rn(n_1+n_2+\mathrm{}+n_r)r=2n3$$
$`(IS)`$
Proof. $`{}_{}{}^{\prime \prime }_{}^{\prime \prime }`$ Since $`CG(1,n)`$ is a conex irreducible curve, we have $`1=2(n1)rn+r+_{i=1}^rn_i\text{(see Theorem }\text{2.3}\text{)}rn_{i=1}^rn_ir=2n3`$.
$`{}_{}{}^{\prime \prime }_{}^{\prime \prime }`$ Since $`2(n1)_{i=1}^rnn_i1>0`$, we conclude that $`C`$ is connected, by Theorem 2.3. Moreover, we easily see from Theorem 2.3 that it is irreducible of dimension $`1`$ for a general choice of the $`𝐏^{n_i},i=1,\mathrm{},r`$.
Each $`\mathrm{\Omega }(𝐏^{n_i},𝐏^n)`$ is the intersection of $`G(1,n)`$ with a linear space in $`𝐏^N`$. Moreover, such a linear space is a hyperplane if and only if we have $`n_i=n2`$. If every dimension of the base spaces is equal to $`n2`$, then we will need exactly $`(2n3)𝐏^{n2}`$’s for obtain a curve in $`G(1,n)`$.
Some relevant properties of incidence scrolls will be indicated. These are elementary but of great importance if we want to develop a theory about incidence scrolls.
1. A hyperplane does not impose conditions on $`G(1,n)`$ because $`\mathrm{\Omega }(𝐏^{n1},𝐏^n)`$ $`=G(1,n)`$. Then we can assume $`n_rn2`$.
2. If $`n_i+n_j<n1`$ then the scroll is degenerated. Under the above assumptions, $`\mathrm{\Omega }(𝐏^{n_i},𝐏^n)\mathrm{\Omega }(𝐏^{n_j},𝐏^n)=\mathrm{\Omega }(𝐏^{n_i},𝐏^{n_i+n_j+1})\mathrm{\Omega }(𝐏^{n_j},𝐏^{n_i+n_j+1})`$.
3. A nondegenerate (irreducible) incidence scroll cannot have double points. In other case, if $`P`$ is a double point, there are two generators of the scroll $`g_1,g_2`$ through $`P`$. Thus, the incidence scroll is not irreducible, containing the plane pencil determined by them. In $`G(1,n)`$, we have $`_{i=1}^r\mathrm{\Omega }(𝐏^{n_i},𝐏^n)=\mathrm{\Omega }(P,𝐏^2)(_{j=1}^s\mathrm{\Omega }(𝐏^{m_j},𝐏^m))`$.
4. The generators of the scroll meet each base space in points of a directrix curve.
5. The degree of the scroll can be calculated by a one-to-one correspondence between any two of the directrix curves. It can also be computed by Giambelli’s formula, i.e., it is the number of lines in $`𝐏^n`$ which intersect $`𝐏^{n_1},𝐏^{n_2},\mathrm{},𝐏^{n_r}`$ and a generic $`𝐏^{n2}`$. If $`deg(R)=d`$, then we have the following equality of Schubert cycles: $`\mathrm{\Omega }(n_1,n)\mathrm{}\mathrm{\Omega }(n_r,n)=d\mathrm{\Omega }(0,2)`$.
From now on, we will talk about a decomposable incidence scroll if the corresponding ruled surface $`X=𝐏()`$ has $``$ decomposable. If $``$ is indecomposable, then we will talk of an indecomposable incidence scroll.
Let $`X=𝐏()`$ be a ruled surface over the curve $`C`$ of genus $`g`$, determined by a decomposable normalized bundle $`𝒪_C𝒪_C(\text{e})`$ such that $`deg(\text{e})=e0`$. Let $`HC_o+\text{b}f`$ be the very ample divisor on $`X`$ with $`m=deg(\text{b})`$ which gives the immersion of the ruled surface as the scroll $`R_g^d𝐏^n`$ such that $`d=2me`$ and $`n=2(mg)e+1+i`$, being $`i`$ the speciality of the scroll ($`i=h^1(𝒪_C(\text{b}))+h^1(𝒪_C(\text{b}+\text{e}))`$). Geometrically, $`X`$ has two disjoint directrix, denoted by $`C_o`$ and $`C_1`$, such that $`C_1C_o\text{e}f`$. Moreover, it is easy to check that these satisfy $`\varphi _{\text{b}+\text{e}}:C_oC_g^{me}𝐏^{meg+i_1}`$ and $`\varphi _\text{b}:C_1C_g^m𝐏^{mg+i_2}`$, being $`i_1=h^1(𝒪_C(\text{b}+\text{e}))`$ and $`i_2=h^1(𝒪_C(\text{b}))`$. Therefore $`i=i_1+i_2`$ and $`𝐏^{meg+i_1}𝐏^{mg+i_2}=\mathrm{}`$.
###### Proposition 2.5
If $`R_g^d𝐏^n`$ is a decomposable incidence scroll with base in general position, then $`𝐏^{meg+i_1}`$ and $`𝐏^{mg+i_2}`$ are base spaces.
Proof. Suppose $``$ a base with $`_{i=1}^r(nn_i1)=2n3`$ such that $`𝐏^{n_i}`$’s, $`i=1\mathrm{},r`$, are in general position and $`meg+i_1<n_in2`$. Since the scroll has a minimum directrix curve $`C_g^{me}𝐏^{meg+i_1}`$, we can take a generic hyperplane $`𝐏^{n1}`$ through $`𝐏^{meg+i_1}`$. Then there are $`m`$ lines in $`𝐏^{n1}`$ which meet $`𝐏^{n_11},𝐏^{n_21},\mathrm{}`$ and $`𝐏^{n_r1}`$. Since each $`𝐏^{n_i1}`$ imposes $`nn_i1`$ independent conditions on $`𝐏^{n1}`$, it is impossible.
$`𝐏^{meg+i_1}`$ imposes $`mg+i_2`$ conditions on $`G(1,n)`$. Then we can consider more subspaces to form the base. There is a $`𝐏^{mg+i_2}`$. For otherwise, suppose $`=\{𝐏^{meg+i_1},𝐏^{n_1},\mathrm{},𝐏^{n_r}\}`$ (in general position) with $`mg+i_2<n_in2`$. Then we can take a generic hyperplane $`𝐏^{n1}`$ such that $`𝐏^{mg+i_2}𝐏^{n1}`$ and arguments similar to the above imply that $``$ is not the base of $`R_g^d𝐏^n`$.
###### Theorem 2.6
Let $`R_g^d𝐏^n`$ be an incidence scroll with base in general position. Then:
$$Risdecomposable𝐏^{n_1}𝐏^{n_2}=\mathrm{}.$$
Proof.$``$” Since $`𝐏^{n_1}𝐏^{n_2}=\mathrm{}`$, the directrix curves $`C_g^{d_1}𝐏^{n_1}`$ and $`C_g^{d_2}𝐏^{n_2}`$ satisfy $`C_g^{d_1}C_g^{d_2}=\mathrm{}`$, i.e., $`R_g^d`$ is decomposable.
$``$” By Proposition 2.5, there are two directrix curves $`C_g^{d_1}𝐏^{n_1}`$ and $`C_g^{d_2}𝐏^{n_2}`$ such that $`C_g^{d_1}C_g^{d_2}=\mathrm{}`$. Suppose $`n_j=d_jg+i_j`$ for $`j=1,2`$. Since $`n=n_1+n_2+1`$ and $`𝐏^{n_1}`$ and $`𝐏^{n_2}`$ are in general position, $`𝐏^{n_1}𝐏^{n_2}=\mathrm{}`$.
If $`R`$ is decomposable, $`𝐏^{meg+i_1}`$ and $`𝐏^{mg+i_2}`$ impose $`2(mg)e+i`$ conditions on $`G(1,n)`$. Since $`2n3(2(mg)e+i)=n21`$, we must consider more subspaces to form the base.
In particular, if $`e2g2`$, then $`h^0(𝒪_X(C_o\text{e}f))=e+2g`$. Moreover, if $`g=0,1`$, then we can show (by Bertini’s theorem) that there are $`e+2g`$ linearly independent directrix curves $`CC_o\text{e}f`$, which generates $`|C_o\text{e}f|`$.
###### Proposition 2.7
Let $`R_g^{2me}𝐏^{2(mg)e+1+i}`$ be a decomposable incidence scroll with base in general position.
1. If $`\text{e}\sim ̸0`$, then there are $`(e+2g)𝐏^{mg+i_2}`$’s in $``$, whenever possible, i.e., when $`mg+i_2+(e+2g)(meg+i_1)4(mg)2e+2i1`$ and the generic curve in $`|C_o\text{e}f|`$ is irreducible.
2. If $`\text{e}0`$, then there are $`3𝐏^{meg+i_1}`$’s in $``$.
Proof. $`(a)`$ We proceed by induction in $`j=e+2g`$. The proposition is true for $`j=1`$, by Proposition 2.5. Supposing the proposition true for $`j`$, we prove it for $`(j+1)𝐏^{mg+i_2}`$’s. Let $`=\{𝐏^{meg+i_1},𝐏^{n_1},`$ $`𝐏^{n_2},\mathrm{},𝐏^{n_r}\}`$ be a base of the scroll in general position with $`mg+i_2=n_1=\mathrm{}=n_j<n_{j+1}\mathrm{}n_r2(mg)e+i1`$. Since the scroll has at least one directrix curve $`C_g^m𝐏^{mg+i_2}`$, which is linearly independent from the other directrix curves contained in base spaces of dimension $`mg+i_2`$, we can take a generic hyperplane $`𝐏^{2(mg)e+i}`$ through $`𝐏^{mg+i_2}`$. Then there are $`me`$ lines in $`𝐏^{2(mg)e+ij}`$ which meet $`𝐏^{meg+i_11},𝐏^{n_1j},`$ $`\mathrm{},𝐏^{n_jj},𝐏^{n_{j+1}j1},`$ $`\mathrm{}`$ and $`𝐏^{n_rj1}`$, which is impossible.
The proof for $`(b)`$ is similar.
## 3 Degeneration
###### Proposition 3.1
Let $`R_g^d𝐏^n`$ be an incidence scroll with base $``$ in general position. Suppose that $`𝐏^{n_i}`$ and $`𝐏^{n_j}`$ lie in a hyperplane $`𝐏^{n1}`$ and have in common $`𝐏^m`$, $`m=n_i+n_jn+1`$. Then the scroll breaks up into:
1. $`R_{g_1}^{d_1}𝐏^n`$ with base $`\dot{}=\{𝐏^m,𝐏^{n_1},\mathrm{},\widehat{𝐏^{n_i}},\mathrm{},\widehat{𝐏^{n_j}},\mathrm{},𝐏^{n_r}\}`$ (which is possibly degenerate);
2. $`R_{g_2}^{d_2}𝐏^{n1}`$ with base $`\ddot{}=\{\mathrm{},𝐏^{n_{i1}1},𝐏^{n_i},𝐏^{n_{i+1}1},\mathrm{},𝐏^{n_{j1}1},𝐏^{n_j},`$ $`𝐏^{n_{j+1}1},\mathrm{}\}`$
which have $`\kappa 1`$ generators in common. Then, $`d=d_1+d_2`$ and $`g=g_1+g_2+\kappa 1`$.
Moreover, if $`m=0`$, then the incidence scroll breaks up into a plane and an incidence scroll $`R_g^{d1}𝐏^{n1}`$ with base $`\ddot{}`$ in general position.
Proof. By assumption we have a scroll formed by the lines which pass through $`𝐏^m`$ and a scroll generated by the lines which intersect the base spaces and lie in $`𝐏^{n1}`$. The three scrolls are represented by curves $`C=_{k=1}^r\mathrm{\Omega }(𝐏^{n_k},𝐏^n)`$, $`\dot{C}=\mathrm{\Omega }(𝐏^m,𝐏^n)(_{ki,j}\mathrm{\Omega }(𝐏^{n_k},𝐏^n))`$ and $`\ddot{C}=\mathrm{\Omega }(𝐏^{n_i},𝐏^{n1})\mathrm{\Omega }(𝐏^{n_j},𝐏^{n1})(_{ki,j}\mathrm{\Omega }(𝐏^{n_k1},𝐏^{n1}))`$ because every one of them satisfies $`(IS)`$.
Write $`A_i=\mathrm{\Omega }(𝐏^{n_i},𝐏^n)),i=1,\mathrm{},r`$, where the $`𝐏^{n_i}`$’s are in general position and $`B_j=\mathrm{\Omega }(𝐏^{n_j},𝐏^n)),j=1,\mathrm{},r`$, where $`𝐏^{n_1}𝐏^{n_2}=𝐏^m`$ and the other subspaces are in general position. Then there is an invertible linear transformation of $`𝐏^N`$ into itself which carries $`G(1,n)`$ into itself and $`A_i`$ into $`B_i`$ (see Kleiman , Proposition 4). Consequently, we conclude that $`A_i`$ and $`B_i,i=1,\mathrm{},r`$, define the same cohomology class. In this way we obtain what the set of lines which simultaneously intersect the subspaces of $``$ can be continuously deformed into the union of two set of lines. A set formed by the lines which pass through a fixed point and lie in a plane and other set formed by the lines which simultaneously intersect the subspaces of $`\ddot{}`$. Since the various subvarieties in a continuous system are all assigned the same cohomology class, we find the following equality of Schubert cycles, $`C=\dot{C}+\ddot{C}=d\mathrm{\Omega }(0,2)`$. So, $`d=d_1+d_2`$.
We have in $`G(1,n)`$ a short exact sequence:
$$0𝒪_C𝒪_{\dot{C}}𝒪_{\ddot{C}}𝒪_{\dot{C}\ddot{C}}0.$$
Therefore $`𝒳(C)=𝒳(\dot{C})+𝒳(\ddot{C})𝒳(\dot{C}\ddot{C})`$ where $`𝒳(C)`$ is the Euler characteristic, whence $`g(C)=g(\dot{C})+g(\ddot{C})+\kappa 1`$.
###### Remark 3.2
Let $`C_g^d=_{k=1}^r\mathrm{\Omega }(𝐏^{n_k},𝐏^n)`$ be a curve in $`G(1,n)`$, which defines an incidence scroll in $`𝐏^n`$. If we want to obtain another curve $`\dot{C}`$ which is an intersection of Schubert varieties with the same genus, then $`\ddot{C}`$ must be a line, there is only one generator in common and $`deg(\dot{C})=d1`$.
Let $`D_g^{d_k}`$ be the directrix curve of $`R_g^d𝐏^n`$ contained in $`𝐏^{n_k}`$ for $`k\{1,\mathrm{},r\}`$. Then the directrix curve of $`R_g^{d1}𝐏^{n1}`$ in the $`𝐏^{n_k1}`$ has degree $`d_k1`$, for $`ki,j`$. This is not so in the other cases because if $`k=i,j`$, then the directrix curve in $`𝐏^{n_k}`$ has degree $`d_k`$ (an easy computation of Schubert cycles).
We shall refer to this particular degeneration as join $`𝐏^{n_i}`$ and $`𝐏^{n_j}`$ and to the inverse as separate $`𝐏^{n_i}`$ and $`𝐏^{n_j}`$.
###### Example 3.3
The following examples are a good illustration of join and separate $`𝐏^{n_i}`$ and $`𝐏^{n_j}`$.
1. $`_1=\{2𝐏^{n1},2𝐏^n,𝐏^{2n3}\}R_{n2}^{4n6}𝐏^{2n1}`$.
If we suppose $`𝐏^n𝐏^{2n3}=𝐏^{2n2}`$, i.e., $`𝐏^n𝐏^{2n3}=𝐏^{n1}`$, then the scroll breaks up into $`R_0^{2n2}𝐏^{2n1}`$ with base $`\{3𝐏^{n1},𝐏^n\}`$ and $`R_0^{2n4}𝐏^{2n3}`$ with base $`\{3𝐏^{n2},𝐏^{n1}\}`$ (we will see it in section 4) and $`n1`$ generators in common. It follows that $`d=2n2+2n4=4n6`$ and $`g=n2`$.
2. $`_2=\{𝐏^{n2},2𝐏^{n1},𝐏^n,𝐏^{2n4}\}R_{n3}^{4n9}𝐏^{2n2}`$; apply 3.2 to a $`𝐏^{n1}`$ and a $`𝐏^n`$ in $`_1`$ and write $`n`$ instead of $`n1`$.
3. $`_3=\{𝐏^{n3},3𝐏^{n1},𝐏^{2n5}\}R_{n4}^{4n12}𝐏^{2n3}`$; apply 3.2 to $`𝐏^{n2}`$ and $`𝐏^n`$ in $`_2`$ and write $`n`$ instead of $`n1`$.
4. $`_4=\{3𝐏^{n1},𝐏^{n+1},𝐏^{2n3}\}R_{n3}^{4n8}𝐏^{2n1}`$; apply 3.2 to $`2𝐏^{n1}`$’s in $`_2`$.
We will use the existence of the previous incidence scrolls in section 5.
## 4 Incidence Rational Scrolls
Let $`X=𝐏()`$ be a ruled surface. We say that $`X`$ is a rational ruled surface if $`C𝐏^1`$. One sees immediately that $`X_0=𝐏^1\times 𝐏^1`$ with its first projection is a rational ruled surface. Each rational ruled surface is isomorphic to the ruled surface obtained from $`X_0`$ by applying a finite number of elementary transform at $`P_1,\mathrm{},P_eX_0`$, with a suitable number $`e0`$. We shall write the above expression as $`X_e=elem_{(P_1,\mathrm{},P_e)}(X_0)`$. The integer $`e`$ is an invariant of $`X_e`$. Moreover, for each $`e0`$ there is exactly one (up to isomorphism) rational ruled surface with invariant $`e`$, given by $`=𝒪_{𝐏^1}𝒪_{𝐏^1}(e)`$, i.e., $`X_e𝐏()`$.
###### Proposition 4.1
Let $`DaC_o+bf`$ be a divisor on $`X_e`$. Then:
1. $`D`$ is very ample $``$ $`D`$ is ample $``$ $`a>0`$ and $`b>ae`$;
2. $`|D|`$ contains an irreducible nonsingular curve $``$ $`|D|`$ contains an irreducible curve $``$ $`a=0,b=1`$; or $`a=1,b=0`$; or $`a>0,b>ae`$; or $`e>0,a>0,b=ae`$.
Let $`H`$ be a very ample divisor on $`X_e`$ and let $`\mathrm{\Phi }_H`$ be the closed immersion defined by $`|H|`$. Then we know that $`\mathrm{\Phi }_H(X_e)`$ is a scroll if and only if $`HC_o+mf`$ with $`m>e`$. To see a more complete theory of rational ruled surfaces we refer the reader to hatshor and Lanteri .
The aim of this chapter is to study all the incidence rational scrolls with base in general position. For this, we can now formulate a theorem as follows.
###### Theorem 4.2
Let $`HC_o+mf`$ be a very ample divisor on $`X_e`$ and let $`\mathrm{\Phi }_H:X_eR_0^{2me}𝐏^{2me+1}`$ be the closed immersion defined by $`|H|`$. Then $`R_0^{2me}`$ is an incidence scroll if and only if it satisfies one of the following conditions:
1. $`HC_o+(e+1)f`$ ($`\mathrm{\Phi }_H(C_o)=𝐏^1`$);
2. $`e=0`$ ($`\mathrm{\Phi }_H(C_o)=C_0^m𝐏^m`$);
3. $`e=1`$ ($`\mathrm{\Phi }_H(C_o)=C_0^{m1}𝐏^{m1}`$).
We will divide its proof into a sequence of propositions.
### Rational Scrolls with a Directrix Line
It is clear that an incidence scroll with $`𝐏^1`$ as base space is a rational scroll. The following proposition gives all rational normal scrolls with a line as directrix. The affirmation is not true if normal is deleted from the hypothesis. For example, $`R_0^3𝐏^3`$, which has a double line, is not an incidence scroll if the base spaces are in general position.
###### Proposition 4.3
In $`𝐏^n`$ the incidence scroll with base $`_n=\{𝐏^1,(n1)𝐏^{n2}\}`$ is the rational normal scroll of degree $`n1`$ with a line as minimum directrix.
Proof. We proceed by induction in $`n`$. The proposition is true in $`𝐏^3`$. In this case, the scroll is the quadric surface in $`𝐏^3`$.
Supposing the proposition true in $`𝐏^n`$, we prove it for $`n+1`$. We can add a hyperplane to the base without affecting the scroll. Separating $`𝐏^1`$ and $`𝐏^{n1}`$, as in 3.2, we obtain a base $`_{n+1}=\{𝐏^1,n𝐏^{n1}\}`$ which defines $`R_0^n𝐏^{n+1}`$ as incidence scroll.
###### Corollary 4.4
For every $`e0`$, there is an embedding of $`X_e`$ as incidence rational scroll.
Proof. Use the very ample divisor $`HC_o+(e+1)f`$.
### Rational Scrolls of General Type
$`X_e`$ is said to be of general type if and only if $`𝒪_{𝐏^1}𝒪_{𝐏^1}`$ or $`𝒪_{𝐏^1}𝒪_{𝐏^1}(1)`$. Moreover, if $`X_e`$ is general of degree $`d`$, then the minimum degree of the directrix curves is equal to $`\frac{d}{2}`$ (respectively, $`\frac{d1}{2}`$) for $`d`$ even (respectively, odd). On $`X_e`$, the family of directrices of degree $`m`$ has dimension $`d_m=2md+1`$ (see Ghione , pp. 89-90).
###### Proposition 4.5
All general rational scrolls are incidence scrolls.
Proof. e=0: We will see that in $`𝐏^{2m+1}`$ the incidence scroll generated by $`=\{3𝐏^m,𝐏^{m+1}\}`$ is in fact $`R_0^{2m}𝐏^{2m+1}`$. This incidence scroll have directrix curves in each $`𝐏^m`$ of degree $`d_m`$ and a directrix curve in $`𝐏^{m+1}`$ of degree $`d_{m+1}`$. Then:
1. $`d_{m+1}`$ is the number of lines in $`𝐏^{2m+1}`$ which meet $`4𝐏^m`$’s. Call this number $`z(m)`$ i.e., $`d_{m+1}=z(m)=\mathrm{\Omega }(m,2m+1)^4`$.
2. if we take a hyperplane $`𝐏^{2m}𝐏^{2m+1}`$ through the base $`𝐏^{m+1}`$, then we can calculate $`d_m`$ as the number of generators of the scroll which meet $`𝐏^{2m}`$ in points of the base $`𝐏^m`$. To obtain this number we consider two cases: either the generators lie in $`𝐏^{2m}`$ or they do not. In the former case, we have the lines in $`𝐏^{2m3}`$ which meet $`4𝐏^{m2}`$’s therein. The number of these lines is $`z(m2)`$. In the latter case, this number is equal to the number of lines through $`𝐏^m𝐏^{m+1}`$, which meet the other $`2𝐏^m`$’s. Thus $`d_m=z(m2)+1`$.
Therefore $`deg(R)=d_m+d_{m+1}1=2d_mz(m)=z(m2)+2d_{m+1}=z(m)=d_m+1`$. But $`z(m)`$ is the number of lines meeting $`4𝐏^m`$’s in $`𝐏^{2m+1}`$, i.e., by Schubert calculus, $`z(m)=m+1`$. Then $`deg(R)=2m`$. Moreover, a scroll of degree $`2m`$ in $`𝐏^{2m}`$ is necessarily rational and normal.
e=1: Since $`R_0^{2m1}𝐏^{2m}`$ is of general type, it has one minimum directrix curve. Then we can take as base space a $`𝐏^{m1}`$ which contains it. Choose $`3𝐏^m`$’s containing 3 generic directrix curves $`C_0^{m+1}`$. This is possible because $`h^0(𝒪_{X_1}(C_o+f))=3`$. By Proposition 4.1 and Bertini’s Theorem we know that $`|C_o+f|`$ has irreducible curves and that the generic is irreducible. $`=\{𝐏^{m1},3𝐏^m\}`$ is base of $`R_0^{2m1}𝐏^{2m}`$. We conclude it from 3.2, joining $`2𝐏^m`$ of case $`e=0`$.
### Rational Scrolls which are not Incidence Scrolls
###### Proposition 4.6
Let $`HC_o+mf`$ be a very ample divisor on $`X_e`$ with $`m>e+1>2`$. If the base must be in general position, then $`\mathrm{\Phi }_H(X_e)=R_0^{2me}𝐏^{2me+1}`$ cannot be obtained by incidences.
Proof. Under the above hypotheses, we have $`m+(e+2)(me)>4m2e1(1<e<m1)`$. Then:
$$\begin{array}{cc}m+(\eta )(me)4m2e1\hfill & \hfill ()\\ m+(\eta +1)(me)>4m2e1\hfill & \end{array}$$
for a suitable number $`\eta `$ such that $`1<\eta <e+2`$. An argument similar to Proposition 2.7 shows that there are $`\eta 𝐏^m`$’s in the base. Then $`=\{𝐏^{me},\eta 𝐏^m,𝐏^{n_1},\mathrm{},𝐏^{n_r}\}`$ with $`m+1n_i2me1`$ is base of an incidence scroll, which is not rational. In other case, joining base spaces, we find a rational scroll in $`𝐏^3`$, obtained by Proposition 4.3, which is impossible. If we apply it successively the inverse degeneration (separate $`𝐏^{n_i}`$ and $`𝐏^{n_j}`$) to all the possible pairs $`(n_i,n_j)`$ such that $`𝐏^{n_i}𝐏^{n_j}=\mathrm{}`$, then we must obtain our scroll, but we obtain the rational normal scroll $`R_0^{2me}𝐏^{2me+1}`$ with a line as minimum directrix curve (i.e., $`m=e+1`$) or the rational scroll of general type (i.e., $`e1`$). Note that we can project to any $`𝐏^{2mek}`$, and so join two base spaces, because the scroll is rational and, in particular, decomposable, i.e., $`𝐏^{n_1}𝐏^{n_2}=\mathrm{}`$. If (\*) is an equality, we apply this argument again, with $``$ replaced by $`\{𝐏^{me},\eta 𝐏^m\}`$. Then no the rational scroll has a base which defines it as incidence scroll.
###### Corollary 4.7
Let $`R=R_0^{2me}𝐏^{2me+1}`$ be the projective model of $`X_e`$ defined by $`HC_o+mf`$. For $`e1`$, $`R`$ is an incidence scroll (with base in general position) if and only if $`mh^0(𝒪_X(C_o))+(me)(h^0(𝒪_X(C_o+ef)))=2(2me+1)3`$.
###### Example 4.8
Consider $`R_0^6𝐏^7`$ with a directrix conic, a three-dimensional family of directrix quartics and a five-dimensional one of directrix quintics. This surface is a projective model of $`X_2`$ by $`HC_o+4f`$. If we apply Corollary 4.7, then we see that it is not an incidence scroll because $`4+4(42)\overline{)}=273`$.
From Propositions 4.3, 4.5 and 4.6, we obtain the proof of Theorem 4.2. In these proofs, we have built a base for each incidence rational scroll in $`𝐏^n`$. It is easy notice that if a rational scroll is an incidence scroll, then there is really only one way to choose a base, i.e., the sequence $`n_1,\mathrm{},n_r`$ is unique. Table $`1`$ contains all the incidence rational scrolls up $`𝐏^8`$.
| TABLE 1. INCIDENCE RATIONAL SCROLLS | | | | | | | | | |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| | $`n_i,i=1,\mathrm{},6`$ | | | | | | | | |
| Scroll | $`1`$ | $`2`$ | $`3`$ | $`4`$ | $`5`$ | $`6`$ | Min. Dir.$`()`$ | Normalized | $`deg\left(\text{b}\right)`$ |
| $`R_0^2𝐏^3`$ | 3 | - | - | - | - | - | $`𝐏^1\left(\mathrm{}^1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}`$ | 1 |
| $`R_0^3𝐏^4`$ | 1 | 3 | - | - | - | - | $`𝐏^1\left(1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}\left(1\right)`$ | 2 |
| $`R_0^4𝐏^5`$ | 1 | - | 4 | - | - | - | $`𝐏^1\left(1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}\left(2\right)`$ | 3 |
| $`R_0^4𝐏^5`$ | - | 3 | 1 | - | - | - | $`C_0^2𝐏^2\left(\mathrm{}^1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}`$ | 2 |
| $`R_0^5𝐏^6`$ | 1 | - | - | 5 | - | - | $`𝐏^1\left(1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}\left(3\right)`$ | 4 |
| $`R_0^5𝐏^6`$ | - | 1 | 3 | - | - | - | $`C_0^2𝐏^2\left(1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}\left(1\right)`$ | 3 |
| $`R_0^6𝐏^7`$ | 1 | - | - | - | 6 | - | $`𝐏^1\left(1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}\left(4\right)`$ | 5 |
| $`R_0^6𝐏^7`$ | - | - | 3 | 1 | - | - | $`C_0^3𝐏^3\left(\mathrm{}^1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}`$ | 3 |
| $`R_0^7𝐏^8`$ | 1 | - | - | - | - | 7 | $`𝐏^1\left(1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}\left(5\right)`$ | 6 |
| $`R_0^7𝐏^8`$ | - | - | 1 | 3 | - | - | $`C_0^3𝐏^3\left(1\right)`$ | $`𝒪_{𝐏^1}𝒪_{𝐏^1}\left(1\right)`$ | 4 |
| Number of minimum directrix curves | | | | | | | | | |
## 5 Incidence Elliptic Scrolls
Let $`X=𝐏()`$ be a ruled surface over a elliptic curve $`C`$. Let d be a divisor on $`C`$. Then:
1. $`deg(\text{d})e+2`$ there is a section $`DC_o+\text{d}f`$ and $`|D|`$ has no base points;
2. $`deg(\text{d})e+3`$ $`|C_o+\text{d}f|`$ is very ample.
If we take $`HC_o+\text{b}f`$ a very ample divisor on $`X`$ with $`m=deg(\text{b})`$, then we obtain the closed immersion $`\mathrm{\Phi }_H:XR_1^{2me}𝐏^{2me1}`$. We denote the scroll briefly by $`R`$.
We want to know if such a scroll is an incidence scroll. For this we need a base $``$. Since such a base is formed by linear spaces which contain directrix curves of the scroll, we begin to study the families of directrix curves in $`R`$. The aim of this section is to prove the following theorem.
###### Theorem 5.1
Let $`HC_o+\text{b}f`$ be a very ample divisor on $`X`$ with $`m=deg(\text{b})`$ and let $`\mathrm{\Phi }_H:XR_1^{2me}𝐏^{2me1}`$ be the closed immersion defined by $`|H|`$. Then $`R_1^{2me}`$ is an incidence scroll (with base in general position) if and only if it satisfies one of the following conditions:
1. $`e=1`$ and $`m=2`$ ($`=\{5𝐏^2\}`$);
2. $`\text{e}0`$ and $`m=4`$ ($`=\{3𝐏^3,2𝐏^5\}`$);
3. $`X`$ decomposable, $`0e3`$ and $`m=e+3`$ ($`=\{𝐏^2,(e+1)𝐏^{e+2},(3e)𝐏^{e+3}\}`$).
### Decomposable incidence elliptic scrolls
Let $`X=𝐏()`$ be a ruled surface over an elliptic curve $`C`$, defined by a normalized bundle $`𝒪_C𝒪_C(\text{e})`$ and let $`HC_o+\text{b}f`$ be a very ample divisor on $`X`$. Suppose $`R`$ an incidence scroll with base in general position.
Consider $`e1`$. Since $`h^0(𝒪_X(C_o))=1,`$ there does not exist another directrix curve of minimum degree. The following directrix curves of the scroll are linearly equivalent to $`C_o\text{e}f`$ and, by $`H`$, are directrix curves of degree $`m`$. Since $`h^0(𝒪_X(C_o\text{e}f))=e+12`$, we know that there are $`(e+1)𝐏^{m1}`$’s in $``$, whenever possible ($`m1+(e+1)(me1)4m2e5`$).
###### Lemma 5.2
Let $`X=𝐏(𝒪_C𝒪_C(\text{e}))`$ be a ruled surface over an elliptic curve $`C`$ with $`e1`$and let $`HC_o+\text{b}f`$ be a very ample divisor on $`X`$. Then $`R`$ is an incidence scroll if and only if $`deg(\text{b})=e+3`$ and $`1e3`$. Moreover, $`=\{𝐏^2,(e+1)𝐏^{e+2},(e3)𝐏^{e+3}\}`$.
Proof. We know that there are a $`𝐏^{me1}`$ and $`(e+1)𝐏^{m1}`$’s in $``$ when $`e=1,2,3`$. For $`e=1`$, the following directrix curves are divisors $`CC_o+(P+Q)f,P,QC`$. Since $`h^0(𝒪_X(C_o+(Q\text{e})f))=3`$, we show that the generic curve is irreducible. For $`m5`$, we can see that one $`𝐏^{m2}`$, $`2𝐏^{m1}`$’s and one $`𝐏^m`$ are in $``$, i.e., we impose $`4m8`$ independent conditions on $`G(1,2m2)`$ so that $`=\{𝐏^{m2},2𝐏^{m1},𝐏^m,𝐏^{2m4}\}`$ is base of an incidence scroll. Such a scroll is of type $`R_{m3}^{4m9}𝐏^{2m2}`$ (see Example 3.3) which is not elliptic. For $`m=4`$, $`=\{𝐏^2,2𝐏^3,2𝐏^4\}`$ is the base of $`R_1^7𝐏^6`$.
For $`e=2`$, we only need another condition, i.e., $`=\{𝐏^{m3},3𝐏^{m1},𝐏^{2m5}\}`$. Such a scroll is $`R_{m4}^{4m12}𝐏^{2m3}`$ (see Example 3.3), which is elliptic if and only if $`m=5`$.
For $`e=3`$, we have seen that if $`R`$ is an incidence scroll, then a $`𝐏^{m4}`$ and $`4𝐏^{m1}`$’s are in $``$. If $`m=6`$, $`=\{𝐏^2,4𝐏^5\}`$ is base of $`R_1^9𝐏^8`$. For $`e3`$ and $`m6`$, we have $`m1+(e+1)(me1)>4m2e5`$. Consequently, we cannot take $`(e+1)𝐏^{m1}`$’s as base spaces because these impose too many conditions on $`G(1,2me1)`$. Since $`m1<2me3`$, we conclude (as in Proposition 4.6) that $`R`$ is not an incidence scroll.
###### Remark 5.3
We have seen that, in $`𝐏^N`$, if the base space of greater dimension is $`𝐏^{N2}`$’s, then it is not necessary to take the number of linearly independent directrix curves. From this, we need to show if $`m1<2me3`$.
We now evaluate the case $`e=0`$. There are two choices of the divisor e. If $`\text{e}\sim ̸0`$, there are two choices of normalized $``$, namely $`𝒪_C𝒪_C(\text{e})`$ and $`𝒪_C𝒪_C(\text{e})`$. There are exactly two choices of $`C_o`$, both with $`C_o^2=0`$, namely $`C_o`$ and $`C_o\text{e}f`$. Since $`h^0(𝒪_X(C_o))=h^0(𝒪_C(C_o\text{e}f))=1`$, these are the only minimum directrix curves. If $`R`$ is an incidence scroll with base in general position, then there are $`2𝐏^{m1}`$’s in $``$. The following directrix curves are in linear systems $`|C_o+(P+\text{e})f|`$ and $`|C_o+(P\text{e})f|,PC`$, where $`h^0(𝒪_X(C_o+(P+\text{e})f))=h^0(𝒪_X(C_o+(P\text{e})f))=2`$. As in the proofs described above we can see that, if $`R`$ is an incidence scroll, there are $`2𝐏^m`$’s as base spaces. Then $`=\{2𝐏^{m1},2𝐏^m,𝐏^{2m3}\}`$ and the incidence scroll has degree $`4m6`$ and genus $`m2`$ (see Example 3.3), which is elliptic if and only if $`m=3`$.
In the other case, $`\text{e}0`$, there is a one-dimensional family of directrix curves of degree $`m`$. If $`R`$ is an incidence scroll, there are $`3𝐏^{m1}`$’s in $``$. The following directrix curves are in a linear system $`|C_o+(P+Q)f|`$ where $`h^0(𝒪_X(C_o+(P+Q)f))=4`$. By Bertini, we know that it has irreducible curves and that so is the generic curve. If $`R`$ is an incidence scroll, then $`m4`$ (no three $`𝐏^{m1}`$’s suffice and another $`𝐏^{m1}`$ imposes too many conditions). For $`m4`$, $`=\{3𝐏^{m1},𝐏^{m+1},𝐏^{2m3}\}`$ generates a scroll of genus $`m3`$ (see Example 3.3), which is elliptic if and only if $`m=4`$.
###### Remark 5.4
If we want to extend the study of incidence elliptic scrolls to incidence scrolls of genus $`g`$, then we must consider the following properties.
1. If the scroll has a unique minimum directrix curve, then the space which contains it is in $``$.
2. If the scroll has two minimum directrix curves, then the correspondent spaces are in $``$.
3. If the scroll has an one-dimensional family of minimum directrix curves, then the base has three spaces which contain three of these curves.
4. If $`mg+(e+2g)(meg)>4m2e4g1`$, then $`R`$ is not an incidence scroll ($`g=0,1`$).
### Indecomposable incidence elliptic scrolls
Let $`X=𝐏()`$ be a ruled surface over an elliptic curve $`C`$, corresponding to an indecomposable $``$. Then $`e=0,1`$ and there is exactly one ruled surface over $`C`$ for each of these values of $`e`$. Let $`HC_o+\text{b}f`$ be a divisor with $`deg(\text{b})=me+3`$, which defines $`\mathrm{\Phi }_H:XR_1^{2me}𝐏^{2me1}`$. We want to know if such a scroll is an incidence scroll. For this, we need to find a suitable number of base spaces in $`𝐏^{2me1}`$ whose dimension is $`2me3`$
If $`e=0`$ then $``$ is an extension of $`𝒪_C`$ by $`𝒪_C`$ i.e. given by an exact sequence $`0𝒪_C𝒪_C0`$. Since $`C_0^2=0`$ and $`X`$ is indecomposable, there is a unique directrix of minimum degree equal to $`m`$. For any $`PC`$, $`h^0(𝒪_X(C_o+Pf))=2`$. Then we have an one-dimensional family $`|C_o+Pf|`$ of directrix curves of degree $`m+1`$, for any $`PC`$. Suppose that $`R`$ is an incidence scroll with base in general position. An easy computation shows that $`=\{𝐏^{m1},2𝐏^m,𝐏^{n_1},\mathrm{},𝐏^{n_r}\}`$ where $`mn_i2m3`$ and $`_{i=1}^r(2m2n_i)=m`$. Then we can obtain another incidence elliptic scroll in $`𝐏^{2m}`$ with $`e=1`$ and base $`\dot{}=\{𝐏^{m1},𝐏^m,𝐏^{m+1},𝐏^{n_1+1},\mathrm{},𝐏^{n_r+1}\}`$ (separate $`𝐏^{m1}`$ and $`𝐏^m`$ in $``$). This happens when $`m=3,n_1=2`$ and $`n_2=3`$, which contradicts the hypothesis $`mn_i`$. Hence there are no indecomposable elliptic scrolls with $`e=0`$ defined by incidences.
If $`e=1`$ then we have an exact sequence $`0𝒪_C𝒪_C(P)0`$. Then $`\mathrm{\Phi }_H(C_o)`$ is a directrix curve of degree $`m+1`$. If $`D`$ is a section with $`D^2=1`$, then $`DC_o+(\delta \text{e})f`$ with $`deg(\delta )=1`$. In fact, the sections $`C_o`$ with $`C_o^2=1`$ form a one-dimensional algebraic family parametrized by $`C`$ such that no two of them are linearly equivalent. If $`R`$ is an incidence scroll, then $`=\{2𝐏^m,𝐏^{n_1},\mathrm{},𝐏^{n_r}\}`$ with $`mn_i2m2`$ and $`_{i=1}^r(2mn_i1)=2m1`$. We can obtain another incidence elliptic scroll in $`𝐏^{2m+1}`$ with $`\text{e}\sim ̸0`$ and $`e=0`$ (separate $`2𝐏^m`$’s in $``$). Consequently, there is only one incidence scroll with $`e=1`$. This happens when $`m=n_i=2`$ and $`r=3`$. So $`\mathrm{\Phi }_H(X)=R_1^5𝐏^4`$ is an incidence scroll with base $`=\{5𝐏^2\}`$.
In Table $`2`$ we have compiled all incidence elliptic scrolls. The interest of this table is that each scroll is a degenerate scroll of the remainder of them.
| TABLE 2. INCIDENCE ELLIPTIC SCROLLS | | | | | | | |
| --- | --- | --- | --- | --- | --- | --- | --- |
| | $`n_i,i=1,\mathrm{},4`$ | | | | | | |
| Scroll | $`2`$ | $`3`$ | $`4`$ | $`5`$ | Min. Dir. | Normalized | $`deg\left(\text{b}\right)`$ |
| $`R_1^5𝐏^4`$ | 5 | - | - | - | $`C_1^3𝐏^2`$ | $`Ext^1(𝒪_C\left(P\right),𝒪_C)`$ | 2 |
| $`R_1^6𝐏^5`$ | 2 | 3 | - | - | $`C_1^3𝐏^2`$ | $`𝒪_C𝒪_C\left(\text{e}\right)(\text{e}\sim ̸0`$) | 3 |
| $`R_1^7𝐏^6`$ | 1 | 2 | 2 | - | $`C_1^3𝐏^2`$ | $`𝒪_C𝒪_C\left(P\right)`$ | 4 |
| $`R_1^8𝐏^7`$ | 1 | - | 3 | 1 | $`C_1^3𝐏^2`$ | $`𝒪_C𝒪_C\left(PQ\right)`$ | 5 |
| $`R_1^8𝐏^7`$ | - | 3 | - | 2 | $`C_1^4𝐏^3`$ | $`𝒪_C𝒪_C`$ | 4 |
| $`R_1^9𝐏^8`$ | 1 | - | - | 4 | $`C_1^3𝐏^2`$ | $`𝒪_C𝒪_C\left(PQR\right)`$ | 6 |
Acknowledgements: Rosa Cid-Mu$`\stackrel{~}{n}`$oz was supported by a fellowship of Xunta de Galicia (Spain).
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# Untitled Document
Market Ecology, Pareto Wealth Distribution and Leptokurtic Returns in
Microscopic Simulation of the LLS Stock Market Model
Sorin Solomon <sup>1</sup> and Moshe Levy <sup>2</sup>
<sup>1</sup> Racah Institute of Physics, Givat Ram 91104,
<sup>2</sup> School of Business Administration, Mount Scopus 91905,
Hebrew University of Jerusalem
Abstract:
The LLS stock market model is a model of heterogeneous quasi-rational investors operating in a complex environment about which they have incomplete information. We review the main features of this model and several of its extensions. We study the effects of investor heterogeneity and show that predation, competition, or symbiosis may occur between different investor populations. The dynamics of the LLS model lead to the empirically observed Pareto wealth distribution. Many properties observed in actual markets appear as natural consequences of the LLS dynamics: truncated Levy distribution of short-term returns, excess volatility, a return autocorrelation ”U-shape” pattern, and a positive correlation between volume and absolute returns.
1. The LLS Model
LLS is a microscopic representation model of the stock market. Its details and some generalizations of it can be found in . In the present account we introduce the basic LLS ideas and the model main results. We consider below a market with only two investment options: a bond and a stock (see for an extension to a multiple stocks case). The model involves a large number of virtual investors characterized each by a current wealth, portfolio structure, probability expectations and risk taking preferences. These personal characteristics come into play in each investor’s decision making process as schematically seen in Fig. 1.
The bond is assumed to be a risk-less asset yielding a return at the end of each time period. The bond is exogenous and investors can buy from it as much as they wish at a given rate.
The stock is a risky asset with overall returns rate $`H(t)`$ composed of two elements: (i). Capital gain (loss): If an investor holds a stock, any rise (fall) in the market price of the stock contributes to an increase (decrease) in the investors’ wealth. (ii). Dividends: The company earns income and distributes dividends.
Each investor $`i`$ is confronted with a decision where the outcome is uncertain: which is the optimal fraction $`X(i)`$ of his/her wealth to invest in stock? According to the standard theory of investment each investor is characterized by a utility function (of its wealth) $`U(W)`$ that reflects his/her personal risk taking preference (here we take for simplicity $`U(W)=\mathrm{ln}W`$ see for a prospect theory extension). The optimal $`X(i)`$ is the one that maximizes the expected value of his/her $`U(W)`$ (we take into account all the unknown factors influencing decision-making (such as liquidity constraints or deviations from rationality) by adding a small random variable (or ”noise”) to the optimal proportion $`X(i)`$). The expected value of $`U(W)`$ depends of course on the expected probabilities for the various values of $`H`$ to be realized in the future. In LLS the investors expectations for the future $`H`$’s are based on extrapolating the past values. More precisely each investor recalls the last k returns on the stock and expects that each of them may take place again with equal probability. The extrapolation range k differs between various investors and it will be in the sequel of this paper the main parameter inducing market inhomogeneity.
At fixed time intervals each investor revises the composition of its portfolio and decides for a new market order. The aggregate of these orders determines the new stock price by the market clearance condition as explained below. Once each investor decides on the proportion of his/her wealth $`X(i)`$ that (s)he wishes to hold in stocks, one can derive the number of stocks $`N(i,p_h)`$ it wishes to hold corresponding to each hypothetical stock price $`p_h`$. Since the total number of shares in the market $`N`$, is fixed there is a particular value of the price $`p`$ for which the sum of the $`N(i,p)`$ equals $`N`$. This value $`p`$ is the new market equilibrium price. Upon updating accordingly the traders’ portfolios, wealth and list of last $`k`$ returns, one is ready for the next market iteration. This process is repeated for each time step, and the market prices are recorded throughout the run. Figure 1: The Flow Chart of the LLS market framework
2. Crashes, Booms and Cycles
The LLS model provides already at the level of a quite homogenous traders population a convincing description of the emergence of cycles of booms and crashes in the stock markets. In a market with one species of investors all having a homogenous memory (extrapolation) range spanning the last k returns of the stock, the stock price alternates regularly between two very different price levels. The explanation for this behavior is as follows:
Assume that the rate of return $`H(t)`$ on the stock at a time $`t`$ is higher than the oldest remembered return ($`H(tk)`$). The addition of $`H(t)`$ and the elimination of $`H(tk)`$ creates then a new distribution of past returns that is better than the previous one. Since the LLS extrapolating investors use the past $`k`$ returns to estimate the distribution of the next period’s return, they will be lead to be more optimistic and increase their investments in the stock. This, in turn, will cause the stock price to rise, which will generate an even higher return.
This positive feedback loop stops only when investors reach the maximum investment proportion (i.e. $`X(i)=100\%`$: we do not allow borrowing or short selling), and can no longer increase their investment proportion in the stock. The dividend contribution to the returns is small compared with this high price at this stage. In the absence of noise the returns on the stock at this plateau converge to a constant growth rate which is just slightly higher than the riskless interest rate (see ). In other words, in the absence of noise the price remains almost constant, growing only because of the interest paid on the bond (more money entering the system and being invested in the stock).
When there is some noise in the system the price fluctuates a little around the asymptotic high level, because of the small random fluctuations in the investment proportions. These fluctuations generate some negative returns (on a downward fluctuation) and some high returns (when the price goes back up). One might suspect that a large downward fluctuation might trigger a reverse positive feedback effect, trader expectations will lower, investment proportions will decrease, the price will drop, generating further negative returns and so on: a crash. This can happen during the ”plateau” period but only after the previous sharp price boom which generated an extremely high return, is forgotten. And, indeed, this is exactly what happens. Since it takes $`k`$ steps to forget the boom, the high price plateaus are a bit longer than the extrapolation span ($`k`$ days to forget the boom $`+O(1)`$ more days until a large enough negative fluctuation occurs).
The crash generates a disastrous return and, until it is forgotten, investment proportions and hence the price remains very low. When the price is low, the dividend becomes significant and the returns on the stock are relatively high (compared with the bond). Once the crash has been forgotten, all the returns that are remembered are therefore high, and the price jumps back up. Thus, the low price plateaus are $`k`$ steps long. This completes one cycle, which is repeated throughout the run. This (quasi-)periodicity is best viewed in the Fourier transform of the price time evolution (Fig. 2) as a series of narrow peaks around the frequency $`2k+O(1)`$ and its harmonics (note however that the dynamics is not perfectly periodic and therefore in spite of its simplicity, according to some mathematical criteria it may fall into the ”complex” category. In the present paper we reserve however the term ”complex” for dynamics that are truly complicated to the degree that they do not admit simple verbal or mathematical description or understanding).
The homogenous stock market described above exhibits booms and crashes. However, the homogeneity of investors leads to unrealistic periodicity. As shown below, when there is more than one investor species the dynamics becomes much more complex and realistic.
Figure 2 : The Fourier transform of the price in a market with one species with extrapolation range $`k=10`$. The market contained 10000 traders that had initially equal wealth invested half in stock and half in bonds.
3. Realistic Features in LLS with Many Species
Our numerical experiments within the LLS framework have found that already a small number of trader species (characterized by different extrapolation ranges $`k`$) leads qualitatively to many of the empirically observed market phenomena.
In reality, we would expect not just a few trader types, but rather an entire spectrum of investors. When the full spectrum of different trader species (fundamentalists and various other types - see for the detailed operational definition) is considered it turns out that ”more is different” : the price dynamics becomes realistic: booms and crashes are not periodic or predictable, and they are also less frequent and dramatic. At the same time, we still obtain many of the usual market anomalies described by the experimental studies (however in the limit of infinite times or infinite number of investors, the dynamics may revert to predictable patterns ).
We list below a few such realistic features:
Return Autocorrelations: Momentum and Mean-Reversion
In the heterogeneous population LLS model trends are generated by the same positive feedback mechanism that generated cycles in the homogeneous case (section 2): high (low) returns tend to make the extrapolating investors more (less) aggressive, this generates more high (low) returns, etc.
The difference between the two cases is that in the heterogeneous case there is a very complicated interaction between all the different investor species and as a result there are no distinct regular cycles but rather, smoother and more irregular trends. There is no single cycle length - the dynamics is a combination of many different cycles corresponding to the many extrapolation ranges $`k`$. This makes the autocorrelation pattern also smoother and more continuous. The return autocorrelations in the heterogeneous LLS model conform to the empirical findings: In the short run the autocorrelation is positive - this is the empirically documented phenomenon known as momentum: high returns during a trading quarter tend to be followed by more high returns in the following months, (and low returns tend to be followed by more low returns). In the longer run the autocorrelation is negative (after a few years of boom, one usually experiences a few ”dry” years), which is known as mean-reversion. For even longer lags the autocorrelation eventually tends to zero . The short run momentum, longer run mean-reversion, and eventual diminishing autocorrelation creates the general ”U-shape” that is found in empirical studies .
Excess Volatility
In markets with a large fundamentalist population (see for their detailed operative definition in the LLS model), the price level is generally determined by the fundamental value of the stock. However, the market extrapolating investors occasionally induce temporary departures of the price away from the fundamental value. These temporary departures from the fundamental value make the price more volatile than the fundamental value.
Following Shiller’s methodology we measured the standard deviations of the detrended price and fundamental value. Averaging over $`100`$ independent simulations we found respectively $`27.1`$ and $`19.2`$, which is an excess volatility of $`41\%`$.
Heavy Trading Volume
In an LLS market with both fundamentalists and market extrapolating investors (over various $`k`$ ranges), shares change hands continuously between the various groups:
When a ”boom” starts, the extrapolating investors observe higher ex-post returns and become more optimistic, while the fundamentalists view the stock as becoming overpriced and become more pessimistic. Thus, at this stage the market extrapolators buy most of the shares from the fundamentalists.
When the stock crashes, the opposite is true: the extrapolators are very pessimistic, but the fundamentalists buy the stock once it falls below the fundamental value. Thus, there is substantial trading volume in this market. The average trading volume in a typical LLS simulation was about $`1,000`$ shares per period, or about $`10\%`$ of the total outstanding shares.
Volume is Positively Correlated with Absolute Returns
The typical scenario in an LLS run is that when a positive trend is induced by the extrapolating investors, the opinions of the fundamentalists and the extrapolating investors change in opposite directions:
\- The extrapolating investors see a trend of rising prices as a positive indication about the future return distribution, while
\- The fundamentalists believe that the higher the price level is (the more overpriced the stock is), the harder it will eventually fall.
The exact opposite holds for a trend of falling prices. Thus, price trends are typically interpreted differently by the two investor types, and therefore induce heavy trading volume. The more pronounced the trend (large price changes), the more likely it is to lead to heavy volume.
In order to verify this relationship quantitatively we regressed volume $`V(t)`$ on the absolute rates of return $`r(t)`$ for $`100`$ independent simulations. We run the regressions:
$$V(t)=a+b|r(t)|+random(t)$$
We found an average value of $`870`$ for $`b`$ with an average t-value of $`5.0`$. Similar results were obtained for time lagged-returns.
4. Predation, Competition and Symbiosis between Trader Species
In section 2 it was explained that a homogenous population of traders that extrapolate the last $`k`$ returns leads to cycles of booms and crashes of period $`2k+O(1)`$. When there are two species with extrapolation ranges $`k_1`$ and respectively $`k_2`$, we observe sharp irregular transitions between eras where one species dominates (cycles of period $`2k_1+O(1)`$) and market eras where the other species dominates (cycles of period $`2k_2+O(1)`$). When the number of trader species is three, there are dramatic qualitative changes: generically, the dynamics becomes complex. We show that complexity is an intrinsic property of the stock market. This suggests an alternative explanation to the widely accepted but empirically questionable random walk hypothesis. We discuss below some of the market ecologies possible with only two species of traders. Of course the picture becomes more complex later, when 3 or more species are introduced.
Market Ecologies with Two Trader Species
When there are two trader species with different extrapolation spans it turns out that the nature of the dynamics is determined by the ratio of the extrapolation spans of the two species. In we performed a qualitative theoretical analysis of this phenomenon and supported it by microscopic simulations. We showed that in market eras in which one species (of extrapolation range $`k_0`$) dictates the dynamics (i.e. boom-crash cycles have periods of length $`2k_0+O(1)`$) the second species (with extrapolation range $`k`$) has generically the following performance:
A : If $`k_0<k<2k_0`$ then $`k`$ is performing very poorly (looses money)
B : If $`2nk_0<k<(2n+1)k_0`$ (with $`n`$ natural number), then $`k`$ is doing relatively well
C : If $`(2n+1)k_0<k<2nk_0,n>1,`$ then $`k`$ does better than in A but worse than in B
D : If $`k<k_0`$, then $`k`$ is doing well .
These facts turned out sufficient to understand the main 3 cases that a 2-species ecology can display:
Case 1: predator - prey dynamics
If one considers one species with an extrapolation range $`k_1=10`$ and a second species with an extrapolation range $`k_2=14`$ it turns out that the resulting ecology dynamics is a predator-prey one. In fact the LLS market dynamics leads in this case to the extinction (total impoverishment) of the $`k_1`$ species: after some time the entire wealth on the market belongs to the species $`k_1=10`$ (Fig. 3). As a consequence, the market price presents clear cycles of booms and crashes of periodicity clustered around $`24=210+O(1)`$.
This is easily understood since according to the property A above the $`k_2=14`$ population is performing poorly when the $`k_1`$ dictates the market periodicity while the population $`10`$ is performing well according D in the hypothetical periods when $`k_2=14`$ dictates the market periodicity. This is only an example of a large class of parameters that lead to predator-prey systems and which may result in the total extinction of one of the species.
Figure 3 : Fraction of the wealth that the species $`k_1=10`$ possesses in Case 1. The traders in the market belonged to 2 species consisting each of 5000 traders. Each trader owed at the beginning 5000 dollars in cache and 5000 shares (worth each 1.4 dolars).
Case 2: competitive species
If one chooses $`k_1=10,k_2=26`$, the species with extrapolation range $`26`$ gains during the periods when the species $`k_1=10`$ dominates (property B) but species $`k_1=10`$ gains when the species $`k_2=26`$ dominates (property D). It is therefore reasonable that one species can not dominate the other indefinitely. Indeed, a look at the fraction of the wealth held by the species with extrapolation range $`k_1=10`$ reveals alternating eras of dominance (Figure 4). This is also reflected in the alternance between price cycles ($`56`$) corresponding to $`k_2=26`$ and price cycles ($`24`$) corresponding to $`k_1=10`$. Clearly this alternance between the 2 species corresponds to a classical competitive ecology, in which two competing species take turns in dominating the ecology. Note however that most of the time it is the population $`k_2`$ which dominates the wealth. This seems to be a generic tendency in the long runs limit.
Figure 4: Fraction of the wealth that the species $`k_1=10`$ possesses in Case 2. The initial conditions were similar to Figure 3.
Case 3 symbiotic species
In the case $`k_1=10,k_2=36`$, similarly to the $`1026`$ market, the investors with extrapolation range $`k_2=36`$ are doing better than those with extrapolation range $`k_1=10`$ when $`k_1=10`$ dictates the dynamics (cf. property C). On the other hand $`k_1=10`$ are doing better when the species $`k_2=36`$ dictates the dynamics (cf. D). Hence, we may speculate that again we will find alternating eras of dominance. Figure 5 shows that this is not the case. The difference between this case and the $`1026`$ case is that here the market remains stuck in a ”metastable” state: the extrapolation range $`36`$ population never gains enough wealth to dictate long cycles. Thus, the system remains in a state of symbiosis throughout the run: the price cycles correspond to the short species extrapolation range span $`k_1=10`$ while $`7080\%`$ of the wealth stays with the long extrapolation span species $`k_2`$.
For very long $`k_2`$ extrapolation ranges, the share of the total wealth detained by $`k_2`$ can be even larger (approaching unity).
Figure 5 -Fraction of the wealth that the species $`k_1=10`$ possesses in case 3.
In conclusion has uncovered a quite lively ecology of the traders populations in the LLS model and
”observed phenomena ranging from complete dominance of one population to alternating eras of domination and to symbiosis. . . . Our results suggest that complexity is an intrinsic property of the stock market. The dynamic and complex behavior of the market need not be explained as an effect of external random information. It is a natural property of the market, emerging from the strong nonlinear interaction between the different investor subgroups of the market . . .”
The main source of endogenous dynamics in the LLS model turns out to be the feedback between the market price fluctuations and the wealth of the investors belonging to various species:
\- On one hand the wealth of the investors determines their influence on the price changes (at the short range): e.g. the richest determine the periodicity of the boom-crash cycles.
\- On the other hand, the variations in the price determine changes in the distribution of wealth, which iterated over longer time intervals, result in changes in the market price cycle periodicity regime.
The entire cycle of rise and fall of a given species can be schematically described as: $``$ The species has by chance a (momentary) winning strategy $``$ Investors belonging to the species gain wealth $``$ Overall wealth of the investors belonging to the species increases $``$ Bids of investors belonging to the species become large $``$ Investor bids influence the market price adversely (self-defeating) $``$ Trading of investors belonging to the species becomes inefficient $``$ Investors lose money $``$ Investors belonging to the species become poor $``$ Species wealth and market relevance vanish $``$ Other species with different strategies become winners $``$ Cycle re-starts (with the new winning strategies).
A few comments are in order:
1. the concept of efficient strategy is only a temporary one as it depends crucially of the state of the market: by its very efficiency at a certain moment, a strategy prepares the seeds of its failure in the future.
2. the biological and cognitive analogies are useful but their limits should be understood:
\- in biology, the species selection mechanism is based on the disappearance of the inefficient individuals.
\- in the learning adaptive agents’ case, the individuals discard loosing strategies for new ones.
In the LLS market framework, while it is possible to include the above effects, they are not necessary: the strategies selection takes place automatically by their carriers (traders belonging to the species) losing or gaining: for the market to be efficient, no a priori intelligence nor explicit criteria for the evaluation and comparison of market performance are required: just the natural (Adam Smith’s ”invisible hand”) market mechanisms.
3. While the adverse influence on the market price implied by the large orders coming from rich agents’ leads automatically to inefficiency in their operations (except for rich agents which follow a buy-and-hold strategy and therefore do not influence (adversely) the market), the mere lack of market influence due to poverty does not guarantee a winning strategy. It is necessary therefore that there are enough strategies and enough agents in the market for insuring its efficiency.
Three Investor Species
One might suspect that the three species dynamics is a natural extension of the two species dynamics. Instead of alternating between two cycle lengths the system may just alternate between the three possible states of dominance. Figures 6-7 shows that this is not at all the case. These figures depict a typical part of the dynamics of a three species market, with extrapolation spans of 10, 141 and 256 respectively. With the introduction of a third species the system has underwent a qualitative change: there is no specific cycle length describing the time series. Instead, we see a mixture of different time scales: the system has become complex. Prediction becomes very difficult, and in this sense the market is much more realistic. Figure 6 shows the power struggle between the three species while the Figure 7 depicts the Fourier transform of the price evolution during this run.
Although the dynamics is complex, it is clear from Figures 6 and 7 that there is an underlying structure, which perhaps may be analyzed by the properties A, B, C, D and their generalizations. For instance it would appear from the Figure 6 that 141 and 256 take turns in dominating while 10 has a chance to a non-vanishing wealth share only occasionally in the transition intervals between 141 and 256 dominated eras. The dynamics generated by only three investor species can be extremely complex, even without any external random influences.
Figure 6: The species wealths in a market with 3 species of extrapolation ranges of respectively 10, 141 and 256 days. Initially the 3 species possessed equal wealth distributed equally between stock and bond. Each species consisted of 1000 traders.
Figure 7: Fourier transform of the stock price time evolution in the market described in Fig 6.
5. Generalized Lotka-Volterra models for markets with multiple species
Inspired by the above facts we devised an effective dynamics that stylized the features uncovered in the LLS model and extended them to a more generic framework. Instead of following in detail the way the market price influence each species and individual $`i`$ , we assumed that this influence can be represented through multiplying their wealth $`w_i(t)`$ by stochastic multiplicative factors $`\lambda _i(t)`$. This is natural in the LLS model in which the investments of the individuals (and consequently their returns) are fractions of their wealth (as implied by the constant relative risk aversion utility functions). The stochastic proportionality between personal returns and personal wealth is consistent with the real data that show that the (annual) individual income distribution is proportional to the individual wealth distribution .
We proposed therefore a model including the above stochastic autocatalytic properties of the capital as well as the cooperative, diffusive and competitive/ predatory interactions between the species. The resulting model turned out to be a straightforward generalization of the Lotka-Volterra system (discrete logistic equation) well known previously in population biology:
$$w_i(t+1)=\lambda _i(t)w_i(t)+\underset{k}{}a_kw_k(t)w_i\underset{k}{}b_kw_k(t)$$
Where the sum is over all the $`N`$ traders participating in the market. There are a few other (mutually non-exclusive) possible interpretations for $`w_i`$ in addition to the individual wealth: the wealth associated with a particular investing strategy, the capitalization associated with a particular company/industry, or the number of investors following a common trend (herd).
Particularly interesting cases were studied subsequently:
1) The linear case where the total wealth diverges to larger and larger values (inflation, production):
$$w_i(t+1)=\lambda _i(t)w_i(t)+\underset{k}{}a_kw_k(t)$$
2) The case in which the binary interactions between individuals are expressible in terms of interactions with the total wealth $`W=_kw_k`$:
$$w_i(t+1)=\lambda _i(t)w_i(t)+aW(t)bw_iW(t)$$
3) The case in which individual wealth is bounded from below by a certain fraction $`c`$ of the average wealth $`\overline{w}=W/N`$ :
$$w_i(t+1)=\lambda _i(t)w_i(t)$$
except if
$$\lambda _i(t)w_i(t)<c\overline{w}$$
when
$$w_i(t+1)=c\overline{w}$$
4) The case of the random multiplicative wealth dynamics
$$w_i(t+1)=\lambda _i(t)w_i(t)$$
with variable number of traders :
\- traders which fall below a certain fixed minimal wealth $`w_{min}`$ drop from the market
\- a number of traders proportional to the total wealth increase:
$$\mathrm{\Delta }N=c(W(t+1)W(t))/w_{min}$$
join the market at each time step i.e. the average wealth remains constant in this process:
$$\overline{w}=w_{min}/c$$
.
One assumes that each of the new traders brings an initial investment equal to $`w_{min}`$ which means that the total amount of added wealth is
$$\mathrm{\Delta }Nw_{min}=c(W(t+1)W(t))$$
i.e. a fraction $`c`$ of the total wealth increase $`(W(t+1)W(t))`$.
In most of these systems were assumed asynchronous: at each time step, only one (randomly chosen) $`w_i`$ was updated. Very striking generic results can be obtained with all these models in certain relevant regimes. We limit ourselves below to the universal scaling properties (power laws).
6. Pareto Law in LLS and Lotka-Volterra models
The efficient market hypothesis and the Pareto law are some of the most striking and basic concepts in economic thinking. It is therefore very important that our models above succeed to connect them in a very essential way.
Let us discuss this in more detail. More than a hundred years ago, Pareto discovered that the number of individuals with wealth (or incomes) with a certain value $`w`$ is proportional to $`w^{1\alpha }`$. This later became known as the Pareto Law. The LLS model treats the individual investor wealth as a crucial quantity, and it views its feedback relation to the market dynamics as the main source driving the endogenous dynamics of the market.
It turns out that in the conditions in which the participants in the market do not have a systematic advantage one over the other (which is in fact expected in an efficient market), a dynamics of the LLS type leads always to a Pareto law. The actual value of the exponent $`\alpha `$ depends on the particular parameters used in the model. Mainly, as explained below $`\alpha `$ is influenced by the social security policy. If one does do not allow any individual to become poorer than a certain fraction c of the current average wealth then, for a wide range of conditions $`\alpha =1/(1c)`$. This is confirmed in Figure 8 which plots the wealth distribution in the LLS model with $`k=3`$ and $`c=0.2`$ (and $`U(W)=\mathrm{ln}W`$).
Figure 8: The wealth distribution of the investors in an LLS model with a poverty line of $`c=20\%`$ of the average wealth. One a double logarithmic scale one obtains a straight line with slope 2.2 corresponding to an $`\alpha `$ of 1.2.
The market consisted of 10000 traders and the measurement was performed as a ”snapshot” after 1 000 000 ”thermalization” market steps. Initially all the traders had equal wealth ($`\$`$1000) equally distributed between bond and stock.
In fact it has been proven \[11-13\] theoretically that any of the effective dynamics of the type 1-4 with $`\lambda _i`$ distribution independent on $`i`$ or $`w_i`$ leads always to a power law Pareto distribution.
In a wide range of models, the generic rule is that
$`\alpha `$ = 1/(1-c)
where $`c`$ is essentially the market global impact factor :
c= exogenous new capital ADDED to the market / increase in stock capitalization due to market price increase
Since the increase in the capitalization is the increase in wealth that the owners incurr upon their investment of new capital, the ratio can be also interpreted as the long range market return factor
c = 1/ (long range market return factor).
Let us explain in short how such results were obtained . The crucial observation is that for large $`w_i`$ values, the non-stationary multiplicative system of interacting $`w_i`$’s is formally equivalent to a statistical mechanical (additive) system in thermal equilibrium when expressed in terms of the variables $`u_i(t)=\mathrm{ln}w_i/\overline{w}(t)`$. For instance the system 3) is mapped into a system of particles diffusing in an energy potential field $`u`$ with a ground level $`u_0=\mathrm{ln}c`$. In thermal equilibrium, all such systems (independent on the details of the interactions between their particles) have an universal probability distribution discovered by Boltzmann more than 100 years ago:
$$P(u)exp(\alpha u)$$
When re-expressed in terms of the original $`w_i`$ variables this gives a Pareto power law distribution:
$$P(w)w^{1\alpha }$$
The exponent $`\alpha `$ can be estimated from the integrals representing the total wealth and the total number of traders . For instance, in the models 3-4, in the limit of $`N\mathrm{}`$ , the result is:
$$\alpha =1/(1w_{min}/\overline{w})$$
i.e.
$$\alpha =1/(1c)$$
Thus, the Pareto law is the exact analogue of the Boltzmann law for stochastic systems that are multiplicative rather then additive.
For finite $`N`$ the $`\alpha `$ is given by the implicit transcendental equation :
$$N=[((1(N/c)^\alpha )/\alpha ]/[((1(N/c)^{\alpha 1})/(\alpha 1)]$$
which for $`Ne^{1/c}`$ gives approximately:
$$\alpha \mathrm{ln}N/(\mathrm{ln}(N/c))<1$$
which incidentally means that in this regime all the wealth belongs to only a few individuals.
In the system 1 defined above, in the appropriate thermodynamic limit The analog result is:
$$\alpha =1/(1c)$$
with
$$c2a/(<\lambda >+\sigma ^2/2)$$
where $`\sigma `$ is the standard deviation of $`\lambda `$
And in the case 2:
$$c2a/(\sigma ^2+a^2)$$
independent on $`b`$ and $`<\lambda >`$.
7. Market Efficiency, Pareto Law and Thermal Equilibrium
The formal equivalence between the non-stationary systems of interacting $`w_i`$’s and the equilibrium statistical mechanics systems governed by the universal Boltzmann distribution has far reaching implications: it relates the Pareto distribution to the efficient market hypothesis: In order to obtain a Pareto power law wealth distribution it is necessary and sufficient that the returns of all the strategies practiced in the market are stochastically the same, i.e. there are no investors that can obtain ”abnormal” returns.
Therefore, the presence of a Pareto wealth distribution is a measure of the market efficiency in analogy to the Boltzmann distribution whose presence is a measure to thermal equilibrium. Indeed physical systems which are not in thermal equilibrium (e.g. are forced by some external field - say by laser pumping) do not fulfill the Boltzmann law. Similarly, markets that are not efficient (e.g. when some groups of investors make systematically more profit than others) do not yield power laws (see Fig 9). Optimal market and power laws are the short time and long time faces of the same medal/phenomenon.
This analogy is consistent with the interpretation of market efficiency as analog to the Second law of Thermodynamics:
\- one can extract energy (only) from systems that are not in thermal equilibrium
\- one can extract wealth (only) from markets that are not efficient.
\- by extracting energy from a non-equilibrium thermal system one gets it closer to an equilibrium one.
\- by extracting wealth from a non-efficient market one brings it closer to an efficient one
-in the process of approaching thermal equilibrium, one also approaches the Boltzmann energy distribution
\- in the process of approaching the efficient market one also approaches the Pareto wealth distribution.
-by having additional knowledge on a thermodynamic system state one can extract additional energy (e.g. Maxwell demons gedanken experiment)
-by having additional knowledge on a financial system one can extract additional wealth.
This double analogy
thermodynamic equilibrium $``$ efficient market
Boltzmann law $``$ Pareto law
holds in the details of their microscopic origins:
\- the convergence to statistical mechanics equilibrium depends on the balance of the probability flow entering and exiting each energy level. This is usually insured microscopically by the fact that the a priori probability for a molecule to gain or loose an energy quanta in a collision is the same for any energy level with the exception of the collisions including molecules in the ground state which can only receive (but not give) energy.
\- in the stochastic models 1-4, the convergence of the wealth to the power-law is insured by the balance of flow of investors from one level of \[log (relative wealth)\] to another. At the individual level, this is enforced by all the individuals having the same (relative) returns probability distribution (except for the individuals possessing the lowest allowed wealth). If this condition is not fulfilled, one does not get a wealth distribution power law.
These facts should guide us in the practical runs in establishing which combinations of strategies (or the strategy selection strategies) are producing a realistic market ”in the Pareto sense”. In Figure 9 one sees the wealth distribution in a model in which there are 2 trader species with slightly different distributions of $`\lambda `$. One sees that even a small violation of the $`\lambda `$ uniformity leads to significant departures from the Pareto law which are inconsistent with the historical experimental facts. The absence of such departures in real life is a strong indication of the market efficiency in the weak stochastic sense (that all investors have stochastically the same relative returns distribution).
Figure 9: Wealth distribution for 2 investor species with different return distributions. Model 3 was used with a lower wealth bound of $`c=20\%`$.
$`\lambda `$ is randomly drawn. For the first species $`\lambda `$ is 1.10 or 0.95 with equal probability. For the second ”more talented” species $`\lambda `$ is 1.11 or 0.96 with equal probability. The 2 species were each composed of 10000 traders with initially equal wealth (1000 dollars each).
The measurement of the wealth distribution was performed after a ”thermalization period” of 100 000 wealth updatings.
8. Leptokurtic Market Returns in LLS
It has been long known that the distribution of stock returns is leptokurtic or ”fat-tailed”. Furthermore, a specific functional form has been suggested for the short-term return distribution (at least in a certain finite range) - the Levy distribution . This feature is present in the LLS model, and is directly related to the Pareto distribution of wealth.
The central limit theorem insures that in a wide range of conditions the distance reached by a random walk of $`t`$ steps of average squared size $`s^2`$ is a Gaussian with standard deviation $`s\sqrt{t}`$. Suppose that at time $`t=0`$ one has $`N`$ positive numbers $`w_i(0);i=1,\mathrm{}.,N`$ of order 1 and sum $`W(0)`$. Suppose that at each time step one of the numbers varies (increases or decreases) by a fraction $`s_i(t)1`$ extracted from a random distribution with average squared $`s^2`$ (and $`0`$ mean). What will be the probability distribution of the sum $`W(t)`$ after $`t`$ steps? According to the central limit theorem this would be the Gaussian
$$P(W,t)=1/(\sqrt{2\pi ts^2})e^{(W(t)W(0))^2/2ts^2}$$
since it consists of $`t`$ steps of average squared size $`s^2`$.
If one interprets $`w_i(t)`$ as the value of the stocks owned by the trader $`i`$ at time $`t`$, then $`W(t)=_iw_i`$ is the total market value of the stock and therefore $`(W(t)W(0))/W(0)`$ is the relative stock return for the time interval $`t`$.
One sees that if the central limit theorem would hold, one would predict a Gaussian stock returns distribution. This is in fact the case for real stocks and time intervals longer than a few weeks. For significantly shorter times $`t`$ however, the distribution of returns is very different from a Gaussian. Even though the exact shape of the returns distribution is not yet established experimentally, it is generally agreed that in certain ranges (typically ”in the tails”- i.e. for large $`w_i`$ values) it fits better a power law rather than a Gaussian.
Such a situation can in principle be explained by the following scenario:
Suppose that at time $`t=0`$ one has an arbitrarily large number of positive numbers $`w_i(0)`$. Suppose moreover that the probability distribution for the sizes of $`w_i(0)`$ is
$$P(w)w^{1\alpha }$$
Suppose that at each time step one of the $`w_i`$’s varies (increases or decreases) by a fraction $`s_i(t)1`$ of average squared size $`s^2`$. What will be the probability distribution of the variation of the $`w_i`$’s sum $`W(t)W(0)`$ after $`t`$ steps?
One is tempted to think that the correct answer is given by
$$P(W,t)=1/(\sqrt{2\pi ts^2})e^{(W(t)W(0))^2/2ts^2}$$
for some $`s`$. However this is wrong. Indeed, assuming such an $`s`$ exists would imply that the probability for the sum variation $`W(t)W(0)`$ to be 10 after a time $`t=1/(2s^2)`$ is:
$$P(W(t)=W(0)+10,t=1/(2s^2))e^{10^2}10^{32}$$
while in reality a lower bound for the probability of getting $`W(t)W(0)=10`$ in just one step it is obviously that given by
$$P(w)w^{1\alpha }$$
I.e. $`P(W(t)=W(0)+10,t=1/(2s^2))`$ is at least of order $`10^{1\alpha }`$ which for $`\alpha <2`$ means it is larger than $`10^3`$ !
This coarse estimations highlights the difference between the Gaussian distributions and the distributions generated by random walks with power distributed step sizes (called Levy distributions ): the presence of $`w_i`$’s of arbitrary size implied by a power law distribution insures that the large returns distribution is dominated by the power law of the individual step sizes rather than the combinatorics of the multiple events characterizing the Gaussian system.
One sees now that the systems 1-4 (and consequently LLS) are exactly of the type one needs to explain returns distribution power tails:
\- on one hand according to section 6, the models 1-4 (and consequently LLS) insure a power distribution of $`w_i`$’s.
\- on the other hand, in the models 1-4 the variation of the stock index W(t) is the sum of the variations of the individual $`w_i(t)`$’s.
\- these variations $`w_i(t+1)w_i(t)`$ are stochastic fractions $`s_i(t)=\lambda _i(t)1`$ of $`w_i`$ as above (the fact that $`\lambda _i(t)1`$ has not $`0`$ mean is taken care by working actually with $`u_i=\mathrm{ln}(w_i/\overline{w})`$).
Therefore, according to the argument above, the effective models 1-4, which reflect the stochastic proportionality in LLS between individual wealth, individual investments and individual gains/losses predict that the price fluctuations in the LLS model will obey a Levy distribution (and in particular fit a power in some range of the ”tail”).
There is a proviso for this argument to hold: the number of individual terms $`N`$ has to be larger than the number of time steps $`t`$. Otherwise the finite size of the sample of $`w_i`$’s will show up in the absence of sizes $`w_i`$ larger than a certain value. In fact for $`tN`$ one recovers (slowly) the Gaussian distribution.
In the LLS case, if the portfolio updatings are performed simultaneously by all the investors, the unit time step corresponds already to a time $`t=N`$. In order to verify the (truncated) Levy distribution and the power ”tail” predictions, one has to look at the dynamics at a finer time scale. We therefore performed LLS runs in which at each time step only one trader $`i`$ reconsiders its portfolio investment proportion $`X(i)`$. In such conditions, one expects to obtain a distribution which fits in a significant range a power law (up to large $`w_i`$ values where the finite $`N`$ effects become important).
This is in fact confirmed by the numerical experiments. While for the global updating steps one gets a Gaussian distribution, for the trader-by-trader procedure one obtains a truncated Levy distribution (Fig. 10).
Figure 10: The returns distribution in the LLS model in which only one trader re-evaluates his/her portfolio per unit time. $`c=0.2`$, $`k=3`$, U= $`\mathrm{ln}W`$.
The market contained 10 000 traders with initially equal wealth and portfolio composition (half in stock and half in bonds).
The number of market returns in intervals of 0.001 were measured during 5 000 000 market steps (after an initial 1 000 000 equilibration period).
Note that in the central region of the short time returns (before the cut-off becomes relevant) the Levy distribution is characterized by an $`\alpha `$ equal to the exponent $`\alpha `$ of the traders’ wealth distribution.
As explained in Section 6, in certain conditions (e.g. model 4) one can interpret $`\alpha `$ as
$`\alpha `$ = 1/(1- 1/(long term market return factor)).
Therefore the analysis above relates the stochastic distribution of the short term returns to the value of the long term returns via the exponent of the Pareto power law of individual incomes/wealths.
Moreover the long term returns are related (e.g. model 3) via the value of the Pareto exponent $`\alpha `$ to the ratio ($`\overline{w}/w_{min}`$) between the average wealth/income and the lowest admissible wealth/income: the value $`\alpha 1.4`$ implies (cf. models 3-4) for both these quantities values of the order of
$$1/c=\alpha /(\alpha 1)3.5.$$
Speculatively, one may try to use the above relation in order to explain the stability of the Pareto constant $`\alpha `$ over the past century (and over the various countries and economies).
Indeed one may relate the implied value $`3.5`$ for both $`\overline{w}/w_{min}`$ and the long term market return to some basic biological invariant which is the average number of dependents / offsprings humans have:
\- if $`w_{min}`$ is the minimal amount necessary to keep alive one person in a certain society (cost of life), then the average income $`\overline{w}`$ will have to equal roughly $`w_{min}`$ times the number of dependents the average household head has to support.
\- at the social level, the total effort/wealth that one generation invests in the economy has to result in an economical growth capable to support a population larger by a factor equal to the average number of descendents.
Acknowledgement We thank T. Lux and D. Stauffer for very intensive and detailed correspondence on various simulation experiments using the LLS model.
References
1. M. Levy, H. Levy, S. Solomon, Microscopic Simulation of Financial Markets, Academic Press, New York, 2000.
2. Levy, M., Levy H. and S. Solomon, ”A Microscopic Model of the Stock Market : Cycles, Booms, and Crashes,” Economics Letters, 45, 1994.
Moss de Oliveira S., H. de Oliveira and D.Stauffer, Evolution, Money, War and Computers, B.G. Teubner Stuttgart-Leipzig 1999.
Solomon, S. , The microscopic representation of complex macroscopic phenomena, Annual Reviews of Computational Physics II, 243-294, D. Stauffer (editor), World Scientific 1995.
Solomon S. Behaviorally realistic simulations of stock markets; Traders with a soul, Computer Physics Communications 121-122 (1999) 161
3. Hellthaler T., ”The Influence of Investor Number on a Microscopic Market”, Int. J. Mod. Phys.C 6, 1995
Kohl R., ”The Influence of the Number of Different Stocks on the Levy, Levy Solomon Model,” Int. J. Mod. Phys. C 8, 1997.
4. Levy, M., Levy H. and S. Solomon, ”Simulation of the Stock Market: The Effects of Microscopic Diversity,” Journal de Physique I, 5, 1995.
5. Anderson P. W., J. Arrow and D. Pines , eds. The Economy as an Evolving Complex System (Redwood City, Calif.: Addison-Wesley, 1988);
6. Egenter E., T. Lux and D. Stauffer, ”Finite Size Effects in Monte Carlo Simulations of Two Stock Market Models,” Physica A 268 ,1999.
7. Fama, E., and K. French, ”Permanent and Temporary Components of Stock Prices,” Journal of Political Economy, 96, 1988.
8. Shiller, Robert J., ”Do Stock Returns Move too Much to be Justified by Subsequent Changes in Dividends?” American Economic Review, 1981.
9. Levy M., Persky N., and Solomon S., ”The Complex Dynamics of a Simple Stock Market Model, ” International Journal of High Speed Computing, 8, 1996.
$`http://www.ge.infm.it/econophysics/papers/solomonpapers/stock\mathrm{\_}ex\mathrm{\_}model.ps.gz`$ see also
Farmer J.D., ”Market Force, Ecology and Evolution,” e-print adap-org/9812005. See also J.D. Farmer and S. Joshi, ”Market Evolution Toward Marginal Efficiency” SFI report 1999.
10. Levy, M., and S. Solomon, ”New Evidence for the Power Law Distribution of Wealth”, Physica, A 242, 1997.
Levy, M. ”Are Rich People Smarter?” UCLA Working Paper 1997.
Levy, M. ”Wealth Inequality and the Distribution of Stock Returns” Hebrew University Working Paper 1999.
11. Solomon S. and M. Levy, ”Spontaneous Scaling Emergence in Generic Stochastic Systems,” International Journal of Modern Physics C , 7(5), 1996.
S. Solomon, Stochastic Lotka-Volterra systems of competing auto-catalytic agents lead generically to truncated Pareto power wealth distribution, truncated Levy distribution of market returns, clustered volatility, booms and crashes, In Computational Finance 97, Eds. A-P. N. Refenes, A.N. Burgess, J.E. Moody, (Kluwer Academic Publishers 1998)
O. Biham, O. Malcai, M. Levy and S. Solomon, Phys Rev E 58, 1352, (1998).
12. Levy M., S. Solomon, ”Power Laws are Logarithmic Boltzmann Laws” International Journal of Modern Physics C , 7 ( 4) 1996.
13. Aharon Blank and Sorin Solomon cond-mat/0003240 Power Laws and Cities Population
14. V. Pareto, Cours d’Economique Politique, 2 (1897).
15. Sorin Solomon, cond-mat/9901250; Generalized Lotka-Volterra (GLV) Models and Generic Emergence of Scaling Laws in Stock Markets; To appear in the proceedings of the ”International Conference on Computer Simulations and the Social Sciences, Paris 2000” Hermes Science Publications.
16. O. Malcai, O. Biham and S. Solomon, Phys. Rev. E, 60, 1299, (1999).
17. B. B. Mandelbrot, Comptes Rendus 232, 1638 (1951)
H. A. Simon and C. P. Bonini, Amer. Econ. Rev. 607 (1958)
Mantegna, R. N., ”Levy Walks and Enhanced Diffusion in the Milan Stock Exchange.” Physica A, 179, 1991.
Mantegna, R. N., and Stanley, H. E., ”Scaling Behavior in the Dynamics of an Economic Index.” Nature, 376, 1995.
18. Paul Levy, Theorie de l’Addition des Variables Aleatoires, Gauthier-Villiers, Paris 1937.
19 Lotka, A.J., (editor) Elements of Physical Biology, Williams and Wilkins, Baltimore, 1925;
V. Volterra , Nature, 118, 558.
Figure Captions
Figure 1: The Flow Chart of the LLS market framework
Figure 2 : The Fourier transform of the price in a market with one species with extrapolation range $`k=10`$. The market contained 10000 traders that had initially equal wealth invested half in stock and half in bonds.
Figure 3 : Fraction of the wealth that the species $`k_1=10`$ possesses in Case 1. The traders in the market belonged to 2 species consisting each of 5000 traders. Each trader owed at the beginning 5000 dollars in cache and 5000 shares (worth each 1.4 dolars).
Figure 4: Fraction of the wealth that the species $`k_1=10`$ possesses in Case 2. The initial conditions were similar to Figure 3.
Figure 5 -Fraction of the wealth that the species $`k_1=10`$ possesses in case 3.
Figure 6: The species wealths in a market with 3 species of extrapolation ranges of respectively 10, 141 and 256 days. Initially the 3 species possessed equal wealth distributed equally between stock and bond. Each species consisted of 1000 traders.
Figure 7: Fourier transform of the stock price time evolution in the market described in Fig 6.
Figure 8: The wealth distribution of the investors in an LLS model with a poverty line of $`c=20\%`$ of the average wealth. One a double logarithmic scale one obtains a straight line with slope 2.2 corresponding to an $`\alpha `$ of 1.2.
The market consisted of 10000 traders and the measurement was performed as a ”snapshot” after 1 000 000 ”thermalization” market steps. Initially all the traders had equal wealth ($`\$`$1000) equally distributed between bond and stock.
Figure 9: Wealth distribution for 2 investor species with different return distributions. Model 3 was used with a lower wealth bound of $`c=20\%`$.
$`\lambda `$ is randomly drawn. For the first species $`\lambda `$ is 1.10 or 0.95 with equal probability. For the second ”more talented” species $`\lambda `$ is 1.11 or 0.96 with equal probability. The 2 species were each composed of 10000 traders with initially equal wealth (1000 dollars each).
The measurement of the wealth distribution was performed after a ”thermalization period” of 100 000 wealth updatings.
Figure 10: The returns distribution in the LLS model in which only one trader re-evaluates his/her portfolio per unit time. $`c=0.2`$, $`k=3`$, U= $`\mathrm{ln}W`$.
The market contained 10 000 traders with initially equal wealth and portfolio composition (half in stock and half in bonds).
The number of market returns in intervals of 0.001 were measured during 5 000 000 market steps (after an initial 1 000 000 equilibration period).
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# Table of Contents
## 1 Introduction and Summary
The 26-dimensional critical bosonic string theory admits Dirichlet $`p`$-branes (D-$`p`$-branes) for all $`p26`$. Each of these D-$`p`$-branes admits a tachyonic mode $`T`$ of mass$`{}_{}{}^{2}=1`$, in units where the tension of the fundamental string is equal to $`(2\pi )^1`$ ($`\alpha ^{}=1`$). It has been conjectured that the potential for the tachyon field has a non-trivial translationally invariant (local) minimum at some value $`T_{vac}`$ where the sum of the tachyon potential and the tension of the original brane vanishes . Thus at $`T=T_{vac}`$ the total energy density vanishes, and hence this configuration can be identified as the vacuum of the closed string theory without any D-branes. It has also been conjectured that although this vacuum does not have any perturbative open string excitations, it contains lump-like soliton configurations which approach the vacuum $`T=T_{vac}`$ asymptotically far away from the core of the soliton and represent D-branes of lower dimension . Similar conjectures have also been made involving the tachyon living on the coincident D-brane anti-D-brane pair, or on a non-BPS D-brane of type IIA and IIB superstring theories .
Various pieces of evidence for these conjectures have been found in both the first , and second quantized string theory, and also using AdS/CFT correspondence . The first quantized description has been successful in verifying the conjectures relating the tachyonic solitons to lower dimensional D-branes, but it can only supply indirect evidence for the equivalence between the (local) minimum of the tachyon potential and the vacuum without a D-brane. On the other hand, the second quantized description $``$ open string field theory $``$ can provide direct evidence for this conjecture by explicitly computing the (negative) value of the tachyon potential at the minimum and comparing it with the tension of the original D-brane system. Although open string field theory contains infinite number of fields, and the problem of finding a translationally invariant stationary point of the potential involves solving the equations of motion of the infinite number of zero momentum modes of these fields, the calculations are made feasible by using the level expansion scheme proposed by Kostelecky and Samuel . The procedure is as follows. Using the correspondence between the modes of the string field and states in the conformal field theory describing the first quantized string, we define the level of a mode of the string field as the difference between the $`\widehat{N}`$ eigenvalue of the first quantized string state representing this mode, and the $`\widehat{N}`$ eigenvalue of the state representing the zero momentum tachyon, where $`\widehat{N}`$ is the total ‘number operator’ of the matter and ghost system. The level truncation scheme to order $`(M,N)`$ then corresponds to an approximation in which we keep in the string field theory action all modes of level $`M`$, and all interaction terms for which the sum of the levels of all the modes appearing in the term is $`N`$. This gives a potential (which, for a static field configuration, is just the negative of the action up to a normalization constant) with finite number of fields and a finite number of terms. Thus we can find its extremum and calculate its value at the extremum. The larger the values of $`(M,N)`$, the larger is the number of modes and the number of terms in the potential, and the better is the accuracy.
The calculation of ref. for the tachyon potential was revisited and extended in in terms of background independent fields. It was shown there that the total negative potential energy at the stationary point cancels the energy of the D-brane represented by the string field theory to an accuracy of $`<`$1.5% at the level (4,8) approximation. This calculation was extended in ref. to level (10,20). At this level the contribution from the tachyon potential was found to cancel the tension of the D-brane to an accuracy of about .1%. Similar calculations have also been performed in open superstring field theory . At the level (2,4) approximation the tachyon potential has been shown to cancel about 90% of the tension of the original brane configuration.<sup>1</sup><sup>1</sup>1At present there seems to be some disagreement between refs. and about the level (2,4) results.
The success of string field theory in verifying the conjecture relating the translationally invariant stationary point of the tachyon potential and the vacuum without any D-brane encourages one to ask whether string field theory can also be used in studying the conjectured relation between the tachyonic lump solutions and lower dimensional D-branes. This study was initiated by Harvey and Kraus . In this paper they started with the level (0,0) contribution to the tachyon potential in open bosonic string field theory on a D-$`p`$-brane, and identified a ‘bounce solution’ in this field theory as the D-$`(p1)`$ brane. At this level the tension associated with this solution turns out to be about 78% of the known value of the D-$`(p1)`$-brane tension. This result receives correction not only from the higher level fields, but also from the momentum dependence of the interaction terms which were neglected in the initial analysis. While there is no systematic expansion scheme for taking into account these momentum dependent corrections, a naive expansion of the interaction term in powers of momentum, keeping only the zeroth and first order terms, reduced the tension of the soliton to about 70% of the conjectured answer. On the other hand, taking into account the correction to the potential to level (2,4) increased the answer back to about 82% of the conjectured answer. A systematic method for taking into account the momentum dependent terms in the interaction was suggested in ref. , but this procedure did not give rise to an appreciable change in the tension of the lump. A similar analysis has also been carried out for solitons in the open superstring field theory . Although the answer turns out to be close to the expected answer, it is likely to be an accidental result, as there is no reason to assume that the corrections due to the momentum dependent terms are small in this case.
The purpose of this paper will be to develop a systematic approximation scheme for studying these solitons in string field theory and calculating their tension. We shall focus on the codimension one lump on a D$`p`$-brane of the bosonic string theory $``$ which is conjectured to be equivalent to a D-$`(p1)`$-brane $``$ but it will become clear that the scheme is general enough to be applicable to the study of higher codimension solitons, as well as to solitons in superstring field theory. In the case of a codimension one soliton, we are dealing with a field configuration on the D$`p`$ brane which depends on only one of the spatial coordinates (say $`x`$) on the brane, and is independent of time, as well as the other $`(p1)`$ spatial coordinates. We study this problem by compactifying the coordinate $`x`$ on a circle of radius $`R`$ instead of letting it span the whole real line. In this case, since all field configurations must be periodic in $`x`$, we can decompose all fields into modes carrying discrete momenta along $`x`$ in units of $`(1/R)`$, and the solitonic field configuration that we are looking for must be obtained as an appropriate superposition of these modes. We can now define the level of any such mode as the difference between the $`L_0`$ eigenvalue of the first quantized string state representing this mode, and that of the zero momentum tachyon state, where $`L_0`$ denotes the zeroth component of the Virasoro generator of the combined matter ghost system.<sup>2</sup><sup>2</sup>2Since for the zero momentum states the eigenvalue of the number operator is the same as the $`L_0`$ eigenvalue of the state, the two prescriptions agree for these states. This allows us to define a level $`(M,N)`$ approximation to the potential exactly as before. Working with the potential up to a given level, we can now look for $`x`$ dependent solutions of the string field equations by extremizing the potential with respect to the modes appearing in the potential to this level.
This is precisely the procedure we follow in this paper for studying the tachyonic lump solution on a D-$`p`$-brane.<sup>3</sup><sup>3</sup>3The discretization of the momentum is reminiscent of the procedure followed in ref. , although the precise relationship between these two approaches is not clear. We study this problem for various radii at various levels of approximation, and compare the tension of the lump with the tension of a D-$`(p1)`$-brane. The results for the tension of the lump turn out to be remarkably close to the known tension of the D-$`(p1)`$-brane. Whereas for $`R=\sqrt{3}`$ and $`\sqrt{15/2}`$ we are able to get a lump tension within 1% of the tension of the D-$`(p1)`$-brane, for larger radii ($`R=\sqrt{12}`$ and $`\sqrt{35/2}`$) we get answers within 3% of the expected answer. We also compare the profile of the tachyon field corresponding to the lump for different values of $`R`$, $``$ obtained by superposition of $`\mathrm{cos}(nx/R)`$ for integer $`n`$ $``$ and find remarkable agreement between the profiles for different values of $`R`$.
At this point we should note that the problem of formation of the tachyonic lump on a circle was addressed using the first quantized approach in ref. . There a renormalization group analysis was used to show that the mass of the tachyonic lump on a D-$`p`$-brane is indeed equal to that of a D-$`(p1)`$-brane.<sup>4</sup><sup>4</sup>4This followed earlier work of ref. on the renormalization group flow of the two dimensional field theory under a perturbation corresponding to switching on a tachyon background proportional to $`\mathrm{cos}(x/R)`$. In view of this result one might ask whether the string field theory analysis carried out in this paper gives any new insight into this problem. To this end, we note, first of all, that the relationship between the renormalization group analysis in the first quantized approach, and the string field theory analysis based on the level truncation scheme, is as yet quite unclear, and hence it is certainly illuminating to independently verify the equivalence of the D-$`(p1)`$-brane, and the tachyonic lump on the D-$`p`$-brane in string field theory. Furthermore, string field theory provides us with a definite picture of the tachyon profile as superposition of $`\mathrm{cos}(nx/R)`$ for different $`n`$ with definite coefficients. In contrast the analysis based on the renormalization group flow only tells us that a perturbation by the leading relevant operator $`\mathrm{cos}(x/R)`$ takes the original D-$`p`$-brane to a D-$`(p1)`$-brane, and does not tell us how the higher harmonics mix with $`\mathrm{cos}(x/R)`$ to produce the soliton. Indeed most of the higher harmonics correspond to irrelevant perturbation, and hence their coefficients vanish in the infra-red.<sup>5</sup><sup>5</sup>5Presumably if we could determine the exact location of the infrared fixed point in the space of coupling constants, then the shape of the lump will be determined in this approach. Furthermore, the rigorous results of ref. have not yet been generalized to superstring theory. Thus we believe that despite the exact results based on the renormalization group analysis of the first quantized theory, the present analysis throws new light on the tachyonic soliton solutions.
The rest of the paper is organized as follows. In section 2 we outline the general procedure of level expansion scheme of the string field theory, discuss the possibility of restricting the string field to a background independent subspace for studying the lump solution, and give details of computation of a few terms in the potential. In section 3 we give in detail the results for the potential, the lump solution and its energy for a specific radius $`R=\sqrt{3}`$. We also compare the profile of the lump at different levels of approximation. In section 4 we give the results for several other radii, both larger ($`\sqrt{15/2}`$, $`\sqrt{12}`$ and $`\sqrt{35/2}`$) and smaller ($`\sqrt{11/10}`$) than $`\sqrt{3}`$, and compare the profile of the lump for each radii with the profile at $`R=\sqrt{3}`$. We conclude in section 5 by discussing possible generalization of this analysis and some speculations.
## 2 Level Expansion and the String Field
In this section we will set up a variant of the level expansion method to deal with the problem of finding the profile and mass of the tachyon lump in string field theory. As reviewed in the introduction, such method is desirable as previous computations of lump masses in string field theory have not been very accurate. After explaining this method we will discuss the background independent expansion of the string field suitable for the problem. Then we discuss two methods for estimating the lump mass. We conclude by showing a few samples of typical calculations needed to evaluate the string field action for the lump.
### 2.1 Modified level expansion
When calculating the tachyon potential in search for a spacetime independent vacuum state, all spacetime fields are set to constants, and the evaluation of the string field action does not require the inclusion of terms with spacetime derivatives. The string field is at zero momentum and is thus built by a superposition of zero momentum states times constants representing the zero momentum modes of the spacetime fields. The states are built by acting on a zero-momentum vacuum with oscillators of the relevant conformal field theory (CFT). In this case the level expansion was defined as follows . Let $`\widehat{N}`$ be the number operator, representing the contribution to $`L_0`$ from the system of matter and ghost oscillators. Let $`N_0`$ (=$`1`$) denote the eigenvalue of $`\widehat{N}`$ for the zero momentum tachyon: $`\widehat{N}|T_0=N_0|T_0`$. For a given state $`|\mathrm{\Phi }_i`$, with number eigenvalue $`N_i`$ ($`\widehat{N}|\mathrm{\Phi }_i=N_i|\mathrm{\Phi }_i`$) we define the level $`l(\mathrm{\Phi }_i)`$ of the state $`|\mathrm{\Phi }_i`$ as
$$l(\mathrm{\Phi }_i)N_iN_0.$$
(2.1)
As defined, level is a dimensionless number. For the case of bosonic string theory the levels are all integers while for NS superstrings they can also be half integral. We now define the level $`(M,N)`$ approximation to the action as follows:
* We keep only those fields with level $`M`$.
* We keep only those terms in the action for which the sum of the levels of all the fields in the term is $`N`$.
In order that the quadratic term of all fields with level $`M`$ are kept in the action, we must have $`N2M`$. While variants are possible, it seems most effective when calculating any physical object to use its level $`(M,2M)`$ approximation, as experience shows that increasing the number of terms in the potential keeping the number of fields fixed does not improve the results very much. While there is yet no theoretical explanation for the convergence of the level expansion, the numerical evidence collected thus far is impressive.
Consider now the problem at hand. While all of our discussion applies to soliton solutions on non-BPS D-branes, and D-brane anti- D-brane pairs of superstring theory, we will consider here explicitly only the case of the unstable D-branes of bosonic string theory. Consider therefore, an unstable bosonic D-brane extending over a number of spatial dimensions. We now wish to select one of these dimensions, call it $`x`$ and construct a tachyon lump such that the solution depends only on the $`x`$-coordinate. (Again our discussion applies to lumps depending on more than one coordinates, but we shall not analyze these cases here.) As the lump is not invariant under translation along $`x`$, we now need to include $`x`$-momentum modes in the string field expansion and $`x`$-derivatives, or $`x`$-momentum dependent terms in the string field action. In order to do this systematically we compactify $`x`$ over a circle of radius $`R`$, namely $`xx+2\pi R`$. This quantizes the $`x`$-momentum as $`p_x=n/R`$ for integer $`n`$. For each of the zero momentum states $`|\mathrm{\Phi }_i`$ we had before, we now have discrete states of the type $`|\mathrm{\Phi }_{i,n}`$ that only differ by the fact that they are built on vacua having $`x`$-momentum $`n/R`$. For such states there is a natural generalization of the level. This is the difference between the $`L_0`$ eigenvalue of the state and that of the zero momentum tachyon, where $`\{L_n\}`$ denote the Virasoro generators of the combined matter and ghost system. This is because (with $`\alpha ^{}=1`$) we have that $`L_0=p_x^2+\widehat{N}`$. For zero momentum this is just the previous definition. Still denoting by $`N_i`$ the number eigenvalue of $`|\mathrm{\Phi }_{i,n}`$ we have
$$l(\mathrm{\Phi }_{i,n})=L_0(\mathrm{\Phi }_{i,n})L_0(T_0)=\frac{n^2}{R^2}+N_iN_0.$$
(2.2)
The level is still a dimensionless number as $`R`$ here is measured in units of $`\sqrt{\alpha ^{}}`$ (which has been set to one). We can now define the level $`(M,N)`$ approximation for the action exactly as before. Since the $`L_0`$ eigenvalue of a state plays a crucial role in the conformal map that inserts the state into the disk representing the interaction terms in the action, this is a natural generalization of the level truncation scheme of ref. . This paper will present evidence that this modified version of the level truncation scheme also works very well.
In calculating in this setup in the level $`(M,2M)`$ approximation for any given radius we will have to include states $`|\mathrm{\Phi }_i|\mathrm{\Phi }_{i,0}`$ and “harmonics” $`|\mathrm{\Phi }_{i,n}`$, and clearly the condition $`l(\mathrm{\Phi }_{i,n})M`$ will give an upper bound on $`n`$ for each $`i`$. This also requires $`l(\mathrm{\Phi }_{i,0})M`$, and thus we have a finite number of modes to be included at a given level of approximation. Each term in the action including modes whose sum of levels does not exceed $`2M`$ is computed exactly. It is manifest that in a cubic string field theory the level $`(M,2M)`$ approximation will only require a finite number of computations<sup>6</sup><sup>6</sup>6This will also be the case for the NS superstring field theory discussed in ref. .
### 2.2 Background Independent String Field
The general setup required to study a lump is similar to that developed in to study the mass of the D-brane. To begin with, we assume that the background space-time is the product of a (2+1) dimensional flat space-time, labelled by a pair of space-like coordinates $`(x,y)`$ and a time like coordinate $`x^0`$, and an arbitrary Euclidean manifold $``$ described by a conformal field theory of central charge 23. We take the spatial direction $`y`$ to be non-compact, but $`x`$ to be compact with radius $`R`$. We let $`X`$, $`Y`$ and $`X^0`$ denote the three scalar fields on the string world-sheet associated with the coordinates $`x`$, $`y`$ and $`x^0`$.
We now consider a D-brane with the following properties. For an open string ending on the D-brane we put Neumann boundary condition on the fields $`X`$ and $`X^0`$ and Dirichlet boundary condition on the field $`Y`$.<sup>7</sup><sup>7</sup>7As in ref. , the extra non-compact direction $`y`$ with Dirichlet boundary condition provides a direction along which the brane can move, and we can calculate the tension of the brane by studying its motion in this direction. We leave the boundary condition on the fields associated with the coordinates on $``$ arbitrary, with the only restriction that all the fields on which we put Neumann boundary condition are associated with compact coordinates. This means that all directions tangential to the D-brane are compact, and hence the D-brane has finite mass. From the point of view of the full space-time, this D-brane describes a D-$`p`$ brane for some $`p1`$, with $`(p1)`$ directions wrapped on an internal $`(p1)`$ cycle of $``$, and one direction wrapped on the circle of radius $`R`$ labelled by $`x`$. On the other hand from the point of view of an observer who only sees the (2+1) dimensional space-time labelled by $`(x,y,x^0)`$, this system corresponds to a D1-brane wrapped on a circle of radius $`R`$. From now on we shall refer to this system as the D1-brane or the D-string; with its tension defined as the total energy per unit length along $`x`$. Of course, an ordinary D-string will be a special case of this system, obtained by putting Dirichlet boundary condition on all the fields associated with the coordinates on $``$.
The dynamics of an open string with ends on this D-brane is described by a boundary conformal field theory of central charge 26, which is a direct sum of the boundary conformal field theories associated with the fields $`X`$, $`Y`$, $`X^0`$ and the manifold $``$. We shall denote by CFT($`X`$), CFT($`Y`$) and CFT($`X^0`$) the boundary conformal field theories (each with central charge 1) associated with the fields $`X`$, $`Y`$ and $`X^0`$ respectively, and by CFT$`()`$ the boundary conformal field theory with central charge 23 associated with the manifold $``$. We also define
$$\text{CFT}^{}=\text{CFT}(Y)\text{CFT}(X^0)\text{CFT}(),$$
(2.3)
so that CFT has central charge 25. We denote by $`L_n^X`$ and $`L_n^{}`$ the Virasoro generators of CFT($`X`$) and CFT respectively. If we denote by $`L_n^{ghost}`$ the Virasoro generators of the ghost system, then the total Virasoro generators of the system will be given by $`L_n=L_n^{ghost}+L_n^X+L_n^{}`$.
The compact direction $`x`$ corresponds to the direction in which we shall eventually form the lump. If we follow the normalization convention of ref., then the tension $`𝒯_1`$ of the D-string described above is related to the coupling constant $`g_o`$ of the open string field theory describing the wrapped D-string by the relation:
$$2\pi R𝒯_1=\frac{1}{2\pi ^2g_o^2}.$$
(2.4)
In this normalization convention, a time independent string field configuration represented by a state $`|\mathrm{\Phi }=\mathrm{\Phi }(0)|0`$ in the Hilbert space of first quantized string theory, will have a potential
$$\text{Potential}=S(\mathrm{\Phi })=\frac{1}{g_o^2}𝒱(\mathrm{\Phi })=2\pi R𝒯_12\pi ^2𝒱(\mathrm{\Phi }),$$
(2.5)
where
$$𝒱(\mathrm{\Phi })=\frac{1}{2}\mathrm{\Phi },Q\mathrm{\Phi }+\frac{1}{3}\mathrm{\Phi },\mathrm{\Phi }\mathrm{\Phi }.$$
(2.6)
Here $`Q`$ denotes the BRST charge, $`,`$ denotes BPZ inner product between two states, and $``$ denotes the $``$-product of Witten’s open bosonic string field theory .
A basis of states in CFT(X) is obtained by acting on $`e^{inX/R}(0)|0`$ with the oscillators $`\alpha _m^X`$ of $`X`$. It follows by a simple counting argument that an alternate basis can be formed out of the Verma module, containing states obtained by acting on $`e^{inX/R}(0)|0`$ with the operators $`L_m^X`$, as long as these states are all linearly independent. This is the case if there are no null states in the spectrum. The condition for the appearance of a null state is given by ,
$$\frac{n^2}{R^2}=\frac{(pq)^2}{4}\frac{n}{R}=\frac{(pq)}{2},$$
(2.7)
where $`p`$ and $`q`$ are integers. Since $`n`$ is an integer, we can avoid null states for $`n0`$ with an appropriate choice of $`R`$. Even if we work with a value of $`R`$ for which there are null states, the choice of basis described above is good below the level where the first null state appears. From now on we shall restrict our analysis to situations where this choice of basis based on Verma module is good. In fact our explicit work in later sections will be based on $`R`$ values that are not rational, and thus there will be no null states for $`n0`$.
For $`n=0`$, however, there are null states and hence the basis of states obtained by applying $`L_m^X`$ on $`|0`$ is not complete. For example, $`L_1^X|0`$ is null, and this requires us to explicitly include the primary state $`\alpha _1^X|0`$ in the basis. There are further null states in the Verma module over $`\alpha _1^X|0`$, and hence there are new primary states at higher level which must be explicitly included in the basis. Let us denote by $`\{|\phi _e^i=\phi _e^i(0)|0_X\}`$ and $`\{|\phi _o^i=\phi _o^i(0)|0_X\}`$ the set of zero momentum primary states which are respectively even and odd under the reflection $`XX`$. The complete basis of zero momentum states in CFT(X) is obtained by acting on $`\{|\phi _e^i\}`$ and $`\{|\phi _o^i\}`$ with $`L_n^X`$’s, and removing the null states.
A generic string field configuration is represented by an arbitrary state in the Hilbert space $``$ of ghost number one in the combined matter, ghost conformal field theory. We now claim that in order to discuss a lump along the $`x`$ coordinate, we can restrict the string field $`|\mathrm{\Phi }`$ to a subspace $`\widehat{}`$ of $``$, built by acting with the oscillators
$$\{L_1^X,L_2^X,\mathrm{};L_2^{},L_3^{},\mathrm{};c_1,c_1,c_2,\mathrm{};b_2,b_3,\mathrm{}\}$$
(2.8)
on the following primary states:
* The zero momentum even primaries $`\phi _e^i(0)|0`$ (and removing the null states), and,
* The Fock vacuum states of the form
$$\mathrm{cos}\left(\frac{n}{R}X(0)\right)|0=\frac{1}{2}\left(e^{inX(0)/R}+e^{inX(0)/R}\right)|0=\frac{1}{2}\left(|\frac{n}{R}+|\frac{n}{R}\right)n0,$$
(2.9)
where $`|0`$ is the SL(2,R) vacuum of the combined matter, ghost conformal field theory. A few points should be made. The Virasoro operator $`L_1^{}`$ is not required for it kills the above primary states (this is not the case for $`L_1^X`$). $`b_1`$ and $`b_0`$ also annihilate the vacuum $`|0`$, and hence have been omitted from the list. We have not included the oscillator $`c_0`$ because we work in the Siegel gauge, where all states must be annihilated by $`b_0`$. Finally we can restrict ourselves to states of even twist . This simply requires that the eigenvalue of the number operator $`\widehat{N}`$ must be odd (same as that for the tachyon).
In order to show that the above is a consistent truncation of the string field, one must show that there is no term in the action that couples a single state in $`(\widehat{})`$ to a state in $`\widehat{}`$ via the quadratic term, or to a pair of states in $`\widehat{}`$ via the interaction term. This is readily done by listing the states in $`(\widehat{})`$. We carry along all ghost oscillators and classify the states by their behavior under the matter operators. In this way we get the following disjoint sets:
* States with nonzero momentum $`k_0`$ along $`X^0`$.
* States obtained by acting with the oscillators in (2.8) on Fock vacua of the type $`\mathrm{sin}(\frac{nX(0)}{R})|0`$, or on a state of the form $`\phi _o^i(0)|0`$.
* States obtained by acting with (2.8) on states that (i) have $`k_0=0`$, (ii) are non-trivial primaries of CFT (of dimension greater than zero, by unitarity), and (iii) are CFT(X) primaries.
It is manifest by momentum conservation that a state in the first set cannot couple to states in $`\widehat{}`$. The symmetry $`XX`$ of CFT(X) insures that a state in the second set also cannot couple to states in $`\widehat{}`$. The same is true for the last set as Virasoro Ward identities can be used to show that a correlator involving two states in $`\widehat{}`$ and a state in the last set is proportional to the one point function of the CFT primary in question. Since this primary must have dimension greater than zero, its one point function vanishes. This completes our justification for the use of $`\widehat{}`$.
Since the choice of basis described above requires the use of the basis $`\{|\phi _e^i\}`$, it will be useful to determine at which level the first zero momentum primary (other than the vacuum state) appears. For this we can compare the full partition function of CFT(X) for states even under $`XX`$
$$Z_{even}(q)Tr_{even}(q^{L_0^X\frac{1}{24}})=\frac{1}{2}q^{\frac{1}{24}}\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{1q^n}+\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{1+q^n}\right)$$
(2.10)
with the Virasoro character for $`(c=1,h=0)`$ ,
$$\chi _{c=1,h=0}(q)=q^{\frac{1}{24}}\underset{n2}{}\frac{1}{1q^n}.$$
(2.11)
It can be easily checked that
$$Z_{even}(q)\chi _{c=1,h=0}(q)=q^{\frac{1}{24}}(q^4+O(q^5)).$$
(2.12)
Thus the first non-trivial primary $`|\phi _e^1`$ even under $`XX`$ appears at level four. Indeed, the available descendents at this level, $`L_4^X|0,L_2^XL_2^X|0`$, do not suffice to represent the nonvanishing $`X`$-even states $`\alpha _3^X\alpha _1^X|0,\alpha _2^X\alpha _2^X|0,(\alpha _1^X)^4|0`$.
### 2.3 Mass of the lump
To begin with, the wrapped D-$`p`$ brane, which we have been calling a D-1 brane wrapped on a circle of radius $`R`$, has mass $`2\pi R𝒯_1`$, where $`𝒯_1`$, as defined earlier, is the tension of this D1-brane. We want to compute the mass of the system in a situation where the tachyon field on this D$`1`$-brane develops a lump along a circle of radius $`R`$ (this direction is represented by the world sheet field $`X`$). If we denote by $`\stackrel{}{T}`$ the multicomponent string field configuration on the D1 brane, restricted to $`\widehat{}`$, then, using eq.(2.5), the rest mass energy plus potential energy of the D1 brane stretched on the circle can be written as
$$E(D1)=𝒯_1(2\pi R)\left(1+2\pi ^2𝒱(\stackrel{}{T})\right)$$
(2.13)
where $`𝒱`$ has been defined in eq.(2.6). Before condensation, $`\stackrel{}{T}=0`$ and $`𝒱(\stackrel{}{T})=0`$, and thus the energy formula correctly reproduces the mass of the D1-brane. Recall that for the nontrivial translationally invariant vacuum $`\stackrel{}{T}_{vac}`$, one expects $`𝒱(\stackrel{}{T}_{vac})=1/(2\pi ^2)`$ and the energy formula correctly gives zero (as the D1 brane has disappeared). Using $`𝒱(\stackrel{}{T}_{vac})=1/(2\pi ^2)`$ we can write the energy formula as
$$E(\text{D1})=𝒯_1\mathrm{\hspace{0.17em}2}\pi R2\pi ^2\left(𝒱(\stackrel{}{T})𝒱(\stackrel{}{T}_{vac})\right)$$
(2.14)
The mass of the tachyonic lump solution, represented by the configuration $`\stackrel{}{T}_{lump}`$, is obtained by replacing $`\stackrel{}{T}`$ by $`\stackrel{}{T}_{lump}`$ on the right hand sides of eqs.(2.13) or (2.14). This tachyonic lump on the D-string (wrapped D-$`p`$-brane) is conjectured to be equivalent to a D0-brane (a wrapped D-$`(p1)`$-brane) of mass $`𝒯_0`$. With $`\alpha ^{}=1`$, the ratio of the tension of a D-$`p`$ brane and a D-$`(p1)`$-brane is $`1/(2\pi )`$; using this we get,
$$𝒯_0=2\pi 𝒯_1.$$
(2.15)
This gives
$$r\frac{E_{lump}}{𝒯_0}=2\pi ^2R\left(𝒱(\stackrel{}{T}_{lump})𝒱(\stackrel{}{T}_{vac})\right).$$
(2.16)
The predicted answer for this ratio is 1.
This prediction can be tested for various values of $`R`$, and independently of the chosen value we must obtain unity, since the mass of a D0-brane on a circle of radius $`R`$ does not depend on $`R`$. At fixed $`R`$ and at any level of approximation in the level expansion it is possible to use (2.16) in two ways. We can use that $`2\pi ^2𝒱(\stackrel{}{T}_{vac})`$ at the exact vacuum is indeed $`1`$ and thus we check how accurately
$$r^{(1)}R\left(2\pi ^2𝒱_{(M,N)}(\stackrel{}{T}_{lump})+1\right)$$
(2.17)
approaches unity. Here $`𝒱_{(M,N)}`$ is the potential calculated at the specified level of approximation. Alternatively we can use the translationally invariant vacuum that is obtained with the same level of approximation used to compute the lump.<sup>8</sup><sup>8</sup>8The value of $`𝒱_{(M,N)}(\stackrel{}{T}_{vac})`$ can be read from refs. . This gives
$$r^{(2)}R\left(2\pi ^2𝒱_{(M,N)}(\stackrel{}{T}_{lump})2\pi ^2𝒱_{(M,N)}(\stackrel{}{T}_{vac})\right)$$
(2.18)
We will find that $`r^{(1)}`$ approaches unity monotonically from above as we increase the level of approximation. On the other hand $`r^{(2)}`$ provides a more accurate answer.
### 2.4 Setup and Sample Computations
Let us now describe explicitly the string field we will be using to analyze the bosonic string lump. The zero momentum tachyon $`|T_0=c_1|0`$ now becomes the lowest in a family of states
$$|T_n=c_1\mathrm{cos}\left(\frac{n}{R}X(0)\right)|0,l(T_n)=\frac{n^2}{R^2}$$
(2.19)
where $`l(T_n)`$ denotes the level of $`T_n`$. For any given computation only a finite number of tachyon modes are required. In the zero-momentum computation the next modes that contribute are $`|U_0=c_1|0`$ and $`|V_0=L_2^{matt}|0.`$ In view of our remarks around (2.9) these states actually give rise to three towers
$`|U_n`$ $`=`$ $`c_1\mathrm{cos}\left({\displaystyle \frac{n}{R}}X(0)\right)|0,l(U_n)=2+{\displaystyle \frac{n^2}{R^2}},`$
$`|V_n`$ $`=`$ $`c_1L_2^X\mathrm{cos}\left({\displaystyle \frac{n}{R}}X(0)\right)|0,l(V_n)=2+{\displaystyle \frac{n^2}{R^2}},`$
$`|W_n`$ $`=`$ $`c_1L_2^{}\mathrm{cos}\left({\displaystyle \frac{n}{R}}X(0)\right)|0,l(W_n)=2+{\displaystyle \frac{n^2}{R^2}}.`$ (2.20)
In addition to these three towers there is one more, where the $`n=0`$ state happens to vanish:
$$|Z_n=c_1L_1^XL_1^X\mathrm{cos}\left(\frac{n}{R}X(0)\right)|0,n1,l(Z_n)=2+\frac{n^2}{R^2}.$$
(2.21)
No new fields or towers arise until level four, and for the purposes of the present paper we shall not carry computations that far. Therefore we will use the string field
$`|\stackrel{}{T}`$ $`=`$ $`t_0|T_0+t_1|T_1+t_2|T_2+\mathrm{}`$
$`+u_0|U_0+u_1|U_1+\mathrm{}`$
$`+v_0|V_0+v_1|V_1+\mathrm{}`$
$`+w_0|W_0+w_1|W_1+\mathrm{}`$
$`+z_1|Z_1+\mathrm{}`$
Which fields and which interactions must be kept for any fixed level computation depends on the chosen radius, and this will be discussed in the following sections. We conclude here with some basic comments about the evaluation of the potential (or the action) for a string field of the above type.
This is simply the evaluation of $`𝒱(\stackrel{}{T})`$ as given in (2.6)
$$𝒱(\stackrel{}{T})=\frac{1}{2}\stackrel{}{T},Q\stackrel{}{T}+\frac{1}{3}\stackrel{}{T},\stackrel{}{T}\stackrel{}{T}.$$
(2.23)
We work in units where $`\alpha ^{}=1`$. The stress tensor for the compact coordinate $`X`$ is $`T_X=\frac{1}{4}XX`$ with $`X(z)X(w)2\mathrm{ln}(zw)`$, $`T(z)e^{ipX(w)}\frac{p^2}{(zw)^2}e^{ipX(w)}`$ and $`e^{ip_1X(z)}e^{ip_2X(w)}=(zw)^{2p_1p_2}e^{ip_1X(z)+ip_2X(w)}`$, where $`z`$ and $`w`$ are coordinates on the real line with $`z>w`$. With these conventions $`L_0|p=L_0e^{ipX(0)}|0=p^2|p`$. In addition, the inner product is normalized as
$$\frac{n}{R}\left|c_1c_0c_1\right|\frac{m}{R}=\delta _{n,m}$$
(2.24)
Consider, for example contributions from the tachyon tower to the action. By momentum conservation all kinetic terms must be diagonal. Using (2.9) we see that the contribution from $`t_n`$ ($`n1`$) to $`V`$ is
$$\frac{1}{2}\frac{t_n}{2}\frac{t_n}{2}\left(\frac{n}{R}|+\frac{n}{R}|\right)c_1c_0L_0c_1\left(|\frac{n}{R}+|\frac{n}{R}\right)$$
(2.25)
By momentum conservation there are two cross terms that do not vanish and give identical contributions. We thus get
$$\frac{1}{4}t_n^2\frac{n}{R}\left|c_1c_0\left(1+\frac{n^2}{R^2}\right)c_1\right|\frac{n}{R}=\frac{1}{4}\left(1\frac{n^2}{R^2}\right)t_n^2$$
(2.26)
For $`t_0`$ the normalization factor differs by a factor of two. All this together gives us that the quadratic terms are
$`𝒱(t_0,t_1,t_2,\mathrm{})^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}t_0^2{\displaystyle \frac{1}{4}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(1{\displaystyle \frac{n^2}{R^2}}\right)t_n^2`$ (2.27)
$`=`$ $`{\displaystyle \frac{1}{2}}t_0^2{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{1}{R^2}}\right)t_1^2{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{4}{R^2}}\right)t_2^2+\mathrm{}`$
We will use the first tachyon harmonic $`t_1`$ to drive the unstable vacuum into the lump solution. Note that $`t_1`$ is tachyonic whenever $`R>1`$. We will choose different values of $`R>1`$ to examine how the lump forms. As $`R`$ increases, more and more tachyon harmonics become tachyonic.
It is not difficult to compute the interactions of the various tachyon harmonics. One can use the oscillator expressions for the states and contract them against the 3-string vertex bra $`V_{123}|`$ . Alternatively one can use the conformal field theory definition
$$\stackrel{}{T},\stackrel{}{T}\stackrel{}{T}h_1T(0)h_2T(0)h_3T(0).$$
(2.28)
where $`T(0)`$ denotes the vertex operator associated to the state $`|\stackrel{}{T}`$. Here $`h_1`$, $`h_2`$ and $`h_3`$ are a set of familiar conformal transformations reviewed in . For illustration purposes consider three tachyon harmonics $`t_n,t_m`$ and $`t_{n+m}`$, with $`nm0`$. Such fields contribute to $`𝒱`$ the following interaction
$$\frac{1}{3}6\frac{t_n}{2}\frac{t_m}{2}\frac{t_{n+m}}{2}2h_1(ce^{\frac{inX}{R}})(0)h_2(ce^{\frac{imX}{R}})(0)h_3(ce^{\frac{i(n+m)X}{R}})(0)$$
(2.29)
The factor $`(1/3)`$ is in the definition of $`𝒱`$. The factor of $`6`$ appears because this is the number of ways three different fields can be assigned to the three punctures in the disk. Then come the fields, and then a factor of two, as there are two momentum conserving combinations giving equal contributions. Evaluation of the above gives
$$\frac{1}{2}t_nt_mt_{n+m}K^{3\frac{1}{R^2}(n^2+m^2+(n+m)^2)},K\frac{3\sqrt{3}}{4}$$
(2.30)
Slightly different combinatorics are required for terms of the form $`t_0t_n^2`$ and $`t_n^2t_{2n}`$. Combining all such terms together we obtain
$`𝒱(t_0,t_1,\mathrm{})^{(3)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}K^3t_0^3+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}t_0t_n^2K^{3\frac{2n^2}{R^2}}+{\displaystyle \frac{1}{4}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}t_n^2t_{2n}K^{3\frac{6n^2}{R^2}}`$ (2.31)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{n1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m>n}{\overset{\mathrm{}}{}}}t_nt_mt_{n+m}K^{3\frac{2}{R^2}(n^2+m^2+nm)}.`$
Equations (2.27) and (2.31) give the complete potential for the tachyon tower.
## 3 Calculating the action in the Level Expansion for $`R=\sqrt{3}`$
In this section we will consider different truncation levels to calculate the lump tension. For this we will write explicitly the action at different levels. Though we will work with a fixed radius $`R=\sqrt{3}`$, all our equations will contain $`R`$ as a variable for further use. Once we know the action we can solve the equations of motion numerically for the one-lump solution by giving a nonzero initial value to $`t_1`$. At the end of the section, we will be able to study the convergence of our level truncation scheme by using both (2.17) and (2.18).
We will do these calculations at levels $`(1/3,2/3)`$, $`(4/3,8/3)`$, $`(2,4)`$, $`(7/3,14,3)`$ and $`(3,6)`$. This will require the fields listed in Table 1 with their respective levels (using (2.19), (2.4) and (2.21)). In order to study the truncation method at various levels, we define $`V(m,n)`$ to be the part of the whole potential satisfying the three following conditions:
1. All terms in $`V(m,n)`$ have level $`n`$.
2. All terms in $`V(m,n)`$ contain only fields of level smaller than or equal to $`m`$.
3. All terms in $`V(m,n)`$ contain at least one field of level $`m`$.
This definition ensures that various $`V(m,n)`$’s are disjoint (i.e. $`V(m,n)`$ and $`V(m^{},n^{})`$ don’t contain common terms for $`(m,n)(m^{},n^{})`$). It now follows that the total potential at level $`(M,N)`$ is given by
$$𝒱_{(M,N)}=\underset{mM}{}\underset{nN}{}V(m,n)$$
(3.1)
We shall now compute $`V(m,n)`$ for $`m3`$ and $`n6`$. Though here we will restrict ourselves to levels $`(M,N)`$ of the form $`(M,2M)`$, eq.(3.1) and the results for $`V(m,n)`$ given below can be used to construct the potential $`𝒱_{(M,N)}`$ for arbitrary level $`(M,N)`$ as long as $`M3`$ and $`N6`$. We shall first list all possible terms appearing in each $`V(m,n)`$ consistent with momentum conservation, separating the quadratic and cubic terms. We then use the methods described in section 2 to explicitly calculate the coefficients of each possible term in the $`V(m,n)`$’s.
The list of interactions that must be computed is generated conveniently with the help of the following function:
$`𝒵(x,y,s)`$ $``$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\{(1t_nx(y^n+y^n)s^{n^2/R^2})(1u_nx(y^n+y^n)s^{2+n^2/R^2})`$ (3.2)
$`\left(1v_nx(y^n+y^n)s^{2+n^2/R^2}\right)\left(1w_nx(y^n+y^n)s^{2+n^2/R^2}\right)`$
$`(1z_{n+1}x(y^{n+1}+y^{n1})s^{2+(n+1)^2/R^2})\mathrm{}\}^1`$
Here the formal variables $`x,y`$ and $`s`$ are used to count number of fields, momentum, and level, respectively. If we write
$$𝒵(x,y,s)=\underset{m,n}{}𝒵(m,n,s)x^my^n.$$
(3.3)
The momentum conserving cubic interactions appear in $`𝒵(3,0,s)`$ and an expansion in $`s`$ gives
$$𝒵(3,0,s)=\underset{l}{}𝒵(l)s^l.$$
(3.4)
Let $`\{\psi ^i\}`$ denote the complete set of modes $`(t_n,u_n,\mathrm{})`$. Then $`𝒵(l)`$ has an expression of the form $`𝒵(l)a_{ijk}\psi ^i\psi ^j\psi ^k`$ where each $`a_{ijk}`$ is an integer. If $`a_{ijk}0`$ the interaction $`\psi ^i\psi ^j\psi ^k`$ must be included in the level $`l`$ contribution to the potential. Thus $`𝒵(l)`$ supplies the complete list of momentum conserving cubic interactions of level $`l`$. When useful, we split by hand the terms in $`𝒵(l)`$ to obtain the possible terms which appear in various $`V(m,l)`$’s.
The list of all terms for the various $`V(m,n)`$’s with $`n6`$ (and $`R=\sqrt{3}`$) are given in Table 2.
### 3.1 The terms in the potential
The explicit interactions corresponding to the various terms appearing in the table will be listed here. With $`K=\frac{3\sqrt{3}}{4}`$, as in eq.(2.30), we have at the lowest level:
$`V(0,0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}t_0^2+{\displaystyle \frac{1}{3}}K^3t_0^3.`$ (3.5)
At first nontrivial level we have:
$`V(1/3,2/3)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{1}{R^2}}\right)t_1^2+{\displaystyle \frac{1}{2}}K^{32/R^2}t_0t_1^2.`$ (3.6)
At level 2 we have:
$`V(4/3,2)`$ $`=`$ $`{\displaystyle \frac{1}{4}}K^{36/R^2}t_1^2t_2`$
$`V(2,2)`$ $`=`$ $`{\displaystyle \frac{K}{32}}t_0^2\left(22u_05(v_0+25w_0)\right).`$ (3.7)
At level 8/3 :
$`V(4/3,8/3)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{4}{R^2}}\right)t_2^2+{\displaystyle \frac{1}{2}}K^{38/R^2}t_0t_2^2`$
$`V(2,8/3)`$ $`=`$ $`K^{12/R^2}t_1^2\left({\displaystyle \frac{11}{32}}u_0+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{R^2}}{\displaystyle \frac{5}{32}}\right)v_0{\displaystyle \frac{125}{64}}w_0\right)`$
$`V(7/3,8/3)`$ $`=`$ $`{\displaystyle \frac{1}{32}}K^{12/R^2}t_0t_1\left(22u_1\left(5+{\displaystyle \frac{16}{R^2}}\right)v_1125w_1+\left({\displaystyle \frac{44}{R^2}}+{\displaystyle \frac{32}{R^4}}\right)z_1\right).`$
At level 4:
$`V(2,4)`$ $`=`$ $`{\displaystyle \frac{1}{2}}u_0^2+{\displaystyle \frac{1}{4}}(v_0^2+25w_0^2)+K\{{\displaystyle \frac{1}{576}}t_0(76u_0^2+179v_0^2+9475w_0^2)`$
$`+{\displaystyle \frac{625}{864}}t_0v_0w_0{\displaystyle \frac{55}{432}}t_0u_0(v_0+25w_0)\}`$
$`V(7/3,4)`$ $`=`$ $`{\displaystyle \frac{1}{64}}K^{16/R^2}t_1t_2\left(22u_1\left(5{\displaystyle \frac{48}{R^2}}\right)v_1125w_1+\left({\displaystyle \frac{44}{R^2}}+{\displaystyle \frac{288}{R^4}}\right)z_1\right).`$
At level 14/3:
$`V(2,14/3)`$ $`=`$ $`{\displaystyle \frac{1}{64}}K^{18/R^2}t_2^2\left(22u_0\left(5{\displaystyle \frac{128}{R^2}}\right)v_0125w_0\right)`$
$`V({\displaystyle \frac{7}{3}},{\displaystyle \frac{14}{3}})`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(1+{\displaystyle \frac{1}{R^2}}\right)\left(2u_1^2+\left(1+{\displaystyle \frac{8}{R^2}}\right)v_1^2+25w_1^2+\left({\displaystyle \frac{8}{R^2}}+{\displaystyle \frac{16}{R^4}}\right)z_1^2+{\displaystyle \frac{24}{R^2}}v_1z_1\right)`$
$`+K^{12/R^2}\{{\displaystyle \frac{19}{288}}t_0u_1^2+{\displaystyle \frac{1}{3456}}(537+{\displaystyle \frac{8864}{R^2}}+{\displaystyle \frac{256}{R^4}})t_0v_1^2+{\displaystyle \frac{28425}{3456}}t_0w_1^2`$
$`{\displaystyle \frac{11}{864}}t_0u_1\left(\left(5+{\displaystyle \frac{16}{R^2}}\right)v_1+125w_1\right)+{\displaystyle \frac{125}{1728}}\left(5+{\displaystyle \frac{16}{R^2}}\right)t_0v_1w_1`$
$`+{\displaystyle \frac{19}{144}}t_1u_0u_1{\displaystyle \frac{11}{864}}t_1u_0\left(\left(5+{\displaystyle \frac{16}{R^2}}\right)v_1+125w_1\right)`$
$`{\displaystyle \frac{11}{864}}\left(5{\displaystyle \frac{32}{R^2}}\right)t_1v_0u_1+{\displaystyle \frac{1}{1728}}\left(537+{\displaystyle \frac{944}{R^2}}{\displaystyle \frac{512}{R^4}}\right)t_1v_0v_1`$
$`+{\displaystyle \frac{25}{1728}}\left(25{\displaystyle \frac{160}{R^2}}\right)t_1v_0w_1{\displaystyle \frac{1375}{864}}t_1w_0u_1+{\displaystyle \frac{25}{1728}}\left(25+{\displaystyle \frac{80}{R^2}}\right)t_1w_0v_1`$
$`+{\displaystyle \frac{28425}{1728}}t_1w_0w_1+{\displaystyle \frac{1}{216}}{\displaystyle \frac{1}{R^2}}\left(384+{\displaystyle \frac{1145}{R^2}}+{\displaystyle \frac{336}{R^4}}+{\displaystyle \frac{64}{R^6}}\right)t_0z_1^2`$
$`+{\displaystyle \frac{11}{864}}{\displaystyle \frac{1}{R^2}}\left(44+{\displaystyle \frac{32}{R^2}}\right)t_0u_1z_1+{\displaystyle \frac{1}{432}}\left({\displaystyle \frac{2359}{R^2}}+{\displaystyle \frac{1672}{R^4}}{\displaystyle \frac{128}{R^6}}\right)t_0v_1z_1`$
$`+{\displaystyle \frac{125}{432}}{\displaystyle \frac{1}{R^2}}\left(11{\displaystyle \frac{8}{R^2}}\right)\left(t_0w_1z_1+t_1w_0z_1\right)`$
$`+{\displaystyle \frac{1}{864}}{\displaystyle \frac{1}{R^2}}(11(44+{\displaystyle \frac{32}{R^2}})t_1u_0z_1+(2158{\displaystyle \frac{2832}{R^2}}+{\displaystyle \frac{512}{R^4}})t_1v_0z_1)\}`$
$`V(3,14/3)`$ $`=`$ $`{\displaystyle \frac{1}{2}}K^{314/R^2}t_1t_2t_3.`$ (3.10)
And at level 6:
$`V(2,6)`$ $`=`$ $`K\{{\displaystyle \frac{1}{144}}u_0^3+{\displaystyle \frac{8321}{93312}}v_0^3{\displaystyle \frac{219775}{10368}}w_0^3{\displaystyle \frac{95}{7776}}u_0^2(v_0+25w_0)`$
$`+{\displaystyle \frac{1969}{15552}}u_0v_0^2+{\displaystyle \frac{104225}{15552}}u_0w_0^2{\displaystyle \frac{22375}{31104}}v_0^2w_0{\displaystyle \frac{47375}{31104}}v_0w_0^2+{\displaystyle \frac{6875}{23328}}u_0v_0w_0\}`$
$`V(7/3,6)`$ $`=`$ $`K^{16/R^2}\{{\displaystyle \frac{19}{576}}t_2u_1^2+{\displaystyle \frac{1}{2304}}(179{\displaystyle \frac{1696}{R^2}}+{\displaystyle \frac{768}{R^4}})t_2v_1^2+{\displaystyle \frac{9475}{2304}}t_2w_1^2`$
$`{\displaystyle \frac{11}{1728}}\left(5{\displaystyle \frac{48}{R^2}}\right)t_2u_1v_1{\displaystyle \frac{1375}{1728}}t_2u_1w_1+{\displaystyle \frac{1}{72}}\left({\displaystyle \frac{625}{48}}{\displaystyle \frac{125}{R^2}}\right)t_2v_1w_1`$
$`+{\displaystyle \frac{1}{144}}\left({\displaystyle \frac{128}{R^2}}+{\displaystyle \frac{723}{R^4}}{\displaystyle \frac{2064}{R^6}}+{\displaystyle \frac{1728}{R^8}}\right)t_2z_1^2+{\displaystyle \frac{11}{432}}\left({\displaystyle \frac{11}{R^2}}+{\displaystyle \frac{72}{R^4}}\right)t_2u_1z_1`$
$`+{\displaystyle \frac{1}{288}}({\displaystyle \frac{67}{R^2}}{\displaystyle \frac{808}{R^4}}+{\displaystyle \frac{1152}{R^6}})t_2v_1z_1+{\displaystyle \frac{125}{864}}({\displaystyle \frac{11}{R^2}}{\displaystyle \frac{72}{R^4}})t_2w_1z_1\}`$
$`V(3,6)`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1+{\displaystyle \frac{9}{R^2}}\right)t_3^2+{\displaystyle \frac{1}{2}}K^{318/R^2}t_0t_3^2.`$ (3.11)
### 3.2 Potentials at various truncation levels and mass calculations
From these formulae one can construct the potentials at various truncation levels using (3.1). As we will use them, we give below the explicit sums for $`𝒱_{(1/3,2/3)}`$, $`𝒱_{(4/3,8/3)}`$, $`𝒱_{(2,4)}`$, $`𝒱_{(7/3,14/3)}`$ and $`𝒱_{(3,6)}`$:
$$\begin{array}{ccc}\hfill 𝒱_{(1/3,2/3)}& =& V(0,0)+V(1/3,2/3)\hfill \\ \multicolumn{3}{c}{}\\ \hfill 𝒱_{(4/3,8/3)}& =& 𝒱_{(1/3,2/3)}+V(4/3,2)+V(4/3,8/3)\hfill \\ \multicolumn{3}{c}{}\\ \hfill 𝒱_{(2,4)}& =& 𝒱_{(4/3,8/3)}+V(2,2)+V(2,8/3)+V(2,4)\hfill \\ \multicolumn{3}{c}{}\\ \hfill 𝒱_{(7/3,14/3)}& =& 𝒱_{(2,4)}+V(7/3,8/3)+V(7/3,4)+V(2,14/3)+V(7/3,14/3)\hfill \\ \multicolumn{3}{c}{}\\ \hfill 𝒱_{(3,6)}& =& 𝒱_{(7/3,14/3)}+V(3,14/3)+V(2,6)+V(7/3,6)+V(3,6)\hfill \end{array}$$
(3.12)
In general, the potential at a given level has many extrema. Two of them will be of particular interest for us:
1. We always find a translationally invariant minimum $`\stackrel{}{T}_{vac}`$ corresponding to the tachyon condensation. At this minimum, all fields with nonzero momentum have zero vev. We will use this solution when calculating the ratio $`r^{(2)}`$ defined in eq.(2.18).
2. If we start the numerical algorithm with initial values near $`t_00.25`$ and $`t_10.4`$ then our numerical algorithm converges to the one-lump solution $`\stackrel{}{T}_{lump}`$ that we are interested in.
The solution $`\stackrel{}{T}_{vac}`$ can be found in refs. . In table 3 we give the solutions $`\stackrel{}{T}_{lump}`$ at various truncation levels. Having found $`\stackrel{}{T}_{vac}`$ and $`\stackrel{}{T}_{lump}`$ we can now calculate the ratio of the lump mass to the D0-brane mass using the two different methods (2.17) and (2.18). The results are given in Table 4. We see that the first method gives a monotonically decreasing lump mass whereas the second method is oscillating but gives a lump mass much closer to the expected mass.
It is instructive to plot the profile of the tachyon field:
$$t(x)=\underset{n}{}t_n\mathrm{cos}\frac{nx}{R},$$
(3.13)
as a function of $`x`$ and compare them at different approximations. In figs.1-4 we have plotted the tachyon profiles at the level (1/3,2/3), (4/3,8/3), (2,4) and (7/3,14/3) approximation respectively, each of them being superimposed on the tachyon profile at the level (3,6) approximation.
For future use, we shall now define two new functions $`F_0`$ and $`G_0`$ as follows:
$`F_0(t_0,t_1,t_2,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R)`$ $`=`$ $`𝒱_{(\frac{7}{3},\frac{14}{3})},`$
$`G_0(t_0,t_1,t_2,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R)`$ $`=`$ $`𝒱_{(3,6)}V(3,6)V(2,14/3)V(7/3,4)`$ (3.14)
$`V(7/3,6)V(3,14/3)`$
where $`𝒱_{(\frac{7}{3},\frac{14}{3})}`$ and $`𝒱_{(3,6)}`$ have been defined in eqs.(3.5)-(3.12). The right hand side of this equation has to be interpreted as a function of the various modes $`t_0,\mathrm{}z_1`$ and $`R`$, without $`R`$ being set to $`\sqrt{3}`$. The function $`F_0`$ and $`G_0`$ defined here will be useful in constructing the potential $`𝒱_{(M,N)}`$ for other values of $`R`$, as will be discussed in the next section.
## 4 Tachyon Lump at Other Radii
In this section we shall discuss the construction of the tachyonic lump solution on circles of radii other than $`\sqrt{3}`$, and compare the results with those obtained for $`R=\sqrt{3}`$. As the basic techniques have already been discussed in the previous two sections, in this section we shall only quote the results.
### 4.1 $`R>\sqrt{3}`$
First we need to decide which values of $`R`$ we shall use to study the lump. Although this choice is arbitrary, there is slight simplification of counting levels if we choose $`R`$ such that the level of $`u_1`$, $`v_1`$, $`w_1`$ and $`z_1`$ coincide with that of one of the harmonics (say $`t_n`$) of the tachyon field. This requires
$$2+\frac{1}{R^2}=\frac{n^2}{R^2},R=\sqrt{\frac{n^21}{2}}.$$
(4.1)
We shall consider the values $`n=4,5,6`$ corresponding to $`R=\sqrt{\frac{15}{2}},\sqrt{12},\sqrt{\frac{35}{2}}`$. In each case we shall be using the level $`(2+\frac{1}{R^2},4+\frac{2}{R^2})`$ approximation to the potential. For this we need to include up to the $`n`$-th harmonic of the tachyon field $`t`$ and the first harmonics of the fields $`u,v,w`$ and $`z`$.
For these additional $`R`$ values all interactions present in $`𝒱_{(\frac{7}{3},\frac{14}{3})}`$ at $`R=\sqrt{3}`$ are still present. We need, however, further interactions as can be checked using the generating function (3.2). These additional interactions can be expressed in terms of the following functions:
$`F_1(t_0,\mathrm{},t_4,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{9}{R^2}}\right)t_3^2{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{16}{R^2}}\right)t_4^2`$
$`+{\displaystyle \frac{1}{2}}K^{318/R^2}t_0t_3^2+{\displaystyle \frac{1}{2}}K^{332/R^2}t_0t_4^2+{\displaystyle \frac{1}{2}}K^{314/R^2}t_1t_2t_3+{\displaystyle \frac{1}{4}}K^{324/R^2}t_2^2t_4+{\displaystyle \frac{1}{2}}K^{326/R^2}t_1t_3t_4`$
$`+\left({\displaystyle \frac{11}{32}}u_1{\displaystyle \frac{125}{64}}w_1+\left({\displaystyle \frac{25}{2R^4}}{\displaystyle \frac{11}{16R^2}}\right)z_1+\left({\displaystyle \frac{11}{4R^2}}{\displaystyle \frac{5}{64}}\right)v_1\right)K^{114/R^2}t_2t_3`$
$`F_2(t_0,\mathrm{},t_5,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{25}{R^2}}\right)t_5^2+{\displaystyle \frac{1}{2}}K^{338/R^2}t_2t_3t_5+{\displaystyle \frac{1}{2}}K^{342/R^2}t_1t_4t_5+{\displaystyle \frac{1}{2}}K^{350/R^2}t_0t_5^2`$
$`+\left({\displaystyle \frac{11}{32}}u_0+\left({\displaystyle \frac{5}{64}}+{\displaystyle \frac{9}{2R^2}}\right)v_0{\displaystyle \frac{125}{64}}w_0\right)K^{118/R^2}t_3^2`$
$`+\left({\displaystyle \frac{11}{32}}u_1+\left({\displaystyle \frac{5}{64}}+{\displaystyle \frac{23}{4R^2}}\right)v_1{\displaystyle \frac{125}{64}}w_1+\left({\displaystyle \frac{49}{2R^4}}{\displaystyle \frac{11}{16R^2}}\right)z_1\right)K^{126/R^2}t_3t_4`$
$`F_3(t_0,\mathrm{},t_6,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R)`$
$`=`$ $`{\displaystyle \frac{1}{4}}\left(1{\displaystyle \frac{36}{R^2}}\right)t_6^2+{\displaystyle \frac{1}{4}}K^{354/R^2}t_3^2t_6+{\displaystyle \frac{1}{2}}K^{356/R^2}t_2t_4t_6+{\displaystyle \frac{1}{2}}K^{362/R^2}t_1t_5t_6`$
$`+{\displaystyle \frac{1}{2}}K^{372/R^2}t_0t_6^2+\left({\displaystyle \frac{11}{32}}u_0+\left({\displaystyle \frac{5}{64}}+{\displaystyle \frac{8}{R^2}}\right)v_0{\displaystyle \frac{125}{64}}w_0\right)K^{132/R^2}t_4^2.`$
We shall now write down our results for level $`(2+\frac{1}{R^2},4+\frac{2}{R^2})`$ approximation for the potential for $`R^2=(n^21)/2`$ in terms of the functions $`F_0,\mathrm{}F_3`$ defined in eqs.(3.2), (LABEL:ef1)-(LABEL:ef3). These are as follows:
$`𝒱_{(32/15,64/15)}(t_0,\mathrm{},t_4,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{15/2})`$ (4.5)
$`=`$ $`F_0(t_0,t_1,t_2,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{15/2})`$
$`+F_1(t_0,\mathrm{},t_4,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{15/2}),`$
$`𝒱_{(25/12,25/6)}(t_0,\mathrm{},t_5,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{12})`$ (4.6)
$`=`$ $`F_0(t_0,t_1,t_2,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{12})`$
$`+F_1(t_0,\mathrm{},t_4,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{12})`$
$`+F_2(t_0,\mathrm{},t_5,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{12}),`$
$`𝒱_{(72/35,144/35)}(t_0,\mathrm{},t_6,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{35/2})`$ (4.7)
$`=`$ $`F_0(t_0,t_1,t_2,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{35/2})`$
$`+F_1(t_0,\mathrm{},t_4,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{35/2})`$
$`+F_2(t_0,\mathrm{},t_5,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{35/2})`$
$`+F_3(t_0,\mathrm{},t_6,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{35/2}).`$
As in the previous section, we can find a tachyonic lump solution by starting with a non-zero seed value of $`t_1`$. The numerical solutions are given in Table 5. The result for the two ratios $`r^{(1)}`$ and $`r^{(2)}`$, defined in eqs.(2.17) and (2.18) are given in Table 6.
In Figs.5-7 we have plotted the tachyon field $`t(x)`$ defined in eq.(3.13) as a function of $`x`$ for each of the three values of $`R`$. For reference we have also plotted on the same graph the function $`t(x)`$ obtained in the level (3,6) approximation for $`R=\sqrt{3}`$. As is seen from these figures, the tachyon profiles for different radii are almost undistinguishable from each other even though they are obtained as superpositions of harmonics of very different wave-lengths.
### 4.2 $`R<\sqrt{3}`$
Finally we would like to study how the shape of the soliton changes when $`R`$ is small. For this we take $`R=\sqrt{1.1}`$ and work at the level (40/11, 80/11) approximation of the potential. One can show that to this level of approximation the potential is given by,
$`𝒱_{(40/11,80/11)}(t_0,t_1,t_2,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{1.1})`$ (4.8)
$`=`$ $`G_0(t_0,t_1,t_2,u_0,v_0,w_0,u_1,v_1,w_1,z_1;R=\sqrt{1.1}),`$
where $`G_0`$ has been defined in eq.(3.2). The tachyonic lump solution for this potential is given in table 5. The results for the two ratios $`r^{(1)}`$ and $`r^{(2)}`$ defined in eqs.(2.17) and (2.18) are given in table 6.
We have displayed in fig. 8 the tachyon profile, superimposed on the tachyon profile for the level (3,6) approximation at $`R=\sqrt{3}`$. As can be seen from this figure, for $`R=\sqrt{11/10}`$ there is not enough room for the full lump solution to fit in, but the profile of the lump at smaller radius follows closely the profile at larger radius near the core.
### 4.3 Size of the lump
We can estimate the size of the lump at different radii in a somewhat systematic way by fitting the lump profile with a gaussian curve of the form:
$$G(x)=a+be^{x^2/(2\sigma ^2)}.$$
(4.9)
We calculate the parameters $`a`$, $`b`$ and $`\sigma `$ using a nonlinear regression algorithm on a set of points chosen on the lump profile in the following way. For $`R\sqrt{3}`$:
* We take 100 points, regularly spaced in $`x`$, in the core of the lump from $`x=\sqrt{3}\pi `$ to $`x=\sqrt{3}\pi `$.
* We take a smaller density of points, regularly spaced in $`x`$, on the rest of the circle (where the profile is essentially flat). Here we have taken 20, 30 and 40 points for $`R=\sqrt{15/2}`$, $`R=\sqrt{12}`$ and $`R=\sqrt{35/2}`$ respectively.
In the case of $`R=\sqrt{11/10}`$, we take 100 points from $`x=\sqrt{11/10}\pi `$ to $`x=\sqrt{11/10}\pi `$
The results of the regression at the different radii are given in table 7. We see that the size of the lump, which can be defined as a multiple of $`\sigma `$, is essentially independent of the radius (it increases by about 1.5 % when R increases from $`\sqrt{3}`$ to $`\sqrt{35/2}`$). Even when there is not enough room for the lump to fit in ($`R=\sqrt{11/10}`$), the lump is only slightly compressed (by about 7 %). A reasonable definition for the size would be $`6\sigma `$, with the solution extending by $`3\sigma `$ both along the positive and the negative $`x`$-axis. With this convention, the lump will have a size of approximately $`9.3\sqrt{\alpha ^{}}`$. This is close to the answer obtained in ref..
## 5 Conclusions and Open Questions
In this paper we have developed and tested the level expansion method in string field theory beyond translationally invariant vacuum solutions. This enabled us to give a systematic method for calculating quantities related to tachyon lumps and to give an accurate description of D-branes as tachyonic lumps in bosonic string field theory. Given the accuracy of our calculations (about 1% typically) we are confident that the profile of the lump that we have found is indeed very close to the exact one. As we have seen, as long as the radius is sufficiently big the lump has a definite radius independent profile. Indeed, when approximated by a gaussian, the lump representing a D-brane has $`\sigma 1.55\sqrt{\alpha ^{}}`$. We also considered the profile of the tachyon lump for $`R=\sqrt{1.1\alpha ^{}}`$, a radius sufficiently small that the large radius profile of the lump does not fit on the circle. We saw that the bottom part of the lump is essentially unchanged.
There are some questions related to the present work that we have not addressed. In particular we have not produced a lump solution in string field theory for $`R=1`$, where the tachyon harmonic $`t_1`$ becomes exactly marginal and the D0 and D1 branes have the same mass. Presumably, for small $`(R^21)`$ one must go fairly high in the level expansion to produce an accurate description. We have also not discussed the case $`R<1`$, where the D0 brane is unstable against decay into the D1 brane, or into the translationally invariant vacuum. We have also not tried to describe several D0 branes, all located at the same position.
We have not discussed issues related to the size of the lump representing a D-brane. While in the conformal field theory description a D-brane is an object with a well defined position, in string field theory it is a fat object, with thickness of the order of the string scale. Since string field theory is a gauge theory one may wonder if the size is an artifact of the chosen gauge. We do not at present know the answer to this question. The simplest way to get some insight into the nature of this extended solution would be to try to find out the energy density. This fails since the string field theory action is nonlocal, and hence there is no known expression for energy density in this theory. It would be interesting to examine some physical question that could help interpret the nature of this size . According to the conjectures of refs., all physical quantities calculated in the background of the lump solution must agree with those calculated in the background of a lower dimensional D-brane.
The methods used in this paper should be able to deal with:
* Neveu-Schwarz string field theory, where tachyon kinks rather than lumps represent lower dimensional D-branes. One way to deal with the boundary conditions on a circle would be to place both a kink and an anti-kink at diametrically opposite points of the circle. Another, probably more efficient way would be to include a Wilson line along the circle in such a way that the tachyon boundary conditions are twisted .
* Higher codimension D-branes. In it was observed that as the codimension is increased the naive use of the tachyon “bounce” gave increasingly worse approximations to the lump mass. We believe that our methods will enable calculations to any desired accuracy. The simplest situation would involve making two of the original brane dimensions into circles and including harmonics in both directions by simple extension of the methods of section 2.2.
* Intersecting D-branes. The simplest setup would be to begin with a D2-brane on a torus and generate a pair of transverse D1 branes intersecting at one point.
We hope that our analysis will ultimately provide a more refined understanding of string field theory and its geometry. One application is already apparent; if we could get a formulation of string field theory around the translationally invariant vacuum where the original D-brane is no longer present, such formulation will have more unbroken symmetries than the current formulation.
It is interesting to note that the level expansion method used here incorporates into the calculational scheme an ultra-violet (UV) cutoff. Since $`l=p^2+\mathrm{}`$, working at fixed $`l`$ implies a upper bound to the momentum (in the spatial directions). From this one is naturally led to propose a level expansion method for quantum string field theory. One approach could be to use the Euclidean version of the theory, and make periodic all directions including time, thus turning, at any fixed level $`M`$, the set of all relevant fields into a set of expansion coefficients $`c_n`$, with $`l(|\varphi _n)M`$. Since we are setting the whole system in a box, we also have a natural infra-red cutoff. The whole quantum path integral $`[dc_n]\mathrm{exp}(S(c_n)/\mathrm{})`$ could then be evaluated.<sup>9</sup><sup>9</sup>9Here $`S`$ should be the truncation, to the given level of approximation, of the full quantum action satisfying the exact quantum Batalin-Vilkovisky master equation. It is not clear that the cubic open string field theory provides such solution. The open-closed string field theory introduced in ref. does provide a well defined quantum action, but due to the non-polynomial nature of the action it is not clear how to carry out the level expansion in this theory. Alternatively, one could make all dimensions except time periodic. In this case the result would be the quantum mechanics of the wave functions $`c_n(t)`$. It would be exciting if the level expansion gave a concrete calculational definition of quantum string field theory, a definition one could in practice feed to a computer in order to calculate observables to any desired degree of precision.
Acknowledgements: We would like to thank R. Gopakumar, S. Minwalla, L. Rastelli and W. Taylor for discussions. A.S. acknowledges the Physics department of Harvard University and the Center for Theoretical Physics at MIT for hospitality during part of this work. The work of N.M. and B.Z. was supported in part by DOE contract #DE-FC02-94ER40818.
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# PSpace Reasoning for Graded Modal LogicsThis papers appeared in the Journal of Logic and Computation, Vol. 10 No. 99-47, pp. 1–22 2000.
## 1 Introduction
Propositional modal logics have found applications in many areas of computer science. Especially in the area of knowledge representation, the description logic (DL) $`𝒜𝒞`$, which is a syntactical variant of the propositional (multi-)modal logic $`𝐊_{}`$ \[Sch91\], forms the basis of a large number of formalisms used to represent and reason about conceptual and taxonomical knowledge of the application domain. The graded modal logic $`\mathrm{𝐆𝐫}(𝐊_{})`$ extends $`𝐊_{}`$ by *graded modalities* \[Fin72\], i.e., counting expressions which allow one to express statements of the form “there are at least (at most) $`n`$ accessible worlds that satisfy…”. This is especially useful in knowledge representation because (a) humans tend to describe objects by the number of other objects they are related to (a stressed person is a person given at least three assignments that are urgent), and (b) qualifying number restrictions (the DL’s analogue for graded modalities \[HB91\]) are necessary for modeling semantic data models \[CLN94\].
$`𝐊_{}`$ is decidable in PSpace and can be embedded into a decidable fragment of predicate logic \[AvBN98\]. Hence, there are two general approaches for reasoning with $`𝐊_{}`$: dedicated decision procedures \[Lad77, SSS91, GS96\], and the translation into first order logic followed by the application of an existing first order theorem prover \[OS97, Sch97\]. To compete with the dedicated algorithms, the second approach has to yield a decision procedure and it has to be efficient, because the dedicated algorithms usually have optimal worst-case complexity. For $`𝐊_{}`$, the first issue is solved and, regarding the complexity, experimental results show that the algorithm competes well with dedicated algorithms \[HS97\]. Since experimental result can only be partially satisfactory, a theoretical complexity result would be desirable, but there are no exact results on the complexity of the theorem prover approach.
The situation for $`\mathrm{𝐆𝐫}(𝐊_{})`$ is more complicated: $`\mathrm{𝐆𝐫}(𝐊_{})`$ is known to be decidable, but this result is rather recent \[HB91\], and the known PSpace upper complexity bound for $`\mathrm{𝐆𝐫}(𝐊_{})`$ is only valid if we assume unary coding of numbers in the input, which is an unnatural restriction. For binary coding no upper bound is known and the problem has been conjectured to be ExpTime-hard \[dHR95\]. This coincides with the observation that a straightforward adaptation of the translation technique leads to an exponential blow-up in the size of the first order formula. This is because it is possible to store the number $`n`$ in $`\mathrm{log}_kn`$-bits if numbers are represented in $`k`$-ary coding. In \[OSH96\] a translation technique that overcomes this problem is proposed, but a decision procedure for the target fragment of first order logic yet has to be developed.
In this work we show that reasoning for $`\mathrm{𝐆𝐫}(𝐊_{})`$ is not harder than reasoning for $`𝐊_{}`$ by presenting an algorithm that decides satisfiability in PSpace, even if the numbers in the input are binary coded. It is based on the tableaux algorithms for $`𝐊_{}`$ and tries to prove the satisfiability of a given formula by explicitly constructing a model for it. When trying to generalise the tableaux algorithms for $`𝐊_{}`$ to deal with $`\mathrm{𝐆𝐫}(𝐊_{})`$, there are some difficulties: (1) the straightforward approach leads to an incorrect algorithm; (2) even if this pitfall is avoided, special care has to be taken in order to obtain a space-efficient solution. As an example for (1), we will show that the algorithm presented in \[dHR95\] to decide satisfiability of $`\mathrm{𝐆𝐫}(𝐊_{})`$ is incorrect. Nevertheless, this algorithm will be the basis of our further considerations. Problem (2) is due to the fact that tableaux algorithms try to prove the satisfiability of a formula by explicitly building a model for it. If the tested formula requires the existence of $`n`$ accessible worlds, a tableaux algorithm will include them in the model it constructs, which leads to exponential space consumption, at least if the numbers in the input are not unarily coded or memory is not re-used. An example for a correct algorithm which suffers from this problem can be found in \[HB91\] and is briefly presented in this paper. Our algorithm overcomes this problem by organising the search for a model in a way that allows for the re-use of space *for each successor*, thus being capable of deciding satisfiability of $`\mathrm{𝐆𝐫}(𝐊_{})`$ in PSpace.
Using an extension of these techniques we obtain a PSpace algorithm for the logic $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$, which extends $`\mathrm{𝐆𝐫}(𝐊_{})`$ by inverse relations and intersection of relations. This solves an open problem from \[DLNN97\].
This paper is an significantly extended and improved version of \[Tob99\].
## 2 Preliminaries
In this section we introduce the graded modal logic $`\mathrm{𝐆𝐫}(𝐊_{})`$, the extension of the multi-modal logic $`𝐊_{}`$ with graded modalities, first introduced in \[Fin72\].
###### Definition 2.1 (Syntax and Semantics of $`\mathrm{𝐆𝐫}(𝐊_{})`$)
Let $`𝒫=\{p_0,p_1,\mathrm{}\}`$ be a set of propositional atoms and $``$ a set of *relation names*. The set of $`\mathrm{𝐆𝐫}(𝐊_{})`$-*formulae* is built according to the following rules:
1. every propositional atom is a $`\mathrm{𝐆𝐫}(𝐊_{})`$-formula, and
2. if $`\varphi ,\psi _1,\psi _2`$ are $`\mathrm{𝐆𝐫}(𝐊_{})`$-formulae, $`n`$, and $`R`$ is a relation name, then $`\neg \varphi `$, $`\psi _1\psi _2`$, $`\psi _1\psi _2`$, $`R_n\varphi `$, and $`[R]_n\varphi `$ are formulae.
The semantics of $`\mathrm{𝐆𝐫}(𝐊_{})`$-formulae is based on *Kripke structures*
$$𝔐=(W^𝔐,\{R^𝔐R\},V^𝔐),$$
where $`W^𝔐`$ is a non-empty set of worlds, each $`R^𝔐W^𝔐\times W^𝔐`$ is an *accessibility relation* on worlds (for $`R`$), and $`V^𝔐`$ is a *valuation* assigning subsets of $`W^𝔐`$ to the propositional atoms in $`𝒫`$. For a Kripke structure $`𝔐`$, an element $`xW^𝔐`$, and a $`\mathrm{𝐆𝐫}(𝐊_{})`$-formula, the model relation $``$ is defined inductively on the structure of formulae:
$`𝔐,x`$ $`p\text{iff}xV^𝔐(p)\text{for}p𝒫`$
$`𝔐,x`$ $`\neg \varphi \text{iff}𝔐,x\vDash ̸\varphi `$
$`𝔐,x`$ $`\psi _1\psi _2\text{iff}𝔐,x\psi _1\text{and}𝔐,x\psi _2`$
$`𝔐,x`$ $`\psi _1\psi _2\text{iff}𝔐,x\psi _1\text{or}𝔐,x\psi _2`$
$`𝔐,x`$ $`R_n\varphi \text{iff}\mathrm{}R^𝔐(x,\varphi )>n`$
$`𝔐,x`$ $`[R]_n\varphi \text{iff}\mathrm{}R^𝔐(x,\neg \varphi )n`$
where $`\mathrm{}R^𝔐(x,\varphi ):=|\{yW^𝔐(x,y)R^𝔐\text{and}𝔐,y\varphi \}|`$
The propositional modal logic $`𝐊_{}`$ is defined as the fragment of $`\mathrm{𝐆𝐫}(𝐊_{})`$ in which for all modal operators $`n=0`$ holds.
A formula is called *satisfiable* iff there exists a structure $`𝔐`$ and a world $`xW^𝔐`$ such that $`𝔐,x\varphi `$.
By $`\text{SAT}(𝐊_{})`$and $`\text{SAT}(\mathrm{𝐆𝐫}(𝐊_{}))`$we denote the sets of satisfiable formulae of $`𝐊_{}`$ and $`\mathrm{𝐆𝐫}(𝐊_{})`$, respectively.
As usual, the modal operators $`R_n`$ and $`[R]_n`$ are dual: $`\mathrm{}R^𝔐(x,\varphi )>n`$ means that in $`𝔐`$ more than $`n`$ $`R`$-successors of $`x`$ satisfy $`\varphi `$; $`\mathrm{}R^𝔐(x,\neg \varphi )n`$ means that in $`𝔐`$ all but at most $`n`$ $`R`$-successors satisfy $`\varphi `$.
In the following we will only consider formulae in *negation normal form* (NNF), a form in which negations have been pushed inwards and occur in front of propositional atoms only. We will denote the NNF of $`\neg \varphi `$ by $`\mathrm{}\varphi `$. The NNF can always be generated in linear time and space by successively applying the following equivalences from left to right:
$`\neg (\psi _1\psi _2)`$ $`\neg \psi _1\neg \psi _2`$ $`\neg R_n\psi `$ $`[R]_n\neg \psi `$
$`\neg (\psi _1\psi _2)`$ $`\neg \psi _1\neg \psi _2`$ $`\neg [R]_n\psi `$ $`R_n\neg \psi `$
## 3 Reasoning for $`\mathrm{𝐆𝐫}(𝐊_{})`$
Before we present our algorithm for deciding satisfiability of $`\mathrm{𝐆𝐫}(𝐊_{})`$, for historic and didactic reasons, we present two other solutions: an incorrect one \[dHR95\], and a solution that is less efficient \[HB91\].
From the fact that $`\text{SAT}(𝐊_{})`$ is PSpace-complete \[Lad77, HM92\], it immediately follows, that $`\text{SAT}(\mathrm{𝐆𝐫}(𝐊_{}))`$ is PSpace-hard. The algorithms we will consider decide the satisfiability of a given formula $`\varphi `$ by trying to construct a model for $`\varphi `$.
### 3.1 An incorrect algorithm
In \[dHR95\], an algorithm for deciding $`\text{SAT}(\mathrm{𝐆𝐫}(𝐊_{}))`$ is given, which, unfortunately, is incorrect. Nevertheless, it will be the basis for our further considerations and thus it is presented here. It will be referred to as the *incorrect* algorithm. It is based on an algorithm given in \[DLNN97\] to decide the satisfiability of the DL $`𝒜𝒞𝒩`$, which basically is the restriction of $`\mathrm{𝐆𝐫}(𝐊_{})`$, where, in formulae of the form $`R_n\varphi `$ or $`[R]_n\varphi `$ with $`n>0`$, necessarily $`\varphi =p\neg p`$ holds.
The algorithm for $`\mathrm{𝐆𝐫}(𝐊_{})`$ tries to build a model for a formula $`\varphi `$ by manipulating sets of constraints with the help of so-called *completion rules*. This is a well-known technique to check the satisfiability of modal formulae, which has already been used to prove decidability and complexity results for other DLs (e. g., \[SSS91, HB91, BBH96\]). These algorithms can be understood as variants of tableaux algorithms which are used, for example, to decide satisfiability of the modal logics $`𝐊_{}`$, $`𝐓_{}`$, or $`\mathrm{𝐒𝟒}_{}`$ in \[HM92\].
###### Definition 3.1
Let $`𝒱`$ be a set of variables. A *constraint system* (c.s.) $`S`$ is a finite set of expressions of the form ‘$`x\varphi `$’ and ‘$`Rxy`$’, where $`\varphi `$ is a formula, $`R`$, and $`x,y𝒱`$.
For a c.s. $`S`$, let $`\mathrm{}R^S(x,\varphi )`$ be the number of variables $`y`$ for which $`\{Rxy,y\varphi \}S`$. The c.s. $`[z/y]S`$ is obtained from $`S`$ by replacing every occurrence of $`y`$ by $`z`$; this replacement is said to be *safe* iff, for every variable $`x`$, formula $`\varphi `$, and relation symbol $`R`$ with $`\{xR_n\varphi ,Rxy,Rxz\}S`$ we have $`\mathrm{}R^{[z/y]S}(x,\varphi )>n`$.
A c.s. $`S`$ is said to contain a *clash*, iff for a propositional atom $`p`$, a formula $`\varphi `$, and $`mn`$:
$$\{xp,x\neg p\}S\text{or}\{xR_m\varphi ,x[R]_n\mathrm{}\varphi \}S.$$
Otherwise it is called *clash-free*. A c.s. $`S`$ is called *complete* iff none of the rules given in Fig. 1 is applicable to $`S`$.
To test the satisfiability of a formula $`\varphi `$, the incorrect algorithm works as follows: it starts with the c.s. $`\{x\varphi \}`$ and successively applies the rules given in Fig. 1, stopping if a clash is occurs. Both the rule to apply and the formula to add (in the $`_{}`$-rule) or the variables to identify (in the $`_{}`$-rule) are selected non-deterministically. The algorithm answers “$`\varphi `$ is satisfiable” iff the rules can be applied in a way that yields a complete and clash-free c.s. The notion of *safe* replacement of variables is needed to ensure the termination of the rule application \[HB91\].
Since we are interested in PSpace algorithms, non-determinism imposes no problem due to Savitch’s Theorem, which states that deterministic and non-deterministic polynomial space coincide \[Sav70\].
To prove the correctness of a non-deterministic completion algorithm, it is sufficient to prove three properties of the model generation process:
1. Termination: Any sequence of rule applications is finite.
2. Soundness: If the algorithm terminates with a complete and clash-free c.s. $`S`$, then the tested formula is satisfiable.
3. Completeness: If the formula is satisfiable, then there is a sequence of rule applications that yields a complete and clash-free c.s.
The error of the incorrect algorithm is, that is does not satisfy Property 2, even though the converse is claimed:
> Claim(\[dHR95\]): Let $`\varphi `$ be a $`\mathrm{𝐆𝐫}(𝐊_{})`$-formula in NNF. $`\varphi `$ is satisfiable iff $`\{x_0\varphi \}`$ can be transformed into a clash-free complete c.s. using the rules from Figure 1.
Unfortunately, the *if*-direction of this claim is not true, which we will prove by a simple counterexample. Consider the formula
$$\varphi =R_2p_1[R]_1p_2[R]_1\neg p_2.$$
On the one hand, $`\varphi `$ is not satisfiable. Assume $`𝔐,xR_2p_1`$. This implies the existence of at least three $`R`$-successors $`y_1,y_2,y_3`$ of $`x`$. For each of the $`y_i`$ either $`𝔐,y_ip_2`$ or $`𝔐,y_i\vDash ̸p_2`$ holds by the definition of $``$. Without loss of generality, there are two worlds $`y_{i_1},y_{i_2}`$ such that $`𝔐,y_{i_j}p_2`$, which implies $`𝔐,x\vDash ̸[R]_1\neg p_2`$ and hence $`𝔐,x\vDash ̸\varphi `$.
On the other hand, the c.s. $`S=\{x\varphi \}`$ can be turned into a complete and clash-free c.s. using the rules from Fig. 1, as is shown in Fig. 2. Clearly this invalidates the claim and its proof.
### 3.2 An alternative syntax
At this stage the reader may have noticed the cumbersome semantics of the $`[R]_n`$ operator, which origins from the wish that the duality $`\mathrm{}\varphi \neg \mathrm{}\neg \varphi `$ of $`𝐊`$ carries over to $`[R]_n\varphi \neg R_n\neg \varphi `$ in $`\mathrm{𝐆𝐫}(𝐊_{})`$. This makes the semantics of $`[R]_n`$ and $`R_n`$ un-intuitive. Not only does the $`n`$ in a diamond operator mean “more than $`n`$” while it means “less *or equal* than $`n`$” for a box operator. The semantics also introduce a “hidden” negation.
To overcome these problems, we will replace these modal operators by a syntax inspired by the counting quantifiers in predicate logic: the operators $`R_n`$ and $`R_n`$ with semantics defined by :
$`𝔐,xR_n\varphi `$ $`\text{iff}\mathrm{}R^𝔐(x,\varphi )n,`$
$`𝔐,xR_n\varphi `$ $`\text{iff}\mathrm{}R^𝔐(x,\varphi )n.`$
This modification does not change the expressivity of the language, since $`𝔐,xR_n\varphi `$ iff $`𝔐,xR_{n+1}\varphi `$ and $`𝔐,x[R]_n\varphi \text{iff}𝔐,xR_n\neg \varphi `$. We use the following equivalences to transform formulae in the new syntax into NNF:
$`\neg R_0\varphi `$ $`p\neg p`$
$`\neg R_n\varphi `$ $`R_{n1}\varphi \text{ iff }n>1`$
$`\neg R_n\varphi `$ $`R_{n+1}\varphi `$
### 3.3 A correct but inefficient solution
To understand the mistake of the incorrect algorithm, it is useful to know how soundness is usually established for the kind of algorithms we consider. The underlying idea is that a complete and clash-free c.s. induces a model for the formula tested for satisfiability:
###### Definition 3.2 (Canonical Structure)
Let $`S`$ be a c.s. The *canonical structure* $`𝔐_S=(W^{𝔐_S},\{R^{𝔐_S}R\},V^{𝔐_S})`$ *induced by* $`S`$ is defined as follows:
$`W^{𝔐_S}`$ $`=\{x𝒱x\text{ occurs in }S\},`$
$`R^{𝔐_S}`$ $`=\{(x,y)𝒱^2RxyS\},`$
$`V^{𝔐_S}(p)`$ $`=\{x𝒱xpS\}.`$
Using this definition, it is then easy to prove that the canonical structure induced by a complete and clash-free c.s. is a model for the tested formula.
The mistake of the incorrect algorithm is due to the fact that it did not take into account that, in the canonical model induced by a complete and clash-free c.s., there are formulae satisfied by the worlds even though these formulae do not appear as constraints in the c.s. Already in \[HB91\], an algorithm very similar to the incorrect one is presented which decides the satisfiability of $`𝒜𝒞𝒬`$, a notational variant of $`\mathrm{𝐆𝐫}(𝐊_{})`$.
The algorithm essentially uses the same definitions and rules. The only differences are the introduction of the $`_{\text{choose}}`$-rule and an adaption of the $`_{}`$-rule to the alternative syntax. The $`_{\text{choose}}`$-rule makes sure that all “relevant” formulae that are implicitly satisfied by a variable are made explicit in the c.s. Here, relevant formulae for a variable $`y`$ are those occuring in modal formulae in constraints for variables $`x`$ such that $`Rxy`$ appears in the c.s. The complete rule set for the modified syntax of $`\mathrm{𝐆𝐫}(𝐊_{})`$ is given in Fig. 3. The definition of *clash* has to be modified as well: A c.s. $`S`$ contains a clash iff
* $`\{xp,x\neg p\}S`$ for some variable $`x`$ and a propositional atom $`p`$, or
* $`xR_n\varphi S`$ and $`\mathrm{}R^S(x,\varphi )>n`$ for some variable $`x`$, relation $`R`$, formula $`\varphi `$, and $`n`$.
Furthermore, the notion of safe replacement has to be adapted to the new syntax: the replacement of $`y`$ by $`z`$ in $`S`$ is called *safe* iff, for every variable $`x`$, formula $`\varphi `$, and relation symbol $`R`$ with $`\{xR_n\varphi ,Rxy,Rxz\}S`$ we have $`\mathrm{}R^{[z/y]S}(x,\varphi )n`$.
The algorithm, which works like the incorrect algorithm but uses the expansion rules from Fig. 3—where $``$ is used as a placeholder for either $``$ or $``$—and the definition of clash from above will be called the *standard algorithm*; it is a decision procedure for $`\text{SAT}(\mathrm{𝐆𝐫}(𝐊_{}))`$:
###### Theorem 3.3 (\[HB91\])
Let $`\varphi `$ be a $`\mathrm{𝐆𝐫}(𝐊_{})`$-formula in NNF. $`\varphi `$ is satisfiable iff $`\{x_0\varphi \}`$ can be transformed into a clash-free complete c.s. using the rules in Figure 3. Moreover, each sequence of these rule-applications is finite.
While no complexity result is explicitly given in \[HB91\], it is easy to see that a PSpace result could be derived from the algorithm using the trace technique, employed in \[SSS91\] to show that satisfiability of $`𝒜𝒞`$, the notational variant for $`𝐊_{}`$, is decidable in PSpace.
Unfortunately this is only true if we assume the numbers in the input to be unary coded. The reason for this lies in the $`_{}`$-rule, which generates $`n`$ successors for a formula of the form $`R_n\varphi `$. If $`n`$ is unary coded, these successors consume at least polynomial space in the size of the input formula. If we assume binary (or $`k`$-ary with $`k>1`$) encoding, the space consumption is exponential in the size of the input because a number $`n`$ can be represented in $`\mathrm{log}_kn`$ bits in $`k`$-ary coding. This blow-up can not be avoided because the completeness of the standard algorithm relies on the generation *and identification* of these successors, which makes it necessary to keep them in memory *at one time*.
## 4 An optimal solution
In the following, we will present the algorithm which will be used to prove the following theorem; it contradicts the ExpTime-hardness conjecture in \[dHR95\].
###### Theorem 4.1
Satisfiability for $`\mathrm{𝐆𝐫}(𝐊_{})`$ is PSpace-complete if numbers in the input are represented using binary coding.
When aiming for a PSpace algorithm, it is impossible to generate all successors of a variable in a c.s. at a given stage because this may consume space that is exponential in the size of the input concept. We will give an optimised rule set for $`\mathrm{𝐆𝐫}(𝐊_{})`$-satisfiability that does not rely on the identification of successors. Instead we will make stronger use of non-determinism to guess the assignment of the relevant formulae to the successors by the time of their generation. This will make it possible to generate the c.s. in a depth first manner, which will facilitate the re-use of space.
The new set of rules is shown in Fig. 4. The algorithm that uses these rules is called the *optimised algorithm*. The definition of *clash* is taken from the standard algorithm. We do not need a $`_{}`$-rule.
At first glance, the $`_{}`$-rule may appear to be complicated and therefor is explained in more detail: like the standard $`_{}`$-rule, it is applicable to a c.s. that contains the constraint $`xR_n\varphi `$ if there are less than $`n`$ $`R`$-successors $`y`$ of $`x`$ with $`y\varphi S`$. The rule then adds a new successor $`y`$ to $`S`$. Unlike the standard algorithm, the optimised algorithm also adds additional constraints of the form $`y(\mathrm{})\psi `$ to $`S`$ for each formula $`\psi `$ appearing in a constraint of the form $`xR_n\psi `$. Since we have suspended the application of the $`_{}`$-rule until no other rule applies to $`x`$, by this time $`S`$ contains all constraints of the form $`xR_n\psi `$ it will ever contain. This combines the effects of both the $`_{\text{choose}}`$\- and the $`_{}`$-rule of the standard algorithm.
### 4.1 Correctness of the optimised algorithm
To establish the correctness of the optimised algorithm, we will show its termination, soundness, and completeness.
To analyse the memory usage of the algorithm it is very helpful to view a c.s. as a graph: A c.s. $`S`$ induces a labeled graph $`G(S)=(N,E,)`$ with
* The set of nodes $`N`$ is the set of variables appearing in $`S`$.
* The edges $`E`$ are defined by $`E:=\{xyRxyS\text{for some }R\}`$.
* $``$ labels nodes and edges in the following way:
+ For a node $`xN`$: $`(x):=\{\varphi x\varphi S\}`$.
+ For an edge $`xyE`$: $`(xy):=\{RRxyS\}`$.
It is easy to show that the graph $`G(S)`$ for a c.s. $`S`$ generated by the optimised algorithm from an initial c.s. $`\{x_0\varphi \}`$ is a tree with root $`x_0`$, and for each edge $`xyE`$, the label $`(xy)`$ is a singleton. Moreover, for each $`xN`$ it holds that $`(x)\text{clos}(\varphi )`$ where $`\text{clos}(\varphi )`$ is the smallest set of formulae satisfying
* $`\varphi \text{clos}(\varphi )`$,
* if $`\psi _1\psi _2\text{or}\psi _1\psi _2\text{clos}(\varphi )`$, then also $`\psi _1,\psi _2\text{clos}(\varphi )`$,
* if $`R_n\psi \text{clos}(\varphi )`$, then also $`\psi \text{clos}(\varphi )`$,
* if $`\psi \text{clos}(\varphi )`$, then also $`\mathrm{}\psi \text{clos}(\varphi )`$.
We will use the fact that the number of elements of $`\text{clos}(\varphi )`$ is bounded by $`2\times |\varphi |`$ where $`|\varphi |`$ denotes the length of $`\varphi `$. This is easily shown by proving
$$\text{clos}(\varphi )=\text{sub}(\varphi )\{\mathrm{}\psi \psi \text{sub}(\varphi )\}$$
where $`\text{sub}(\varphi )`$ denotes the set of all sub-formulae of $`\varphi `$. The size of $`\text{sub}(\varphi )`$ is obviously bounded by $`|\varphi |`$.
#### 4.1.1 Termination
First, we will show that the optimised algorithm always terminates, i.e., each sequence of rule applications starting from a c.s. of the form $`\{x_0\varphi \}`$ is finite. The next lemma will also be of use when we will consider the complexity of the algorithm.
###### Lemma 4.2
Let $`\varphi `$ be a formula in NNF and $`S`$ a c.s. that is generated by the optimised algorithm starting from $`\{x_0\varphi \}`$.
* The length of a path in $`G(S)`$ is limited by $`|\varphi |`$.
* The out-degree of $`G(S)`$ is bounded by $`|\text{clos}(\varphi )|\times 2^{|\varphi |}`$.
###### Proof.
For a variable $`xN`$, we define $`\mathrm{}(x)`$ as the maximum depth of nested modal operators in $`(x)`$. Obviously, $`\mathrm{}(x_0)|\varphi |`$ holds. Also, if $`xyE`$ then $`\mathrm{}(x)>\mathrm{}(y)`$. Hence each path $`x_1,\mathrm{},x_k`$ in $`G(S)`$ induces a sequence $`\mathrm{}(x_1)>\mathrm{}>\mathrm{}(x_k)`$ of natural numbers. $`G(S)`$ is a tree with root $`x_0`$, hence the longest path in $`G(S)`$ starts with $`x_0`$ and its length is bounded by $`|\varphi |`$.
Successors in $`G(S)`$ are only generated by the $`_{}`$-rule. For a variable $`x`$ this rule will generate at most $`n`$ successors for each $`R_n\psi (x)`$. There are at most $`|\text{clos}(\varphi )|`$ such formulae in $`(x)`$. Hence the out-degree of $`x`$ is bounded by $`|\text{clos}(\varphi )|\times 2^{|\varphi |}`$, where $`2^{|\varphi |}`$ is a limit for the biggest number that may appear in $`\varphi `$ if binary coding is used. ∎
###### Corollary 4.3 (Termination)
Any sequence of rule applications starting from a c.s. $`S=\{x_0\varphi \}`$ of the optimised algorithm is finite.
###### Proof.
The sequence of rules induces a sequence of trees. The depth and the out-degree of these trees is bounded in $`|\varphi |`$ by Lemma 4.2. For each variable $`x`$ the label $`(x)`$ is a subset of the finite set $`\text{clos}(\varphi )`$. Each application of a rule either
* adds a constraint of the form $`x\psi `$ and hence adds an element to $`(x)`$, or
* adds fresh variables to $`S`$ and hence adds additional nodes to the tree $`G(S)`$.
Since constraints are never deleted and variables are never identified, an infinite sequence of rule application must either lead to an arbitrary large number of nodes in the trees which contradicts their boundedness, or it leads to an infinite label of one of the nodes $`x`$ which contradicts $`(x)\text{clos}(\varphi )`$. ∎
#### 4.1.2 Soundness and Completeness
The following definition will be very helpful to establish soundness and completeness of the optimised algorithm:
###### Definition 4.4
A c.s. $`S`$ is called *satisfiable* iff there exists a Kripke structure $`𝔐=(W^𝔐,\{R^𝔐R\},V^𝔐)`$ and a mapping $`\alpha :𝒱W^𝔐`$ such that the following properties hold:
1. If $`y,z`$ are distinct variables such that $`Rxy,RxzS`$, then $`\alpha (y)\alpha (z)`$.
2. If $`x\psi S`$ then $`𝔐,\alpha (x)\psi `$.
3. If $`RxyS`$ then $`(\alpha (x),\alpha (y))R^𝔐`$.
In this case, $`𝔐,\alpha `$ is called a *model* of $`S`$.
It easily follows from this definition, that a c.s. $`S`$ that contains a clash can not be satisfiable and that the c.s. $`\{x_0\varphi \}`$ is satisfiable if and only if $`\varphi `$ is satisfiable.
###### Lemma 4.5 (Local Correctness)
Let $`S,S^{}`$ be c.s. generated by the optimised algorithm from a c.s. of the form $`\{x_0\varphi \}`$.
1. If $`S^{}`$ is obtained from $`S`$ by application of the (deterministic) $`_{}`$-rule, then $`S`$ is satisfiable if and only if $`S^{}`$ is satisfiable.
2. If $`S^{}`$ is obtained from $`S`$ by application of the (non-deterministic) $`_{}`$\- or $`_{}`$-rule, then $`S`$ is satisfiable if $`S^{}`$ is satisfiable. Moreover, if $`S`$ is satisfiable, then the rule can always be applied in such a way that it yields a c.s. $`S^{}`$ that is satisfiable.
###### Proof.
$`SS^{}`$ for any rule $``$ implies $`SS^{}`$, hence each model of $`S^{}`$ is also a model of $`S`$. Consequently, we must show only the other direction.
1. Let $`𝔐,\alpha `$ be a model of $`S`$ and let $`x\psi _1\psi _2`$ be the constraint that triggers the application of the $`_{}`$-rule. The constraint $`x\psi _1\psi _2S`$ implies $`𝔐,\alpha (x)\psi _1\psi _2`$. This implies $`𝔐,\alpha (x)\psi _i`$ for $`i=1,2`$. Hence $`𝔐,\alpha `$ is also a model of $`S^{}=S\{x\psi _1,x\psi _2\}`$.
2. Firstly, we consider the $`_{}`$-rule. Let $`𝔐,\alpha `$ be a model of $`S`$ and let $`x\psi _1\psi _2`$ be the constraint that triggers the application of the $`_{}`$-rule. $`x\psi _1\psi _2S`$ implies $`𝔐,\alpha (x)\psi _1\psi _2`$. This implies $`𝔐,\alpha (x)\psi _1`$ or $`𝔐,\alpha (x)\psi _2`$. Without loss of generality we may assume $`𝔐,\alpha (x)\psi _1`$. The $`_{}`$-rule may choose $`\chi =\psi _1`$, which implies $`S^{}=S\{x\psi _1\}`$ and hence $`𝔐,\alpha `$ is a model for $`S^{}`$.
Secondly, we consider the $`_{}`$-rule. Again let $`𝔐,\alpha `$ be a model of $`S`$ and let $`xR_n\varphi `$ be the constraint that triggers the application of the $`_{}`$-rule. Since the $`_{}`$-rule is applicable, we have $`\mathrm{}R^S(x,\varphi )<n`$. We claim that there is a $`wW^𝔐`$ with
$$(\alpha (x),w)R^𝔐,𝔐,w\varphi ,\text{and}w\{\alpha (y)RxyS\}.$$
($``$)
Before we prove this claim, we show how it can be used to finish the proof. The world $`w`$ is used to “select” a choice of the $`_{}`$-rule that preserves satisfiability: Let $`\{\psi _1,\mathrm{},\psi _n\}`$ be an enumeration of the set $`\{\psi xR_n\psi S\}`$. We set
$$S^{}=S\{Rxy,y\varphi \}\{y\psi _i𝔐,w\psi _i\}\{y\mathrm{}\psi _i𝔐,w\vDash ̸\psi _i\}.$$
Obviously, $`𝔐,\alpha [yw]`$ is a model for $`S^{}`$ (since $`y`$ is a fresh variable and $`w`$ satisfies $`()`$), and $`S^{}`$ is a possible result of the application of the $`_{}`$-rule to $`S`$.
We will now come back to the claim. It is obvious that there is a $`w`$ with $`(\alpha (x),w)R^𝔐`$ and $`𝔐,w\varphi `$ that is not contained in $`\{\alpha (y)Rxy,y\varphi S\}`$, because $`\mathrm{}R^𝔐(x,\varphi )n>\mathrm{}R^S(x,\varphi )`$. Yet $`w`$ might appear as the image of an element $`y^{}`$ such that $`Rxy^{}S`$ but $`y^{}\varphi S`$.
Now, $`Rxy^{}S`$ and $`y^{}\varphi S`$ implies $`y^{}\mathrm{}\varphi S`$. This is due to the fact that the constraint $`Rxy^{}`$ must have been generated by an application of the $`_{}`$-rule because it has not been an element of the initial c.s. The application of this rule was suspended until neither the $`_{}`$\- nor the $`_{}`$-rule are applicable to $`x`$. Hence, if $`xR_n\varphi `$ is an element of $`S`$ now, then it has already been in $`S`$ when the $`_{}`$-rule that generated $`y^{}`$ was applied. The $`_{}`$-rule guarantees that either $`y^{}\varphi `$ or $`y^{}\mathrm{}\varphi `$ is added to $`S`$. Hence $`y^{}\mathrm{}\varphi S`$. This is a contradiction to $`\alpha (y^{})=w`$ because under the assumption that $`𝔐,\alpha `$ is a model of $`S`$ this would imply $`𝔐,w\mathrm{}\varphi `$ while we initially assumed $`𝔐,w\varphi `$. ∎
From the local completeness of the algorithm we can immediately derive the global completeness of the algorithm:
###### Lemma 4.6 (Completeness)
If $`\varphi \text{SAT}(\mathrm{𝐆𝐫}(𝐊_{}))`$ in NNF, then there is a sequence of applications of the optimised rules starting with $`S=\{x_0\varphi \}`$ that results in a complete and clash-free c.s.
###### Proof.
The satisfiability of $`\varphi `$ implies that also $`\{x_0\varphi \}`$ is satisfiable. By Lemma 4.5 there is a sequence of applications of the optimised rules which preserves the satisfiability of the c.s. By Lemma 4.3 any sequence of applications must be finite. No generated c.s. (including the last one) may contain a clash because this would make it unsatisfiable. ∎
Note that since we have made no assumption about the order in which the rules are applied (with the exception that is stated in the conditions of the $`_{}`$-rule), the selection of the constraints to apply a rule to as well as the selection which rule to apply is “don’t-care” non-deterministic, i.e., if a formula is satisfiable, then this can be proved by an arbitrary sequence of rule applications. Without this property, the resulting algorithm certainly would be useless for practical applications, because any deterministic implementation would have to use backtracking for the selection of constraints and rules.
###### Lemma 4.7 (Soundness)
Let $`\varphi `$ be a $`\mathrm{𝐆𝐫}(𝐊_{})`$-formula in NNF. If there is a sequence of applications of the optimised rules starting with the c.s. $`\{x_0\varphi \}`$ that results in a complete and clash-free c.s., then $`\varphi \text{SAT}(\mathrm{𝐆𝐫}(𝐊_{}))`$.
###### Proof.
Let $`S`$ be a complete and clash-free c.s. generated by applications of the optimised rules. We will show that the canonical model $`𝔐_S`$ together with the identity function is a model for $`S`$. Since $`S`$ was generated from $`\{x_0\varphi \}`$ and the rules do not remove constraints from the c.s., $`x_0\varphi S`$. Thus $`𝔐_S`$ is also a model for $`\varphi `$ with $`𝔐_S,x_0\varphi `$.
By construction of $`𝔐_S`$, Property 1 and 3 of Definition 4.4 are trivially satisfied. It remains to show that $`x\psi S`$ implies $`𝔐_S,x\psi `$, which we will show by induction on the norm $``$ of a formula $`\psi `$. The norm $`\psi `$ for formulae in NNF is inductively defined by:
$$\begin{array}{ccccc}p\hfill & :=& \neg p\hfill & :=& 0\text{for }p𝒫\hfill \\ \psi _1\psi _2\hfill & :=& \psi _1\psi _2\hfill & :=& 1+\psi _1+\psi _2\hfill \\ R_n\psi \hfill & & & :=& 1+\psi \hfill \end{array}$$
This definition is chosen such that it satisfies $`\psi =\mathrm{}\psi `$ for every formula $`\psi `$.
* The first base case is $`\psi =p`$ for $`p𝒫`$. $`xpS`$ implies $`xV^{𝔐_S}(p)`$ and hence $`𝔐_S,xp`$. The second base case is $`x\neg pS`$. Since $`S`$ is clash-free, this implies $`xpS`$ and hence $`xV^{𝔐_S}(p)`$. This implies $`𝔐_S,x\neg p`$.
* $`x\psi _1\psi _2S`$ implies $`x\psi _1,x\psi _2S`$. By induction, we have $`𝔐_S,x\psi _1`$ and $`𝔐_S,x\psi _2`$ holds and hence $`𝔐_S,x\psi _1\psi _2`$. The case $`x\psi _1\psi _2S`$ can be handled analogously.
* $`xR_n\psi S`$ implies $`\mathrm{}R^S(x,\psi )n`$ because otherwise the $`_{}`$-rule would be applicable and $`S`$ would not be complete. By induction, we have $`𝔐_S,y\psi `$ for each $`y`$ with $`y\psi S`$. Hence $`\mathrm{}R^{𝔐_S}(x,\psi )n`$ and thus $`𝔐_S,xR_n\psi `$.
* $`xR_n\psi S`$ implies $`\mathrm{}R^S(x,\psi )n`$ because $`S`$ is clash-free. Hence it is sufficient to show that $`\mathrm{}R^{𝔐_S}(x,\psi )\mathrm{}R^S(x,\psi )`$ holds. On the contrary, assume $`\mathrm{}R^{𝔐_S}(x,\psi )>\mathrm{}R^S(x,\psi )`$ holds. Then there is a variable $`y`$ such that $`RxyS`$ and $`𝔐_S,y\psi `$ while $`y\psi S`$. For each variable $`y`$ with $`RxyS`$ either $`y\psi S`$ or $`y\mathrm{}\psi S`$. This implies $`y\mathrm{}\psi S`$ and, by the induction hypothesis, $`𝔐_S,y\mathrm{}\psi `$ holds which is a contradiction. ∎
The following theorem is an immediate consequence of Lemma 4.3, 4.6, and 4.7:
###### Corollary 4.8
The optimised algorithm is a non-deterministic decision procedure for $`\text{SAT}(\mathrm{𝐆𝐫}(𝐊_{}))`$.
### 4.2 Complexity of the optimised algorithm
The optimised algorithm will enable us to prove Theorem 4.1. We will give a proof by sketching an implementation of this algorithm that runs in polynomial space.
###### Lemma 4.9
The optimised algorithm can be implemented in PSpace
###### Proof.
Let $`\varphi `$ be the $`\mathrm{𝐆𝐫}(𝐊_{})`$-formula to be tested for satisfiability. We can assume $`\varphi `$ to be in NNF because the transformation of a formula to NNF can be performed in linear time and space.
The key idea for the PSpace implementation is the *trace technique* \[SSS91\], i.e., it is sufficient to keep only a single path (a trace) of $`G(S)`$ in memory at a given stage if the c.s. is generated in a depth-first manner. This has already been the key to a PSpace upper bound for $`𝐊_{}`$ and $`𝒜𝒞`$ in \[Lad77, SSS91, HM92\]. To do this we need to store the values for $`\mathrm{}R^S(x,\psi )`$ for each variable $`x`$ in the path, each $`R`$ which appears in $`\text{clos}(\varphi )`$ and each $`\psi \text{clos}(\varphi )`$. By storing these values in binary form, we are able to keep information *about* exponentially many successors in memory while storing only a single path at a given stage.
Consider the algorithm in Fig. 5, where $`_\varphi `$ denotes the set of relation names that appear in $`\text{clos}(\varphi )`$. It re-uses the space needed to check the satisfiability of a successor $`y`$ of $`x`$ once the existence of a complete and clash-free “subtree” for the constraints on $`y`$ has been established. This is admissible since the optimised rules will never modify this subtree once is it completed. Neither do constraints in this subtree have an influence on the completeness or the existence of a clash in the rest of the tree, with the exception that constraints of the form $`y\psi `$ for $`R`$-successors $`y`$ of $`x`$ contribute to the value of $`\mathrm{}R^S(x,\psi )`$. These numbers play a role both in the definition of a clash and for the applicability of the $`_{}`$-rule. Hence, in order to re-use the space occupied by the subtree for $`y`$, it is necessary and sufficient to store these numbers.
Let us examine the space usage of this algorithm. Let $`n=|\varphi |`$. The algorithm is designed to keep only a single path of $`G(S)`$ in memory at a given stage. For each variable $`x`$ on a path, constraints of the form $`x\psi `$ have to be stored for formulae $`\psi \text{clos}(\varphi )`$. The size of $`\text{clos}(\varphi )`$ is bounded by $`2n`$ and hence the constraints for a single variable can be stored in $`𝒪(n)`$ bits. For each variable, there are at most $`|_\varphi |\times |\text{clos}(\varphi )|=𝒪(n^2)`$ counters to be stored. The numbers to be stored in these counters do not exceed the out-degree of $`x`$, which, by Lemma 4.2, is bounded by $`|\text{clos}(\varphi )|\times 2^{|\varphi |}`$. Hence each counter can be stored using $`𝒪(n^2)`$ bits when binary coding is used to represent the counters, and all counters for a single variable require $`𝒪(n^4)`$ bits. Due to Lemma 4.2, the length of a path is limited by $`n`$, which yields an overall memory consumption of $`𝒪(n^5+n^2)`$. ∎
Theorem 4.1 now is a simple Corollary from the PSpace-hardness of $`𝐊_{}`$, Lemma 4.9, and Savitch’s Theorem \[Sav70\].
## 5 Extensions of the Language
It is possible to extend the language $`\mathrm{𝐆𝐫}(𝐊_{})`$ without loosing the PSpace property of the satisfiability problem. In this section we extend the techniques to obtain a PSpace algorithm for the logic $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$, which extends $`\mathrm{𝐆𝐫}(𝐊_{})`$ by intersection of accessibility relations and inverse relations. These extension are mainly motivated from the world of Description Logics, where they are commonly studied. In this context, the logic $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$ can be perceived as a notational variant of the Description Logic $`𝒜𝒞𝒬`$.
###### Definition 5.1 (Syntax and Semantics of $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$)
Let $`𝒫=\{p_0,p_1,\mathrm{}\}`$ be a set of proposition letters and let $``$ be a set of *relation names*. The set $`\overline{}:=\{R^1|R\}`$ is called the set of $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-relations.
The set of $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-*formulae* is the smallest set such that
1. every proposition letter is a $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-formula and,
2. if $`\varphi ,\psi _1,\psi _2`$ are formulae, $`n`$, and $`R_1,\mathrm{},R_k`$ are (possibly inverse) $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-relations, then $`\neg \varphi `$, $`\psi _1\psi _2`$, $`\psi _1\psi _2`$, $`R_1\mathrm{}R_k_n\varphi `$, and $`R_1\mathrm{}R_k_n\varphi `$ are $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-formulae.
The semantics are extended accordingly:
$`𝔐,x`$ $`R_1\mathrm{}R_k_n\varphi \text{iff}\mathrm{}(R_1\mathrm{}R_k)^𝔐(x,\varphi )n`$
$`𝔐,x`$ $`R_1\mathrm{}R_k_n\varphi \text{iff}\mathrm{}(R_1\mathrm{}R_k)^𝔐(x,\varphi )n`$
where
$$\mathrm{}(R_1\mathrm{}R_k)^𝔐(x,\varphi )=|\{yW^𝔐(x,y)R_1^𝔐\mathrm{}R_k^𝔐\text{and}𝔐,y\varphi \}|,$$
and, for $`R`$, we define
$$(R^1)^𝔐:=\{(y,x)(x,y)R^𝔐\}.$$
We will use the letters $`\omega ,\sigma `$ to range over intersections of $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-relations. By abuse of notation we will sometimes identify an intersection of relations $`\omega `$ with the set of relations occurring in it and write $`R\omega `$ iff $`\omega =R_1\mathrm{}R_k`$ and there is some $`1ik`$ with $`R=R_i`$. To avoid dealing with relations of the form $`(R^1)^1`$ we use the convention that $`(R^1)^1=R`$ for any $`R`$.
Obviously every $`\mathrm{𝐆𝐫}(𝐊_{})`$ formula is also a $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$ formula. Using standard bisimluation arguments one can show that $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$ is strictly more expressive than $`\mathrm{𝐆𝐫}(𝐊_{})`$.
### 5.1 Reasoning for $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$
We will use similar techniques as in the previous section to obtain a PSpace-algorithm for $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$. The definition of a constraint system remains unchanged, but we additionally require that, for any $`R`$, a c.s. $`S`$ contains the constraint ‘$`Rxy`$’ iff it contains the constraint ‘$`R^1yx`$’. For a c.s. $`S`$, an intersection of $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-relations $`\omega =R_1\mathrm{}R_k`$, and a formula $`\varphi `$, let $`\mathrm{}\omega ^S(x,\varphi )`$ be the number of variables $`y`$ such that $`\{R_1xy,\mathrm{},R_kxy,y\varphi \}S`$.
We modify the definition of *clash* to deal with intersection of relations as follows. A c.s. $`S`$ contains a clash iff
* $`\{xp,x\neg p\}S`$ for some variable $`x`$ and a proposition letter $`p`$, or
* $`x\omega _n\varphi S`$ and $`\mathrm{}\omega ^S(x,\varphi )>n`$ for some variable $`x`$, intersection of $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-relations $`\omega `$, formula $`\varphi `$ and $`n`$.
The set of rules dealing with the extended logic is shown in Figure 6. We require the algorithm to maintain a binary relation $`_S`$ between the variables in a c.s. $`S`$ with $`x_Sy`$ iff $`y`$ was inserted by the $`_{}`$-rule to satisfy a constraint for $`x`$. When considering the graph $`G(S)`$, the relation $`_S`$ corresponds to the successor relation between nodes. Hence, when $`x_Sy`$ holds we will call $`y`$ a successor of $`x`$ and $`x`$ a predecessor of $`y`$. We denote the transitive closure of $`_S`$ by $`_S^+`$. For a set of variables $`𝒳`$ and a c.s. $`S`$, we denote the subset of $`S`$ in which no variable from $`𝒳`$ occurs in a constraint by $`S𝒳`$. The $`_{}`$-, $`_{}`$\- and $`_{\text{choose}}`$-rule are called “non-generating rules” while the $`_{}`$-rule is called a “generating rule”. The algorithm which uses these rules will be called the $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-algorithm.
The $`_{}`$-rule, while looking complicated, is a straightforward extension of the $`_{}`$-rule for $`\mathrm{𝐆𝐫}(𝐊_{})`$, which takes into account that we also need to guess additional *relations* between the old variable $`x`$ and the freshly introduced variable $`y`$. The $`_{\text{choose}}`$-rule requires more explanation.
For $`\mathrm{𝐆𝐫}(𝐊_{})`$, the optimised algorithm generates a c.s. $`S`$ in a way that, whenever $`xR_n\psi S`$, then, for any $`y`$ with $`RxyS`$, either $`y\psi S`$ of $`y\mathrm{}\psi S`$. This was achieved by suspending the generation of any successors $`y`$ of $`x`$ until $`S`$ contained all constraints of the from $`x\varphi `$ it would ever contain. In the presence of inverse relations, this is no longer possible because $`y`$ might be generated as a predecessor of $`x`$ and hence before it was possible to know which $`\psi `$ might be relevant. There are at least two possible ways to overcome this problem. One is, to guess, for every $`x`$ and *every* $`\psi \text{clos}(\varphi )`$, whether $`x\psi `$ or $`x\mathrm{}\psi `$. In this case, since the termination of the optimised algorithm as shown in Lemma 4.3 relies on the fact that the modal depth strictly decreases along a path in the induced graph $`G(S)`$, termination would no longer be guaranteed. It would have to be enforced by different means.
Here, we use another approach. We can distinguish two different situations where $`\{x\omega _n\psi ,Rxy\}S`$ for some $`R\omega `$, and $`\{y\psi ,y\mathrm{}\psi \}S=\mathrm{}`$, namely, whether $`y`$ is a predecessor of $`x`$ ($`y_Sx`$) or a successor of $`x`$ ($`x_Sy`$). The second situation will never occur. This is due to the interplay of the $`_{}`$-rule, which is suspended until all known relevant information has been added for $`x`$, and the $`_{\text{choose}}`$-rule, which deletes certain parts of the c.s. whenever new constraints have to be added for predecessor variables.
The first situation is resolved by non-deterministically adding either $`y\psi `$ or $`y\mathrm{}\psi `$ to $`S`$. The subsequent deletion of all constraints involving variables from $`\{zy_S^+z\}`$, which corresponds to all subtrees of $`G(S)`$ rooted at successors of $`y`$, is necessary to make this rule “compatible” with the trace-technique we want to employ in order to obtain a PSpace-algorithm. The correctness of the trace-approach relies on the property that, once we have established the existence of a complete and clash-free “subtree” for a node $`x`$, we can remove this tree from memory because it will not be modified by the algorithm. In the presence of inverse relations this can be no longer taken for granted as can be shown by the formula
$$\varphi =R_1_0qR_1_1(pq)R_2_1R_2^1_0R_1_1p$$
Figure 7 shows the beginning of a run of the algorithm for $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$. After a number of steps, a successor $`y`$ of $`x`$ has been generated and the expansion of constraints has produced a complete and clash-free subtree for $`y`$. Nevertheless, the formula $`\varphi `$ is not satisfiable. The expansion of $`R_2_1R_2^1_0R_1_1p`$ will eventually lead to the generation of the constraint $`x\mathrm{}R_1_1p=R_1_0p`$, which clashes with $`yp`$. If the subtree for $`y`$ would already have been deleted from memory, this clash would go undetected. For this reason, the $`_{\text{choose}}`$-rule deletes all successors of the modified node, which, while duplicating some work, makes it possible to detect these clashes even when tracing through the c.s. A similar technique has been used in \[HST99\] to obtain a PSpace-result for a Description Logic with inverse roles.
### 5.2 Correctness of the Algorithm
As for $`\mathrm{𝐆𝐫}(𝐊_{})`$, we have to show termination, soundness, and correctness of the algorithm for $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$.
#### 5.2.1 Termination
Obviously, the deletion of constraints in $`S`$ makes a new proof of termination necessary, since the proof of Lemma 4.3 relied on this fact. Please note, that the Lemma 4.2 still holds for $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$.
###### Lemma 5.2 (Termination)
Any sequence of rule applications starting from a c.s. $`S=\{x_0\varphi \}`$ of the $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$ algorithm is finite.
###### Proof.
The sequence of rule applications induces a sequence of trees. As before, the depth and out-degree of this tree is bounded in $`|\varphi |`$ by Lemma 4.2. For each variable $`x`$, $`(x)`$ is a subset of the finite set $`\text{clos}(\varphi )`$. Each application of a rule either
* adds a constraint of the form $`x\psi `$ and hence adds an element to $`(x)`$, or
* adds fresh variables to $`S`$ and hence adds additional nodes to the tree $`G(S)`$, or
* adds a constraint to a node $`y`$ and deletes all subtrees rooted at successors of $`y`$.
Assume that algorithm does not terminate. Due to the mentioned facts this can only be because of an infinite number of deletions of subtrees. Each node can of course only be deleted once, but the successors of a single node may be deleted several times. The root of the completion tree cannot be deleted because it has no predecessor. Hence there are nodes which are never deleted. Choose one of these nodes $`y`$ with maximum distance from the root, i.e., which has a maximum number of ancestors in $`_S`$. Suppose that $`y`$’s successors are deleted only finitely many times. This can not be the case because, after the last deletion of $`y`$’s successors, the “new” successors were never deleted and thus $`y`$ would not have maximum distance from the root. Hence $`y`$ triggers the deletion of its successors infinitely many times. However, the $`_{\text{choose}}`$-rule is the only rule that leads to a deletion, and it simultaneously leads to an increase of $`(y)`$, namely by the missing concept which caused the deletion of $`y`$’s successors. This implies the existence of an infinitely increasing chain of subsets of $`\text{clos}(\varphi )`$, which is clearly impossible. ∎
#### 5.2.2 Soundness and Completeness
###### Lemma 5.3 (Soundness)
Let $`\varphi `$ be a $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-formula in NNF. If the completion rules can be applied to $`\{x_0\varphi \}`$ such that they yield a complete and clash-free c.s., then $`\varphi \text{SAT}(\mathrm{𝐆𝐫}(𝐊__{}^1))`$.
###### Proof.
Let $`S`$ be a complete and clash-free c.s. obtained by a sequence of rule applications from $`\{x_0\varphi \}`$. We show that the canonical structure $`𝔐_S`$ is indeed a model of $`\varphi `$, where the canonical structure for $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$ is defined as in Definition 3.2. Please note, that we need the condition “$`RxyS`$ iff $`R^1yxS`$” to make sure that all information from the c.s. is reflected in the canonical structure.
By induction over the norm of formulae $`\psi `$ as defined in the proof of Lemma 4.7, we show that, for a complete and clash-free c.s. $`S`$, $`x\psi S`$ implies $`𝔐_S,x\psi `$. The only interesting cases are when $`\psi `$ starts with a modal operator.
* $`x\omega _n\psi S`$ implies $`\omega ^S(x,\psi )n`$ because $`S`$ is complete. Hence, there are $`n`$ distinct variables $`y_1,\mathrm{},y_n`$ with $`y_i\psi S`$ and $`Rxy_iS`$ for each $`1in`$ and $`R\omega `$. By induction, we have $`𝔐_S,y_i\psi `$ and $`(x,y_i)\omega ^{𝔐_S}`$ and hence $`𝔐_S,x\omega _n\psi `$.
* $`x\omega _n\psi S`$ implies, for any $`R\omega `$ and any $`y`$ with $`RxyS`$, $`y\psi S`$ or $`y\mathrm{}\psi S`$. For any predecessor of $`x`$, this is guaranteed by the $`_{\text{choose}}`$-rule, for any successor of $`x`$ by the $`_{}`$-rule which is suspended until no non-generating rule rules can applied to $`x`$ or any predecessor of $`x`$ together with the reset-restart mechanism that is triggered by constraints “moving upwards” from a variable to its predecessor.
We show that $`\mathrm{}\omega ^{𝔐_S}(x,\psi )\mathrm{}\omega ^S(x,\psi )`$: assume $`\mathrm{}\omega ^{𝔐_S}(x,\psi )>\mathrm{}\omega ^S(x,\psi )`$. This implies the existence of some $`y`$ with $`(x,y)R^{𝔐_S}`$ for each $`R\omega `$ and $`𝔐_S,y\psi `$ but $`y\psi S`$. This implies $`y\mathrm{}\psi S`$, which, by induction yields $`𝔐_S,y\mathrm{}\psi `$ in contradiction to $`𝔐_S,y\psi `$.
Since constraints for the initial variable $`x_0`$ are never deleted from $`S`$, we have that $`x_0\varphi S`$ and hence $`𝔐_S,x_0\varphi `$ and $`\varphi \text{SAT}(\mathrm{𝐆𝐫}(𝐊__{}^1))`$. ∎
The following lemma combines the local and global completeness proof for the $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-algorithm
###### Lemma 5.4 (Completeness)
If $`\varphi \text{SAT}(\mathrm{𝐆𝐫}(𝐊__{}^1))`$ in NNF, then there is a sequence of the $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-rule starting with $`S=\{x_0\varphi \}`$ that results in a complete and clash-free c.s.
###### Proof.
Let $`𝔐`$ be a model for $`\psi `$ and $`\overline{}_\varphi `$ the set of relations that occur in $`\varphi `$ together with their inverse. We use $`𝔐`$ to guide the application of the non-deterministic completion rules by incremently defining a function $`\alpha `$ mapping variables from the c.s. to elements of $`W^𝔐`$. The function $`\alpha `$ will always satisfy the following conditions:
$$\begin{array}{cc}1.\hfill & \text{if }x\psi S\text{ then }𝔐,\alpha (x)\psi \hfill \\ 2.\hfill & \text{if }RxyS\text{ then }\{RRxyS\}=\{R(\alpha (x),\alpha (y))R^𝔐\}\overline{}_\varphi \hfill \\ 3.\hfill & \text{if }y,z\text{ are distinct variables such that }\{R_1xy,R_2xz\}S,\text{ then }\alpha (y)\alpha (z)\hfill \end{array}\}()$$
Claim: Whenever $`()`$ holds for a c.s. $`S`$ and a function $`\alpha `$ and a rule is applicable to $`S`$ then it can be applied in a way that maintains $`()`$.
* The $`_{}`$-rule: if $`x\psi _1\psi _2S`$, then $`𝔐,\alpha (x)(\psi _1\psi _2)`$. This implies $`𝔐,\alpha (x)\psi _i`$ for $`i=1,2`$, and hence the rule can be applied without violating $`()`$.
* The $`_{}`$-rule: if $`x\psi _1\psi _2S`$, then $`𝔐,\alpha (x)(\psi _1\psi _2)`$. This implies $`𝔐,\alpha (x)\psi _1`$ or $`𝔐,\alpha (x)\psi _2`$. Hence the $`_{}`$-rule can add a constraint $`x\chi `$ with $`\chi \{\psi _1,\psi _2\}`$ such that $`()`$ still holds.
* The $`_{\text{choose}}`$-rule: obviously, either $`𝔐,\alpha (y)\psi `$ or $`𝔐,\alpha (y)\mathrm{}\psi `$ for any variable $`y`$ in $`S`$. Hence, the rule can always be applied in a way that maintains $`()`$. Deletion of nodes does not violate $`()`$.
* The $`_{}`$-rule: if $`x\omega _n\varphi ^{}S`$, then $`𝔐,\alpha (x)\omega _n\varphi ^{}`$. This implies $`\mathrm{}\omega ^𝔐(\alpha (x),\varphi ^{})n`$. We claim that there is an element $`tW^𝔐`$ such that
$$\begin{array}{c}(\alpha (x),t)R^𝔐\text{ for each }R\omega ,\text{ and }𝔐,t\psi ,\text{ and }\hfill \\ t\{\alpha (y)RxyS\}\hfill \end{array}\}()$$
We will come back to this claim later. Let $`\psi _1,\mathrm{},\psi _k`$ be an enumeration of the set $`\{\psi x\sigma _mS\}`$ The $`_{}`$-rule can add the constraints
$`S^{}`$ $`=\{y\psi _i𝔐,t\psi _i\}\{y\mathrm{}\psi _i𝔐,t\vDash ̸\psi _i\}`$
$`S^{\prime \prime }`$ $`=\{RxyR\overline{}_\varphi ,(\alpha (x),t)R^𝔐\}\{RyxR\overline{}_\varphi ,(t,\alpha (x))R^𝔐\}`$
as well as $`\{y\varphi ^{}\}`$ to $`S`$. If we set $`\alpha ^{}:=\alpha [yt]`$, then the obtained c.s. together with $`\alpha ^{}`$ satisfies $`()`$.
Why does there exists an element $`t`$ that satisfies $`()`$? Let $`sW^𝔐`$ be an arbitrary element with $`(\alpha (x),s)\omega ^𝔐`$ and $`𝔐,s\psi `$ that appears as an image of an arbitrary element $`y`$ with $`RxyS`$ for some $`R\overline{}_\varphi `$. Condition 2 of $`()`$ implies that $`RxyS`$ for any $`R\omega `$ and also $`y\psi S`$ must hold as follows:
Assume $`y\psi S`$. This implies $`y\mathrm{}\psi S`$: either $`y_Sx`$, then in order for the $`_{}`$-rule to be applicable, no non-generating rules and especially the $`_{\text{choose}}`$-rule is not applicable to $`x`$ and its ancestor, which implies $`\{y\psi ,y\mathrm{}\psi \}S\mathrm{}`$. If not $`y_Sx`$ then $`y`$ must have been generated by an application of the $`_{}`$-rule to $`x`$. In order for this rule to be applicable no non-generating rule may have been applicable to $`x`$ or any of its ancestors. This implies that at the time of the generation of $`y`$ already $`x\omega _n\psi S`$ held and hence the $`_{}`$-rule ensures $`\{y\psi ,y\mathrm{}\psi \}S\mathrm{}`$.
In any case $`y\mathrm{}\psi S`$ holds and together with Condition 1 of $`()`$ this implies $`𝔐,s\vDash ̸\psi `$ which contradicts $`𝔐,s\psi `$.
Together this implies that, whenever an element $`s`$ with $`(\alpha (x),s)\omega ^𝔐`$ and $`𝔐,s\psi `$ is assigned to a variable $`y`$ with $`RxyS`$, then it must be assigned to a variable that contributes to $`\mathrm{}\omega ^S(x,\psi )`$. Since the $`_{}`$-rule is applicable there are less than $`n`$ such variables and hence there must be an unassigned element $`t`$ as required by $`()`$.
This concludes the proof of the claim. The claim yields the lemma as follows: obviously, $`()`$ holds for the initial c.s. $`\{x_0\varphi \}`$, if we set $`\alpha (x_0):=s_0`$ for an element $`s_0`$ with $`𝔐,s_0\varphi `$ (such an element must exist because $`𝔐`$ is a model for $`\varphi `$). The claim implies that, whenever a rule is applicable, then it can be applied in a manner that maintains $`()`$. Lemma 5.2 yields that each sequence of rule applications must terminate, and also each c.s. for which $`()`$ holds is necessarily clash-free. It cannot contain a clash of the form $`\{xp,x\neg p\}S`$ because this would imply $`𝔐,\alpha (x)p`$ and $`𝔐,\alpha (x)\vDash ̸p`$. It can neither contain a clash of the form $`x\omega _n\psi S`$ and $`\mathrm{}\omega ^S(x,\psi )>n`$ because $`\alpha `$ is an injective function on $`\{yRxyS\}`$ and preserves all relations in $`\overline{}_\varphi `$. Hence $`\mathrm{}\omega ^S(x,\psi )>n`$ implies $`\mathrm{}\omega ^𝔐(x,\psi )>n`$, which cannot be the case since $`𝔐,\alpha (x)\omega _n\psi `$. ∎
As a corollary of Lemma 5.2, 5.3, and 5.4 we get:
###### Corollary 5.5
The $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-algorithm is a non-deterministic decision procedure for $`\text{SAT}(\mathrm{𝐆𝐫}(𝐊__{}^1))`$.
### 5.3 Complexity of the Algorithm
As for the optimised algorithm for $`\mathrm{𝐆𝐫}(𝐊_{})`$, we have to show that the $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-algorithm can be implemented in a way that consumes only polynomial space. This is done similarly to the $`\mathrm{𝐆𝐫}(𝐊_{})`$-case, but we have to deal with two additional problems: we have to find a way to implement the “reset-restart” caused by the $`_{\text{choose}}`$-rule, and we have to store the values of the relevant counters $`\omega ^S(x,\psi )`$. It is impossible to store the values for each possible intersection of relations $`\omega `$ because the are exponentially many of these. Fortunately, storing only the values for those $`\omega `$ which actually appear in $`\varphi `$ is sufficient.
###### Lemma 5.6
The $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-algorithm can be implemented in PSpace.
###### Proof.
Consider the algorithm in Figure 8, where $`\mathrm{\Omega }_\varphi `$ denotes all intersections of relations that occur in $`\varphi `$. As the algorithm for $`\mathrm{𝐆𝐫}(𝐊_{})`$, it re-uses the space used to check for the existence of a complete and clash-free “subtree” for each successor $`y`$ of a variable $`x`$. Counter variables are used to keep track of the values $`\mathrm{}\omega ^S(x,\psi )`$ for all relevant $`\omega `$ and $`\psi `$. This can be done in polynomial space. Resetting a node and restarting the generation of its successors is achieved by resetting all successor counters. Please note, how the predecessor of a node is taken into account when initialising the counter variables.
Since the length of paths in a c.s. is polynomial bounded in $`|\varphi |`$ and all necessary book-keeping information can be stored in polynomial space, this proves the lemma. ∎
Obviously, $`\text{SAT}(\mathrm{𝐆𝐫}(𝐊__{}^1))`$ is PSpace-hard, hence Lemma 5.6 and Savitch’s Theorem \[Sav70\] yield:
###### Theorem 5.7
Satisfiability for $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$ is PSpace-complete if the numbers in the input are represented using binary coding.
As a simple corollary, we get the solution of an open problem in \[DLNN97\]:
###### Corollary 5.8
Satisfiability for $`𝒜𝒞𝒩`$ is PSpace-complete if the numbers in the input are represented using binary coding.
###### Proof.
The DL $`𝒜𝒞𝒩`$ is a syntactic restriction of the DL $`𝒜𝒞𝒬`$, which, in turn, is a syntactical variant of $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$. Hence, the $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$-algorithm can immediately be applied to $`𝒜𝒞𝒩`$-concepts. ∎
## 6 Conclusion
We have shown that by employing a space efficient tableaux algorithm satisfiability of the logic $`\mathrm{𝐆𝐫}(𝐊_{})`$ can be decided in PSpace, which is an optimal result with respect to worst-case complexity. Moreover, we have extended the technique to the logic $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$, which extends $`\mathrm{𝐆𝐫}(𝐊_{})`$ both by inverse relations and intersection of relations. This logic is a notational variant of the Description Logic $`𝒜𝒞𝒬`$, for which the complexity of concept satisfiability has also been open. This settles the complexity of the DL $`𝒜𝒞𝒩`$ for which the upper complexity bound with binary coding had also been an open problem \[DLNN97\]. While the algorithms presented in this work certainly are only optimal from the viewpoint of worst-case complexity, they are relatively simple and will serve as the starting-point for a number of optimisations leading to more practical implementations. They also serve as tools to establish the upper complexity bound of the problems and thus shows that tableaux based reasoning for $`\mathrm{𝐆𝐫}(𝐊_{})`$ and $`\mathrm{𝐆𝐫}(𝐊__{}^1)`$ can be done with optimum worst-case complexity. This establishes a kind of “theoretical benchmark” that all algorithmic approaches can be measured against.
#### Acknowledgments.
I would like to thank Franz Baader, Ulrike Sattler, and an anonymous referee for valuable comments and suggestions. Part of this work was supported by the DFG, Project No. GR 1324/3-1.
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# Modelling fluctuations of financial time series: from cascade process to stochastic volatility model
## 1 Introduction
As shown by most recent empirical studies on huge amount of data, the market price changes are characterized by several “universal” features : price increments are not correlated, volatilities are strongly (power-law) correlated and price increment probability density function (pdf) shapes depend on the time scale. From quasi Gaussian at rather large time scales, these pdf are characterized by fat tails at fine scales. Many authors in the recently emerged field of “econophysics” as well as in classical empirical finance, aim at proposing simple, discrete or continuous time models that are able to account for these observations. Among all the proposed models, one can distinguish several streams, from the simplest Brownian process, that constitutes the main tool used by practitians, to the class of “heteroskedastic” nonlinear processes as proposed in Refs. . To account for the letpokurtic nature of the small scale pdf, Mandelbrot and Fama proposed the Levy stable paradigm that has been recently improved in the “truncated Levy” version . More recently, an interesting comparison between market price variations and the fluctuations of the fluid velocity field in fully developed turbulence has been suggested . Besides the real pertinence of such an analogy that has been widely commented , this work opens the door to another important paradigm to model financial time series, namely multifractal processes. The multifractal processes<sup>1</sup><sup>1</sup>1people sometimes refer to “multi-affine” processes or processes that display “multi-scaling”, and the deeply connected mathematics of large deviations and multiplicative cascades, are well known to be useful to describe the intermittent nature of fully developed turbulence . Recent empirical findings suggest that in finance, this framework is also likely to be pertinent as far as the time scale dependence of the statistical properties of price variations is concerned.
Our purpose in this paper is twofold. First, we make a brief review of multifractals in order to specify what is a multifractal process. We try to provide several complementary points of view and to understand what are the main ingredients for “multi-scaling”. We also comment about the criticisms raised by several authors about multifractality in finance. Our second goal is to propose a simple multifractal “stochastic volatility” model that captures very well all the above mentionned features of financial fluctuations. This model, that has been originally introduced in Ref. , is compared to real data and some models proposed elsewhere. We discuss its possible multivariate extension in order to use it in management applications. The paper is organized as follows. The review on multifractal processes is made in section 2. We introduce notations, the related notions of multi-scaling, scale-invariance, cascade process and self-similarity kernel. We illustrate our purpose using empirical estimates for some high frequency financial data. In section 3 we review some findings of Ref. concerning the magnitude correlations for cascade models and suggest a link with $`1/f`$ processes as recently observed in financial time series. In section 4 we introduce the multifractal random walk defined in Ref. as a stochastic volatility model. We discuss its main properties and propose a natural multivariate generalization. Our discussion is illustrated by numerical simulations. In section 5 we propose estimators for the few parameters of our model and compute them for some intraday and daily time series. In section 6 we discuss some related works about multifractality in finance. Conclusions and some prospects are reported in section 7.
## 2 Multifractal processes and cascade models
In this section we briefly discuss the related notions of multifractality and multiplicative cascade. Most of the ideas and concepts that we recall below have been introduced in the field of fully developed turbulence where people aim at accounting for the so-called “intermittency phenomenon” (for a review of this subject see e.g., ).
### 2.1 Multifractality of financial time series
Let us consider the variations of a stochastic process $`X(t)`$ at a time scale $`l`$. For that purpose, one can consider the increments of the process, $`\delta _lX(t)=X(t+l)X(t)`$ or more generally its wavelet transform
$$T(t,l)=l^1\psi \left(\frac{t^{}t}{l}\right)X(t^{})𝑑t^{}$$
where $`\psi (t)`$ is the so-called analyzing wavelet, i.e, a function well localized in both Fourier and direct spaces<sup>2</sup><sup>2</sup>2One nice property of wavelet transform is that it can be inverted, i.e., one can recover the original signal from its wavelet coefficients. Another interesting feature is that there exist orthonormal wavelet bases. Such bases are very useful for signal synthesis and modelling, as it is illustrated for cascade processes in Ref. . Let us denote $`M(q,l)`$ the order $`q`$ absolute moment of $`\delta _lB(t)`$ or $`T(t,l)`$, (in this paper $`E(.)`$ will be used for the mathematical expectation and we will always suppose that the considered processes has stationary increments)
$$M(q,l)=E(|\delta _lX(t)|^q).$$
(1)
We will say that the process is scale-invariant, if the scale behavior of the absolute moment $`M(q,l)`$ is a power law. Let us call $`\zeta _q`$ the exponent of this power law, i.e.,
$$M(q,l)C_ql^{\zeta _q},$$
(2)
where $`C_q`$ is a prefactor that will be interpreted below. The process is called monofractal if $`\zeta _q`$ is a linear fonction of $`q`$ and multifractal if $`\zeta _q`$ is nonlinear. Note that, from the concavity of the moments of a random variable, it is easy to show that $`\zeta _q`$, as defined from the scaling behavior (2) in the limit $`l0^+`$, is necessarily a convex fonction of $`q`$. The same argument leads to the conclusion that such scaling behavior with a nonlinear $`\zeta _q`$ cannot hold for all scales $`l`$. Thus, for a multifractal process there exists at least one characteristic time $`T`$ (hereafter referred to as the integral time) above which the behavior (2) is no longer valid. Multifractality has been introduced in the context of fully developed turbulence in order to describe the spatial fluctuations of the fluid velocity at very high Reynolds number . As suggested by recent studies , multifractality is likely to be a pertinent concept to account for the prices fluctuations in financial time-series. This is illustrated in Fig. 1 where the $`\zeta _q`$ function is estimated for the future S&P500 index over the period 1988-1999. The original intraday time-series has been sampled at a $`10`$ mn rate (Fig. 1(a)) in order to obtain equi-sampled data. We consider the associated continuously compounded return time-series, (i.e., the logarithm of the index value) that has been detrended and de-seasonalized<sup>3</sup><sup>3</sup>3The amplitude of the return variations in each intraday period is normalized according to the estimated U-shaped intraday r.m.s.. The $`\zeta _q`$ spectrum in Fig. 1(c) is obtained using linear regression fit of “log-log” representations of the behavior of the $`q`$-th order moment versus the time scale as illustrated in Fig. 1(b). In this figure, the scales span an interval from 10 minutes to approximately 1 year. Moment estimates at larger time scales are very poor because of the finite size of the overall record. From the linear behavior of such curves, one clearly sees that the scale-invariance hypothesis is satisfied over around 3 decades. In Fig. 1(b) we have plotted $`\mathrm{log}_2\frac{M(q,l)}{M(1,l)^q}`$ versus $`\mathrm{log}_2(l)`$. The fact that such plots are not constant reflects the nonlinearity the $`\zeta _q`$ spectrum. The future S&P500 can thus be considered, at least at this description level, as a multifractal signal. Let us notice that we have computed, in Fig. 1(d), the $`\zeta _q`$ values for $`q`$-th order moments that include negative values of $`q`$. This can be achieved using a wavelet based technique that has been introduced in Refs. . This spectrum turns out to be well fitted by a parabolic shape $`\zeta _q=0.53q0.015q^2`$. The non linear parabolic component of $`\zeta _q`$ has been plotted in the inset of Fig. 1(d).
### 2.2 Multifractal processes, self-similar processes and multiplicative cascades.
Multifractality (in the sense defined above) is a notion that is often related to an underlying multiplicative cascading process. In the context of deterministic functions the situation is rather clear since the analyticity of the $`\zeta _q`$ spectrum is deeply connected to the self-similarity properties of the function . Roughly speaking, a function is self-similar if it can be written as a multiplicative cascade in an appropriate space-scale (or time-scale) representation . In that context, the so-called multifractal formalism is valid, i.e., one can relate the $`\zeta _q`$ spectrum to the $`D(h)`$ singularity spectrum that provides information about the statistical distribution of singularity (Hölder) exponents $`h`$. The things are somehow more complex for stochastic processes. One of the goals of this paper is to provide some simple elements about this subject.
In the mathematical literature, a process $`X(t)`$ is called self-similar of exponent $`H`$ if $`\lambda >0`$, $`\lambda ^HX(\lambda t)`$ is the same process as $`X(t)`$. According to this definition, the Brownian motion is self-similar with an exponent $`H=1/2`$. This definition is however too restrictive for our purpose since it excludes multifractal processes. Indeed, let us consider $`P_l(\delta X)`$ the probability density function (pdf) of $`\delta _lX(t)`$<sup>4</sup><sup>4</sup>4Note that from stationarity of the increments, the law of $`\delta _lX(t)`$ is the same as the law of $`X(l)`$ if one assumes that $`X(0)=0`$.. If $`X(t)`$ is self-similar with an exponent $`H`$, then it is easy to prove that
$$P_l(\delta X)=\lambda ^HP_{\lambda l}(\lambda ^H\delta X).$$
(3)
Then, the moments at scale $`l`$ and $`L=\lambda l`$ are related by
$$M(q,l)=C_q\left(\frac{l}{L}\right)^{qH},$$
(4)
with $`C_q=M(q,L)`$. Thus one has a “monofractal” process with $`\zeta _q=qH`$. In order to account for multifractality, one has to generalize this classical definition of self-similarity. This can be done by introducing a weaker notion, as originally proposed in the field of fully developed turbulence by B. Castaing and co-authors . According to Castaing’s definition of self-similarity, a process is self-similar if the increment pdf’s at scales $`l`$ and $`L=\lambda l`$ ($`\lambda >1`$) are related by the relationship :
$$P_l(\delta X)=G_{l,L}(u)e^uP_L(e^u\delta X)𝑑u,$$
(5)
where the self-similarity kernel $`G_{l,L}`$ depends only on $`l/L`$. Let us note that this definition generalizes Eq. (3) that corresponds to the “trivial” case $`G_{l,L}(u)=\delta (uH\mathrm{ln}(l/L))`$. This equation basically states that the pdf $`P_l`$ can be obtained through a “geometrical convolution” between the kernel $`G_{l,L}`$ and the pdf $`P_L`$. A simple argument shows that the logarithm of the Fourier transform of the kernel $`G_{l,L}`$ can be written as $`F_{l,L}(k)=\mathrm{ln}\widehat{G}_{l,L}(k)=F(k)\mathrm{ln}(l/L)`$ <sup>5</sup><sup>5</sup>5It essentially results from the fact that $`G_{l,L}`$ depends only on $`l/L`$ and satisfies the semi-group composition law $`G_{l_1,l_3}=G_{l_1,l_2}G_{l_2,l_3}`$ where $`l_1l_2l_3`$ and $``$ is the convolution product .. Thus, from Eq. (5), one can easily show that the $`q`$ order absolute moments at scales $`l`$ and $`L`$ are related by:
$$M(q,l)=\widehat{G}_{l,L}(iq)M(q,L)=M(q,L)\left(\frac{l}{L}\right)^{F(iq)},$$
(6)
and then $`C_q=M(q,L)`$ and $`\zeta _q=F(iq)`$. A nonlinear $`\zeta _q`$ spectrum implies that $`F`$ is nonlinear and thus that $`G`$ is different from a Dirac delta function<sup>6</sup><sup>6</sup>6Note that from the above mentionned semi-group property, the Levy theorem implies that $`G`$ is necessarily an infinitely divisible law. For example, the simplest non linear case is the so-called log-normal model that corresponds to a parabolic $`\zeta _q`$ function and thus to a function $`G`$ that is Gaussian.
The equation (5) can be interpreted as follows: the pdf at scale $`l`$, $`P_l`$ is written as a weighted superposition of the rescaled versions of the pdf at scale $`L`$, $`P_L`$, the self-similarity kernel $`G_{l,L}`$ being the associated distribution of weights. In the case of a monofractal process as described by Eq. (3), a single value of $`u`$ is sufficient in the equation (5) since $`P_l`$ and $`P_L`$ have the same shape and differ only by the scale factor $`e^u=(l/L)^H=\lambda ^H`$. This explains the Dirac function for the kernel $`G`$. This situation can be easily generalized by considering other shapes for the kernel $`G_{l,L}`$. In that case, the shapes of the pdf $`P_l`$ across scales are no longer the same: when going to small scales, fat tails emerge and the pdf become strongly leptokurtic (see Refs. or Fig. 4).
Let us now make the link with multiplicative cascades. This can be easily done if one consider discrete scales $`l_n=2^nL`$. Let us suppose that the local variation of the process $`\delta _{l_n}X`$ at scale $`l_n`$ is obtained from the variation at scale $`L`$ as
$$\delta _{l_n}X(t)=\left(\underset{i=1}{\overset{n}{}}W_i\right)\delta _LX(t)$$
(7)
where $`W_i`$ are i.i.d. random positive factors. This is the cascade paradigm. Realizations of such processes can be constructed using orthonormal wavelet bases as discussed in Ref. . If one defines the magnitude $`\omega (t,l)`$ at time $`t`$ and scale $`l`$ as the logarithm of “local volatility” :
$$\omega (t,l)=\frac{1}{2}\mathrm{ln}(|\delta _lX(t)|^2),$$
(8)
then the previous cascade equation becomes a simple random walk equation, at fixed time $`t`$, versus the logarithm of scales:
$$\omega (t,l_{n+1})=\omega (t,l_n)+\mathrm{ln}(W_{n+1}).$$
If the noise $`\mathrm{ln}W_i`$ is normal $`N(\mu ,\lambda ^2)`$, the pdf of $`\omega `$, $`P_l(\omega )`$, thus satisfies a simple diffusion equation with a Gaussian kernel:
$$P_{l_n}(\omega )=\left(N(\mu ,\lambda ^2)^np_L\right)(\omega )$$
(9)
where $``$ is the convolution product. Going back to the original variable $`\delta X`$, the previous equation corresponds exactly to Castaing’s formulation of self-similarity (5) with the log-normal propagator:
$$G_{l_n,L}=N(\mu ,\lambda ^2)^n=N(n\mu ,n\lambda ^2).$$
Conversely, let us consider a process that satisfies Castaing’s equation with a normal kernel $`G`$. This means that one can write,
$$\delta _lX(t)W\delta _{2l}X(t)$$
(10)
where $``$ means the equality in law of the two random variables and $`W`$ is a log-normal random variable which mean $`\mu `$ and variance $`\lambda ^2`$ do not depend on $`l`$. By iterating this equation $`n`$ times, one thus recover, at least heuristically, the cascade equation (7). Thus, the cascade picture across scales, constitutes a kind of paradigm of non-trivial self-similar processes. As explained in Ref. , the problem with such processes is that they involve representations (e.g., orthonormal wavelet bases) that are constructed on a discrete set of scales (e.g., dyadic scales $`l_n=2^n`$) and in turn cannot be invariant under continuous scale dilations.
## 3 Magnitude correlations and 1/f spectra
We have seen in the previous section that multifractality can be interpreted as a diffusion of the magnitude of the variations of the return from large time scales to small time scales. In the financial framework, magnitudes at all scales are nothing but a logarithmic representation of local volatilities. In this section we would like to address the problem of volatility correlations. The “heteroskedastic” nature of financial time-series is now a well established empirical fact. Volatility possesses long-range positive correlations: periods of strong activity alternate with quiet periods. A lot of models have been proposed to account for this phenomenon from the famous GARCH models to various stochastic volatility models. Let us proceed with the multifractal and cascade picture and study what kind of correlations are associated to these models. This problem has already been considered by Arneodo, Muzy and Sornette in Ref. (see also Refs ). These authors have shown that a log-normal cascade model on the dyadic tree associated to the orthonormal wavelet representation leads naturally to magnitude correlation functions $`C_\omega (l,\tau )=\text{Cov}(\omega (t,\tau ),\omega (t+l,\tau ))`$ that behave as $`\lambda ^2\mathrm{ln}(l/T)`$ for $`T>l>\tau `$. This behavior has been shown to provide good fits of the empirical estimates of the correlation functions from real data . In Fig. 2 is reported the magnitude correlation function $`C_\omega (l,\tau )`$ (we choose $`\tau =10`$ min) of the S&P500 time series studied in Fig. 1. One can see that, when plotted versus the logarithm of the time lag, $`\mathrm{ln}(l)`$, the correlation function decreases linearly with a slope $`\lambda ^20.025`$. The intercept of such straight line provides an estimator of the integral time $`T`$ that is, in our case, approximately $`T3`$ years (note that because we get an estimate of $`\mathrm{ln}(T)`$ the error on the value of $`T`$ is very large). We have checked that those results are stable when changing the reference time $`\tau `$ for return calculation.
As it will be illustrated in section 5, for various financial time series, the “cascade ansatz” is very pertinent to describe the volatility correlations. Let us notice that the very slow (logarithmic) decrease of the correlation functions for time lags below the integral time $`T`$, is reminiscent of the ultrametric nature of the tree naturally associated to the time-scale (or time-frequency) representation . Moreover, let us remark that, as far as power spectrum is concerned, Gaussian processes with such correlation functions can be seen as “$`1/f`$” processes. Indeed, if the correlation function is given by the above expression, the power spectrum can be shown to reduce to
$$S(f)=2\lambda ^2f^1_0^{Tf}x^1\mathrm{sin}(x)𝑑x.$$
(11)
In the high frequency limit $`f+\mathrm{}`$, we then have $`S(f)\lambda ^2\pi f^1`$. Another intuitive way to understand this property, comes from the fact that the logarithmic decay of the correlation function can be understood as the limit $`H0`$ in the power-law correlation function $`k^{2H}`$ of a fractional gaussian noise of exponent $`H`$. This property will be explicitely used in the discussion of section 6.3. Let us finally remark that $`1/f`$ spectra have been observed in a wide range of applications . Recently, Bonanno et al. suggested the possible pertinence of such processes to account for the fluctuations of the number of trades of different stocks.
## 4 A simple solvable multifractal model
As emphasized previously, multiplicative cascade models represent the paradigm of multifractal processes in that they contain the main ingredient leading to multifractality, i.e, the scale evolution of the magnitudes, from coarse to fine scales, is a random walk. Besides the problems of continuous scale invariance and stationarity of standard hierarchical constructions of such processes, they cannot be formulated using a stochastic evolution equation as one would expect for a model for financial time series. In this section we propose a “stochastic volatility” model that has been introduced in Ref. , that does not possess any of these drawbacks: it has stationary increments, it has log-normal multifractal properties and is invariant under continuous dilations. The key idea underlying this model is that the stochastic volatility possesses, as for cascading processes, a “$`1/f`$” spectrum, or, more precisely, a correlation function with a logarithmic behavior.
### 4.1 The multifractal random walk
Let us briefly recall the construction of the multifractal random walk (MRW) proposed in . A discretized version of the model $`X_{\mathrm{\Delta }t}`$ (using a time discretization step $`\mathrm{\Delta }t`$) is built by adding up $`t/\mathrm{\Delta }t`$ random variables :
$$X_{\mathrm{\Delta }t}(t)=\underset{k=1}{\overset{t/\mathrm{\Delta }t}{}}\delta X_{\mathrm{\Delta }t}[k],$$
where the process $`\{\delta X_{\mathrm{\Delta }t}[k]\}_k`$ is a noise whose variance is stochastic, i.e.,
$$\delta X_{\mathrm{\Delta }t}[k]=ϵ_{\mathrm{\Delta }t}[k]e^{\omega _{\mathrm{\Delta }t}[k]},$$
(12)
where $`\omega _{\mathrm{\Delta }t}[k]`$ is the logarithm of the stochastic variance. More specifically, we will choose $`ϵ_{\mathrm{\Delta }t}`$ to be a gaussian white noise independent of $`\omega `$ and of variance $`\sigma ^2\mathrm{\Delta }t`$. The choice for the process $`\omega _{\mathrm{\Delta }t}`$ introduced in , is dictated by the cascade picture. It corresponds to a gaussian stationary process whose covariance can be written
$$\text{Cov}(\omega _{\mathrm{\Delta }t}[k],\omega _{\mathrm{\Delta }t}[l])=\lambda ^2\mathrm{ln}\rho _{\mathrm{\Delta }t}[|kl|]$$
where $`\rho _{\mathrm{\Delta }t}`$ is chosen in order to mimic the correlation structure observed in cascade models with an integral time $`T`$:
$$\rho _{\mathrm{\Delta }t}[k]=\{\begin{array}{cc}\frac{T}{(|k|+1)\mathrm{\Delta }t}\hfill & \text{for}|k|T/\mathrm{\Delta }t1\hfill \\ 1\hfill & \text{otherwise}\hfill \end{array}$$
Hereafter, we will refer to the process $`\omega (t)`$ as the “magnitude process”. In order the variance of $`X_{\mathrm{\Delta }t}(t)`$ to converge when $`\mathrm{\Delta }t0`$, one must choose the mean of the process $`\omega _{\mathrm{\Delta }t}`$ such that
$$E\left(\omega _{\mathrm{\Delta }t}[k]\right)=\text{Var}\left(\omega _{\mathrm{\Delta }t}[k]\right)=\lambda ^2\mathrm{ln}(T/\mathrm{\Delta }t),$$
for which we find $`\text{Var}(X_{\mathrm{\Delta }t}(t))=\sigma ^2t`$. Let us review the multifractal properties of MRW.
### 4.2 $`\zeta _q`$ spectrum: computation of the moments
The $`q`$th-order moment of the increments of the MRW can be computed. Since, by construction, the increments of the model are stationary, the law of $`X_{\mathrm{\Delta }t}(t+l)X_{\mathrm{\Delta }t}(t)`$ does not depend on $`t`$ and is the same law as $`X_{\mathrm{\Delta }t}(l)`$. In Ref. , it is proven that the moments of $`X(l)X_{\mathrm{\Delta }t0^+}(l)`$ can be expressed as
$$E(X(l)^{2p})=\frac{\sigma ^{2p}(2p)!}{2^pp!}_0^l𝑑u_1\mathrm{}_0^l𝑑u_p\underset{i<j}{}\rho (u_iu_j)^{4\lambda ^2},$$
(13)
where $`\rho `$ is defined by
$$\rho (t)=\{\begin{array}{cc}T/|t|\hfill & \text{for}|t|T\hfill \\ 1\hfill & \text{otherwise}\hfill \end{array}.$$
Using this expression in the above integral, a straightforward scaling argument leads to
$$M(2p,l)=K_{2p}\left(\frac{l}{T}\right)^{p2p(p1)\lambda ^2},$$
(14)
where we have denoted the prefactor
$$K_{2p}=T^p\sigma ^{2p}(2p1)!!_0^1𝑑u_1\mathrm{}_0^1𝑑u_p\underset{i<j}{}|u_iu_j|^{4\lambda ^2}.$$
(15)
Note that $`K_{2p}`$ is nothing but the moment of order $`2p`$ of the random variable $`X(T)`$ or equivalently of $`\delta _TX(t)`$. From the above expression, we thus obtain
$$\zeta _{2p}=p2p(p1)\lambda ^2$$
and by analytical continuation, the corresponding full $`\zeta _q`$ spectrum is thus the parabola
$$\zeta _q=(qq(q2)\lambda ^2)/2.$$
(16)
Let us remark that one can show that $`K_q=+\mathrm{}`$ if $`\zeta _q<0`$ (i.e., $`q>2+1/\lambda ^2`$) and thus the pdf of $`\delta _lX(t)`$ have fat tails . In order to control the order of the first divergent moment (without changing $`\lambda `$), one could simply choose for the $`ϵ_{\mathrm{\Delta }t}`$’s a law with fat tails. Indeed, the prefactor $`\sigma ^{2p}(2p1)!!`$ in Eq. (15) comes directly from the fact that the $`ϵ_{\mathrm{\Delta }t}`$’s have been chosen to be Gaussian. Using instead fat tail laws (e.g., t-student laws) would allow us to control the divergence of this prefactor.
In Fig. 3 we have estimated the scaling behavior of the absolute moments $`M(q,l)`$ for a discrete simulation of a MRW (Fig. 3(a)). In order to simulate the sampling of a time continuous MRW, we have generated a discretized MRW using $`\mathrm{\Delta }t<<1`$ and then subsampled it at the sample period 1. Using this procedure, we have generated a $`2^{17}`$ long time-series using the parameters $`\mathrm{\Delta }t=1/16`$, $`T=2^{15}`$, $`\sigma ^2=1`$ and $`\lambda ^2=0.03`$. In Fig. 2(b), we have plotted, in double logarithmic representation $`M(q,l)`$ versus $`l`$ for different values of $`q`$. In these representations, the linear behavior of each moment indicates that the scaling hypothesis is verified. The estimation of $`\zeta _q`$ (made by estimating the slope of each of such curve) is reported in Fig 2(c). As expected this spectrum is a parabola that is in very good agreement with expression (16).
It is clear that the same power law scaling does not stand when $`l`$ goes to $`+\mathrm{}`$. Since $`\rho (l)=1`$ for large $`l`$ (as compared to $`T`$), we get
$`E(X(l)^{2p})`$ $`_{l>>T}`$ $`{\displaystyle \frac{\sigma ^{2p}(2p)!}{2^pp!}}{\displaystyle _0^l}𝑑u_1\mathrm{}{\displaystyle _0^l}𝑑u_p`$
$``$ $`Cl^p`$
Thus, there is a cross-over from the parabolic multifractal behavior at time scales $`lT`$ which is described by Eq. (16) to the Brownian-like behavior at larger time scales ($`l>>T`$)
$$\zeta _q=q/2.$$
In Eq. (6), we have shown that there exists a deep link between the self-similarity kernel and the $`\zeta _q`$ spectrum. This suggests that the probability distribution functions of our model satisfy Castaing’s equation when going from large to small time scales with a gaussian kernel $`G_{l,T}`$. Thus, as far as the increment pdf at different time scales are concerned, they will satisfy an evolution equation from “quasi-Gaussian” at very large scale ($`l>>T`$) to fat tailed pdf’s at small scales. This transformation of the pdf’s is illustrated in Fig. 4(a) where are plotted, in logarithmic scale, the standardized pdf’s (the variance has been set to one) for different time scales in the range $`[1,4T]`$. The pdf’s have been estimated for 500 realizations of size $`2^{17}`$ of MRW with parameters $`\lambda ^2=0.03`$ and $`T=2^{13}`$. In solid line, we have superimposed the Castaing’s transformation obtained from the coarse scale pdf (at scale $`T`$) using the appropriate normal self-similarity kernel. If Fig. 4(b) we have reproduced similar analysis for the S&P500 future variations. Besides statistical convergence limitations, one can observe the same features as in Fig. 4(a).
### 4.3 Volatility and magnitude correlation functions
#### 4.3.1 Volatility correlation functions
As recalled in the introduction, increments of financial time series are well known to be uncorrelated (for time lags large enough) while their amplitude (“local volatilities”) possesses power-law correlations. Let us show that our model satisfies these two properties at all time scales smaller than the “integral time” $`T`$. By construction, the increment correlation function,
$$(X_{\mathrm{\Delta }t}(t+\tau )X_{\mathrm{\Delta }t}(t))(X_{\mathrm{\Delta }t}(t_1+\tau _1)X_{\mathrm{\Delta }t}(t_1))$$
($`|t_1t|>\tau `$), is zero in our model. Let us study the correlation function of the squared increments. Since the increments are stationary, we can choose arbitrarily $`t_1=0`$. Thus we need to compute, in the limit $`\mathrm{\Delta }t0`$, the following correlation function, that corresponds to a lag $`l`$ between increments of size $`\tau `$
$$C(l,\tau )=(X_{\mathrm{\Delta }t}(l+\tau )X_{\mathrm{\Delta }t}(l))^2X_{\mathrm{\Delta }t}(\tau )^2.$$
(17)
From the results of Ref. and in the case $`0l<T`$, $`0\tau +l<T`$, we get,
$`C(l,\tau )=\sigma ^4{\displaystyle _l^{l+\tau }}𝑑u{\displaystyle _0^\tau }𝑑v\rho (uv)^{4\lambda ^2}.`$
A direct computation shows that
$`{\displaystyle _l^{l+\tau }}du{\displaystyle _0^\tau }dv|uv]^{4\lambda ^2}=`$
$`{\displaystyle \frac{1}{(14\lambda ^2)(24\lambda ^2)}}((l+\tau )^{24\lambda ^2}+(l\tau )^{24\lambda ^2}2l^{24\lambda ^2}),`$
and consequently
$$C(l,\tau )=K(|l+\tau |^{24\lambda ^2}+|l\tau |^{24\lambda ^2}2|l|^{24\lambda ^2})$$
(18)
where
$$K=\frac{\sigma ^4T^{4\lambda ^2}}{(14\lambda ^2)(24\lambda ^2)}.$$
Let us note that in the usual case $`0\tau <<l`$, one gets
$$C(l,\tau )\sigma ^4\tau ^2\left(\frac{l}{T}\right)^{4\lambda ^2}$$
(19)
i.e., for fixed $`\tau `$, the volatility correlation function scales as
$$C(l)l^{2\nu }$$
(20)
with $`\nu =2\lambda ^2`$. From the estimates $`\lambda ^20.0250.05`$ for financial assets (see section 5), one thus obtains $`\nu 0.050.1`$, values very close to the ones observed empirically in many works.
#### 4.3.2 Power of returns and magnitude correlation functions
Let us now show that magnitude correlation functions behave as expected, i.e, decrease very slowly as a logarithmic behavior.
For that purpose, the previous computation of the correlation function can be extended to the power of returns $`|X_{\mathrm{\Delta }t}(l+\tau )X_{\mathrm{\Delta }t}(\tau )|^p`$. Several empirical works have concerned the study of such “generalized volatilities” and people often noticed variations of amplitude of the correlation and of the power-law exponent $`\nu _p`$ when varying the order $`p`$ . In Ref. , it is shown that the quantity,
$$C_p(l,\tau )=|X_{\mathrm{\Delta }t}(l+\tau )X_{\mathrm{\Delta }t}(l)|^p|X_{\mathrm{\Delta }t}(\tau )|^p,$$
(21)
behaves, when $`\tau `$ is small enough, as
$$C_p(l,\tau )K_p^2\left(\frac{\tau }{T}\right)^{2\zeta _p}\left(\frac{l}{T}\right)^{\lambda ^2p^2}$$
(22)
where the constant $`K_p`$ has been defined previously. Using analytical continuation of the behavior of $`C_p`$ in the limit $`p=ϵ0`$, we can obtain, from previous expression, the behavior of the magnitude correlation function $`C_\omega (l,\tau )`$:
$$C_\omega (l,\tau )ϵ^2\left(C_ϵ(l,\tau )M(ϵ,\tau )^2\right)\lambda ^2\mathrm{ln}(\frac{l}{T}).$$
(23)
The magnitude correlation function, for $`\tau `$ small enough, has thus the same behavior as the correlation function of the underlying magnitude process $`\omega _{\mathrm{\Delta }t}`$. This result is checked in Fig. 4 where we have plotted the magnitude correlation function for $`\tau =32\mathrm{\Delta }t`$ as a function of $`\mathrm{ln}(l)`$. This correlation function has been estimated using a single realization of the process of $`2^{17}`$ sampled points, i.e, $`16`$ integral scales. The linear behavior we obtain is exactly the same one as predicted from Eq. (23) and Fig. 2. Measures of the slope and the intercept of such straight line provide a good estimate of respectively $`\lambda ^2`$ and $`T`$.
### 4.4 Extension to a multivariate process
In order to account for the fluctuations of financial portfolios and to consider management applications of our approach, it is important to build a multivariate, i.e. a vector valued, version of the previous multifractal random walk. Since only gaussian random variables are involved in the construction of section 4, this generalization can be done by considering two uncorrelated gaussian random vectors $`ϵ_{\mathrm{\Delta }t}(t)`$ and $`\omega _{\mathrm{\Delta }t}(t)`$ whose covariance matrices are denoted respectively $`𝚺`$ and $`𝚲`$. Hereafter, we will refer to these matrices as respectively the “Markowitz matrix” $`𝚺`$ and the “multifractal matrix” $`𝚲`$. One can then define the multivariate multifractal random walk (MMRW) $`\stackrel{}{X}(t)`$ as:
$$X_i(t+\mathrm{\Delta }t)X_i(t)=ϵ_i(t)e^{\omega _i(t)},$$
(24)
with $`\text{Cov}(ϵ_i(t),ϵ_j(t+\tau ))=\delta (\tau )𝚺_{ij}`$ and $`\text{Cov}(\omega _i(t),\omega _j(t+\tau ))=𝚲_{ij}\mathrm{ln}(T_{ij}/|\mathrm{\Delta }t+|\tau |)`$ (note that the previously defined coefficients $`\sigma ^2`$ and $`\lambda ^2`$ for an asset $`i`$ correspond respectively to the diagonal elements $`𝚺_{ii}`$ and $`𝚲_{ii}`$). Let us briefly review some of the properties of this model, postponing its detailed analysis to a forthcoming publication . A quantity that will be of central interest is the $`k`$-point joint moment of order $`q_1,q_2,\mathrm{},q_k`$ that can be defined as:
$$M_{i_1,\mathrm{},i_k}(q_1,\mathrm{},q_k)=E\left(|X_{i_1}(l)|^{q_1}\mathrm{}|X_{i_k}(l)|^{q_k}\right).$$
(25)
When $`k=2`$, by denoting $`i_1=i`$, $`i_2=j`$, $`q_1=p`$ and $`q_2=q`$, let us define the joint scaling exponent spectrum as:
$$M_{i,j}(p,q)=C_{i,j}(p,q)l^{\zeta _{i,j}(p,q)}.$$
(26)
This spectrum can be computed analytically. If the matrix $`𝚺`$ is diagonal (the $`ϵ_i`$’s are uncorrelated), a straightforward calculation shows that the scaling exponent $`\zeta _{i,j}(p,q)`$ is the following:
$$\zeta _{ij}(p,q)=\zeta _i(p)+\zeta _j(q)𝚲_{ij}pq,$$
(27)
where $`\zeta _i(q)`$ is the $`\zeta _q`$ spectrum for the component $`X_i(t)`$. Thus, for uncorrelated $`\omega _i`$’s, one has $`\zeta _{ij}(p,q)=\zeta _i(p)+\zeta _j(q)`$ while for the extreme case $`\omega _i=\omega _j`$, the exponent becomes $`\zeta _{ij}(p,q)=\zeta _i(p+q)=\zeta _j(p+q)`$. The computation of the scaling exponent is trickier for general Markowitz and multifractal matrices. Under some mild conditions that are necessary for the existence of a non trivial limit $`\mathrm{\Delta }t0^+`$, one can show that the previous scaling law remains valid even for non diagonal matrix $`𝚺`$ .
In order to define a simple way to get an estimate of the multifractal covariance coefficient $`𝚲_{ij}`$, let us define the moment ratio:
$$R_{ij}(q,l)=\frac{E(|X_i(l)|^q|X_j(l)|^q)}{E(|X_i(l)|^q)E(|X_j(l)|^q)}l^{\kappa _{ij}(q)}$$
(28)
From Eq. (27), the value of $`\kappa _{ij}(q)`$ is simply
$$\kappa _{ij}(q)=𝚲_{ij}q^2.$$
(29)
Thus, the non-diagonal element in the multifractal matrix $`𝚲`$ corrresponds to the nonlinear behavior of the exponent spectrum $`\kappa (q)`$ of the moment ratio $`R`$. Along the same line as for the computation of the magnitude auto-correlation in previous section, one can get the correlation function of magnitudes $`\omega _i(t,l)`$ and $`\omega _j(t,l)`$ from the limit $`q0`$ of $`R_{ij}(q,l)`$:
$$\text{Cov}(\omega _i(t,l),\omega _j(t,l))𝚲_{ij}\mathrm{ln}(l)+C,$$
(30)
where $`C`$ is a constant related to $`T_{ij}`$ . Thus the scale behavior of the magnitude covariance provides an estimate of the multifractal correlation coefficient $`𝚲_{ij}`$. This is the generalization of the classical result in multifractal analysis that relates the intermittency coefficient $`\lambda ^2=𝚲_{ii}`$ to the scale behavior of the variance of the magnitude.
Let us remark that the covariance of the variations of the assets $`i`$ and $`j`$ can be obtained by a direct calculation:
$$\text{Cov}(X_i(l),X_j(l))=𝚺_{ij}e^{\frac{1}{2}(𝚲_{ii}+𝚲_{jj}2𝚲_{ij})}l.$$
(31)
This covariance between $`X_i(t)`$ and $`X_j(t)`$ thus depends not only on $`𝚺`$, the “Markowitz” covariance matrix, but also on the multifractal matrix $`𝚲`$. This expression, allows us to get an estimate of the value of $`𝚺_{ij}`$ once the values of $`𝚲`$ are known.
Finally, let us mention that the idea of “multivariate multifractality” has been recently introduced in Ref. where the authors propose a phenomenological generalization of Castaing’s equation to the multivariate setting. Evidences that financial assets are characterized by non trivial multifractal matrices are also provided. We are currently working to obtain further empirical evidences towards such conclusions. Moreover, a precise link between the present model and the extended Castaing’s approach of Ref. is under progress.
## 5 Parameter estimation for real financial data
We have seen that the MRW is characterized mainly by 3 parameters: $`\sigma ^2`$, the white noise variance, $`T`$ the integral scale and $`\lambda ^2`$ the magnitude variance. We have shown that this model is able to reproduce all the main features of the future S&P 500 time series. Natural estimators of those parameters can be defined from the results of previous section. The parameter $`\lambda ^2`$ can be obtained from the shape of the $`\zeta _q`$ spectrum that is itself estimated using the scaling behavior of the absolute moments $`M(q,l)`$. This parameter can also be estimated thanks to the magnitude correlation function $`C_\omega (l,\tau )`$ that behaves as $`\lambda ^2\mathrm{ln}(l/T)`$. From the intercept of such correlation function as a function of $`\mathrm{ln}(l)`$, we can define an estimator of the integral scale $`T`$. Finally, the parameter $`\sigma ^2`$ can be obtained using the classical relationship $`\text{Var}(\delta _lX_{\mathrm{\Delta }t}(t))=\sigma ^2l`$. In this section, we report estimates of the multifractal parameters $`\lambda ^2`$ and $`T`$ for some financial time series. We do not have the ambition to provide fine estimates of those parameters. Our aim is rather to get an idea of realistical values of the parameters of the model for real assets. A precise discussion of the properties of various estimators from a statistical point of view is out of the scope of this paper and will be addressed in a forthcoming publication. Note that similar empirical study has already been performed in Ref. . We have studied some high frequency future time series that are sampled at a 10 min rate over the 7 years period from 1991 to 1997. We have also processed a set of daily index values for 8 different countries over the period from 1973 to 1997. The results are reported in table 1.
We remark that the values of the multifractal parameter $`\lambda ^2`$ are all very close to $`\mathrm{2.5\hspace{0.33em}10}^2`$ (excepted for the hong-kong index). The integral time $`T`$ values are centered around 3 years but with a large spread. Let us notice that we get an estimate of $`\mathrm{ln}(T)`$ and thus the error on the estimate of $`T`$ can be very large. We do not report here the values of the errors and confidence intervals for the proposed estimators that will be studied elsewhere.
## 6 Discussion about other approaches and findings
In this section, we make some comments about related studies that concern multifractals and finance.
### 6.1 Turbulence and finance
The analogy between turbulence and finance has been originally proposed by Ghashghaie et al. . These authors proposed to describe the pdf’s of FX price changes at different time scales in the same way physicists describe the pdf’s of velocity variations at different space separations in fully developed turbulence. This approach naturally leads to the notions of cascading process, Castaing’s formula and multifractality as described in section 2. This work suggests that the key mechanism at the origin of these observations, is an information cascade according to which short-term traders are influenced by long-term traders. This cascade is the analog of the Richardson’s kinetic energy cascade in turbulence where small eddies result from the breakdown of larger ones and so on . If the observations reported in Ref. strongly support this point of view, its quantitative understanding in terms of “microscopic” mechanisms remains an open question. In this section we would like to comment about some criticisms that have been raised about the analogy between turbulence and finance. The first one concerns the power spectrum behavior in both situations . In turbulence, Kolmogorov theory predicts a $`k^{5/3}`$ power spectrum that is confirmed in experimental situations. In finance, since price fluctuations are almost uncorrelated, they are characterized by a $`k^2`$ spectrum. For a general multifractal process, the exponent $`\beta `$ of the power spectrum behavior can be shown to be related to the value of $`\zeta _2`$: $`\beta =1+\zeta _2`$. Thus, from the cascading process point of view, nothing prevents the exponent $`\beta `$ from being equal to the exponent of the Brownian motion, i.e., $`\beta =2`$. In other words, as examplified by the MRW, a cascading process can have uncorrelated increments. We could also remark, that in turbulence $`\beta =5/3`$ has a dimensional origin, i.e., it is the exponent of the spatial spectrum of velocity fluctuations within an Eulerian description. If one adopts a Lagrangian description and one is interested by temporal fluctuations of a fluid particle velocity, then the dimensional value of the power spectrum exponent is $`\beta =2`$. Thus the value of this exponent is not a pertinent argument to reject the analogy with turbulence. Another difference that has been raised in concerns the behavior of the probability of return to origin $`P_l(0)`$ that has been shown to possess a scaling regime in finance while its behavior is more complex for a turbulent velocity field. First of all, let us point out that whatever the quantity studied (probability of return or absolute moments), it is well known that there is no observed well-defined scaling regime in turbulence: the classical “log-log” plots always display some curvature across scales. This curvature is Reynolds number dependent and several studies suggest that it vanishes, i.e. the field is scale-invariant, only in the limit of infinite Reynolds number . However, within the cascade paradigm and using Castaing’s equation, the scaling behavior of the probability of return to origin is easy to show. Indeed, by setting $`\delta X=0`$ in (5), one obtains, from the definition of $`\zeta _q`$ and the self-similarity kernel:
$$P_l(0)=P_T(0)G_{l,T}(u)e^u𝑑u=P_T(0)\left(\frac{l}{T}\right)^{\zeta _1}.$$
(32)
The exponent for the probability of return to origin is thus simply $`\zeta _1`$. For the log-normal stochastic volatility model introduced in section 4, we thus get
$$P_l(0)l^{\frac{1+3\lambda ^2}{2}}.$$
(33)
To conclude, neither the power spectrum exponent, nor the scaling behavior of the probability of return to origin can be used as argument against the existence of a cascading process at the origin of the fluctuations of financial time series.
### 6.2 Subordinated processes. Multifractal time
Subordinated processes are Markov processes in a time variable $`\mu (t)`$ that is itself an (increasing) random process . Such processes have been introduced in finance by Mandelbrot and Taylor to account for the existence of Levy stable laws as the result of a Brownian motion in some stochastic time. Today, the idea of modelling financial return fluctuations as a Brownian motion in a “fractal time, “trading time” or “financial time” can be found in many approaches. In Refs. , the multifractal nature of these fluctuations has been modelled by a (fractional) Brownian motion subordinated with a multifractal stochastic measure. In this section, without any concern for rigor, we would like to make a link between our stochastic volatility approach and the multifractal time approach of Mandelbrot and co-authors. Let us first remark that if we drop the noise $`ϵ`$ in Eq. (12) and keep only the stochastic volatility $`\sigma (t)`$, we can construct a stochastic measure $`\mu (dt)`$ that satisfies $`\mu (dt)=e^{\omega (t)}dt`$. Using exactly the same kind of computation as in section 4, one can show that this measure is stationary and its multifractal spectrum $`\tau (q)`$ is
$$\tau (q)=\zeta _{2q}1,$$
(34)
as usually defined by
$$\mu ([0,t])^qt^{\tau (q)+1}.$$
(35)
Let us note that the existence and the construction of such a measure that is stationary and possesses a continuous scale invariance, was at the heart of the construction in Refs. and was still an open problem. According to these studies, one can thus construct a multifractal process by simply considering the subordinated process $`S(t)=B(\mu ([0,t]))`$ where $`B(t)`$ is the standard Brownian motion. The $`\zeta _q`$ spectrum of such process would be exactly the same as the stochastic volatility process defined in section 4. This is not so surprising since, formally, a differential form for the process $`S(t)`$ would be $`dS=\frac{d\mu }{dt}dB(\mu (t))`$. If one assumes that a white noise that is subordinated remains a white noise, one thus obtains $`dS=e^{\omega (t)}dB(t)`$ that is the equation that defines the MRW of section 4. The questions of well-definiteness of this construction, its statistical properties and the precise mathematical justification of such results, will be addressed in a forthcoming work.
### 6.3 Some remarks about Bouchaud, Potters and Meyer’s model
Besides multifractal and cascade pictures, our present approach has been inspired by a recent paper by Bouchaud, Potters and Meyer . These authors have proposed a model that is very similar to ours: the stochastic volatility $`\sigma (t)`$ instead of being log-normal ($`e^{\omega (t)}`$) is a normal ($`\omega (t)`$) random process with long-range (power-law) correlations. By a simple analytical computation, they have shown that the q-order cumulants of such a process satisfy a simple scaling behavior but the moments display apparent multiscaling caused by a “competition” between the different cumulant behavior on a finite scale range. They thus conclude that a distinction between multifractality and such “apparent multifractality” is a difficult task for finite size time series. As far as multifractal analysis and modelling of financial time series are concerned, this work is very interesting and the previous assertion is undoubtedly difficult to infirm. However, let us remark that in order to illustrate their purpose, Bouchaud et al. choose a “stochastic volatility” $`\sigma (t)=e^{\omega (t)}`$ instead of their “monofractal” model $`\sigma (t)=|\omega (t)|`$. The reason invoked by the authors is that the log-normal is “a more realistic time series as compared with real data…without changing the feature of the above model, i.e. the very slow decay of the volatility correlations”. They thus claim that the scaling features of both models are the same and thus that the multifractality observed for the simulations of the “log-normal” volatility model is only apparent as predicted by their theory for the “normal” volatility model. The results reported in section 4 can be used to show that this interpretation is not correct. Let us indeed reconsider both results of Ref. and section 4. According to the “normal” volatility model, the moment of order $`q=2p`$ is written in terms of cumulants and behaves as :
$$M(2p,l=N\mathrm{\Delta }t)=A_{2p,0}N^{(1\nu )p}+\mathrm{}+A_{2,..,2}N^p$$
(36)
where the constants $`A_{q_1,..,q_k}`$ depend only $`q_i`$, $`\lambda ^2`$ the variance of $`\omega `$ and $`\nu `$ the exponent for the correlation function of $`\omega `$: $`C_\omega (l)l^\nu `$. According to this equation, if $`N`$ is small enough, $`M(2p,l)l^{(1\nu )p}`$ while, for $`N`$ very large, $`M(2p,l)l^p`$. The transition scale $`N^{}(q=2p)`$ above which the scaling exponent is $`\zeta _q=q/2`$ can be estimated if we define it as the scale where the contribution to the moment of order $`q`$ of the cumulant of order 4 and 2 are equal. Using the expression in Ref. for second and fourth cumulants, $`C_2`$ and $`C_4`$, we can show that at scale $`N^{}(q=2p)`$, we have $`(2p1)!!C_2^pp(p1)(2p1)!!C_2^{p2}C_4/6`$. From the value of $`C_4`$, we obtain ($`q>1`$):
$$N^{}(q=2p)=\left(p(p1)\nu ^22^{2(\nu 1)}\underset{m=1}{\overset{+\mathrm{}}{}}m^{2(\nu 1)}\right)^{\frac{1}{2\nu }}$$
(37)
This function only depends on $`q`$ and $`\nu `$ and is increasing as $`q+\mathrm{}`$. Thus, the larger the $`q`$ value, the wider the range of scales on which apparent multifractality exists. However, a numerical computation of the values of $`N^{}`$ for $`\nu =0.2`$ shows that the value $`N^{}100`$ is reached only for the moment of order $`p=15`$. The greatest moment value attained in practical situations is $`q6`$, for which $`N^{}<1`$ ! That means that for all moments less that $`10`$, the model of Bouchaud Potters and Meyer predicts the trivial spectrum $`\zeta _q=q/2`$ without any cross-over phenomenon. For their numerical simulations, they have used a log-normal model. However, within the log-normal ansatz, the conclusions of Ref. are questionnable since, when $`\nu `$ is small enough, this model is very close to the model introduced in section 4. Let us indeed consider as in that $`C_\omega (l)\lambda ^2\left(\mathrm{\Gamma }(\nu )\mathrm{cos}(\frac{\pi \nu }{2})\right)^{\frac{1}{2}}l^\nu `$ with $`\nu `$ very small<sup>7</sup><sup>7</sup>7Notice that there is no reason to consider the same value of $`\nu `$ for the normal and log-normal models. The results of this paper suggest that, in finance, the value for the correlation exponent in the log-normal model is very close to zero and significantly smaller than 0.2. By expanding this expression, we obtain
$$C_\omega (l)=\frac{\lambda ^2}{\nu }\lambda ^2\mathrm{ln}(l)+O(\nu \mathrm{ln}(l)).$$
(38)
If we set $`T=e^{\nu ^1}`$, then for $`1l<<T`$, this equation becomes
$$C_\omega (l)=\lambda ^2\mathrm{ln}(T/l)$$
(39)
that is the same correlation function as introduced for the multifractal model in section 4. Let us notice that for $`\nu =0.1`$ we have $`T2.10^4`$, $`T3.10^5`$ for $`\nu =0.08`$ and $`T5.10^8`$ for $`\nu =0.05`$ !. In this model $`T`$ is increasing very fast as $`\nu `$ goes to zero. We can thus conclude, that the model numerically studied in Ref. can be seen as multifractal from one point of view: whatever that scaling range $`[1,T]`$, there exists $`\nu `$ small enough ($`\nu 1/\mathrm{ln}(T)`$) such that the model displays multiscaling with log-normal $`\zeta _q`$ spectrum in this scale range. According to these remarks, we thus think that the multifractal picture is more realistic to describe multiscaling in financial time series.
## 7 Summary and prospects
In this paper we have reviewed what are the main features of multifractal processes. We have shown that the Multifractal Random Walk is a very attractive alternative to classical cascade processes in the sense that it is stationary, continuously scale-invariant and formulated using a simple stochastic evolution equation. As a model for financial engineering, MRW are interesting for many reasons. First, as illustrated in details for the S&P 500 intraday time series, this model is able to reproduce the main empirical properties observed for financial time series. Moreover, as Brownian motion and other stable walks, it is a “scale-free” model in the sense that it does not have to fit a particular time-scale since it is scale-invariant. This kind of stability with respect to time “aggregation” is a serious advantage as compared to classical ARCH-like models which parameters strongly depend on the time-scale one is interested in. Moreover, as discussed in section 4.4, a simple multivariate formulation of MRW can be proposed. To our knowledge, it is the first example of an extension of the notions of multifractality to a vector field. The empirical results reported in Ref. suggest that MMRW can be pertinent for portfolio theory. We are currently working on applications of MRW to classical problems of finance like management problems and option pricing theory.
From a theoretical point of view, MRW can be seen as the simplest model that contains the main ingredients for multifractality. In that respect, it can be very helpful to elucidate, in many fields where multiscaling is observed, what are the generic mechanisms that are involved leading to “non-trivial” self-similarity properties. Various “microscopic” models, as proposed in finance or other fields, could be considered within this perspective. It could also be interesting to recast our approach within a field theoretical formulation involving some renormalization procedure. From a mathematical point of view, this problem is deeply linked to the existence of a limit stochastic process when the sampling time $`\mathrm{\Delta }t`$ goes to zero. The convergence of the moments is not sufficient to prove this non trivial assertion. Such a limit could be very useful to develop a new stochastic calculus within which, for example, one could formulate the model of multifractal time of Mandelbrot and co-authors very naturally (see section 6.2). Finally, in a forthcoming work, we will discuss the generalization of such approach to other laws than the (log-)normal.
Acknowledgement
We acknowledge Matt Lee and Didier Sornette for the permission to use their financial data. We are also very grateful to Alain Arneodo, Jean-Philippe Bouchaud and Didier Sornette for interesting discussions and remarks.
All the computations in this paper have been made using the free GNU licensed sofware LastWave .
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# 1 Introduction
## 1 Introduction
In recent years, reaction–diffusion systems have been studied by many people, using different methods. Among them are the field theoretic methods, which allow for perturbative approaches to build up correlations in low dimensions . As mean field techniques can not be used for low dimensional systems, people are motivated to study stochastic models in low dimensions, which can be solved exactly. Moreover, solving one dimensional systems should in principle be easier. Applying a similarity transformation on an integrable model, one may construct stochastic models, their integrability may be not obvious. Recently, Some people have studied such transformations \[3–8\].
Exact results for some models in a one–dimensional lattice have been obtained, for example in . In These cases, the time evolution of the system is determined by a master equation . Models with no diffusion received less attention in the literature \[11–14\]: It is said that unless the system has long–range reactions , the time dependence involves exponential relaxation rather than power law behaviour typical of the fast diffusion reactions.
In , a 10–parameter family of stochastic models has been studied. In these models, the $`k`$–point equal time correlation functions $`n_in_j\mathrm{}n_k`$ satisfy linear differential equations involving no higher–order correlations. These linear equations for the average density $`n_i`$ has been solved. But, these set of equations may not be solved easily for higher order correlation functions. The spectrum is also partially obtained. The model which we address in this article is a special case of that 10–parameter stochastic model.
In this work, we report the exact solution for a system with two–particle annihilation and decoagulation. This model may be considered as a biased voting model, in the sense that there are two different opinions. If the two persons on two adjacent sites have different opinions, they may interact so that their opinions become the same. The bias parameter corresponds to the dominance of the left (or right) sight. In the absence of bias, this system is equivalent to the zero–temperature Glauber model . This system is related to free fermion system, through a similarity transformation, and hence is solvable. Note that the system itself is not a free fermion system and can not be solved by applying only Jordan–Wigner transformation.
When there is right–left symmetry, the average density decays to its final value in the form of power law ($`t^{\frac{1}{2}}`$). But in the general case (biased model) it decays in the form of an exponential. Moreover, the profile of the deviation of the average density from its final value is not uniform but exponential in terms of the site number. In fact, the parameter representing the right–left asymmetry, in some sense, determines the dominance of the right sites over the left sites, or vice versa.
The spectrum of the Hamiltonian of the system is found. It is shown that the steady state is two–fold degenerate. The probability of finding the system in each of these two states is determined by the initial average density, and is time–independent. It is shown that at large times, any $`n`$–point function is equal to the 1–point function, which is position–independent.
$$n_i(\mathrm{})n_j(\mathrm{})\mathrm{}n_k(\mathrm{})=n_i(\mathrm{})=\frac{1}{L}\underset{m}{}n_m(0)$$
(1)
This is due to the fact that the system has two steady states; either completely full, or completely empty, as it will be shown. This means that the mean–field approach does not work and this system is highly correlated.
The scheme of the paper is as follows. In section 2, similarity transformations relating stochastic systems to other (stochastic or non–stochastic) systems are investigated. In section 3, a solvable model is obtained through a similarity transformation on a free–fermion system. The spectrum of the system is also obtained in this section. In section 4, the 1–point function is calculated and its large–time behavior is investigated. In section 5, the two–point function and its limiting behavior is obtained. In section 6, The null vectors of the Hamiltonian are obtained and from that the steady state of the system is obtained in terms of its one–point function at $`t=0`$. Finally, in section 7 we consider the next–to–leading term of the one–point function at large times, and from this obtain the way the system relaxes to its final state.
## 2 Similarity transformations as a method for obtaining solvable stochastic models
Here some standard material is introduced, just to fix notation. The master equation for $`P(\sigma ,t)`$ is
$$\frac{}{t}P(\sigma ,t)=\underset{\tau \sigma }{}\left[\omega (\tau \sigma )P(\tau ,t)\omega (\sigma \tau )P(\sigma ,t)\right],$$
(2)
where $`\omega (\tau \sigma )`$ is the transition rate from the configuration $`\tau `$ to $`\sigma `$. Introducing the state vector
$$|P(t)=\underset{\sigma }{}P(\sigma ,t)|\sigma ,$$
(3)
where the summation runs over all possible states of the system, one can write the above equation in the form
$$\frac{}{t}|P=|P,$$
(4)
where the matrix elements of $``$ are
$`\sigma ||\tau =\omega (\tau \sigma ),\tau \sigma ,`$ (5)
$`\sigma ||\sigma ={\displaystyle \underset{\tau \sigma }{}}\omega (\sigma \tau ).`$ (6)
The basis $`\{\sigma |\}`$ is dual to $`\{|\sigma \}`$, that is
$$\sigma |\tau =\delta _{\sigma ,\tau }.$$
(7)
The operator is $``$ is called a Hamiltonian, and it is not necessarily hermitian. But, it has some properties. Conservation of probability,
$$\underset{\sigma }{}P(\sigma ,t)=1,$$
(8)
shows that
$$S|=0,$$
(9)
where
$$S|=\underset{\beta }{}\beta |.$$
(10)
So, the sum of each column of $``$, as a matrix, should be zero. As $`S|`$ is a left eigenvector of $``$ with zero eigenvalue, $``$ has at least one right eigenvector with zero eigenvalue. This state corresponds to the steady state distribution of the system and it does not evolve in time. If the zero eigenvalue is degenerate, the steady state is not unique. The transition rates are non–negative, so the off–diagonal elements of the matrix $``$ are non–negative. Therefore, if a matrix $``$ has the following properties,
$$\begin{array}{c}S|=0,\hfill \\ \sigma ||\tau 0,\hfill \end{array}$$
(11)
then it can be considered as the generator of a stochastic process. The real part of the eigenvalues of any matrix with the above conditions should be less than or equal zero.
The dynamics of the state vectors (4) is given by
$$|P(t)=\mathrm{exp}(t)|P(0),$$
(12)
and the expectation value of an observable $`𝒪`$ is
$$𝒪(t)=\underset{\sigma }{}𝒪(\sigma )P(\sigma ,t)=S|𝒪\mathrm{exp}(t)|P(0).$$
(13)
If $``$ is integrable, one can solve the problem, that is, one can calculate the expectation values. Suppose now, that a Hamiltonian is integrable but is not stochastic. There arises a question, whether or not there exist exist a similarity transformation which transforms it to a stochastic integrable Hamiltonian. Consider an integrable Hamiltonian $`\stackrel{~}{}`$. The similarity transformation
$$:=\stackrel{~}{}^1$$
(14)
leaves its eigenvalues invariant. Consider a special case: The system consists of a one dimensional lattice, with nearest–neighbor interaction,
$$\stackrel{~}{}=\underset{i=1}{\overset{L}{}}\stackrel{~}{}_{i,i+1}.$$
(15)
Suppose, also, that the system is translation–invariant:
$$\stackrel{~}{}_{ii+1}=\underset{i1}{\underset{}{1\mathrm{}1}}\stackrel{~}{H}\underset{Li1}{\underset{}{1\mathrm{}1}},$$
(16)
and we are using periodic boundary conditions. A simple class of similarity transformations is then
$$=\mathrm{\Gamma }_1\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_L.$$
(17)
The simplest case is when all $`\mathrm{\Gamma }_i`$’s are the same. In this case, if one can find $`\mathrm{\Gamma }`$ such that
$$H=\mathrm{\Gamma }\mathrm{\Gamma }\stackrel{~}{H}\mathrm{\Gamma }^1\mathrm{\Gamma }^1$$
(18)
is stochastic, then $``$ defined through (14) would be stochastic. A more general class of similarity transformations is obtained through
$$\mathrm{\Gamma }_i:=\mathrm{\Gamma }(g)^i,$$
(19)
where $`g`$ should have the property
$$[gg,\stackrel{~}{H}]=0.$$
(20)
In this case, one obtains
$$H=(\mathrm{\Gamma }\mathrm{\Gamma }g)\stackrel{~}{H}(\mathrm{\Gamma }\mathrm{\Gamma }g)^1.$$
(21)
Define $`s|`$ to be the sum of all bra–states corresponding to a single site. We then have
$$S|=\underset{L}{\underset{}{s|\mathrm{}s|}}.$$
(22)
For $`H`$ to be stochastic, its off–diagonal elements should be non–negative, and we must have
$$s|s|H=0.$$
(23)
This shows that
$$\alpha |\beta |:=s|\mathrm{\Gamma }s|\mathrm{\Gamma }g,$$
(24)
should be an eigenvector with zero eigenvalue of $`\stackrel{~}{H}`$, that is, $`\stackrel{~}{H}`$ should have a decomposable left eigenvector. So, in order that this prescription of constructing integrable stochastic model works, one must begin with a Hamiltonian $`\stackrel{~}{H}`$, the left eigenvector with zero eigenvalue of which is decomposable. The real part of all other eigenvalues of $`\stackrel{~}{H}`$ should, of course, be non–positive.
## 3 A one–parameter solvable system on the basis of a free–fermion system
Consider the Hamiltonian
$`\stackrel{~}{}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{L}{}}}\{{\displaystyle \frac{1+\eta }{2}}[s_{i+1}^+s_i^{}n_i(1n_{i+1})]`$ (27)
$`+{\displaystyle \frac{1\eta }{2}}[s_{i+1}^{}s_i^+n_{i+1}(1n_i)]`$
$`+\lambda [s_{i+1}^{}s_i^{}n_in_{i+1}]\},`$
where $`s^+,s^{},\mathrm{and}n`$ are
$$s^+:=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),s^{}:=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right),n:=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),$$
(28)
and the subscript $`i`$ represents the site, on which the operator acts. This Hamiltonian describes the following processes
$`A\mathrm{}\mathrm{}A`$ $`\text{with the rate}{\displaystyle \frac{1+\eta }{2}}`$ (29)
$`\mathrm{}AA\mathrm{}`$ $`\text{with the rate}{\displaystyle \frac{1\eta }{2}}`$ (30)
$`AA\mathrm{}\mathrm{}`$ $`\text{with the rate}\lambda .`$ (31)
This model has been recently studied. In the case $`\lambda =0`$, the above model describes an asymmetric exclusion process. For $`\lambda =1`$, the Hamiltonian is bilinear in terms of creation $`s^+`$ and annihilation $`s^{}`$ operators. This problem has been solved via a Jordan–Wigner transformation . In the notation of the previous section the matrix form of $`\stackrel{~}{H}`$ is
$$\stackrel{~}{H}:=\left(\begin{array}{cccc}\lambda & 0& 0& 0\\ 0& \frac{1+\eta }{2}& \frac{1\eta }{2}& 0\\ 0& \frac{1+\eta }{2}& \frac{1\eta }{2}& 0\\ \lambda & 0& 0& 0\end{array}\right)$$
(32)
This matrix has two eigenvalues, 0 and -1, both of them are two–folded degenerate. One of the zero left eigenvectors can be decomposed into a tensor product. Doing the above mentioned procedure, this Hamiltonian can be transformed to another stochastic one. For this case, One can show that the matrix $`g`$ is the identity matrix, and the similarity transformation for all sites become the same. This has been done in .
One of the left eigenvectors corresponding to the eigenvalue -1 has also the desired property. To use the prescription described in the previous section to construct a stochastic Hamiltonian, we define a new Hamiltonian,
$$\stackrel{~}{H}^{}:=\stackrel{~}{H}1,$$
(33)
and apply the similarity transformation on this new Hamiltonian. One of the zero left eigenvectors of $`\stackrel{~}{H}^{}`$ is
$$(1000)=(10)(10).$$
(34)
The similarity transformation should map $`s|s|`$ to this vector:
$$(11)\mathrm{\Gamma }(11)\mathrm{\Gamma }g=\alpha (10)(10)$$
(35)
So,
$$\begin{array}{c}(11)\mathrm{\Gamma }=\alpha \nu (10)\hfill \\ (11)\mathrm{\Gamma }g=\frac{\alpha }{\nu }(10).\hfill \end{array}$$
(36)
Scaling $`\mathrm{\Gamma }`$ and $`g`$ does not alter the Hamiltonian $`H`$. So we can remove $`\alpha `$ and $`\nu `$ by scaling the matrices $`\mathrm{\Gamma }`$ and $`g`$. Then the above relation gives some constraints on the elements of $`\mathrm{\Gamma }`$ and $`g`$. the condition of positivity of rates, fixes $`g`$ and $`\mathrm{\Gamma }`$:
$$\mathrm{\Gamma }=\frac{1}{2}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right)$$
(37)
$$g=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)$$
(38)
The two site Hamiltonian, then, takes the following form
$$H=\left(\begin{array}{cccc}0& \frac{1\eta }{2}& \frac{1+\eta }{2}& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& \frac{1+\eta }{2}& \frac{1\eta }{2}& 0\end{array}\right)$$
(39)
and $``$ is
$``$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{L}{}}}\{{\displaystyle \frac{1\eta }{2}}[n_is_{i+1}^++(1n_i)s_{i+1}^{})]`$ (42)
$`+{\displaystyle \frac{1+\eta }{2}}\left[s_i^+n_{i+1}+s_i^{}(1n_{i+1})\right]`$
$`[n_i(1n_{i+1})(1n_i)n_{i+1}]\}.`$
This Hamiltonian describes the following processes
$`A\mathrm{}AA`$ $`{\displaystyle \frac{1\eta }{2}}`$ (43)
$`A\mathrm{}\mathrm{}\mathrm{}`$ $`{\displaystyle \frac{1+\eta }{2}}`$ (44)
$`\mathrm{}AAA`$ $`{\displaystyle \frac{1+\eta }{2}}`$ (45)
$`\mathrm{}A\mathrm{}\mathrm{}`$ $`{\displaystyle \frac{1\eta }{2}}.`$ (46)
The Hamiltonian (42) is not quadratic in $`s^+`$ and $`s^{}`$. So, one can not map this Hamiltonian to a free fermion system, using a Jordan–Wigner transformation. But the Hamiltonaian $`\stackrel{~}{}`$ is integrable and can be mapped to a free fermion system by a Jordan–Wigner transformation. Consider the following Jordan–Wigner transformation
$`a_j`$ $`:=`$ $`Q_{j1}s_j^{}`$ (47)
$`a_j^{}`$ $`:=`$ $`Q_{j1}s_j^+`$ (48)
$`Q_j`$ $`:=`$ $`{\displaystyle \underset{i=1}{\overset{j}{}}}(s_i^3).`$ (49)
It can be easily shown that the number operator at each site $`n_i`$ is, in terms of new generators,
$$n_i:=\frac{1+s_i^3}{2}=a_i^{}a_i$$
(50)
Using this transformation, The Hamiltonian $`\stackrel{~}{}`$ takes the following form
$$\stackrel{~}{}=\underset{i=1}{\overset{L}{}}\left[\frac{1\eta }{2}a_i^{}a_{i+1}+\frac{1+\eta }{2}a_{i+1}^{}a_i+a_{i+1}a_ia_i^{}a_i\right],$$
(51)
$`a_i`$ and $`a_i^{}`$ fulfill the fermionic anti–commutation relations
$$\begin{array}{c}\{a_i,a_j\}=\{a_i^{},a_j^{}\}=0\hfill \\ \{a_i,a_j^{}\}=\delta _{ij}.\hfill \end{array}$$
(52)
Note that it is in the limit $`L\mathrm{}`$ that the Jordan–Wigner transformation we are using, works. Otherwise, there are some boundary terms in (51) as well. So, all the results we obtain hereafter, are valid only in this limit. Now, introducing the Fourier transformation
$`a_j`$ $`:=`$ $`{\displaystyle \frac{1}{\sqrt{L}}}{\displaystyle \underset{k}{}}b_k\mathrm{exp}\{{\displaystyle \frac{2\pi ijk}{L}}\}`$ (53)
$`a_j^{}`$ $`:=`$ $`{\displaystyle \frac{1}{\sqrt{L}}}{\displaystyle \underset{k}{}}b_k^{}\mathrm{exp}\{{\displaystyle \frac{2\pi ijk}{L}}\},`$ (54)
and substituting it in (52), it is seen that
$`\{b_k,b_l\}=\{b_k^{},b_l^{}\}=0`$ (55)
$`\{b_k,b_l^{}\}=\delta _{kl}.`$ (56)
As a result, the Hamiltonian $`\stackrel{~}{}`$ takes the form
$`\stackrel{~}{}`$ $`=`$ $`{\displaystyle \underset{k}{}}\left[{\displaystyle \frac{1\eta }{2}}\mathrm{exp}({\displaystyle \frac{2\pi ik}{L}})+{\displaystyle \frac{1+\eta }{2}}\mathrm{exp}({\displaystyle \frac{2\pi ik}{L}})\right]b_k^{}b_k+b_kb_k\mathrm{exp}({\displaystyle \frac{2\pi ik}{L}})b_k^{}b_k`$ (57)
$`=`$ $`{\displaystyle \underset{k}{}}\left[ϵ_kb_k^{}b_ki\mathrm{sin}({\displaystyle \frac{2\pi k}{L}})b_kb_k\right]`$ (58)
where
$$ϵ_k:=1+\mathrm{cos}(\frac{2\pi k}{L})i\eta \mathrm{sin}(\frac{2\pi k}{L})$$
(59)
One can now, easily obtain the time dependence of $`b_k`$ and $`b_k^{}`$, using (55) and $`{\displaystyle \frac{\mathrm{d}O}{\mathrm{d}t}}=[O,H]`$
$`b_k(t)=b_k(0)e^{ϵ_kt}`$ (60)
$`b_k^{}(t)=e^{ϵ_kt}\{b_k^{}(0)i\mathrm{cot}({\displaystyle \frac{\pi k}{L}})[e^{(ϵ_k+ϵ_k)t}1]b_k(0)\}`$ (61)
Now we return to our problem: determining the expectation values of a system evolving with the Hamiltonian $``$. The expectation value of a quantity $`𝒪`$ is
$$\begin{array}{cc}𝒪(t)\hfill & =S|𝒪\mathrm{exp}(t)|P(0)\hfill \\ & =S|\mathrm{exp}(t)𝒪\mathrm{exp}(t)|P(0).\hfill \end{array}$$
(62)
Substituting $`=\stackrel{~}{}^1L\mathrm{𝟏}`$, where $`\mathrm{𝟏}`$ stands for the identity matrix, yields
$$𝒪(t)=\mathrm{\Omega }|\stackrel{~}{𝒪}(\stackrel{~}{t})^1|P(0)$$
(63)
where
$$\stackrel{~}{𝒪}:=^1𝒪,$$
(64)
$$\stackrel{~}{𝒪}(\stackrel{~}{t}):=e^{t\stackrel{~}{}}𝒪e^{t\stackrel{~}{}}$$
(65)
and
$$\mathrm{\Omega }|:=\left(\begin{array}{cc}1& 0\end{array}\right)\left(\begin{array}{cc}1& 0\end{array}\right)\mathrm{}\left(\begin{array}{cc}1& 0\end{array}\right).$$
(66)
The main expectation values of interest are the correlation functions of $`n_i`$‘s. To determine these, we use
$$\begin{array}{c}\mathrm{\Gamma }^1n\mathrm{\Gamma }=\frac{1}{2}(1s^+s^{})\hfill \\ (\mathrm{\Gamma }g)^1n\mathrm{\Gamma }g=\frac{1}{2}(1+s^++s^{})\hfill \end{array}$$
(67)
So
$$^1n_i=\frac{1}{2}[1(1)^i(s^++s^{})].$$
(68)
Now, we want to calculate the expectation value of $`𝒪`$ where
$$𝒪:=n_{i_m}\mathrm{}n_{i_2}n_{i_1},i_1i_2\mathrm{}i_m.$$
(69)
We have
$$\stackrel{~}{𝒪}=\frac{1}{2^m}\left[1(1)^{i_m}(s_{i_m}^++s_{i_m}^{})\right]\mathrm{}\left[1(1)^{i_1}(s_{i_1}^++s_{i_1}^{})\right].$$
(70)
Using the Jordan–Wigner transformation, one arrives at
$$\stackrel{~}{𝒪}=\frac{1}{2^m}\left[1(1)^{i_m}Q_{i_m1}(a_{i_m}^{}+a_{i_m})\right]\mathrm{}\left[1(1)^{i_1}Q_{i_11}(a_{i_1}^{}+a_{i_1})\right].$$
(71)
It is easy to check that $`\mathrm{\Omega }|Q_i=(1)^i\mathrm{\Omega }|`$. So in calculating $`𝒪`$, one can use $`𝒪^{}`$ instead of $`\stackrel{~}{𝒪}`$:
$$𝒪^{}:=\frac{1}{2^m}\left[1+(a_{i_m}^{}+a_{i_m})\right]\mathrm{}\left[1+(a_{i_1}^{}+a_{i_1})\right].$$
(72)
Instead of $`𝒪^{}`$, It is enough to set $`𝒪^{\prime \prime }`$ in the expectation value of $`𝒪`$, where
$$𝒪^{\prime \prime }:=\frac{1}{2^m}\left(1+a_{i_m}^{}\right)\mathrm{}\left(1+a_{i_1}^{}\right).$$
(73)
To prove this, one should use $`\mathrm{\Omega }|\stackrel{~}{}=L\mathrm{\Omega }|`$ and $`\mathrm{\Omega }|a_i(0)=0`$.
## 4 The one–point function
As the first example, consider the one–point function $`n_m(t)`$:
$$n_m(t)=\frac{1}{2}\mathrm{\Omega }|[1+a_m^{}(\stackrel{~}{t})]^1|P(0).$$
(74)
Using the Fourier transformation (53), the time dependence of $`b_k^{}`$ (60), and remembering $`\mathrm{\Omega }|b_k(0)=0`$, we obtain
$$n_m(t)=\frac{1}{2}+\frac{1}{2\sqrt{L}}\underset{k}{}e^{\frac{2\pi ikm}{L}}\mathrm{\Omega }|b_k^{}(0)^1|P(0)e^{ϵ_kt}$$
(75)
Now, we use the inverse Fourier–, and Jordan–Wigner–transformations, and arrive at
$$n_m(t)=\frac{1}{2}+\frac{1}{2L}\underset{k,j}{}e^{\frac{2\pi ik(jm)}{L}}S|s_j^+^1|P(0)e^{ϵ_kt}(1)^{j1}.$$
(76)
This can be written in a simpler form, using
$$s_j^+^1=(1)^{j1}\frac{2n_j1+s_j^{}s_j^+}{2}$$
(77)
and
$$s|(2n1)=s|(s^{}s^+).$$
(78)
One then arrives at
$$n_m(t)=\frac{1}{2}+\frac{1}{2L}\underset{k,j}{}e^{\frac{2\pi ik(jm)}{L}}S|(2n_j(0)1)|P(0)e^{ϵ_kt}.$$
(79)
Using $`S|P(0)=_\sigma P(\sigma ,0)=1`$, one arrives at
$$n_m(t)=\underset{j}{}\mathrm{\Lambda }_{mj}(t)n_j(0),$$
(80)
where
$$\mathrm{\Lambda }_{mj}(t):=\frac{1}{L}\underset{k}{}e^{\frac{2\pi ik(jm)}{L}}e^{ϵ_kt}.$$
(81)
Now, consider the limit $`t\mathrm{}`$. In this limit, the only contribution in the above summation comes from the term $`k=0`$. So,
$$\underset{t\mathrm{}}{lim}n_m(t)=\frac{1}{L}\underset{j}{}n_j(0),$$
(82)
which shows that in the limit $`t\mathrm{}`$, the expectation value of the number of particles in any site tends to the average of the initial value of this quantity. In the last section, we will find the next leading term of $`n_i(t)`$, for large times.
Now, we want to calculate the expectation value of the number of particles in the site $`j`$ in the limit $`L\mathrm{}`$. First, we calculate $`\mathrm{\Lambda }_{mj}`$ in this limit. To do so, we define $`z:=\mathrm{exp}(i\frac{2\pi k}{L})`$. We then (in this limit) arrive at
$$\mathrm{\Lambda }_{mj}(t)=e^t\frac{\mathrm{d}z}{2\pi iz}z^{jm}\mathrm{exp}[t(\frac{1\eta }{2}z+\frac{1+\eta }{2}z^1)].$$
(83)
Changing the variable $`z`$ to $`w:=z\sqrt{\frac{1+\eta }{1\eta }}`$, the above integral takes the form
$$\mathrm{\Lambda }_{mj}(t)=e^t\left(\frac{1\eta }{1+\eta }\right)^{\frac{jm}{2}}\frac{\mathrm{d}w}{2\pi iw}w^{jm}\mathrm{exp}[t\sqrt{\frac{1\eta ^2}{2}}(w+w^1)],$$
(84)
or, using the change of variable $`w:=e^{i\theta }`$,
$$\mathrm{\Lambda }_{mj}(t)=e^t\left(\frac{1\eta }{1+\eta }\right)^{\frac{mj}{2}}_0^{2\pi }\frac{\mathrm{d}\theta }{2\pi }e^{i(jm)\theta +t\sqrt{1\eta ^2}\mathrm{cos}\theta }.$$
(85)
The above integral is an integral representation of the modified Bessel function:
$$\mathrm{\Lambda }_{mj}(t)=\left(\frac{1\eta }{1+\eta }\right)^{\frac{mj}{2}}\mathrm{I}_{mj}(t\sqrt{1\eta ^2})e^t.$$
(86)
$`n_m(t)`$, in the limit $`L\mathrm{}`$, is then
$$n_m(t)=\underset{j}{}\left(\frac{1\eta }{1+\eta }\right)^{\frac{mj}{2}}\mathrm{I}_{mj}(t\sqrt{1\eta ^2})e^tn_j(0).$$
(87)
## 5 The two–point function
The other quantity which we want to calculate is $`n_m(t)n_l(t)`$. Without loss of generality, one may assume $`(m>l)`$. To calculate this, we use (73), which gives
$`n_m(t)n_l(t)={\displaystyle \frac{1}{4}}\mathrm{\Omega }|[1+a_m^{}(\stackrel{~}{t})][1+a_l^{}(\stackrel{~}{t})]^1|P(0)`$ (88)
$`={\displaystyle \frac{1}{4}}+{\displaystyle \frac{1}{2}}[n_m(t)+n_l(t)]+{\displaystyle \frac{1}{4}}\mathrm{\Omega }|a_m^{}(\stackrel{~}{t})a_l^{}(\stackrel{~}{t})^1|P(0).`$ (89)
The main thing is to calculate the last term. To do this, we first use the Fourier transformation of $`a_i^{}`$’s,
$$\mathrm{\Omega }|a_m^{}(\stackrel{~}{t})a_l^{}(\stackrel{~}{t}))^1|P(0)=\frac{1}{L}_{k,p}e^{i2\pi \frac{(km+pl)}{L}}\mathrm{\Omega }|b_k^{}(\stackrel{~}{t})b_p^{}(\stackrel{~}{t})^1|P(0),$$
(90)
and then substitute the time dependence of $`b_k^{}`$s.
$`\mathrm{\Omega }|a_m^{}(\stackrel{~}{t})a_l^{}(\stackrel{~}{t}))^1|P(0)`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{k,p}{}}e^{i2\pi \frac{(km+pl)}{L}+(ϵ_k+ϵ_p)t}\mathrm{\Omega }|b_k^{}(0)`$ (92)
$`\left[b_p^{}(0)+i\mathrm{cot}({\displaystyle \frac{\pi p}{L}})\left(1e^{(ϵ_p+ϵ_p)t}\right)b_p(0)\right]^1|P(0).`$
Now we use inverse Fourier transformation for the $`b_k^{}b_p^{}`$ term. The other term is easily summed. We arrive at,
$`\mathrm{\Omega }|a_m^{}(\stackrel{~}{t})a_l^{}(\stackrel{~}{t})^1|P(0)`$ $`=`$ $`{\displaystyle \underset{r,s}{}}\mathrm{\Lambda }_{mr}(t)\mathrm{\Lambda }_{ls}(t)\mathrm{\Omega }|a_r^{}a_s^{}^1|P(0)`$ (94)
$`+{\displaystyle \frac{i}{L}}{\displaystyle \underset{k}{}}e^{i2\pi \frac{k(lm)}{L}}\mathrm{cot}({\displaystyle \frac{\pi k}{L}})\left(1e^{(ϵ_k+ϵ_k)t}\right),`$
or,
$`\mathrm{\Omega }|a_m^{}(\stackrel{~}{t})a_l^{}(\stackrel{~}{t})^1|P(0)`$ $`=`$ $`{\displaystyle \underset{k,p,r,s}{}}\mathrm{\Lambda }_{mr}(t)\mathrm{\Lambda }_{ls}(t)(2n_r1)(2n_s1)_0\mathrm{sgn}(rs)`$ (96)
$`+{\displaystyle \frac{i}{L}}{\displaystyle \underset{k}{}}e^{i2\pi \frac{k(lm)}{L}}\mathrm{cot}({\displaystyle \frac{\pi k}{L}})\left(1e^{(ϵ_k+ϵ_k)t}\right),`$
where we have used the definition of $`\mathrm{\Lambda }_{ij}`$, and $`\mathrm{}_0`$ means the expectation value at the initial time. Adding all terms in (88) together, one arrives at
$`n_m(t)n_l(t)`$ $`=`$ $`{\displaystyle \frac{1}{4}}(n_m(t)+n_l(t))+{\displaystyle \underset{r,s}{}}\mathrm{\Lambda }_{mr}(t)\mathrm{\Lambda }_{ls}(t)(n_r{\displaystyle \frac{1}{2}})(n_s{\displaystyle \frac{1}{2}})_0\mathrm{sgn}(rs)`$ (98)
$`+{\displaystyle \frac{i}{4L}}{\displaystyle \underset{k}{}}e^{i2\pi \frac{k(lm)}{L}}\mathrm{cot}({\displaystyle \frac{\pi k}{L}})\left(1e^{(ϵ_k+ϵ_k)t}\right).`$
The last term is independent of initial conditions, So we can calculate it for a special case, e.g. $`|P(0)=|0`$. Then, the final result is
$$n_m(t)n_l(t)=\frac{1}{2}(n_m(t)+n_l(t))+\underset{r,s}{}\mathrm{\Lambda }_{mr}(t)\mathrm{\Lambda }_{ls}(t)\mathrm{sgn}(rs)n_rn_s\frac{n_r+n_s}{2}_0$$
(99)
For large times, it is seen that
$$\underset{t\mathrm{}}{lim}n_m(t)n_l(t)=n(\mathrm{})$$
(100)
## 6 Null vectors of the Hamiltonian, the steady state of the system, and the $`n`$–point function
Now we want to study the null eigenvectors of the Hamiltonian $``$. It is easy to see that the Hamiltonian (42) has at least two null eigenvectors, which means that the steady state is not unique. These states are one in which all sites are occupied, and one in which no site is occupied. One can check this easily by acting the Hamiltonian (42) on these states. It was shown that the Hamiltonian $`\stackrel{~}{}`$ may be written as
$`\stackrel{~}{}`$ $`=`$ $`{\displaystyle \underset{k}{}}\left[ϵ_kb_k^{}b_ki\mathrm{sin}({\displaystyle \frac{2\pi k}{L}})b_kb_k\right]`$ (101)
$`=`$ $`{\displaystyle \underset{k>0}{}}\left[ϵ_kb_k^{}b_k+ϵ_kb_k^{}b_k2i\mathrm{sin}({\displaystyle \frac{2\pi k}{L}})b_kb_k\right]+ϵ_0b_0^{}b_0.`$ (102)
This Hamiltonian is obviously block diagonal. In each four dimensional block, one can choose a basis $`\{|0,b_k^{}|0,b_k^{}|0,b_k^{}b_k^{}|0\}`$. The eigenvalues of this four dimensional block are $`0,ϵ_k,ϵ_k,ϵ_k+ϵ_k`$, or $`𝒩_kϵ_k+𝒩_kϵ_k`$, where $`𝒩`$’s are zero or one. The eigenvalues of the Hamiltonian are, therefore,
$$\stackrel{~}{E}\{𝒩\}=\underset{k}{}𝒩_kϵ_k.$$
(103)
From this, one can obtain the eigenvalues of $``$ as
$`E\{𝒩\}`$ $`=`$ $`{\displaystyle \underset{k}{}}(𝒩_kϵ_k+1)`$ (104)
$`=`$ $`{\displaystyle \underset{k}{}}(𝒩_k+1)ϵ_k.`$ (105)
Here we have used
$$\underset{k}{}(1+ϵ_k)=0.$$
(106)
Now, it is easy to see that $`E`$ is zero iff $`1𝒩_k=0,k0`$. This shows that the null eigenvector is two–fold degenerate. As the final state is two–fold degenerate, and it is known that the totally full state $`(|\mathrm{\Omega }:=\left(\begin{array}{c}1\\ 0\end{array}\right)\left(\begin{array}{c}1\\ 0\end{array}\right)\mathrm{}\left(\begin{array}{c}1\\ 0\end{array}\right))`$ and totally empty state $`(|0:=\left(\begin{array}{c}0\\ 1\end{array}\right)\left(\begin{array}{c}0\\ 1\end{array}\right)\mathrm{}\left(\begin{array}{c}0\\ 1\end{array}\right))`$ are null eigenvectors of the system, we have
$$|P(\mathrm{})=\alpha |0+\beta |\mathrm{\Omega },$$
(107)
where $`\alpha +\beta =1`$. Using (82), and
$`S|n_i|0=0,`$ (108)
$`S|n_i|\mathrm{\Omega }=1,`$ (109)
it is seen that
$$\beta =n(\mathrm{})=\frac{1}{L}\underset{j}{}n_j(0)=:\rho _0$$
(110)
and
$$\alpha =1\rho _0.$$
(111)
From this, one obtains
$$|P(\mathrm{})=[1\rho _0]|0+\rho _0|\mathrm{\Omega }.$$
(112)
Using this, it is easy to find all $`m`$-point functions in the limit $`t\mathrm{}`$. We have
$`n_{i_m}\mathrm{}n_{i_1}(\mathrm{})`$ $`=`$ $`S|n_{i_m}\mathrm{}n_{i_1}|P(\mathrm{})`$ (113)
$`=`$ $`\rho _0`$ (114)
$`=`$ $`n(\mathrm{})`$ (115)
$``$ $`n(\mathrm{})^m`$ (116)
This clearly shows that the mean–field approximation does not work here.
## 7 Relaxation of the system toward its steady state
It was shown that in the limit $`t\mathrm{}`$, the expectation of the number of particles in any site tends to the average of the initial value of this quantity. Now, we want to study the behaviour of the system at large times. Starting from (87), and representing $`n_j(0)`$ by its Fourier transform, we have
$$n_m(t)=e^t_0^{2\pi }\frac{\mathrm{d}u}{2\pi }\underset{j}{}\left(\frac{1\eta }{1+\eta }\right)\mathrm{I}_{mj}(t\sqrt{1\eta ^2})e^{iuj}[f(u)+2\pi \overline{n}\delta (u)],$$
(117)
where $`[f(u)+2\pi \overline{n}\delta (u)]`$ is the Fourier transform of $`n_j(0)`$, and $`\overline{n}`$ is the average density. We have extracted this part of the Fourier transform, so that the remaining is a smooth function of $`u`$. Then, $`f(u)`$ denotes the Fourier transform of the deviation $`n_j(0)\overline{n}`$. The summation on $`j`$ is easily done, using
$$\underset{n}{}x^n\mathrm{I}_n(y)=e^{(y/2)(x+1/x)}.$$
(118)
So, one arrives at
$$n_m(t)=\overline{n}+_0^{2\pi }\frac{\mathrm{d}u}{2\pi }e^{tϵ(u)}f(u)e^{imu},$$
(119)
where $`ϵ(u)`$ is the same as $`ϵ_k`$ with $`k=Lu/(2\pi )`$. The above integral is simplified for large times, using the steepest descent method. Using the change of variable $`z:=e^{iu}`$, the integral becomes
$$n_m(t)=\overline{n}+\frac{\mathrm{d}z}{2\pi iz}e^{t[2+(1\eta )z+(1+\eta )z^1]/2}\stackrel{~}{f}(z)z^m,$$
(120)
where the integration contour is the unit circle. The multiplier of $`t`$ in the exponent is stationary at
$$z_1=\sqrt{\frac{1+\eta }{1\eta }},$$
(121)
and
$$z_2=\sqrt{\frac{1+\eta }{1\eta }}.$$
(122)
As the real part of this multiplier is larger at $`z=z_1`$, the integral gets its main contribution from this point. (This point is not on the integration contour. But, assuming $`\stackrel{~}{f}`$ to be analytic, one deforms the integration contour so that it passes from $`z_1`$, and then uses the steepest descent method.) We arrive at
$$n_m(t)\overline{n}\frac{1}{\sqrt{t}}\left(\frac{1\eta }{1+\eta }\right)^{m/2}e^{t(\sqrt{1\eta ^2}1)}.$$
(123)
The effect of the Fourier transform $`\stackrel{~}{f}`$, and the second derivative of the multiplier of $`t`$ in the exponent is a multiplier independent of $`m`$ and $`t`$. Two general features, independent of the initial condition, are seen from the above relation. First, the decay to the final state is not in the form of a power law, but in the form of an exponential. It becomes a power law only in the symmetric case $`\eta =0`$. Second, if $`\eta >0`$, the expectation at the rightmost sites tends rapidly to its final value. That is, the profile of the deviation from the final value is decreasing with respect to $`m`$. This is so, since in this case the two–site reaction is favorable to the state where the left site changes so that it becomes identical to the right site. This means that cases where the right–site changes is less probable than cases where the left site changes. So, the right site arrives earlier to its final state. This expression seems to be unbounded for either $`m\mathrm{}`$ or $`m\mathrm{}`$. For any fixed $`t`$, this is true. But it simply means that in order that this term represents the leading term for some $`m`$, $`t`$ must be greater than some $`T`$, which does depend on $`m`$.
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# The impact of LEP 2 data on possible anomalous enhancements of 𝑅𝑒ϵ'/ϵ
## Abstract
It has been shown in the past that the real part of the $`ϵ^{}/ϵ`$ratio is particularly sensitive to anomalous gauge couplings that modify the Standard Model Lagrangian. Due to the loose bounds on these couplings coming from low energy processes and to the poor sensitivity of hadron colliders to couplings such as $`\mathrm{\Delta }g_1^Z`$, it has been argued that anomalous couplings could still produce an enhancement of $`Reϵ^{}/ϵ`$ bringing this observable closer to the experimental value obtained by KTeV, NA31 and NA48. The impact of the new measurements done at LEP 2 in these years is discussed and new severe constraints to this hypothesis are determined.
xxx
Since 1995, it has been noted that anomalous triple gauge couplings (TGC) involving $`WW\gamma `$ and $`WWZ`$ vertices could modify significantly the Standard Model (SM) prediction concerning $`ϵ^{}/ϵ`$. More precisely, strong penguin diagrams and isospin breaking due to quark masses completely dominate the SM prediction of $`Reϵ^{}/ϵ`$for low values of $`m_t`$. However, for $`m_t`$ of the order of 170 GeV, the effects of electroweak penguin diagrams are sizeable and therefore $`Reϵ^{}/ϵ`$becomes sensitive to anomalies in the bosonic sector of the SM.
The possibility that anomalous TGC could modify and, particularly, enhance $`Reϵ^{}/ϵ`$has been considered with interest, especially because several experimental measurements still point towards a rather high value of $`Reϵ^{}/ϵ`$with respect to the SM expectation. In fact, present theoretical errors on $`Reϵ^{}/ϵ`$, coming mainly from the uncertainties in the models for hadronic matrix elements, prevent us from drawing conclusions about the effectiveness of SM in the CP violating sector. On the other hand, it is interesting to quantify to what extent the hypothesis of anomalous enhancement of $`Reϵ^{}/ϵ`$is corroborated by the intense experimental investigation on the bosonic sector of the SM carried out at LEP 2 since 1996.
It is customary to express general deviations from the SM in the bosonic sector in the framework of effective theories . In this case, the most general Lagrangian invariant under $`U_{em}(1)`$ that contributes to $`WWV`$ vertices ($`V=Z,\gamma `$) is:
$$\begin{array}{cc}i_{eff}^{WWV}=\hfill & g_{WWV}[g_1^VV^\mu (W_{\mu \nu }^{}W^{+\nu }W_{\mu \nu }^+W^\nu )+\kappa _VW_\mu ^+W_\nu ^{}V^{\mu \nu }+\hfill \\ & +\frac{\lambda _V}{M_W^2}V^{\mu \nu }W_\nu ^{+\rho }W_{\rho \nu }^{}+ig_5^V\epsilon _{\mu \nu \rho \sigma }((^\rho W^\mu )W^{+\nu }W^\mu (^\rho W^{+\nu }))V^\sigma +\hfill \\ & +ig_4^VW_\mu ^{}W_\nu ^+(^\mu V^\nu +^\nu V^\mu )\frac{\stackrel{~}{\kappa }_V}{2}W_\mu ^{}W_\nu ^+\epsilon ^{\mu \nu \rho \sigma }V_{\rho \sigma }\frac{\stackrel{~}{\lambda }_V}{2M_W^2}W_{\rho \mu }^{}W_\nu ^{+\mu }\epsilon ^{\nu \rho \alpha \beta }V_{\alpha \beta }],\hfill \end{array}$$
(1)
where $`W^{\pm \mu }`$ are the $`W`$ boson fields and $`V=\gamma ,Z`$. Defining $`g_{WW\gamma }=e`$ and $`g_{WWZ}=ecot\theta _W`$, this Lagrangian allows anomalous values for the C- and P-conserving couplings $`\kappa _\gamma `$, $`\kappa _Z`$, $`g_1^Z`$, $`g_1^\gamma `$ (equal to 1 in SM). Moreover, new contributions coming from operators absent in the standard theory are present. These are $`\lambda _V`$, which also conserves both $`C`$ and $`P`$-parity; $`g_5^V`$ ($`C`$ and $`P`$ violating but $`CP`$ conserving) and the $`CP`$ violating terms $`g_4^V`$, $`\stackrel{~}{\kappa }_V`$ and $`\stackrel{~}{\lambda }_V`$. In terms of the $`W`$ magnetic dipole and electric quadrupole they contribute as
$`\mu _W=(g_1^\gamma +\kappa _\gamma +\lambda _\gamma ){\displaystyle \frac{e}{2m_W}}s;`$ $`Q_W={\displaystyle \frac{e(\kappa _\gamma \lambda _\gamma )}{m_W^2}},`$ (2)
where $`eg_\gamma ^1`$ is the $`W`$ charge and $`s`$ its spin. To simplify the notations, it is convenient to introduce the deviation from the SM couplings
$`\mathrm{\Delta }g_1^Vg_1^V1,\mathrm{\Delta }\kappa _V\kappa _V1.`$
CP-violating terms are constrained by the neutron and electron electric dipole moments . CP-conserving anomalous couplings affect low-energy rare decays , $`W`$ production processes at high scales (LEP 2 and Tevatron) and electroweak corrections to the $`W`$,$`Z`$, $`\gamma `$ propagator (“oblique corrections”). Oblique corrections have been extensively tested at the scale of $`m_Z`$ at LEP 1 and SLC. In order to evade the tight constraints already obtained from LEP at the $`Z^0`$-pole, present LEP 2 measurements are expressed as limits to couplings contributing to the effective Lagrangian (1) after imposing $`SU(2)U(1)`$ gauge invariance and retaining only the lowest dimension operators . This approach implies relations amongst the various TGC and reduce the number of independent couplings to three: $`\mathrm{\Delta }g_1^Z`$, $`\lambda _\gamma `$and $`\mathrm{\Delta }\kappa _\gamma `$.
The CP-violating $`\mathrm{\Delta }S=1`$ interaction responsible for $`K\pi \pi `$ is affected, beyond strong penguins, by electroweak penguins and (possibly) TGC contributions changing both the $`I=0`$ and $`I=2`$ amplitude. At the scale of $`\mu =m_W`$, the SM effective Hamiltonian ($`H_{eff}`$) for $`\mathrm{\Delta }S`$=1, modified in order to cope with anomalous TGC, has been computed in . It has been shown that the CP-violating couplings and $`\mathrm{\Delta }\kappa _Z`$ do not contribute to leading order, being suppressed by factors of $`O((m_{d,s}^2,m_K^2)/m_W^2)`$. Running $`H_{eff}`$ at the scale of $`\mu =1`$ GeV implies the knowledge of the boundary conditions of the Wilson coefficients in SM and in the occurrence of anomalous couplings . In the calculation of a cut-off $`\mathrm{\Lambda }`$ of 1 TeV has been used for terms proportional to $`\mathrm{\Delta }g_1^Z`$and $`\mathrm{\Delta }\kappa _\gamma `$. The results obtained in depend, in general, on the anomalous couplings $`\mathrm{\Delta }g_1^Z`$,$`\mathrm{\Delta }\kappa _\gamma `$, $`\lambda _\gamma `$and $`g_5^Z`$, on the imaginary part of $`V_{td}V_{ts}^{}`$ and on the assumptions about the hadronic matrix elements.
The dependence of $`Reϵ^{}/ϵ`$to anomalous TGC can be expressed in the following way:
$$Re\left(\frac{ϵ^{}}{ϵ}\right)Im(V_{td}V_{ts}^{})(1+0.96\mathrm{\Delta }\kappa _\gamma +0.16\lambda _\gamma 4.08\mathrm{\Delta }g_1^Z+0.44g_5^Z).$$
(4)
where the overall normalisation factor $``$ depends mainly on the knowledge of the hadronic matrix elements and, according to the calculation of , $`8.66`$. Due to the clear dominance of the term involving $`\mathrm{\Delta }g_1^Z`$, Tevatron results do not modify significantly the low-energy constraints and are not considered here. In the present analysis, we take the current allowed range of $`Im(V_{td}V_{ts}^{})`$ from and we assume
$$Im(V_{td}V_{ts}^{})=(1.14\pm 0.11)10^4$$
(5)
It is well known that the status of the experimental measurement of $`Reϵ^{}/ϵ`$is not completely satisfactory. First results were obtained by NA31 and were not confirmed by E731 . Preliminary results from the KTeV collaboration and from NA48 strongly support a non-zero value of $`Reϵ^{}/ϵ`$(fig.1). On the theoretical side, current calculations of $`Reϵ^{}/ϵ`$in SM point, in general, towards lower values than those suggested by current experiments, even if the statistical significance is of the order of 2-2.5$`\sigma `$ (for a review see ) and make the hypothesis of an anomalous enhancement of $`Reϵ^{}/ϵ`$due to TGC rather attractive. It is interesting to note, however, that some very recent updates of the theoretical calculations of the hadronic matrix elements could imply a higher value of $`Reϵ^{}/ϵ`$which, if confirmed, would ease the agreement of the SM predictions with the current experimental measurements without invoking new physics. On the other hand, for what concerns TGC the hadronic matrix elements mainly affects the overall normalisation factor $``$ being the relative contributions of the couplings (last factor of eq.(4) ) practically independent to them. Therefore, the allowed ranges derived in the following assuming $`8.66`$ can be updated in a straightforward manner by proper rescaling of $``$ in eq.(4).
Figure 2 shows the possible enhancement of $`Reϵ^{}/ϵ`$as a function of $`\mathrm{\Delta }g_1^Z`$, without including current LEP 2 data and assuming all TGC but $`\mathrm{\Delta }g_1^Z`$at their SM values. The allowed range of $`\mathrm{\Delta }g_1^Z`$has been computed including just the low energy constraints coming from rare $`B`$ and $`K`$ decay, as in . The vertical width of the dark band represents the variation of $`Reϵ^{}/ϵ`$corresponding to a change of $`\pm 2\sigma `$ of $`Im(V_{td}V_{ts}^{})`$. The light band indicates the allowed range from NA31. All the limits are computed at 95% C.L. The plot represents approximately the experimental situation at the beginning of the high energy run of LEP.
The current situation is depicted in figure 3. Here the world average has been computed inflating the error until $`\chi ^2/ndf=1`$, in order to deal with the discrepancy between the old E731 measurement and the analyses of KTeV, NA31 and NA48. Simple rejection of E731 results in a world average of $`(21.6\pm 2.7)10^4`$. The allowed range for $`\mathrm{\Delta }g_1^Z`$combines the 95% C.L. limits coming from ALEPH, DELPHI, OPAL and L3 summing up the accumulated statistics from 1996 to 1999 taken at centre-of-mass energies ranging from 161 GeV to 202 GeV. Here, statistical and systematic uncertainties are included.
The remarkable improvement in the experimental tests of anomalous values of $`\mathrm{\Delta }g_1^Z`$strongly constraints the hypothesis of anomalous enhancement, bringing the allowed range of $`Reϵ^{}/ϵ`$to
$$\mathrm{8.0\; 10}^4<Reϵ^{}/ϵ<\mathrm{13.7\; 10}^4(95\%C.L.),$$
(6)
to be compared with the pre-LEP 2 measurement of
$$\mathrm{5.4\; 10}^4<Reϵ^{}/ϵ<\mathrm{30.5\; 10}^4(95\%C.L.).$$
(7)
These limits have been extracted assuming that only $`\mathrm{\Delta }g_1^Z`$is different from zero. In fact, single $`W`$ production at LEP 2 strongly constraints contributions to the effective Lagrangian coming from operators proportional to $`\mathrm{\Delta }\kappa _\gamma `$. The accumulated LEP 2 statistics allows a simultaneous fit of $`\mathrm{\Delta }g_1^Z`$and $`\mathrm{\Delta }\kappa _\gamma `$(again, all other couplings have been fixed to their SM value). The corresponding allowed range for $`Reϵ^{}/ϵ`$is
$$\mathrm{7.9\; 10}^4<Reϵ^{}/ϵ<\mathrm{15.6\; 10}^495\%C.L.$$
(8)
Contributions from other couplings have very limited impact on $`Reϵ^{}/ϵ`$by virtue of eq.(4). It has to be noted that results for the TGC parameters have been quoted without use of a form factor $`\mathrm{\Lambda }`$, describing the scale at which new physics should become manifest. Inclusion of such a parameter with $`\mathrm{\Lambda }=1`$ TeV, for sake of consistency with , and form factors of the type used, for example, in would increase the limits obtained for $`Reϵ^{}/ϵ`$by no more than 4%.
In conclusion, we have shown that present LEP 2 data highly constrain the hypothesis of a non-standard enhancement of $`Reϵ^{}/ϵ`$coming from anomalies in the bosonic sector of the SM. In particular, the sensitivity of the process $`e^+e^{}W^+W^{}`$ to $`\mathrm{\Delta }g_1^Z`$allows a reduction of the possible contribution to $`Reϵ^{}/ϵ`$from this coupling by a factor of about 4.
Acknowledgements
I’m greatly indebted to X.-G. He, C. Matteuzzi and R. Sekulin for useful discussions on this subject. Interesting suggestions of R. Felici and V. Verzi are gratefully acknowledged.
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# 1 Introduction
## 1 Introduction
$`A`$ formulation of statistical mechanics on the basis of Shannon’s information theory <sup>1)</sup> has been proposed by Jaynes in 1957.<sup>2)</sup> This formulation utilizes the least biased statistical inference about a physical system on the basis of a limited amount of information available.
In the present note we discuss quantum statistical mechanics,<sup>3-9)</sup> which incorporates Jaynes’s proposal, and we examine what kind of picture appears if one distinguishes the statistical aspects and dynamical aspects of quantum statistical mechanics to the maximum extent. We work exclusively on quantum statistcal mechanics, though in the course of our discussion we comment on the recent progress in the Boltzmann approach to classical statistical mechanics .<sup>10-12)</sup>
We would like to briefly summarize the basic aspects of our analysis. We examine a quantum mechanical mixed state, which is slightly away from thermal equilibrium, in the representation where the total energy and particle numbers are diagonal. The average value of any macroscopic observable $`\widehat{O}`$ is given by $`\widehat{O}(t)=Tr\widehat{O}\widehat{\rho }(t)`$. For the system slightly away from equilibrium, $`\widehat{O}(t)`$ is time dependent in general, and thus the density matrix $`\widehat{\rho }(t)`$ is time dependent. Following von Neumann, we introduce the (dynamical) entropy defined by
$$H=kTr\widehat{\rho }\mathrm{ln}\widehat{\rho }$$
(1.1)
which is a generalization of Gibbs entropy to a quantum system. It is well known that $`H`$ thus defined is constant in time.
The basic observation in this note is that we can define another quantity
$$S=k\underset{n}{}p_n\mathrm{ln}p_n$$
(1.2)
in terms of $`p_nn|\widehat{\rho }|n`$ in the representation where the total energy and particle numbers are diagonal. By definiton we have $`p_n(t)=p_n(0)`$ and thus $`S`$ is also constant in time. We propose to identify $`S`$ thus defined, which can be regarded as basically statistical quantity, as the amount of uncertainty of Shannon which was introduced into statistical mechanics by Jaynes.<sup>2)</sup> Note that $`H`$ in (1.1) agrees with $`S`$ in (1.2) only when $`\widehat{\rho }`$ is diagonalizable simultaneously with the total Hamiltonian and the total particle number operator. This clear distinction between $`H`$ in (1.1) and $`S`$ in (1.2), to our knowledge, has not been made in the past.<sup>13)</sup>
In the present formulation of quantum statistical mechanics, we regard the least biased estimate on the basis of limited amount of information discussed by Jaynes as a least biased estimate of initial conditions on the diagonal elements $`p_n`$. Note that the diagonal elements of $`\widehat{\rho }(t)`$ are constant in time and thus they cannot approach equilibrium values by any dynamical motion but by statistical inference. (Our initial conditions are thus final conditions also.) The initial state is thus specified by $`\{p_n\}`$ and a set of macroscopic observables other than the total energy and particle number. From the days of Boltzmann, it is well known that the second law of thermodynamics is mechanically represented only by means of assumptions regarding initial conditions,<sup>12)</sup> or by choosing the typical states in the classical Boltzmann analysis.<sup>11)</sup> In quantum statistical mechanics, as formulated here, it is clear that Shannon’s statistical inference does not precede physical dynamics.
As for the main problem of quantum statistical mechanics as to what kind of dynamical properties of the system ensure the approach to eventual thermal equilibrium, we propose to analyze a set of macroscopic observables, each of which is characterized by an intrinsic time scale $`\tau `$. As for the entropy law, some form of coarse graining is necessary not only in the Gibbs approach to quantum statistical approach<sup>8)</sup> but also in the Boltzmann approach to classical statistical mechanics.<sup>10)</sup> In our approach it turns out to be more convenient to take a coarse graining in the “time direction” or a suitable time averaging.<sup>14,9)</sup> The entropy law of Clausius in the present formulation is expressed as the approach of the macroscopic observables of the system on a suitable time average to those of the almost equilibrium state, whose physical entropy is estimated by the maximum value of Shannon’s $`S`$.
## 2 Shannon’s Least Biased Inference
We consider the variable $`x`$ which takes the $`n`$ values $`\{x_1,x_2,\mathrm{}.,x_n\}`$ and define the probability $`p_i`$ for the variable $`x`$ to assume the value $`x_i`$. The non-negative probability $`p_i`$ is constrained by the condition
$$\underset{i=1}{\overset{n}{}}p_i=1,$$
(2.1)
which means that the total probability is unity. We also consider a smooth function $`f(x)`$ of the variable $`x`$, such as $`f(x)=x`$.
We then ask what we can say about the set of probabilities $`\{p_1,p_2,,\mathrm{}.,p_n\}`$, if only available information is the average value $`<f>`$ of $`f(x)`$ defined by
$$<f>\underset{i=1}{\overset{n}{}}p_if(x_i).$$
(2.2)
Clearly it is impossible to determine all $`p_i`$ uniquely for a large value of $`n`$ since we know only $`<f>`$.
Shannon introduced the notion of amount of uncertainty $`S(p_1,p_2,\mathrm{}.,p_n)`$ for the set of variables $`\{p_1,p_2,\mathrm{}.,p_n\}`$, and he proposed to determine each $`p_i`$ by allowing the maximum amount of uncertainty, or equivalently, the least bias for the chosen solution of $`\{p_1,p_2,\mathrm{}.,p_n\}`$. On the basis of a composition law, Shannon derived the amount of uncertainty<sup>1)</sup> ( see also Appendix in ref. 2)
$$S(p_1,p_2,\mathrm{}.,p_n)=k\underset{i=1}{\overset{n}{}}p_i\mathrm{ln}p_i$$
(2.3)
with a positive constant $`k`$.
We now apply the above theory of inference to statistical mechanics. Consider a closed system for which one knows that the total energy is confined within a small range
$$E\frac{1}{2}\mathrm{\Delta }EE_nE+\frac{1}{2}\mathrm{\Delta }E$$
(2.4)
with the constraint
$$\underset{n}{}E_np_n=E$$
(2.5)
for sufficiently small $`\mathrm{\Delta }E`$ (with a fixed particle number $`N`$ and a fixed volume). The maximum of the Shannon’s amount of uncertainty then gives rise to the probability for $`E_n`$
$`p_n=\mathrm{exp}[\beta E_n]/Z,`$
$`Z={\displaystyle \underset{E\frac{1}{2}\mathrm{\Delta }EE_nE+\frac{1}{2}\mathrm{\Delta }E}{}}\mathrm{exp}[\beta E_n].`$ (2.6)
The parameter $`\beta =\beta (E)`$, which is introduced as a multiplier, is defined by
$$E=\frac{\mathrm{ln}Z}{\beta }.$$
(2.7)
This formulation, which exhibits the temperature explicitly, is more convenient than the conventional formulation in microcanonical ensemble.<sup>15)</sup> This formulation, if one assumes static equilibrium, is reduced to the microcanonical ensemble in the limit of small $`\mathrm{\Delta }E`$: The equal a priori probabilities are obtained as
$$p_n\frac{\mathrm{exp}[\beta E]}{\mathrm{\Delta }W(E,N)\mathrm{exp}[\beta E]}=\frac{1}{\mathrm{\Delta }W(E,N)}$$
(2.8)
and the thermodynamic relation
$`F`$ $`=`$ $`{\displaystyle \frac{1}{\beta }}\mathrm{ln}Z`$ (2.9)
$``$ $`{\displaystyle \frac{1}{\beta }}\mathrm{ln}\{\mathrm{\Delta }W(E,N)\mathrm{exp}[\beta E]\}`$
$`=`$ $`E{\displaystyle \frac{1}{\beta }}S/k=ETS.`$
Here we used the Boltzmann entropy
$$S=k\mathrm{ln}\mathrm{\Delta }W(E,N)k\underset{n}{}p_n\mathrm{ln}p_n$$
(2.10)
where $`\mathrm{\Delta }W(E,N)`$ stands for the number of quantum states within the above energy range. Our determination of the time independent quantities $`p_n`$ is thus consistent with the microcanonical ensemble for sufficiently small $`\mathrm{\Delta }E`$.
In the next section, we discuss a formulation of quantum statistical mechanics which incorporates the above statistical inference, and yet a logically consistent formulation in the sense that statistical inference does not precede dynamical time development. Note that we rely on the statistical inference since we deal with an enormous number of states $`E_n`$ in (2.4).
## 3 Statistical Inference and Quantum Statistical Mechanics
We discuss how to describe near-equilibrium states and their approach to thermal equilibrium in a framework which incorporates the least biased statistical inference. We assume that we analyze a physical system which is completely characterized by its total energy and particle number, if the thermal equilibrium should be realized for a fixed volume. In the representation where the total energy and the particle number are diagonal, we have the density matrix<sup>16)</sup> $`\widehat{\rho }(t)`$ which satisfies<sup>8)</sup>
$`Tr\widehat{\rho }(t)={\displaystyle \underset{n}{}}\widehat{\rho }(t)_{nn}={\displaystyle \underset{n}{}}p_n=1,`$
$`Tr\widehat{\rho }\widehat{}={\displaystyle \underset{n}{}}E_np_n=E,`$
$`Tr\widehat{\rho }\widehat{\nu }={\displaystyle \underset{n}{}}\nu _np_n=N.`$ (3.1)
We here assume for simplicity the presence of only one kind of particles. We first note that $`p_n`$ is time independent $`p_n=n|\widehat{\rho }(t)|n=n|e^{i\widehat{}t/\mathrm{}}\widehat{\rho }(0)e^{i\widehat{}t/\mathrm{}}|n=n|\widehat{\rho }(0)|n`$, since the total Hamiltonian is diagonal in the present representation. We assume that either none of the energy levels are degenerate, or if some of them are degenerate, the density matrix $`\widehat{\rho }(t)`$ is diagonalized by a (constant) unitary transformation beforehand in each sector which contains the degenerate energy levels. Consequently, all the possible off-diagonal elements of the density matrix $`\widehat{\rho }(t)`$ are time dependent. Our proposal is to define the Shannon’s amount of uncertainty, which is based on the information available, by
$$S=k\underset{n}{}p_n\mathrm{ln}p_n.$$
(3.2)
This $`S`$, which carries no characteristic properties of quantum theory, may be assigned a purely statistical meaning, and it is time independent.
The dynamical entropy of von Neumann
$$H=kTr\widehat{\rho }\mathrm{ln}\widehat{\rho }$$
(3.3)
which is a quantum generalization of Gibbs entropy, is also time independent since the time development of $`\widehat{\rho }`$ is a unitary transformation. In ref. 9 (and also in ref. 5), this entropy $`H`$, which in principle contains the effects of quantum coherence, is called the information entropy. In this paper we stick to the classical notion of information and thus to the amount of uncertainty defined in (3.2). The advantage of choosing (3.2) as the amount of uncertainty becomes clear later. The average value of any operator in the Schrödinger representation is defined by
$$\widehat{O}(t)=Tr\widehat{\rho }(t)\widehat{O}(0).$$
(3.4)
In the framework of Shannon’s least biased statistical inference used by Jaynes, the maximum value of the amount of uncertainty, which is identified with $`S`$ in (3.2) in the present formulation, is considered with the constraints (3.1). For a closed system we are considering, the constraints are actually replaced by the conditions (2.4) and (2.5) in Section 2. One then obtains the standard result (2.6) for the constant diagonal elements of $`\widehat{\rho }(t)`$
$`p_n=\mathrm{exp}[\beta (E)E_n]/Z,`$
$`Z={\displaystyle \underset{E\frac{1}{2}\mathrm{\Delta }EE_nE+\frac{1}{2}\mathrm{\Delta }E}{}}\mathrm{exp}[\beta (E)E_n]`$ (3.5)
but the time dependent off-diagonal elements of $`\widehat{\rho }(t)`$ are left completely unspecified, namely, we remain maximally noncommittal with regard to missing information.<sup>17)</sup>
On the other hand, the maximum of the von Neumann entropy $`H`$ with conditions (2.4) and (2.5) would give rise to<sup>6,7,8)</sup>
$`n|\widehat{\rho }(t)|n`$ $`=`$ $`\mathrm{exp}[\beta (E)E_n]/Z,`$
$`n|\widehat{\rho }(t)|m`$ $`=`$ $`0,nm.`$ (3.6)
This density matrix $`\widehat{\rho }`$ is completely diagonal; in other words, if one imposes the maximum condition on the von Neumann entropy, we arrive at the conventional microcanonical ensemble without any freedom of time development. In fact, the conventional analysis of statistical mechanics utilizes this property of the entropy $`H`$ and attempts to prove the Boltzmann’s H-theorem for (a coarse grained form of) $`H`$ in (3.3) as an indicator of the general tendency toward thermal equilibrium.<sup>8,9,18)</sup> In contrast, we here attempt to characterize the approach to thermal equilibrium from a different perspective. Also, we will later suggest that the dynamical density matrix does not approach the static form (3.6) even in thermal equilibrium, since the oscillations with microscopic time scales, which should be described by the dynamical density matrix, always exist in the system.
For the general situation, the maximum uncertainty inference of Shannon as formulated here does not specify the density matrix completely, and the best estimate of the average of a general macroscopic operator $`\widehat{O}`$ is
$$\widehat{O}_0\underset{n}{}p_n\widehat{O}(0)_{nn}=\underset{n}{}\widehat{O}(0)_{nn}\mathrm{exp}[\beta (E)E_n]/Z$$
(3.7)
where $`p_n`$ is defined in (3.5) and the sum is taken over the available eigenstates in the representation where the total Hamiltonian and particle number are diagonal. This quantity is time independent by definition and agrees with the conventional average in thermal equilibrium. We will later show that, only after a suitable time averaging, the true average (3.4) is well approximated by the conventional thermal average (3.7) if the system satisfies certain dynamical properties. It is clear that the least biased inference of Shannon, as formulated here, is purely statistical and does not provide any information about dynamical time development.
We next note an inequality between $`S`$ and $`H`$ (see also refs. 5 and 19)
$$HS.$$
(3.8)
This relation is shown by using the standard technique of the statistical mechanics<sup>20)</sup>. This inequality, $`HS`$, valid for any $`\widehat{\rho }(t)`$ suggests that we can impose the maximum amount of uncertainty condition on $`S`$ without any dynamical constraint on $`H`$. (In contrast, if one should use (3.3) as Shannon’s amount of uncertainty, the statistical inference would precede physical dynamics, since the statistical inference would then determine the time dependent off-diagonal elements of $`\widehat{\rho }(t)`$ as well.) Besides, $`HS`$ shows that we start with an initial state with smaller entropy.<sup>12)</sup>
In analogy with the definition of a quantum state in terms of a complete set of commuting hermitian operators,<sup>22)</sup> we assume that our density matrix $`\widehat{\rho }(t)`$, which is a generalization of the Schrödinger wave function, is well specified by a set of macroscopic observables $`\{\widehat{O}\}`$ , and total energy and particle number; in the present case, the operators $`\{\widehat{O}\}`$ do not commute with the total Hamiltonian by our assumption. As in classical Boltzmann statistical mechanics,<sup>11)</sup> where one works exclusively on macroscopic variables, one may define macroscopic observables $`\widehat{O}_\tau `$ in the present framework by the condition
$$Tr\widehat{\rho }(t)\widehat{O}_\tau =\frac{1}{\tau }_t^{t+\tau }𝑑tTr\widehat{\rho }(t)\widehat{O}_\tau .$$
(3.9)
This condition is written in full detail as
$`\widehat{O}_\tau (t)`$ $`=`$ $`{\displaystyle \underset{m,n}{}}\rho _{mn}(0)n|\widehat{O}_\tau |m\mathrm{exp}[i(E_nE_m)t/\mathrm{}]`$ (3.10)
$`=`$ $`{\displaystyle \frac{1}{\tau }}{\displaystyle _t^{t+\tau }}𝑑t{\displaystyle \underset{m,n}{}}\rho _{mn}(0)n|\widehat{O}_\tau |m\mathrm{exp}[i(E_nE_m)t/\mathrm{}]`$
and thus the macroscopic observables are not sensitive to the microscopic time (shorter than $`\tau `$) dependence of the density matrix.<sup>23)</sup> Each macroscopic observable in our definition is labeled by a characteristic time scale $`\tau `$, which may in general depend on the temperature contained in the diagonal components $`\{p_n\}`$ of the density matrix.<sup>24)</sup> An operator with large $`\tau `$ gives a macroscopic observable which agrees with our intuitive understanding: For example, $`\tau =\mathrm{}`$ for the total energy or particle number.
We then define the non-equilibrium state operationally by the relation
$`\mathrm{\Delta }\widehat{O}_\tau (t)`$ $``$ $`Tr\widehat{\rho }(t)[\widehat{O}_\tau {\displaystyle \underset{n}{}}p_n<n|\widehat{O}_\tau |n>]`$ (3.11)
$`=`$ $`Tr\widehat{\rho }(t)\widehat{O}_\tau {\displaystyle \underset{n}{}}p_n<n|\widehat{O}_\tau |n>0`$
for some macroscopic observables $`\widehat{O}_\tau `$ other than the total energy and particle number. Here $`p_n`$ is defined in (3.5). If we do not find any sensible macroscopic observable $`\widehat{O}_\tau `$ which satisfies the above relation ( after a suitable time averaging described later ), the system is in thermal equilibrium.
Our system described by the density matrix $`\widehat{\rho }(t)`$ then develops with time following Schrödinger equation with a fixed value of the von Neumann entropy $`H`$.
## 4 Second Law in Quantum Statistical Mechanics
Our next task is to specify what kind of dynamical properties of a many particle system ensure that the system with initial conditions defined by our statistical inference will in the long run approach the almost equilibrium state. We first note that the time average of our $`\widehat{\rho }(t)`$ over a sufficiently long period approaches arbitrarily close to the equilibrium $`\widehat{\rho }_0`$ with diagonal elements $`p_n`$ in (3.5). In this sense, our system by its construction satisfies the Boltzmann’s ergodic postulate; the time average behavior of a system is the same as its equilibrium behavior. Namely a suitable time averaging of (3.11) gives rise to
$$\overline{\mathrm{\Delta }\widehat{O}_\tau }0$$
(4.1)
which is the (necessary) condition for equilibrium.
We need to sharpen the time averaged behavior (4.1) to be a sufficient dynamical condition for equilibrium. We first recall a quantum version of Poincare’s recurrence theorem.<sup>25)</sup> The theorem states that observables in a (finite) many particle system with discrete energy spectrum are almost periodic, namely, after a suitable time lapse the system comes back to arbitrarily close to the original configuration. This means that the system has a rather well defined dynamical property and we here deal with those finite systems.<sup>26)</sup> We however assume that the recurrence time for a many particle system is sufficiently long by the time scale of our laboratory. The reccurence theorem provides a partial justification for the independent specification of diagonal elements of $`\widehat{\rho }`$ in (3.5) and the non-diagonal elements in (3.11), since an apparently thermal equilibrium state, which is related to (3.5), can come back to any original starting configuration.
To ensure the thermal equilibrium, we need to avoid the persistent synchronized collective oscillation even if the averaged behavior (4.1) is satisfied. We expect that the probability for a great number of oscillators in $`\widehat{\rho }(t)`$ to synchronize persistently is negligibly small for a system of a many particle system. In the generic situation, our system is expected to give a negligible time correlation between $`\mathrm{\Delta }\widehat{O}_\tau (t)`$ for a large time difference $`|t_1t_2|`$, which is however very small compared to the recurrence time.
The averaged behavior (4.1) is now sharpened to be a stronger dynamical postulate as ( see also (3.11))
$`|\overline{\mathrm{\Delta }\widehat{O}_\tau }(t,\mathrm{\Delta }t_c)|`$ $``$ $`|{\displaystyle \frac{1}{\mathrm{\Delta }t_c}}{\displaystyle _t^{t+\mathrm{\Delta }t_c}}𝑑t\mathrm{\Delta }\widehat{O}_\tau (t)|`$ (4.2)
$`=`$ $`|{\displaystyle \frac{1}{\mathrm{\Delta }t_c}}{\displaystyle _t^{t+\mathrm{\Delta }t_c}}𝑑tTr\widehat{\rho }(t)\widehat{O}_\tau {\displaystyle \underset{n}{}}p_n<n|\widehat{O}_\tau |n>|`$
$``$ $`|\widehat{O}_\tau _0|`$
for any macroscopic observable $`\widehat{O}_\tau `$ and a fixed finite $`\mathrm{\Delta }t_c`$; $`\widehat{O}_\tau _0=_np_n<n|\widehat{O}_\tau |n>|`$ is defined in (3.7) with $`p_n`$ in (3.5). For $`\mathrm{\Delta }t_c=\mathrm{}`$, the left-hand side of this relation vanishes by our construction. We impose this condition for a finite $`\mathrm{\Delta }t_c`$ and assume that $`\overline{\mathrm{\Delta }\widehat{O}_\tau }(t,\mathrm{\Delta }t_c)`$ is not sensitive to the absolute value of $`t`$ and a small variation of $`\mathrm{\Delta }t_c`$. Physically, this condition means the existence of a well-defined relaxation time for a set of macroscopic observables. The actual magnitude of $`\mathrm{\Delta }t_c`$, which is expected to be microscopically quite long and of the order of macroscopic time scale, will generally depend on the specific system we are analyzing. In terms of the parameter $`\tau `$ in (3.9), we expect (for an operator with a finite $`\tau `$)
$$\tau <\mathrm{\Delta }t_c<\mathrm{}.$$
(4.3)
To observe the relaxation in terms of the macroscopic observables $`\widehat{O}_\tau `$, we need to have $`\tau <\mathrm{\Delta }t_c`$.
Note that the condition (4.2) does not contradict the reccurence theorem:<sup>25)</sup> In practice, after the initial relaxation, it might be that one can take $`\mathrm{\Delta }t_c\tau `$, namely, one cannot recognize the sizable deviation from thermal equilibrium by the time resolution of macroscopic observables. However, if the reccurence occurs, one need to wait for the time $`\mathrm{\Delta }t_c`$ for the system to relax again.
Since
$$\widehat{O}_\tau (t)=\underset{m,n}{}\rho _{mn}(0)n|\widehat{O}_\tau |m\mathrm{exp}[i(E_nE_m)t/\mathrm{}]$$
(4.4)
we have only the near diagonal components after the above time averaging (4.2); namely, only the terms with
$$|E_nE_m|2\pi \mathrm{}/\mathrm{\Delta }t_c$$
(4.5)
survive the time averaging. Note that $`2\pi \mathrm{}/\mathrm{\Delta }t_c`$ is very small compared to $`\mathrm{\Delta }E`$ in (3.5) for a many particle system we are interested in.<sup>27)</sup> Also the real time-independent diagonal components are subtracted in $`\mathrm{\Delta }\widehat{O}_\tau (t)`$. In the remaining terms we have a sum of complex amplitudes $`\rho _{mn}(0)n|\widehat{O}_\tau |m`$ with nearly equal frequencies $`(E_nE_m)/\mathrm{}`$: Our condition (4.2) is that the sum of oscillating quantities are either absent or destructively interfere
$$|\underset{0<|E_nE_m|2\pi \mathrm{}/\mathrm{\Delta }t_c}{}\rho _{mn}(0)n|\widehat{O}_\tau |m\mathrm{exp}[i(E_nE_m)t/\mathrm{}]||\widehat{O}_\tau _0|$$
(4.6)
for a general value of $`t`$.
As a concrete example of our analysis, we illustrate the problem of a gas confined in the left-half of a box and then removing the partition, although this probelm corresponds to the case of far from equilibrium.<sup>28)</sup> We make a statistical inference to fix $`p_n`$ on the basis of the information about all the possible energy spectra of the entire box, the average energy and particle number. The macroscopic observable $`\widehat{O}_\tau (\stackrel{}{x})`$ may be chosen as the particle number density (in a suitably smeared sense) inside the entire box. We define $`\widehat{O}_\tau _0(\stackrel{}{x})`$ by using the result of the above inference and (3.7). The observable
$$Tr\widehat{\rho }(t)\mathrm{\Delta }\widehat{O}_\tau (\stackrel{}{x})=Tr\widehat{\rho }(t)\widehat{O}_\tau (\stackrel{}{x})\widehat{O}_\tau _0(\stackrel{}{x}),$$
(4.7)
which has a positive peak in the left-half of the box and a negative peak in the right-half of the box at $`t=0`$, is described by choosing the off-diagonal time dependent elements of $`\widehat{\rho }(t)`$ suitably. Our inference agrees with the conventional answer of statistical mechanics, if the possible macroscopic oscillation in $`Tr\widehat{\rho }(t)\mathrm{\Delta }\widehat{O}_\tau (\stackrel{}{x})`$ diminishes soon.
To analyze the effects of time averaging, one may consider a simpler example of the average position of particles (instead of the particle number density)
$$\widehat{O}=\stackrel{}{X}=\underset{i}{}\stackrel{}{x}_i/N$$
(4.8)
of the gas confined into the left-half of the box at $`t=0`$, which is one of the indicators of the macroscopic motion of particles. One may make a crude estimate of the typical frequency contained in $`\stackrel{}{x}_i`$ by considering the matrix element $`n+1|\stackrel{}{x}_i|n`$ for a free particle as
$$\omega (\mathrm{}\pi (n+1))^2/(2mL^2\mathrm{})(\mathrm{}\pi n)^2/(2mL^2\mathrm{})(\mathrm{}\pi n)/(mL^2)v/L$$
(4.9)
where $`v=(\mathrm{}\pi n)/mL`$ is the typical velocity, which is determined by the temperature appearing in the diagonal elements $`p_n`$. Here $`L`$ is the size of the box. It is then unlikely that $`\widehat{O}(t)`$ contains the sizable components with frequency
$$\omega =|E_nE_m|/\mathrm{}2\pi /\mathrm{\Delta }t_c$$
(4.10)
to survive the time averaging for $`\mathrm{\Delta }t_cL/v`$ (or equivalently, $`2\pi /\mathrm{\Delta }t_c2\pi v/L`$), except for the static diagonal elements. For this choice of $`\widehat{O}`$ and the crude estimate, the gas which was in thermal equilibrium and confined into the left-half of the box at $`t=0`$ satisfies the condition (4.2) for $`\mathrm{\Delta }t_cL/v`$. The relaxation time in the present example is determined by the typical transport time scale $`\tau L/v`$, provided that a soliton-like persistent collective motion does not occur.
### 4.1 Entropy law of Clausius
The physical entropy of the final thermodynamic state defined by this time averaging (4.2), which is characterized by $`\mathrm{\Delta }t_c`$, is estimated by
$$\overline{H}(\mathrm{\Delta }t_c)kTr\overline{\widehat{\rho }}\mathrm{ln}\overline{\widehat{\rho }}$$
(4.11)
with
$$\overline{\widehat{\rho }}=(1/\mathrm{\Delta }t_c)_t^{t+\mathrm{\Delta }t_c}\widehat{\rho }(t)𝑑t$$
(4.12)
for a generic value of $`t`$. This value is expected to be close to the maximum of Shannon’s statistical amount of uncertainty $`S`$, which is the maximum value of any sensible definition of entropy because of (3.8): Note that
$$S\overline{H}(\mathrm{\Delta }t_c)$$
(4.13)
since $`\overline{\widehat{\rho }}_{nn}=\widehat{\rho }(t)_{nn}=p_n`$ and $`Tr\overline{\widehat{\rho }}=1`$, which are sufficient to prove (3.8). We interprete this approach of $`\overline{H}(\mathrm{\Delta }t_c)`$ to $`S`$ as a manifestation of the entropy law in the present formulation of quantum statistical mechanics.<sup>29)</sup> We reiterate that we used two ingredients to formulate the entropy law in our approach : The first is the statistical input related to the least biased inference on the basis of a limited amount of information available (3.5), and the second is the dynamical input related to the time averaging in (3.9) and (4.2).
The von Neumann entropy $`H`$ is in contrast rigidly defined by the basic dynamics, and it does not allow any arbitrary manipulation such as taking a time averaging of $`\widehat{\rho }(t)`$. The dynamical entropy $`H`$ of von Neumann stays constant throughout the unitary time development of the system regardless of our time averaging procedure. In fact, the above time averaging (4.2) to define the physical thermodynamic state resolves the discrepancy of the dynamical von Neumann $`H`$ and the physical statistical entropy (4.11) we defined. Physically, the von Neumann entropy $`H`$ is sensitive to the dynamical motion of all the time scales in the system, whereas the thermodynamic entropy (4.11) is not sensitive to the motion with time scales shorter than $`\mathrm{\Delta }t_c`$; moreover, the condition (4.2) states that the macroscopic motion with time scales larger than $`\mathrm{\Delta }t_c`$ in the system is negligible.
We here note an interesting analogy of the present formulation of physical entropy with the renormalization group in field theory. The parameter $`\mathrm{\Delta }t_c`$ (or to be precise, $`\mathrm{}/\mathrm{\Delta }t_c`$) characterizes the energy scale of the theory, and the entropy $`\overline{H}(\mathrm{\Delta }t_c)kTr\overline{\widehat{\rho }}\mathrm{ln}\overline{\widehat{\rho }}`$ in (4.11) corresponds to the renormalized running coupling constant. The ultra-violet limit $`\mathrm{\Delta }t_c0`$ gives rise to the von Neumann entropy $`H`$, which corresponds to the bare coupling constant, the fundamental quantity defined by the basic dynamics , namely, quantum mechanics in the present case. But physics is not sensitive to the bare coupling constant. The infrared limit of $`\overline{H}`$ for $`\mathrm{\Delta }t_clarge`$ gives rise to the measurable quantity, the maximum of Shannon’s $`S`$, corresponding to the coupling constant $`\alpha =1/137`$ in QED defined in the Thomson limit. In this analogy, the entropy law of Clausius corresponds to a statement of the existence of a stable infrared fixed point for finite $`\mathrm{\Delta }t_c`$.
This picture also suggests that the dynamical density matrix $`\widehat{\rho }(t)`$ does not approach the static diagonal form in (3.6) even in thermal equilibrium, since the oscillations with microscopic time scales always exist. The microscopic time dependence of the equilibrium density matrix is also expected in the conventional Gibbs ensemble theory, if one defines the thermal equilibrium by macroscopic observables $`\widehat{O}_\tau `$ with well defined time scale $`\tau `$. Our definition of the macroscopic observables $`\widehat{O}_\tau `$ in (3.10) shows that the macroscopic observables are not sensitive to the possible microscopic time (shorter than $`\tau `$) dependence of the equilibrium density matrix.
Our picture is expected to have an implication on the linear response theory for time dependent current correlations<sup>30)</sup> with frequencies which are comparable to $`1/\tau `$, since the possible time dependence of the equilibrium density matrix is not ignored for such a case. In the applications of the fluctuation-dissipation theorem, one usually examines the correlation functions of operators averaged with the static diagonal equilibrium density matrix $`e^{\beta \widehat{}}/Z`$. The possible microscopic time dependence of equilibrium density matrix may lead to interesting implications on the analysis of the foundations of linear response theory.
## 5 Discussion
The purpose of this note has been to discuss quantum statistical mechanics which incorporates Shannon’s statistical inference, and to analyze a resulting picture for the second law. Our statistical inference does not resolve the basic issue why the statistical inference works for a many particle system, but our statistical inference ,unlike the equal a priori probabilities, allows an analysis of the dynamical aspects of the second law.
Our picture for the general tendency toward thermal equilibrium is quite different from the one in the conventional formulation of quantum mechanical H-theorem:<sup>8,9,18)</sup> We note that the conventional coarse-grained approach performs not only the statistical operation by setting the diagonal elements in each subsector of the density matrix to be equal by assuming equal a priori probabilities, but also the dynamical operation by setting the off-diagonal time dependent elements to be zero to let the entropy increase.
We now briefly comment on the two ambitious approaches to the second law on the basis of purely dynamical considerations, to be compared with our more conservative statistical analysis. The first is the approach of Tasaki<sup>31)</sup> and the second is the approach of Van Hove:<sup>32)</sup>
Tasaki analyzed a possible quantum mechanical derivation of a canonical ensemble starting from a pure quantum state. The basic idea is to consider two quantum systems, a bath $`|\alpha `$ and a subsystem $`|a`$, and one then introduces a suitable interaction between them so that one obtains a pure quantum state for the combined system,
$$|\psi =\underset{a,\alpha }{}M_{\alpha ,a}|\alpha |a$$
(5.1)
which has a vanishing von Neumann entropy. One then traces out the bath system in the density matrix
$$\widehat{\rho }_{tot}=|\psi \psi |\widehat{\rho }_{sub}=\underset{a,b}{}(M^{}M)_{ab}|ba|.$$
(5.2)
Because of the quantum entanglement of two systems, one then obtains a mixed state for the subsystem which has a non-vanishing entanglement entropy<sup>33)</sup>
$$H_{sub}=kTr\widehat{\rho }_{sub}\mathrm{ln}\widehat{\rho }_{sub}>0.$$
(5.3)
To reproduce a canonical ensemble for the subsystem, Tasaki assumes the “hypothesis of equal weights for eigenstates” for the combined system. He also argues the cancellation of oscillating non-diagonal components of the density matrix for the subsystem at sufficiently large $`t`$ by using Chebyshev’s inequality.<sup>31)</sup> If one can show that the “hypothesis of equal weights for eigenstates” holds for a rather general class of dynamical systems, the quantum entanglement entropy would provide a physical explanation of the statistical entropy. At this moment, it appears that the generality of the hypothesis has not been established.
In the analysis of entanglement entropy, a clear distinction between the von Neumann’s $`H`$, which is equal for both of the bath and the subsystem<sup>33)</sup> $`H_{bath}=H_{sub}`$ and thus not extensive, and the Shannon’s $`S`$, which could be vastly different for the bath and the subsystem and thus could be extensive, is expected to be essential.
Van Hove<sup>32)</sup> analyzed the possible approach of a general mixed state to microcanonical states by a dynamical time development in the limit $`t\mathrm{}`$. To be specific, he starts with a mixed state
$$\widehat{\rho }(0)𝑑\alpha |\alpha |c_\alpha |^2\alpha |$$
(5.4)
where $`\widehat{H}_0|\alpha =ϵ_\alpha |\alpha `$ , and the unitary time development generated by
$$U(t)=\mathrm{exp}[i(\widehat{H}_0+\lambda \widehat{V})t/\mathrm{}].$$
(5.5)
His major claim is ( see Eq. (1.4) in the first paper in ref. 32)
$$Tr\widehat{O}U(t)\widehat{\rho }(0)U^1(t)\widehat{O}_{microcanonical}$$
(5.6)
for $`t\mathrm{}`$. Since the operator $`U(t)`$ is unitary for whatever large but finite $`t`$, both of the von Neumann’s $`H`$ and Shannon’s $`S`$ remain at the initial value different from the microcanonical value. Van Hove however considers a singular limit $`\lambda ^2tconstant`$ for $`t\mathrm{}`$ (i.e., $`\lambda t1/\lambda \mathrm{}`$ for $`t\mathrm{}`$). In such a limit, if literally taken, the $`S`$-matrix
$$\widehat{S}=\underset{t_\pm \pm \mathrm{}}{lim}e^{i\widehat{H}_0t_+/\mathrm{}}e^{i(\widehat{H}_0+\lambda \widehat{V})(t_+t_{})/\mathrm{}}e^{i\widehat{H}_0t_{}/\mathrm{}}$$
(5.7)
(and consequently perturbation theory he uses) is not defined, since the condition of an adiabatic switch-on and switch-off of the interaction is not satisfied. See also ref. 34. If one can provide a mathematical basis for the singular limit, one would be able to derive the microcanonical ensemble from a general mixed state by a unitary time development. At this moment, to our knowledge, such a mathematical basis appears to be missing.
Finally, we mention recent activities on an alternative approach to the second law. Jarzynski<sup>35)</sup> found the following amusing identity
$`{\displaystyle \frac{Z_1}{Z_0}}`$ $`=`$ $`{\displaystyle \frac{1}{Z_0}}{\displaystyle 𝑑\mu (z^{})e^{\beta H_1(z^{})}}`$ (5.8)
$`=`$ $`{\displaystyle \frac{1}{Z_0}}{\displaystyle 𝑑\mu (z^{})e^{\beta H_0(z)\beta (H_1(z^{})H_0(z))}}`$
$`=`$ $`{\displaystyle \frac{1}{Z_0}}{\displaystyle 𝑑\mu (z)e^{\beta H_0(z)\beta W(z)}}=e^{\beta W(z)}`$
where we defined the mapping
$$z=(q(t_0),p(t_0))z^{}=(q^{}(t_1),p^{}(t_1))$$
(5.9)
as the canonical transformation generated by the time dependent Hamiltonian (a quantum version is also known<sup>36)</sup>)
$$H_{\lambda (t)},t_0tt_1$$
(5.10)
with
$$\lambda (t_0)=0,\lambda (t_1)=1.$$
(5.11)
The Liouville theorem $`d\mu (z^{})=d\mu (z)`$ is essential in the above identity. We also defined the work done during the time development by $`W(z)H_1(z^{})H_0(z)`$.
If one defines the Helmholtz free energy for the system described by Hamiltonian $`H_1`$ at temperature $`\beta `$ by $`F_1(\beta )`$ , one obtains from (5.8)
$$\mathrm{\Delta }F(\beta )=F_1(\beta )F_0(\beta )=\frac{1}{\beta }\mathrm{ln}e^{\beta W(z)}.$$
(5.12)
and, by noting the mathematical inequality $`\mathrm{exp}[A(z)]\mathrm{exp}[A(z)]`$,
$$\mathrm{\Delta }F(\beta )W$$
(5.13)
which resembles the basic thermodynamic inequality,<sup>6)</sup> an alternative expression of the second law. The identity (5.8) as it stands is equivalent to the Liouville theorem and thus contains no information about the thermal entropy generation. In fact, it is known (from an explicit analysis of harmonic oscillators, for example) that the equality sign in (5.13) does not hold for an infinitely slow adiabatic work<sup>35)</sup>, which is by itself consistent: But to analyze the equality sign in (5.13), one need to analyze the approach of a system in (5.8), once driven out of equilibrium by an external work, to thermal equilibrium again. The analysis of Jarzynski is thus complementary to the analysis in the present note, and certainly it does not replace our analysis. See recent works<sup>37,38</sup> related to the above identity. See also Lenard<sup>39)</sup> for an analysis of similar inequality associated with a canonical ensemble.<sup>40)</sup>
In conclusion, we have illustrated a physical picture of the second law when one makes a clear distinction between statistical aspects and dynamical aspects in statistical mechanics, which is made possible if one uses $`S`$ in (3.2). Our basic view as presented here is rather conservative, namely, it is based on the premise that the entropy law of Clausius is not a direct consequence of microscopic dynamical laws alone. In the context of classical Boltzmann approach, this view appears to be shared with experts just to quote “ It follows that the macroscopic dynamics cannot be a consequence of the microscopic dynamics alone”.<sup>10)</sup> The remaining basic issue in the present approach is to specify precisely the class of many-particle Hamiltonians which ensure (4.2).
I thank H. Tasaki for numerous clarifying comments, and A. Shimizu and M. Ueda for helpful comments at the initial stage of this work.
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# On semi-infinite cohomology of finite dimensional algebras
## 1. Introduction
Semi-infinite cohomology of associative algebras was studied, in particular, by S. Arkhipov (see \[Ar1\], \[Ar2\], \[Ar3\]). Recall that the definition of semi-infinite cohomology in \[Ar1\] works in the following set-up. We are given an associative graded algebra $`A`$, two subalgebras $`B,NA`$ such that $`A=BN`$ as a vector space, satisfying some additional assumptions. In this situation the space of semi-infinite Ext’s, $`Ext^{\mathrm{}/2+}(X,Y)`$ is defined for $`X,Y`$ in the bounded derived category of graded $`A`$-modules. The definition makes use of explicit complexes.
In this note we show that under some additional assumptions semi-infinite Ext groups $`Ext^{\mathrm{}/2+}(X,Y)`$ has a categorical interpretation. More precisely, given a category $`𝒜`$ and subcategory $`𝒜`$ one can define for $`X,Y𝒜`$ the set of morphisms from $`X`$ to $`Y`$ ”through $``$”; we denote this space by $`Hom_𝒜_{}(X,Y)`$. We then show that if $`𝒜`$ is the bounded derived category of $`A`$-modules, and $``$ is the full triangulated subcategory generated by $`B`$-projective $`A`$-modules, then, under certain assumptions one has
(1)
$$Ext^{\mathrm{}/2+i}(X,Y)=Hom_𝒜_{}(X,Y[i]).$$
Notice that the right hand side of (1) makes sense for a wide class of pairs $`(A,B)`$ (an associative algebra, and a subalgebra), and $`X,YD^b(Amod)`$; in particular we do not need $`A,B`$ to be graded. Thus one may consider (1) as providing a generalization of the definition of semi-infinite Ext’s to this set up. However, we should warn the reader that under our working assumptions, but not in general, $``$ also equals the full triangulated subcategory generated by $`B`$-injective modules, or by modules (co)induced from a ”complemental” subalgebra $`NA`$, so one has at least four different obvious generalizations of the definition of the right-hand side of (1).
In fact, a description of semi-infinite cohomology similar to (1) in a general situation (in particular, in the case of enveloping algebras of infinite-dimensional Lie algebras) requires additional ideas, and is the subject of a forthcoming joint work with Arkhipov and Positselskii.
An example of the situation considered in this paper is provided by a small quantum group at a root of unity \[L\], or by the restricted enveloping algebra of a simple Lie algebra in positive characteristic. Computation of semi-infinite cohomology in the former case is due to S. Arkhipov \[Ar1\] (the answer suggested as a conjecture by B. Feigin). This example was a motivation for the present work. We informally explain the relation of our Theorem 1 to the answer for semi-infinite cohomology of small quantum groups in Remark 5 below (we plan to derive it from Theorem 1 elsewhere).
Acknowledgements. I thank S. Arkhipov for stimulating interest, and L. Positselskii for helpful comments. I thank the Clay Mathematical Institute and NSF grant DMS-0071967 for financial support.
## 2. Categorical preliminaries: morphisms through a functor
Let $`𝒜`$, $``$ be (small) categories, and $`\mathrm{\Phi }:𝒜`$ be a functor. For $`X,YOb(𝒜)`$ define the set of ”morphisms from $`X`$ to $`Y`$ through $`\mathrm{\Phi }`$” as $`\pi _0`$ of the category of diagrams
(2)
$$X\mathrm{\Phi }(\mathrm{?})Y,\mathrm{?}.$$
This set will be denoted by $`Hom_{𝒜_\mathrm{\Phi }}(X,Y)`$. Thus elements of $`Hom_{𝒜_\mathrm{\Phi }}(X,Y)`$ are diagrams of the form (2), with two diagrams identified if there exists a morphism between them. Composing the two arrows in (2) we get a functorial map
(3)
$$Hom_{𝒜_\mathrm{\Phi }}(X,Y)Hom_𝒜(X,Y).$$
If $`𝒜`$, $``$ are additive and $`\mathrm{\Phi }`$ is an additive functor, then addition of diagrams of the form (2) is defined by
$$(X\stackrel{𝑓}{}\mathrm{\Phi }(Z)\stackrel{𝑔}{}Y)+(X\stackrel{f^{}}{}\mathrm{\Phi }(Z^{})\stackrel{g^{}}{}Y)=(X\stackrel{f\times f^{}}{}\mathrm{\Phi }(ZZ^{})\stackrel{gg^{}}{}Y);$$
it induces an abelian group structure on $`Hom_{𝒜_\mathrm{\Phi }}(X,Y)`$. Proposition 3 in \[ML\], VIII.2 shows that for $`Z`$ the tautological map
$$Hom(X,\mathrm{\Phi }(Z))_{}Hom(\mathrm{\Phi }(Z),Y)Hom_{𝒜_\mathrm{\Phi }}(X,Y)$$
is compatible with addition.
We have the composition map
$$Hom_𝒜(X^{},X)\times Hom_{𝒜_\mathrm{\Phi }}(X,Y)\times Hom_𝒜(Y,Y^{})Hom_{𝒜_\mathrm{\Phi }}(X^{},Y);$$
in particular, for $`𝒜`$, $``$, $`\mathrm{\Phi }`$ additive, $`Hom_{𝒜_\mathrm{\Phi }}(X,Y)`$ is an $`End(X)End(Y)`$ bimodule.
Given $`\mathrm{\Phi }:𝒜`$, $`\mathrm{\Phi }^{}:𝒜^{}^{}`$ and $`F:𝒜𝒜^{}`$, $`G:^{}`$ with $`F\mathrm{\Phi }\mathrm{\Phi }^{}G`$ we get for $`X,Y𝒜`$ a map
(4)
$$Hom_{𝒜_\mathrm{\Phi }}(X,Y)Hom_𝒜_{^{}}^{}(F(X),F(Y)).$$
If the left adjoint functor $`\mathrm{\Phi }^{}`$ to $`\mathrm{\Phi }`$ is defined on $`X`$, then we have
$$Hom_{𝒜_\mathrm{\Phi }}(X,Y)=Hom_𝒜(\mathrm{\Phi }(\mathrm{\Phi }^{}(X)),Y),$$
because in this case the above category contracts to the subcategory of diagrams of the form $`X\stackrel{can}{}\mathrm{\Phi }(\mathrm{\Phi }^{}(X))Y`$, where $`can`$ stands for the adjunction morphism. If the right adjoint functor $`\mathrm{\Phi }^!`$ is defined on $`Y`$, then
$$Hom_{𝒜_\mathrm{\Phi }}(X,Y)=Hom_𝒜(X,\mathrm{\Phi }(\mathrm{\Phi }^!(Y)))$$
for similar reasons. In particular, if $`\mathrm{\Phi }`$ is a full imbedding then (3) is an isomorphism provided either $`X`$ or $`Y`$ lie in the image of $`\mathrm{\Phi }`$.
In all examples below $`𝒜`$ will be a triangulated category, and $`\mathrm{\Phi }:𝒜`$ will be an imbedding of a (strictly) full triangulated subcategory. Given $`𝒜`$ we will tacitly assume $`\mathrm{\Phi }`$ to be the imbedding, and write $`Hom_𝒜_{}`$ (”morphisms through $``$”) instead of $`Hom_{𝒜_\mathrm{\Phi }}`$.
###### Example 1.
Let $`M`$ be a Noetherian scheme, and $`𝒜=D^b(Coh_M)`$ be the bounded derived category of coherent sheaves on $`M`$; let $`I:𝒜`$ be the full subcategory of complexes whose cohomology ia supported on a closed subset $`i:NM`$. Then the functor $`II^!=i_{}i^!`$ takes values in a larger derived category of quasi-coherent sheaves (i.e. ind-coherent sheaves), and $`II^{}=i_{}i^{}`$ takes values in the Grothendick-Serre dual category, the derived category of pro-coherent sheaves (introduced in Deligne’s appendix to \[H\]). Still we have
$$Hom_𝒜_{}(X,Y)=Hom(X,i_{}(i^!(Y)))=Hom(i_{}(i^{}(X)),Y).$$
In particular, if $`X=𝒪_M`$ is the structure sheaf, we get
(5)
$$Hom_𝒜_{}(𝒪_M,Y[i])=H_N^i(Y),$$
where $`H_N^{}(Y)`$ stands for cohomology with support on $`N`$ (local cohomology) \[H\].
## 3. Recollection of the definition of $`Ext^{\mathrm{}/2+}`$
All algebras below will be associative and unital algebras over a field.
We recall a variant of definition of semi-infinite Ext’s (available under certain restrictions on the algebra and subalgebras) suited for our purpose (see e.g. \[FS\], §2.4, pp 180-183, for this definition in the particular case of small quantum groups; the general case is analogous).
We make the following assumptions. A $``$-graded algebra $`A`$ and graded subalgebras $`A^0`$, $`A^0`$, $`A^0`$ $`A`$ are fixed and satisfy the following conditions:
(1) $`A^0`$, $`A^0`$ are graded by, respectively, $`^0`$, $`^0`$, and $`A^0=A^0A^0`$ is the component of degree 0 in $`A^0`$ and in $`A^0`$.
(2) The maps $`A^0_{A^0}A^0A`$ and $`A^0_{A^0}A^0A`$ provided by the multiplication map are isomorphisms.
(3) $`A`$ is finite dimensional; $`A^0`$ is semisimple, and $`A^0`$ is self-injective (i.e. the free $`A^0`$-module is injective).
By a ”module” we will mean a finite dimensional graded module, unless stated otherwise. By $`Amod`$ we denote the category of (graded finite dimensional) $`A`$-modules.
Recall that a bounded below complex of graded modules is called convex if the weights ”go down”, i.e. for any $`n`$ the sum of weight spaces of degree more than $`n`$ is finite dimensional. A bounded below complex of graded modules is called concave if the weights ”go up” in the similar sense.
###### Lemma 1.
i) Any $`A`$-module admits a right convex resolution by $`A`$-modules, which are injective as $`A^0`$-modules. It also admits a right concave resolution by $`A`$-modules, which are $`A^0`$-injective.
ii) Any finite complex of $`A`$-modules is a quasiisomorphic subcomplex of a bounded below convex complex of $`A^0`$-injective $`A`$-modules. It is also a quasiisomorphic subcomplex of a bounded below concave complex of $`A^0`$-injective $`A`$-modules.
###### Proof.
To deduce (ii) from (i) imbed given finite complex $`C^{}Com(Amod)`$ into a complex of $`A`$-injective modules $`I^{}Com^0(Amod)`$ (notice that condition (2) above implies that an $`A`$-injective module is also $`A^0`$ and $`A^0`$ injective), and apply (i) to the module of cocycles $`Z^n=I^n/d(I^{n1})`$ for large $`n`$.
To check (i) it suffices to find for any $`MAmod`$ an imbedding $`MI`$, where $`I`$ is $`A^0`$ injective, and if $`n`$ is such that all graded components $`M_i`$ for $`i<n`$ vanish, then $`M_n\stackrel{~}{}I_n`$. (This would prove the second part of the statement; the first one is obtained from the first one by renotation.) It suffices to take $`I=CoInd_{A^0}^A(Res_{A^0}^A(M))`$. It is indeed $`A^0`$-injective, because of the equality
(6)
$$Res_{A^0}^A(CoInd_{A^0}^A(M))=CoInd_{A^0}^{A^0}(M)),$$
which is a consequence of assumption (2) above. o
We set $`D=D^b(Amod)`$.
###### Definition 1.
(cf. \[FS\], §2.4) The assumptions (1–3) are enforced. Let $`X,YD`$. Let $`J_{}^X`$ be a convex bounded below complex of $`A^0`$-injective (= projective) modules quasiisomorphic to $`X`$, and $`J_{}^Y`$ be a concave bounded below complex of $`A^0`$-injective modules quasiisomorphic to $`Y`$. Then one defines
(7)
$$Ext^{\mathrm{}/2+i}(X,Y)=H^i(Hom^{}(J_{}^X,J_{}^Y)).$$
###### Remark 1.
Independence of the right-hand side of (7) on the choice of resolutions $`J_{}^X`$, $`J_{}^Y`$ follows from the argument below. Since particular complexes used in \[Ar1\] to define $`Ext^{\mathrm{}/2+}`$ satisfy our assumptions, we see that this definition agrees with the one in loc. cit.
###### Remark 2.
Notice that $`Hom`$ in the right-hand side of (7) is $`Hom`$ in the category of graded modules. As usual, it is often convenient to denote by $`Ext^{\mathrm{}/2+i}(X,Y)`$ the graded space which in present notations is written down as $`\underset{n}{}Ext^{\mathrm{}/2+i}(X,Y(n))`$, where $`(n)`$ refers to shift of grading by $`n`$.
###### Remark 3.
The next standard Lemma shows that conditions on the resolutions $`J_{}^X`$, $`J_{}^Y`$ used in the (7) can be formulated in terms of the subalgebra $`A^0`$ alone (or, alternatively, in terms of $`A^0`$ alone); this conforms with the fact that the left-hand side of (11) in Theorem 1 below depends only on $`A^0`$. However, existence of a ”complemental” subalgebra $`A^0`$ is used in the construction of a resolution $`J_{}^X`$ with required properties.
###### Lemma 2.
An $`A`$-module is $`A^0`$-injective iff it is has a filtration with subquotients of the form $`CoInd_{A^0}^A(M)`$, $`MA^0mod`$.
###### Proof.
The ”if” direction follows from semisimplicity of $`A^0`$, and equality (6) above. To show the ”only if” part let $`M`$ be an $`A^0`$-injective $`A`$-module. Let $`M^{}`$ be its graded component of minimal degree; then the canonical morphism
(8)
$$MCoInd_{A^0}^{A^0}(M^{})$$
is surjective. If $`M`$ is actually an $`A`$-module, then the projection $`MM^{}`$ is a surjection of $`A^0`$-modules, hence yields a morphism
(9)
$$MCoInd_{A^0}^A(M^{}).$$
(6) shows that $`Res_{A^0}^A`$ sends (9) into (8); in particular (9) is surjective. Thus the top quotient of the required filtration is constructed, and the proof is finished by induction. o
###### Remark 4.
In two special cases $`Ext^{\mathrm{}/2+i}(X,Y)`$ coincides with a traditional derived functor. First, suppose that $`Res_{A^0}^A(X)`$ has finite injective (equivalently, projective) dimension; then one can use a finite complex $`J_{}^X`$ in (7) above. It follows immediately, that in this case we have
$$Ext^{\mathrm{}/2+i}(X,Y)Hom(X,Y[i]).$$
On the other hand, suppose that $`Res_{A^0}^A(Y)`$ has finite injective dimension, so that the complex $`J_{}^Y`$ in (7) can be chosen to be finite. To describe semi-infinite Ext’s in this case we need another notation. Let $`A^{}`$ denote the co-regular $`A`$-bimodule; for $`MAmod`$ let $`M\stackrel{ˇ}{}=M^{}=Hom_A(M,A^{})`$ denote the corresponding right $`A`$-module, and we use the same notation for the corresponding functor on the derived categories. Let also $`S:D^b(Amod)D^+(Amod)`$ be given by $`S(Y)=RHom_A(A^{},Y)`$. Notice that $`A^{}`$ is $`A^0`$-projective by self-injectivity of $`A^0`$; thus Lemma 2 shows that $`Ext_A^i(A^{},N)=0`$ for $`i>0`$ if $`N`$ is $`A^0`$-injective. In particular, $`S(Y)D^b(Amod)`$ if $`Y|_{A^0}`$ has finite injective dimension. We claim that in this case we have
$$Ext^{\mathrm{}/2+i}(X,Y)X\stackrel{ˇ}{}\stackrel{𝐿}{}_AS(Y).$$
This isomorphism an immediate consequence of the next Lemma. We also remark that if $`A`$ is a Frobenius algebra, then $`SId`$.
###### Lemma 3.
Let $`M,NAmod`$ be such that $`M`$ is $`A^0`$-projective, while $`N`$ is $`A^0`$-injective. Then we have
a) $`Ext_A^i(M,N)=0`$; $`Ext_A^i(A^{},N)=(R^iS)(N)=0`$, $`Tor_i^A(M\stackrel{ˇ}{},S(N))=0`$ for $`i0`$.
b) The natural map
(10)
$$M\stackrel{ˇ}{}_AS(N)=Hom_A(M,A^{})_AHom_A(A^{},N)Hom_A(M,N)$$
is an isomorphism.
###### Proof.
The first equality in (a) follows from Lemma 2, and the second one was checked above. Self-injectivity of $`A^0`$ shows that $`M\stackrel{ˇ}{}`$ is $`A^0`$-projective, and a variant of Lemma 2 ensures that it is filtered by modules induced from $`A^0`$. Thus it sufficies to show that $`S(N)`$ is $`A^0`$-projective. This follows from isomorphisms
$$Hom_A(A^{},CoInd_{A^0}^A(N_0))=Hom_{A^0}(A^{},N_0)Hom_{A^0}((A^0)^{},N_0)_{A^0}A^0.$$
Let us now deduce (b) from (a). Notice that (a) implies that both sides of (10) are exact in $`N`$ (and also in $`M`$), i.e. send exact sequences $`0N^{}NN^{\prime \prime }0`$ with $`N^{}`$, $`N^{\prime \prime }`$ being $`A^0`$-injective into exact sequences. Also (10) is evidently an isomorphism for $`N=A^{}`$. For any $`A^0`$-injective $`N`$ there exists an exact sequence
$$0N(A^{})^n\stackrel{\mathit{\varphi }}{}(A^{})^m$$
with image and cokernel of $`\varphi `$ being $`A^0`$-injective. Thus both sides of (10) turn it into an exact sequence, which shows that (10) is an isomorphism for any $`A^0`$-injective $`N`$. o
## 4. Main result
###### Theorem 1.
Let $`D_{\mathrm{}/2}D`$ be the full tringulated subcategory of $`D`$ generated by $`A^0`$-injective (=projective) modules. For $`X,YD^b(Amod)`$ we have a natural isomorphism
(11)
$$Hom_{D_{D_{\mathrm{}/2}}}(X,Y[i])Ext^{\mathrm{}/2+i}(X,Y).$$
The proof of Theorem 1 is based on the following
###### Lemma 4.
i) Every graded $`A^0`$-injective $`A`$-module admits a concave right resolution consisting of $`A`$-injective modules.
ii) A finite complex of graded $`A^0`$-injective $`A`$-modules is quasiisomorphic to a concave bounded below complex of $`A`$-injective modules.
###### Proof.
(ii) follows from (i) as in the proof of Lemma 1. (Recall that, according to Hilbert, if a bounded below complex of injectives represents an object $`XD^b`$ which has finite injective dimension, then for large $`n`$ the module of cocycles is injective.)
To prove (i) it is enough for any $`A^0`$-injective module $`M`$ to find an imbedding $`MI`$, where $`I`$ is $`A`$-injective, and $`M_n\stackrel{~}{}I_n`$ provided $`M_i=0`$ for $`i<n`$. (Notice that cokernel of such an imbedding is $`A^0`$-injective, because $`I`$ is $`A^0`$-injective by condition (2).) We can take $`I`$ to be $`CoInd_{A_0}^A(Res_{A_0}^A(M))`$. Then $`I`$ is indeed injective, because $`M`$ is $`A^0`$-injective by semi-simplicity of $`A^0`$, and condition on weights is clearly satisfied. o
###### Proposition 1.
a) Let $`J_{}`$ be a convex bounded below complex of $`A`$-modules. Let $`J_{}^n`$ be the $`n`$-th stupid truncation of $`J_{}`$ (thus $`J_{}^n`$ is a quotient complex of $`J_{}`$).
Let $`Z`$ be a finite complex of $`A^0`$-injective $`A`$-modules. Then we have
(12)
$$Hom_D(X,Z)\stackrel{~}{}\underset{}{\mathrm{lim}}Hom_D(J_{}^n,Z).$$
In fact, for $`n`$ large enough we have
$$Hom_D(X,Z)\stackrel{~}{}Hom_D(J_{}^n,Z).$$
###### Proof.
Let $`I_{}`$ be a concave bounded below complex of $`A`$-injective modules quasiisomorphic to $`Z`$ (which exists by Lemma 4(ii)). Then the left-hand side of (12) equals $`Hom_{Hot}(J_{},I_{})`$ where $`Hot`$ stands for the homotopy category of complexes of $`A`$-modules. Conditions on weights of our complexes ensure that there are only finitely many degrees for which the corresponding graded components both in $`J_{}`$ and $`I_{}`$ are nonzero; thus any morphism between graded vector spaces $`J_{}`$, $`I_{}`$ factors through the finite dimensional sum of corresponding graded components. In particular, $`Hom^{}(J_{}^n,I_{})\stackrel{~}{}Hom^{}(J_{},I_{})`$ for large $`n`$, and hence
$$Hom_{D(Amod)}(J_{}^n,I_{})=Hom_{Hot}(J_{}^n,I_{})\stackrel{~}{}Hom_{Hot}(J_{},I_{})$$
for large $`n`$. o
Proof of the Theorem. We keep notations of Definition 1. It follows from the Proposition that
$$Hom_{D_{D_{\mathrm{}/2}}}(X,Y[i])=\underset{}{\mathrm{lim}}_nHom_D((J_{}^X)^n,Y[i]).$$
The right-hand side of (11) (defined in (7)) equals $`H^i(Hom^{}(J_{}^X,J_{}^Y))`$. Conditions on weights of $`J_{}^X`$, $`J_{}^Y`$ show that for large $`n`$ we have
$$Hom^{}((J_{}^X)^n,J_{}^Y)\stackrel{~}{}Hom^{}(J_{}^X,J_{}^Y).$$
Lemma 2 implies that $`Ext_A^i(M_1,M_2)=0`$ for $`i>0`$ if $`M_1`$ is $`A^0`$-projective, and $`M_2`$ is $`A^0`$-injective. Thus
$$Hom_D((J_{}^X)^n,Y[i])=H^i(Hom^{}(J_{}^X,J_{}^Y)).$$
The Theorem is proved. o
###### Remark 5.
This remark concerns with the example provided by a small quantum group. So let $`𝔤`$ be a simple Lie algebra over $``$, $`q`$ be a root of unity of order $`l`$, and let $`A=u_q=u_q(𝔤)`$ be the corresponding small quantum group \[L\]. Let $`A^0=b_qu_q`$ and $`A^0=b_q^{}u_q`$ be respectively the upper and the lower triangular subalgebras. Then the above conditions (1–3) are satisfied.
Let $`𝕀`$ denote the trivial $`u_q`$-module. The cohomology $`Ext_{u_q}^{}(𝕀,𝕀)`$, and the semi-infinite cohomology $`Ext^{\mathrm{}/2+}(𝕀,𝕀)`$ were computed respectively in \[GK\] and \[Ar1\]. Let us recall the results of these computations.
Assume for simplicity that $`l`$ is prime to twice the maximal multiplicity of an edge in the Dynkin diagram of $`𝔤`$. Let $`𝒩𝔤`$ be the cone of nilpotent elements, and $`𝔫𝒩`$ be a maximal nilpotent subalgebra. Then the Theorem of Ginzburg and Kumar asserts that
(13)
$$Ext^{}(𝕀,𝕀)𝒪(𝒩),$$
the algebra of regular functions on $`𝒩`$. Also, a Theorem of Arkhipov (conjectured by Feigin) asserts that
(14)
$$Ext^{\mathrm{}/2+}(𝕀,𝕀)H_𝔫^d(𝒩,𝒪),$$
where $`d`$ is the dimension of $`𝔫`$, and $`H_𝔫`$ denotes cohomology with support on $`𝔫`$; one also has $`H_𝔫^i(𝒩,𝒪)=0`$ for $`id`$ (here the choice of $`𝔫`$ is assumed to be compatible with the choice of an upper triangular subalgebra $`b_qu_q`$ via isomorphism (13) in a natural sense).
The aim of this remark is to point out a formal similarity between (14) and equality (5) in Example 1 above. Namely, the Ginzburg-Kumar isomorphism (13) yields a functor $`F:D^b(u_qmod)Coh(𝒩)`$, $`F(X)=Ext^{}(𝕀,X)`$, such that $`F(𝕀)=𝒪_𝒩`$ is the structure sheaf. It is easy to see that if $`XD^b(u_qmod)`$ has finite projective (equivalently, injective) homological dimension over $`b_q`$, then the support of $`F(X)`$ lies in $`𝔫`$ (here by support we mean set-theoretic rather than scheme-theoretic support, so the coherent sheaf $`F(X)`$ may be annihilated by some power of the ideal of $`𝔫`$). Thus if we assume for a moment that the functor $`F`$ can be lifted to a triangulated functor $`\stackrel{~}{F}^{}:D^b(u_qmod)D^b(Coh(𝒩))`$, then (4) and Theorem 1 would yield a morphism from the left-hand side to the right-hand side of (14). Here we say that $`\stackrel{~}{F}^{}`$ is a lifting of $`F`$ if $`FR\mathrm{\Gamma }\stackrel{~}{F}^{}`$, where $`R\mathrm{\Gamma }()=\underset{i}{}H^i()`$ for $`D^b(Coh(N))`$.
It is easy to see that such a functor $`\stackrel{~}{F}^{}`$ does not exist. A meaningful version of the argument is as follows. Let $`𝐎`$ be the differential graded algebra $`RHom_{u_q}(𝕀,𝕀)`$ (thus $`𝐎`$ is a well-defined object of the categroy of differential graded algebras with inverted quasiisomorphisms); the Ginzburg-Kumar theorem (13) shows that the cohomology algebra $`H^{}(𝐎)𝒪(𝒩)`$. Let $`DGmod(𝐎)`$ be the triangulated category of differential graded modules over $`𝐎`$ with inverted quasiisomorphisms. Let $`DDGmod(𝐎)`$ be the full subcategory of DG-modules whose cohomology is a finitely generated module over $`H^{}(𝐎)=𝒪(𝒩)`$, and let $`D_{\mathrm{}/2}D`$ be the full triangulated subcategory of DG-modules, whose cohomology is a coherent sheaf on $`𝒩`$ supported (set-theoretically) on $`𝔫`$.
We have a functor $`\stackrel{~}{F}:D^b(u_qmod)D`$ given by $`\stackrel{~}{F}:XRHom(𝕀,X)`$. It is easy to see that $`\stackrel{~}{F}`$ sends complexes of finite homological dimension over $`b_q`$ to $`D_{\mathrm{}/2}`$; and that $`\stackrel{~}{F}(𝕀)=𝐎`$. Thus, by Theorem 1, (4) provides a morphism
$$Ext^{\mathrm{}/2+}(𝕀,𝕀)Hom_{D_{D_{\mathrm{}/2}}}^{}(𝐎,𝐎).$$
One can then show that this morphism is an isomorphism; and also that the DG-algebra $`𝐎`$ is formal (quasi-isomorphic to the DG-algebra $`H^{}(𝐎)`$ with trivial differential), which implies that
$$Hom_{D_{D_{\mathrm{}/2}}}^{}(𝐎,𝐎)H_𝔫^{}(𝒩,𝒪)$$
(notice that the latter isomorphism is not compatible with homological gradings). This yields the isomorphism (14).
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# THE PROCESS 𝑒⁺𝑒⁻→𝜔𝜋⁰→𝜋⁰𝜋⁰𝛾 UP TO 1.4 GEV.
## 1 Introduction
The process of $`e^+e^{}`$ annihilation into hadrons in the 1–2 GeV energy region is an important source of information about excited states of light vector mesons $`\rho `$, $`\omega `$ and $`\varphi `$. Parameters of these states are still not well established mainly due to poor accuracy of existing experimental data. Recently new data in this energy region became available. The cross sections of the reactions $`e^+e^{}3\pi `$ and $`e^+e^{}4\pi `$ were measured at VEPP-2M collider. Accuracy was also significantly improved in the recent measurements of $`\tau 2\pi \nu _\tau `$ and $`\tau 4\pi \nu _\tau `$ spectral functions related to corresponding $`e^+e^{}`$ annihilation cross sections by the CVC hypothesis .
The reaction $`e^+e^{}\omega \pi ^0`$ is one of the dominant processes in the energy range from 1 to 2 GeV. The PDG values for the mass and width of the $`\rho ^{}(1450)`$ meson are based mainly on phenomenological studies of this process . The most precise measurements of the $`e^+e^{}\omega \pi ^0`$ cross section and $`\tau \omega \pi \nu _\tau `$ spectral function were done in . All these experiments addressed the $`4\pi `$ final state into which other intermediate states, for example $`a_1\pi `$, contribute significantly. In this work the $`e^+e^{}\omega \pi ^0`$ reaction was studied in the $`\pi ^0\pi ^0\gamma `$ final state where other contributions are much smaller. This allows to avoid systematic errors inherent to the $`e^+e^{}4\pi `$ channel due to non-trivial background subtraction, which must take into account interference effects. The process $`e^+e^{}\omega \pi ^0\pi ^0\pi ^0\gamma `$ was studied for the first time in where 20% statistical accuracy was achieved.
## 2 Detector and experiment
SND is a general purpose non-magnetic detector . Its main part is a three-layer scintillation electromagnetic calorimeter consisted of 1630 NaI(Tl) crystals with solid angle coverage about 90% of 4$`\pi `$. The energy resolution of the calorimeter for photons is $`\sigma _E/E=4.2\%/\sqrt[4]{E(\text{GeV})}`$, the angular resolution is about 1.5. The directions of charged particles are measured by two coaxial cylindrical drift chambers covering 95% of $`4\pi `$ solid angle.
The analysis presented in this work is based on data recorded in 1997–1999 in two separate energy regions: 920–980 MeV and 1040–1380 MeV. Analysis of the $`\varphi `$-resonance region (980–1040 MeV) was published earlier . In the first energy region the total integrated luminosity of 1.5 $`\text{pb}^1`$ was collected at 6 energy points. The region above the $`\varphi `$ meson was scanned with a 10 MeV step. Total integrated luminosity accumulated in this region is about 9 $`\text{pb}^1`$. The luminosity was measured with a systematic uncertainty of 3% using $`e^+e^{}e^+e^{}`$ and $`e^+e^{}\gamma \gamma `$ reactions.
## 3 Event selection
For primary selection of
$$e^+e^{}\omega \pi ^0\pi ^0\pi ^0\gamma $$
(1)
events the following criteria were applied:
* five or more photons and no charged tracks are found in an event;
* the energy deposition in the calorimeter is more than $`0.7E`$;
* the total momentum of an event measured by the calorimeter is less than $`0.15E`$;
where $`E`$ is a center of mass energy of $`e^+e^{}`$ pair. The main sources of background surviving these cuts are QED processes
$$e^+e^{}2\gamma ,\mathrm{\hspace{0.17em}3}\gamma .$$
(2)
with extra photons either from the the beam background or splitting of electromagnetic showers in the calorimeter.
For each event satisfying primary selection criteria the kinematic fitting assuming $`e^+e^{}\pi ^0\pi ^0\gamma 5\gamma `$ hypothesis was performed. As a result, two parameters: $`\chi _{\pi \pi \gamma }`$ — the $`\chi ^2`$ of the fit and $`M_{\pi \gamma }`$ — the $`\pi ^0`$ recoil mass closest to that of $`\omega `$ meson were evaluated.
The $`\chi _{\pi \pi \gamma }`$ distribution for the events with $`|M_{\pi \gamma }782|<50`$ and the $`\pi ^0`$ recoil mass spectrum for the events with $`\chi _{\pi \pi \gamma }<30`$ are plotted in figs.2 and 2. Good agreement between experimental distributions and simulation of the process (1) shows that there are no other significant contributions in the selected event sample. For final event selection the following cuts were applied:
$$\chi _{\pi \pi \gamma }<30,|M_{\pi \gamma }782|<50,$$
(3)
which reject the QED background (2) almost completely. In order to estimate residual background the experimental $`\pi ^0`$ recoil mass spectrum was fitted by the sum of simulated spectrum of the process (1) and a linear background. As a result the estimated total number of background events did not exceed 1.5% of all selected events. This value was taken as a systematic error of the measured cross section of the process (1) related to the residual background.
Experimental data were collected in 40 energy points. In the energy region under study the cross section changes slowly so it was possible to reduce the number of energy points combining the neighboring ones. Resulting energies, their standard deviations, integrated luminosities ($`IL`$) and numbers of selected events ($`N`$) are listed in the Table 1. The detection efficiency for the process (1) was determined by simulation based on a formula from which takes into account finite width of $`\omega `$-meson. In the energy range from 1050 to 1400 MeV the efficiency was found to be constant and equal to 40%. Near the $`e^+e^{}\omega \pi ^0`$ threshold the number of events with the $`\pi ^0`$ recoil mass below $`\omega `$-meson mass increases sharply and the efficiency goes down. The energy dependence of the detection efficiency is presented in the Table 1.
## 4 Fitting of the cross section
The visible cross section $`\sigma _{vis}=N/IL`$ is related with the Born cross section of $`e^+e^{}\omega \pi ^0\pi ^0\pi ^0\gamma `$ process as
$$\sigma _{vis}(E)=\epsilon (E)\sigma _0(E)(1+\delta (E)),$$
(4)
where $`\delta (E)`$ is a radiative correction calculated according to . Radiative corrections for the different energies and obtained cross section values are listed in the Table 1. The cross section at $`E=m_\varphi `$ was taken from . Only statistical errors are shown in the table. The systematic error includes the error of the luminosity measurement (3%), the detection efficiency error (4%), possible background contribution (1.5%), and the error of radiative correction (1%). The total systematic error was estimated to be 5%.
Our results in comparison with the most precise CMD-2 , CLEO , and DM2 measurements are shown in Fig.3. The cross sections from measured in the $`\pi ^+\pi ^{}\pi ^0\pi ^0`$ channel were recalculated using the PDG value of $`B(\omega \pi ^0\gamma )`$. The cross section from was obtained from the $`\tau \omega \pi \nu _\tau `$ spectral function assuming the CVC hypothesis . The CLEO results are in good agreement with ours while the CMD-2 measurements are about 10% lower, although the difference observed is smaller than the 15% systematic error quoted in . There is a significant difference between the results of DM2 and CLEO . For the cross section fitting we used our data together with the data from CLEO.
The energy dependence of the process (1) cross section was written as a sum of contributions from $`\rho (770)`$ and its excitations $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$:
$$\sigma _0(E)=\frac{4\pi \alpha ^2}{E^3}\left(\frac{g_{\rho \omega \pi }}{f_\rho }\right)^2\left|\frac{m_\rho ^2}{D_\rho }C_{\rho \omega \pi }+A_1\frac{m_\rho ^{}^2}{D_\rho ^{}}C_{\rho ^{}\omega \pi }+A_2\frac{m_{\rho ^{\prime \prime }}^2}{D_{\rho ^{\prime \prime }}}C_{\rho ^{\prime \prime }\omega \pi }\right|^2P_f(E).$$
(5)
Here $`\alpha `$ is a fine structure constant and $`g_{\rho \omega \pi }`$ is a $`\rho \omega \pi `$ coupling constant. The $`f_\rho `$ coupling constant was calculated from the $`\rho e^+e^{}`$ decay width: $`\mathrm{\Gamma }_{\rho ee}=4\pi m_\rho \alpha ^2/3f_\rho ^2`$. The expression $`m_\rho ^2/D_\rho `$ represents $`\rho `$-meson Breit-Wigner amplitude with $`D_\rho =m_\rho ^2E^2iE\mathrm{\Gamma }_\rho (E)`$, where $`m_\rho `$ and $`\mathrm{\Gamma }_\rho (E)`$ denote the $`\rho `$ meson mass and energy-dependent total width respectively. The real parameters $`A_i=g_{\rho _i\omega \pi }/g_{\rho \omega \pi }f_\rho /f_{\rho _i}`$ are the ratios of the coupling constants of different $`\rho `$ states. The factor $`P_f(E)`$ describes the energy dependence of the final state phase space. In the case of infinitely narrow $`\omega `$ resonance $`P_f(E)=1/3q_\omega ^3B(\omega \pi ^0\gamma )`$, where $`q_\omega `$ is an $`\omega `$-meson momentum. This approximation is good for the energy range above 1050 MeV, but at energies close to the $`e^+e^{}\omega \pi ^0`$ threshold it is more adequate to use precise formula taking into account finite width of the $`\omega `$ meson . The Blatt-Weisskopf factors $`C_{\rho _i\omega \pi }`$ restricting fast growth of the partial widths, were taken in the form :
$$C_{\rho \omega \pi }=\frac{1}{1+(Rq_\omega (E))^2},C_{\rho _i\omega \pi }=\frac{1+(Rq_\omega (m_{\rho _i}))^2}{1+(Rq_\omega (E))^2},i=\rho ^{},\rho ^{\prime \prime },$$
(6)
The range parameter $`R`$ was assumed to be the same for $`\rho `$, $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ mesons. The energy dependence of the $`\rho (770)`$ total width was expressed as:
$$\mathrm{\Gamma }_\rho (E)=\mathrm{\Gamma }_\rho (m_\rho )\left(\frac{m_\rho }{E}\right)^2\left(\frac{q_\pi (E)}{q_\pi (m_\rho )}\right)^3C_{\rho \pi \pi }^2+\frac{g_{\rho \omega \pi }^2}{12\pi }q_\omega ^3(E)C_{\rho \omega \pi }^2,$$
(7)
where $`q_\pi `$ is the pion momentum, $`C_{\rho \pi \pi }`$ is the Blatt-Weisskopf factor for $`\rho \pi ^+\pi ^{}`$ decay :
$$C_{\rho \pi \pi }(E)=\sqrt{\frac{1+(Rq_\pi (m_\rho ))^2}{1+(Rq_\pi (E))^2}}.$$
(8)
There is no generally accepted description of the shapes of the broad excited states $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$. In constant total width were assumed, while in the total widths varied with energy as a sum of the partial widths of all main decay modes.
We considered the both approaches in order to understand model dependence of the fit parameters.
The fit results obtained for three classes of models are listed in the Table 2. The fit parameters are sensitive to the variation of the range parameters R in the Blatt-Weisskopf factors. In the Table 2 we show the intervals of the fit parameter variation when $`R`$ ranges from 0 to 2 GeV<sup>-1</sup>. The typical errors of the parameters obtained for each specific model are listed in the fourth column of the Table 2.
In the model 1 the constant widths of $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ were assumed. The mass and width of the $`\rho ^{\prime \prime }`$ meson were fixed to their PDG values : $`m_{\rho ^{\prime \prime }}=1700\text{ MeV}`$, $`\mathrm{\Gamma }\rho ^{\prime \prime }=235\text{ MeV}`$. The obtained value of the $`\rho ^{\prime \prime }`$ amplitude $`A_2`$ is compatible with zero, so the final fit result for the model 1 is given for $`A_2`$ fixed to zero. Let us discuss the choice of an upper boundary for the $`R`$ parameter. Small values $`g_{\rho \omega \pi }<13`$ GeV<sup>-1</sup> found for $`R>2`$ GeV<sup>-1</sup> are in conflict with the QCD sum rules estimation 16 GeV<sup>-1</sup> and experimental value of 14.4 GeV<sup>-1</sup>, obtained from the $`\omega 3\pi `$ decay width assuming that $`\omega \rho \pi `$ mechanism dominates in this decay. Only for $`R0`$ the extracted mass and width of the $`\rho ^{}`$ meson are compatible with their PDG values. However in this case the fit yields the lowest confidence level of $`P(\chi ^2)=3\%`$.
Models 2 and 3 take into account energy dependence of the total widths of $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ mesons. Since branching ratios of $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ decays are practically unknown the energy dependence of the sum of all rapidly growing partial widths of their multihadron decays was approximated as the energy dependence of $`\rho _i\omega \pi `$ width:
$`\mathrm{\Gamma }_{\rho _i}(E)`$ $`=`$ $`\mathrm{\Gamma }_{\rho _i}(m_{\rho _i})[(1B_{\rho _i\pi \pi })\left({\displaystyle \frac{q_\omega (E)}{q_\omega (m_{\rho _i})}}\right)^3C_{\rho _i\omega \pi }^2`$ (9)
$`+`$ $`B_{\rho _i\pi \pi }\left({\displaystyle \frac{m_{\rho _i}}{E}}\right)^2\left({\displaystyle \frac{q_\pi (E)}{q_\pi (m_{\rho _i})}}\right)^3C_{\rho _i\pi \pi }^2],`$
For the $`\rho ^{\prime \prime }2\pi `$ branching fraction we use the theoretical estimation: $`B_{\rho ^{\prime \prime }\pi \pi }=10\%`$ . For the $`\rho ^{}`$ meson the value $`B_{\rho ^{}\pi \pi }=50\%`$ was chosen which, we think, reflects experimental situation more correctly than the theoretical prediction $`25\%`$ . The only difference between the models 1 and 2 which both consider only one excited $`\rho `$ state is that the model 2 assumes its total width energy dependent. But the fit results for these two models are quite different. Particularly the mass found for the model 2 is close to the PDG value of the $`\rho ^{\prime \prime }`$-meson mass. On the other hand there is a definite signal of the $`\rho ^{}`$ meson with a mass of 1320–1400 MeV and 400–500 MeV width in the pion form factor data . Therefore we also considered the model 3 with two excited $`\rho `$ states, in which the mass and width of the $`\rho ^{}`$ meson were fixed to $`m_\rho ^{}=1400\text{ MeV}`$ and $`\mathrm{\Gamma }_\rho ^{}=500\text{ MeV}`$. This model yields the best fit to the experimental data: $`P(\chi ^2)13\%`$. Fitting curves corresponding to the models 1 and 3 with $`R=0`$ are shown in figure 3.
The main conclusions from the analysis of the fit results are the following. Fitting of the same experimental data by models with fixed and energy-dependent total widths of the excited states yields quite different parameters of these states. This is caused by strong energy dependence of the phase space for the main decay modes of $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ mesons and this effect should be taken into account in the fitting of experimental data. Satisfactory description of the experimental cross section was obtained in the model with two excited states with the masses $`m_\rho ^{}=1400`$ MeV and $`m_{\rho ^{\prime \prime }}1600`$ MeV in which contribution of the higher state dominates. However this result contradicts the theoretical expectation , where $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ are considered as $`2S`$ and $`1D`$ $`q\overline{q}`$ states respectively and the larger contribution of the lower $`2S`$ excitation was predicted. Thus, with the new experimental data the situation in the isovector sector remains unclear. The main problem for data analysis is the absence of consistent phenomenological description of the shapes of broad resonances with strong energy dependence of partial widths.
## 5 Summary
In this work the cross section of the $`e^+e^{}\omega \pi ^0\pi ^0\pi ^0\gamma `$ reaction was measured from the threshold up to 1.4 GeV with a 5% systematic accuracy. This is the most precise measurement of the $`e^+e^{}\omega \pi ^0`$ cross section in this energy range. Our data are in a good agreement with the CLEO measurement of $`\tau \omega \pi \nu _\tau `$ spectral function . The combined fit to our and CLEO data was performed in the vector meson dominance model taking into account the contributions of $`\rho `$, $`\rho ^{}`$ and $`\rho ^{\prime \prime }`$ states. The experimental data can be reasonably well approximated assuming existence of only one excited state but its mass is strongly model-dependent and varies from 1460 to 1700 MeV for different descriptions of the resonance shape. The best fit to the data was obtained for the model with two excited states in which the higher state with $`m_{\rho ^{\prime \prime }}1600`$ MeV dominates.
This work is supported in part by the Russian Fund for Basic Researches (grants No. 99-02-16815, 99-02-17155) and STP “Integration” (grant No. 274).
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# Flavor Gauge Bosons at the Tevatron
## 1 Introduction
The origin of the Standard Model’s (SM) familiar $`SU(3)_c\times SU(2)_L\times U(1)_Y`$ gauge symmetry remains theoretically unclear. In the limit where we neglect all gauge interactions and fermion masses, the fermion sector of the model possesses a large SU(45) global symmetry corresponding to the fact that in this limit there are 45 chiral fermion fields that are indistinguishable. The gauge interactions of the SM are by necessity subgroups of this maximal symmetry but in principle a larger subgroup of this symmetry might be gauged and broken to the SM groups at high energies.
Such gauged flavor symmetries have been invoked in a number of scenarios to play a role in the dynamical generation of fermion masses. For example they may play the part of extended technicolor interactions in technicolor models or top condensation models , feeding the electroweak symmetry (EWS) breaking fermion condensate down to provide masses for the lighter standard model fermions. Strongly interacting flavor gauge interactions may also be responsible for the condensation of the fermions directly involved in EWS breaking. For example, top condensation has been postulated to result from a Topcolor gauge group and in the model of from family gauge interactions. There has been renewed interest in these models recently with the realization that variants, in which the top mixes with singlet quarks, can give rise to both the EW scale and an acceptable top mass via a seesaw mass spectrum . These top seesaw models have the added benefit of a decoupling limit which allows the presence of the singlet fields to be suppressed in precision EW measurements bringing these dynamical models in line with the data. Flavor universal variants of the top-seesaw idea have been proposed in Ref. , where the dynamics is driven by family or large flavor gauge symmetries.
The naive gauging of flavor symmetries at low scales (of order a few TeV) often gives rise to unacceptably large flavor changing neutral currents (FCNC) since gauge and mass eigenstates need not coincide. For instance, gauge symmetries that give rise to direct contributions to $`K^0\overline{K}^0`$ mixing are typically constrained to lie above 500 TeV in mass scale. There are, however, models that survive these constraints. Gauge groups that only act on the third family are less experimentally constrained - Topcolor is such an example. Models in which the chiral flavor symmetries of the SM fermions are gauged preserving the SM $`U(3)^5`$ flavor symmetry can respect the SM GIM mechanism and do not give rise to tree level FCNCs . In addition, there are also strong constraints on gauged flavor models where the dynamics responsible for the breaking of the flavor symmetry does not respect custodial isospin . We shall restrict ourselves to models where the top mass is the sole source of custodial isospin breaking. In particular we we will study a model where the SU(3) chiral family symmetry of the quarks is gauged and another where the full SU(9) family-color multiplicity of the quarks is gauged, corresponding to the models of . In the spirit of these models it is also interesting to consider chiral flavor symmetries that include the leptons which might be expected to give interesting contributions to Drell-Yan production. The obvious extension has a gauged SU(12) flavor symmetry but we show in the final section that an analysis of low energy atomic parity violation experiments places constraints on the gauge bosons of such models of order 10 TeV and they are thus outside the reach of the Tevatron.
Since these new flavor interactions may exist at relatively low scales (a few TeV) and may play an integral part in either EWS breaking or fermion mass generation it is interesting to study current experimental bounds on the corresponding gauge bosons. In a previous paper we investigated the limits from Z-pole precision measurements . Although the limits obtained vary across models, the typical lower bound on the mass scale is $`2`$ TeV. Here we study the potential of direct searches at the Fermilab Tevatron collider. In particular we study effects in dijet production (in the spirit of the analysis in ) and single top production. When possible, we first establish bounds from the existing Run I data (they are typically 1-2 TeV). We then project the sensitivity of the Tevatron in Run II and show the bounds are more than competitive with the precision data bounds. If these gauge symmetries do have a role to play in EWS breaking then they must presumably be broken at scales close to the EW scale and these bounds therefore represent a significant probe of the interesting parameter space.
## 2 Constraints on Models
We present three models of flavoron physics. While this list is not exhaustive, we believe these examples cover a broad range of signals at the Tevatron collider. In what follows, only the couplings to standard model fermions will be specified. Explicit models include additional fermions, necessary for either flavor gauge symmetry breaking and/or anomaly cancellation, which typically have masses of order of the flavor gauge symmetry breaking scale.
### 2.1 Chiral Quark Family Symmetry
The gauging of the chiral family symmetry of the left handed quarks has been motivated in technicolor , top condensate and flavor universal see saw models . The minimal representative model has a gauged SU(3) family symmetry, in addition to the SM interactions, acting on the three left handed quark<sup>1</sup><sup>1</sup>1One can also imagine the same symmetry acting on leptons . Here we only consider the quarks since they lead to signals at hadron colliders. doublets $`Q=((t,b)_L^i,(c,s)_L^i,(u,d)_L^i)`$ where $`i`$ is a QCD index which commutes with the family symmetry . We assume that some massive sector completely breaks the SU(3) family gauge group to an global SU(3) family symmetry, giving the family gauge bosons (“familons”) masses of order $`M_F=g_FV`$ where $`V`$ is the mass scale associated with the symmetry breaking. There is no mixing between the flavor and standard model gauge bosons. Note that with this gauge symmetry and symmetry breaking pattern, the (approximate) SM $`U(3)^5`$ global symmetry responsible for the GIM mechanism remains and the model is free of tree level FCNCs . The interactions of the massive flavorons are summarized by the couplings
$$=ig_FA^{\mu a}\overline{Q}\gamma _\mu T^aQ,$$
(1)
where $`T^a`$ are the generators of SU(3) symmetry acting on the three families of left-handed quarks.
The SU(3) coupling $`g_F`$ cannot be too large or this interaction will cause a chiral symmetry breaking condensate between the left-handed ordinary fermions and right-handed fermions which must be present in the theory to eliminate gauge anomalies. This would result in TeV-scale fermion masses and a scale for electroweak symmetry breaking which is too high. We may estimate the upper bound on $`g_F`$ by approximating, at low energies, the interactions of the massive flavor gauge bosons by a Nambu–Jona-Lasinio (NJL) model with the four-fermion interaction
$$_{\mathrm{eff}}=\frac{2\pi \kappa _F}{M_F^2}\left(\underset{f}{}\overline{Q}\gamma _\mu T^aQ\right)^2,$$
(2)
where $`\kappa g_F^2/4\pi `$. Applying the usual NJL analysis<sup>2</sup><sup>2</sup>2Note that, defining the theory in terms of a momentum-space cutoff $`\mathrm{\Lambda }`$, a four fermion interaction $`G\overline{\psi }\psi \overline{\psi }\psi `$ has a critical coupling $`G_c=2\pi ^2/\mathrm{\Lambda }^2`$ . , we see that $`\kappa _F`$ cannot exceed
$$\kappa _{crit}=\frac{2N\pi }{(N^21)}=2.36,$$
(3)
where $`N=3`$ for chiral quark flavor symmetry.
In Ref. we obtained bounds on flavor gauge boson masses from electroweak precision measurements. The lower bound obtained for a critically coupled familon is $`M_F>1.9`$ TeV, at $`95\%`$ C.L. Here we will investigate the reach of direct searches. First, we consider the bounds from the existing Tevatron data. As is the case for the universal coloron model , stringent limits will come from the study of the angular behavior of the dijet cross section . The contributions arising in the chiral quark family model are the consequence of the exchange of the familon gauge boson in the various possible channels. The resulting modification of the quark scattering matrix elements are given in Section A.2 of the Appendix.
In Fig.1 we plot the ratio of the dijet mass distribution for $`|\eta |<0.5`$ to the mass distribution with $`0.5<|\eta |<1.0`$, with $`\eta `$ the jet pseudo-rapidity. This ratio, as noted for instance in Refs., is very sensitive to new physics producing effects concentrated in the central region, and in general affecting the angular distribution of dijets. Also it is expected that in this ratio there is a large cancellation of uncertainties coming from softer QCD effects. The data points are from the D0 data in Ref., and the error bars show the statistical and systematic errors added in quadrature. The histogram corresponds to the QCD prediction, obtained to next-to-leading order (NLO) with the use of JETRAD (see for details). The familon contribution is known only at leading order (LO). Thus, in order to estimate their NLO dijet spectrum, we compute the fractional excess with respect to LO QCD and then multiply it by the NLO QCD result. We consider various familon masses, with the coupling set at its critical value.
In order to obtain a lower mass limit we follow the procedure described in Ref.. We construct the Gaussian likelihood function
$$P(x)=\frac{1}{2\pi ^2det(S)}\mathrm{exp}\left(\frac{1}{2}[dt(x)]^TS^1[dt(x)]\right),$$
(4)
where the vector $`d`$ contains the data points in the various mass bins, $`t(x)`$ is the vector of theoretical predictions for a given mass and coupling $`x=\kappa _F/M_F^2`$; and $`S`$ is the covariant matrix. To obtain $`95\%`$ confidence level limits, we require
$$Q(x_{\mathrm{max}})_0^{x_{\mathrm{max}}}P(x)𝑑x=0.95Q(\mathrm{}),$$
(5)
with $`x_{\mathrm{max}}`$ the value defining the mass bound. Making use of the the Run I data we then obtain mass bounds for the familon
$$M_F>1.55\mathrm{TeV},95\%\mathrm{C}.\mathrm{L}.,$$
(6)
where we have considered a critically coupled familon. This is consistent with, but somewhat weaker than the $`95\%`$ C.L. limit obtained in Ref., $`M_F>1.9`$ TeV at critical coupling.
During Run II however, measurements of the dijet spectrum at an upgraded Tevatron will yield bounds better than those derived from Z-pole observables. For instance if we consider the nominal luminosity of $`2fb^1`$, and assume a $`30\%`$ reduction in the systematic errors, the bound on the familon mass for Run II would be $`M_F>2.2`$ TeV. An extended Tevatron run or the achievement of higher intensities could therefore result in a mass reach well above that of electroweak precision measurements and cover a large fraction of the interesting parameter space of this model.
In addition to the dijet signal, the chiral quark family model leads to another potentially interesting signal at hadron colliders: anomalous single top production. This occurs due to the existence of non-diagonal couplings to the family gauge bosons. Although these do not lead to $`|\mathrm{\Delta }S|=2`$ signals, because of GIM cancellation, there are flavor changing couplings of quarks. The fact that the family symmetry commutes with $`SU(2)_L`$ implies that there will be tree level familon exchanges such as $`d\overline{b}u\overline{t}`$, where “family number” is preserved. The diagrams relevant for single top production at the Tevatron are s-channel $`d\overline{b}u\overline{t}`$, and t-channel $`u\overline{d}t\overline{b}`$ (dominant) and $`u\overline{b}t\overline{d}`$. Other diagrams also are obtained by the replacements $`ds`$ and $`uc`$. For instance, the s-channel matrix element squared is
$$|(d\overline{b}u\overline{t})|^2=(4\pi )^2\kappa ^2u(um_t^2)\left|\frac{1}{2}P_s\right|^2,$$
(7)
Neglecting $`m_b`$, the t-channel contributions are obtained by replacing $`P_s`$ by the $`P_t`$, where $`P_s`$ and $`P_t`$ are the familon propagators in the corresponding channel as defined in (A.5).
If the coupling is close to critical, these processes will generate important contributions to the single top production cross section. In Fig. 2 we show the familon induced single top production cross section at $`\sqrt{s}=1.8`$ TeV as a function of the familon mass. The horizontal line is the $`95\%`$ C.L. upper limit on single top production as obtained by the CDF collaboration . The most constraining bound, $`\sigma (p\overline{p}tX)<15.4`$ pb translates into the familon mass bound
$$M_F>1.02\mathrm{TeV}95\%\mathrm{C}.\mathrm{L}..$$
(8)
This is somewhat weaker than the bound (6) obtained from the Run I dijet data, but may be improved if a study exploiting the kinematic differences between the SM and the flavoron signals is undertaken.
In Run II, the Tevatron will be sensitive to the SM single top production via $`W`$-gluon fusion as well as the s-channel $`W^{}`$ exchange. The latter process can be separated from the former by making use of double b-tagging, since the b quark produced in association with the top is hard, unlike in $`W`$-gluon fusion. In order to estimate the sensitivity of the Tevatron in Run II to the flavoron contribution to single top production, we take only the dominant flavoron diagram, t-channel mediated $`u\overline{d}t\overline{b}`$. We compare this contribution to the s-channel SM assuming these will be separately observed with the use of double b-tagging . In Fig. 3 we show the $`p_T`$ distribution of the $`b`$ quark produced in association with the top quark for t-channel familon exchange and s-channel $`W^{}`$ exchange. We see that, for example, for $`M_F=2`$TeV the total ($`t\overline{b}+\overline{t}b`$) cross section is about $`50\%`$ larger from familon exchange than in the SM, with the added feature that the $`p_T`$ distribution is harder. We conclude that the sensitivity of Run II could go as far as $`(22.5)`$ TeV for 2 fb<sup>-1</sup>, or perhaps higher depending on the sensitivity to be achieved to the SM s-channel process. Thus, anomalous single top production could be the most constraining channel on the $`SU(3)`$ chiral quark model in Run II at the Tevatron.
### 2.2 SU(9) Chiral Flavor Symmetry
We next consider a natural extension of gauging the quark family symmetry, gauging the full SU(9) symmetry of both the color and family multiplicity of the left handed quarks. Such a symmetry can be implemented as an extended technicolor gauge symmetry (in the spirit of ) or in quark universal seesaw models (as in ). The SU(9) symmetry commutes with the standard weak $`SU(2)_L`$ gauge group and acts on the left handed quarks
$$Q_L=((t,b)^r,(t,b)^b,(t,b)^g,(c,s)^r,\mathrm{}(u,d)^g)_L$$
(9)
with $`r,g,b`$ the three QCD colors. The quark couplings to the SU(9) gauge bosons is given by
$$=ig_FB^{a\mu }\overline{Q}_L\mathrm{\Lambda }^a\gamma _\mu Q_L,$$
(10)
with $`\mathrm{\Lambda }^a`$ the generators of SU(9). These include
$$\frac{1}{\sqrt{3}}\left(\begin{array}{ccc}T^a& 0& 0\\ 0& T^a& 0\\ 0& 0& T^a\end{array}\right),\frac{1}{\sqrt{6}}\left(\begin{array}{ccc}T^a& 0& 0\\ 0& T^a& 0\\ 0& 0& 2T^a\end{array}\right),\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}T^a& 0& 0\\ 0& T^a& 0\\ 0& 0& 0\end{array}\right),$$
(11)
where $`T^a`$ are the 8 3x3 QCD generators. SU(9) further contains
$$\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}0& T^a& 0\\ T^a& 0& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{12}}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}0& iT^a& 0\\ iT^a& 0& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{12}}\left(\begin{array}{ccc}0& i& 0\\ i& 0& 0\\ 0& 0& 0\end{array}\right)$$
(12)
plus the two other similar sets mixing the remaining families. Finally there are two diagonal generators
$$\frac{1}{\sqrt{12}}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{36}}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right)$$
(13)
The model must also contain interactions which give rise to color for the right handed quarks. For this reason, we include an $`SU(3)_{pc}`$ proto-color group that acts on the right handed quarks, which will be combined with the $`SU(3)_C`$ subgroup of $`SU(9)_L`$ to yield ordinary color. We normalize the proto-color gauge bosons couplings such that they have the same generators as the SU(9) bosons
$$=\frac{i}{\sqrt{3}}g_{pc}A^{\mu a}\overline{q}_R\gamma _\mu T^aq_R.$$
(14)
At the flavor breaking scale we assume some massive sector breaks the $`SU(9)_L\times SU(3)_{pc}`$ gauge symmetry down to ordinary color $`SU(3)_C`$ and a global $`SU(3)_F`$ group acting on the three families of quarks. The global $`SU(3)_F`$ symmetry is sufficient to insure the absence of tree-level FCNCs .
For simplicity, we will assume the symmetry breaking sector has an $`SU(9)_L\times SU(9)_{flavor/color}`$ chiral flavor symmetry, under which the symmetry breaking vev transforms as a $`(9,\overline{9})`$. The majority of the SU(9) gauge bosons will then have mass $`M_F=g_FV`$. Eight of the $`SU(9)_L`$ gauge bosons mix with the right handed proto-color group, giving rise to ordinary color and eight massive gluons. The proto-gluons and color-octet flavorons mix through the mass matrix
$$(A^\mu ,B^\mu )\left(\begin{array}{cc}g_{pc}^2& g_{pc}g_F\\ g_{pc}g_F& g_F^2\end{array}\right)V^2\left(\begin{array}{c}A_\mu \\ B_\mu \end{array}\right)$$
(15)
which diagonalizes to
$$(X^\mu ,G^\mu )\left(\begin{array}{cc}g_{pc}^2+g_F^2& 0\\ 0& 0\end{array}\right)V^2\left(\begin{array}{c}X_\mu \\ G_\mu \end{array}\right)$$
(16)
where
$$\left(\begin{array}{c}A^\mu \\ B^\mu \end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\varphi & \mathrm{sin}\varphi \\ \mathrm{sin}\varphi & \mathrm{cos}\varphi \end{array}\right)\left(\begin{array}{c}G_\mu \\ X_\mu \end{array}\right)$$
(17)
with
$$\mathrm{sin}\varphi =\frac{g_{pc}}{\sqrt{g_{pc}^2+g_F^2}},\mathrm{cos}\varphi =\frac{g_F}{\sqrt{g_{pc}^2+g_F^2}},$$
(18)
and $`G^\mu `$ and $`X^\mu `$ are the gluon and color-octet flavoron respectively.
The low energy QCD coupling, with the standard generator normalization is given by
$$g_c=\frac{g_Fg_{pc}}{\sqrt{3(g_{pc}^2+g_F^2)}}$$
(19)
which implies that $`\kappa _F3\alpha _s(2\mathrm{TeV})`$. The interactions of the SM fermions with the massive color octet (with mass $`M_F^{}=\sqrt{g_{pc}^2+g_F^2}V=M_F/c_\varphi `$) are given by
$$g_c\mathrm{tan}\varphi X^{a\mu }\overline{q}_R\gamma _\mu T^aq_R+g_c\mathrm{cot}\varphi X^{a\mu }\overline{q}_L\gamma _\mu T^aq_L,$$
(20)
where $`\mathrm{cot}\varphi =g_F/g_{pc}`$.
As in the case of $`SU(3)_F`$, the coupling $`g_F`$ cannot be too large, or it would likely induce an EWS breaking condensate at the flavor scale. Assuming that at low energies the massive gauge boson interactions with the SM fermions can be approximated by a NJL model (ignoring the effects of the mixing of eight of the generators with proto-color in this estimate), then the critical coupling for chiral symmetry breaking in that approximation is
$$\kappa _{crit}=\frac{2N\pi }{(N^21)}=0.71.$$
(21)
As was the case in the previous two models, the most conspicuous signals are in the dijet spectrum. In the Appendix A.3 we list all the relevant matrix elements for dijet production due to the two gauge bosons with masses $`M_F`$ and $`M_F^{}`$.
In Fig.4 we plot the contributions of these gauge bosons to the cross section ratio as a function of their mass, assuming for simplicity $`M_F=M_F^{}`$. Although, in principle, one could expect the effect to be smaller than for the $`SU(3)`$ chiral familon due to the fact that the critical coupling in eqn.(21) is considerably smaller than that of the $`SU(3)`$ case, the $`SU(9)`$ flavorons contribute to a large number of diagrams leading to dijets. In fact, as can be seen in Fig.4, the effects will be stronger in this model. We follow the procedure described earlier to obtain a mass constraint from the Tevatron Run I data. The mass of the $`SU(9)`$ flavorons is bounded by
$$M_F>1.9\mathrm{TeV},95\%\mathrm{C}.\mathrm{L}.$$
(22)
This is very similar to the $`95\%\mathrm{C}.\mathrm{L}.`$ limit obtained in Ref. from electroweak precision measurements. On the other hand, at $`\sqrt{s}=2`$TeV and with an integrated luminosity of 2 fb<sup>-1</sup>, Run II at the Tevatron will put a limit of $`M_F>2.7\mathrm{TeV}`$, where we assume a 30% reduction in systematic errors. This covers a large fraction of the interesting parameter space of this model.
Just as in the chiral quark family model, in the $`SU(9)`$ model there are also important contributions to anomalous single top production. The fact that some of the $`SU(9)`$ gauge bosons carry color tends to enhance the interactions when compared to the $`SU(3)`$ chiral quark model. On the other hand, the critical coupling in this model is considerably smaller than that in the $`SU(3)`$ case, as can be seen by comparing eqns. (21) with (3). The net effect is a reduction in the single top signal shown in Fig 3, by a factor of
$$\left(\frac{\kappa _F^{SU(9)}}{\kappa _F^{SU(3)}}\right)^2\times \left(\frac{14}{9}\right)0.15$$
(23)
at critical coupling. Since the cross section falls approximately as $`1/M_F^4`$, this will result in a familon mass bound that is smaller than the one to be obtained in the $`SU(3)`$ model by a factor of about $`\sqrt[4]{0.15}0.60`$. Thus, since our expectations for Run II in the single top channel in the $`SU(3)`$ model put the reach somewhere around $`M_F>(22.5)`$TeV, we conclude that the reach of this channel for the $`SU(9)`$ flavoron is still below the Run I mass limit eqn. (22) that we extracted from the dijet data. Although more detailed studies of the single top signal (for instance including all possible single top final states) are possible, we can safely conclude that this channel will not be competitive with the dijet signal in the $`SU(9)`$ model at the Tevatron.
### 2.3 SU(12) Chiral Flavor Symmetry
The final model we consider is one in which we gauge the full SU(12) flavor symmetry of all the left handed SM fermion doublets
$$Q_L=((t,b)^r,(t,b)^b,(t,b)^g,(\nu _\tau ,\tau ),(c,s)^r,\mathrm{}(\nu _e,e))_L.$$
(24)
This is similar to the $`SU(9)`$ model, but it also includes a proto-hypercharge interaction that, after the $`SU(12)`$ breaking, gives rise to the SM $`U(1)_Y`$. The flavor gauge interactions act as
$$=ig_FB^{a\mu }\overline{Q}_L\mathrm{\Lambda }^a\gamma _\mu Q_L,$$
(25)
with $`\mathrm{\Lambda }^a`$ the generators of SU(12), which may be conveniently broken down into the following groupings
$$\frac{1}{\sqrt{3}}\left(\begin{array}{ccc}P^a& 0& 0\\ 0& P^a& 0\\ 0& 0& P^a\end{array}\right),\frac{1}{\sqrt{6}}\left(\begin{array}{ccc}P^a& 0& 0\\ 0& P^a& 0\\ 0& 0& 2P^a\end{array}\right),\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}P^a& 0& 0\\ 0& P^a& 0\\ 0& 0& 0\end{array}\right)$$
(26)
where $`P^a`$ are the 15 4x4 Pati-Salam generators consisting of 8 3x3 blocks that are QCD, 6 step operators between the quarks and leptons and the diagonal generator $`1/\sqrt{24}`$ diag$`(1,1,1,3)`$. SU(12) further contains
$$\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}0& P^a& 0\\ P^a& 0& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{16}}\left(\begin{array}{ccc}0& 1& 0\\ 1& 0& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{2}}\left(\begin{array}{ccc}0& iP^a& 0\\ iP^a& 0& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{16}}\left(\begin{array}{ccc}0& i& 0\\ i& 0& 0\\ 0& 0& 0\end{array}\right)$$
(27)
plus the two other similar sets mixing the remaining families. Finally there are two diagonal generators
$$\frac{1}{\sqrt{16}}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right),\frac{1}{\sqrt{48}}\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 2\end{array}\right)$$
(28)
In order to ensure the SM gauge groups emerge at low energies we must again introduce a proto-color group as in the SU(9) model above. The first 8 generators of SU(12) in (26) are the same as those in the SU(9) model (11) and hence the discussion of the mixing between the proto-color and the 8 SU(12) gauge bosons follows the discussion in the SU(9) model exactly. In addition, in the SU(12) model we must also include a proto-hypercolor gauge boson. Since the Pati-Salam diagonal generator in the first set of generators in (26) is the traditional generator for the hypercharge boson’s coupling to left handed fermions, the proto-hypercharge gauge boson only has to couple to the right handed fermions. The result of the mixing of these two gauge bosons is the massless SM hypercharge gauge boson plus a massive gauge boson coupling to both left and right handed fermions.
If the interactions of flavoron gauge bosons in Eq. (25) at low energies can be modeled by a NJL lagrangian with coupling $`4\pi \kappa /2!M_F^2`$, the critical coupling for chiral symmetry breaking is calculated to be
$$\kappa _{crit}=\frac{2N\pi }{(N^21)}=0.53,$$
(29)
somewhat smaller than in the $`SU(9)`$ case in Eq. (21). Note that combined with the lower constraint from the ability to reproduce the QCD coupling ($`\kappa _F3\alpha _s(2\mathrm{TeV})0.3`$) there is a relatively small window of allowed couplings.
Although this model results in various signals at the Tevatron — such as quark scatterings similar to those of the $`SU(9)`$ model as well as anomalous contributions to Drell-Yan production arising from the flavoron couplings to leptons — the energy scale of this scenario is severely constrained by data from experiments of Atomic Parity Violation (APV) in Cesium. The parity-violating part of the electron-nucleon interaction can be written as
$$_{eq}=\frac{G_F}{\sqrt{2}}\underset{q=u,d}{}\{C_{1q}(\overline{e}\gamma _\mu \gamma _5e)(\overline{q}\gamma ^\mu q)+C_{2q}(\overline{e}\gamma _\mu e)(\overline{q}\gamma ^\mu \gamma _5q)\},$$
(30)
where the coefficients $`C_{1q}`$ and $`C_{2q}`$ are given in the SM by
$$C_{1q}^{SM}=(T_3^q2Q_q\mathrm{sin}^2\theta ),C_{2q}^{SM}=T_3^q(14\mathrm{sin}^2\theta ),$$
(31)
and $`T_3^q`$ is the third component of the quark isospin. The atomic weak charge is then defined as
$$Q_W=2\{C_{1u}(2Z+N)+C_{1d}(N+2Z)\},$$
(32)
with $`Z`$ and $`N`$ the number of protons and neutrons respectively. The APV experiment finds the atomic charge of Cesium to be $`Q_W=72.06\pm 0.28\pm 0.34`$, whereas the SM prediction is $`Q_W=73.09\pm 0.03`$. This translates into a deviation from the SM prediction of
$$\mathrm{\Delta }Q_W=1.33\pm 0.44.$$
(33)
We can write the deviations of $`Q_W`$ as
$$\mathrm{\Delta }Q_W=376\mathrm{\Delta }C_{1u}422\mathrm{\Delta }C_{1d}.$$
(34)
The SU(12) model gives rise to various contributions to $`\mathrm{\Delta }Q_W`$. However, by far the largest of these corresponds to a step operator from the generators in Eq. (26) which connect quarks to leptons. These result in the non-diagonal effective coupling
$$\frac{g_F^2}{8M_F^2}(\overline{e}_L\gamma _\mu d_L)(\overline{d}_L\gamma ^\mu e_L).$$
(35)
After Fierzing and decomposing into the proper vector and axial pieces, the contribution in (35) gives rise to an effect in the weak charge of Cesium given by
$$\mathrm{\Delta }Q_W^F=80.4\kappa _F\frac{(1\mathrm{TeV})^2}{M_F^2}=42.6\frac{(1\mathrm{TeV})^2}{M_F^2},$$
(36)
where $`M_F`$ is understood to be measured in TeV, and the last equality is obtained by using $`\kappa _F=\kappa _{\mathrm{crit}.}`$ as defined in (29). Thus not only is this a large contribution to $`Q_W(Cs)`$, but it also has the opposite sign of Eq. (33). For instance, the $`3\sigma `$ bound would be $`M_F>12`$ TeV. More conservatively, we can estimate the sensitivity of the APV measurement by taking the error in Eq. (33) as the possible size of the effect. This translates into $`M_F>9.8`$TeV. From the model building point of view this is an undesirably large mass scale and raises the issue of fine-tuning. In any event, it is clear that the APV experiment forces the mass scale in the $`SU(12)`$ model to be very high and out of reach of the Tevatron.
Finally, we point out that the constraint on the $`SU(12)`$ model resulting from Eq. (36) is more general since it cannot be completely evaded by lowering the coupling below its critical value. As we mentioned earlier, in order to obtain the correct QCD coupling, $`\kappa _F`$ must satisfy $`\kappa _F3\alpha _s(2\mathrm{TeV})`$. Then, its minimum value of approximately $`0.3`$ translates into the bound $`M_F>7.4`$TeV.
## 3 Conclusions
We have studied the Tevatron collider bounds on two models of broken, gauged, chiral flavor symmetries; an SU(3) chiral family symmetry and an SU(9) chiral flavor symmetry of the SM quarks. These symmetries have been proposed as playing a significant role in theories of EWS breaking and fermion mass generation and are blessed with a GIM mechanism that suppresses FCNCs allowing the gauge bosons to be relatively light. The strongest Tevatron signals result in dijet production and single top production. We summarize the current limits, from precision data and Run I, on the critically coupled gauge boson masses in Table 1 - they are comparable. The Run II expectations are also displayed and should become the leading constraints on the models.
In the $`SU(3)`$ model both, dijet and anomalous single top production, are likely to be important signals. On the other hand, in the $`SU(9)`$ model the dijet cross section receives a large enhancement due to the fact that some of the flavor gauge bosons carry color, resulting in more diagrams contributing (see Appendix A.3). However, since the critical coupling is considerably smaller than in the $`SU(3)`$ case, the single top signal – even after taking into account the color enhancement – is reduced. Thus, the single top channel is crucial in order to separate these two models as the possible origin of a hypothetical deviation in the dijet sample.
For comparison we also display in Table I the equivalent limits for the Universal Coloron model of \- in this model the chiral $`SU(3)_L\times SU(3)_R`$ color group of the quarks is gauged and broken to the QCD group leaving axially coupling massive colorons. This model is considerably more strongly constrained in part because of its large critical coupling and because the dijet channel is a particularly good probe of extra color like interactions. It is notable that in the models we have explored the gauge bosons are potentially lighter, as one might hope if they played a role in EWS breaking, and that the Tevatron can hope to probe interesting regions of parameter space.
Finally we have pointed out a further low energy precision constraint on models where the flavor symmetry is enlarged to include the lepton sector. In particular an SU(12) gauged chiral flavor model gives contributions in low energy atomic parity violation experiments that place the bound on the gauge boson masses out of the Tevatron’s reach.
xxxx
| | EPM | Run I | Run II |
| --- | --- | --- | --- |
| Universal Coloron | 3 | 4.3 | 7 |
| $`SU(3)_F`$ | 1.9 | 1.55 | 2.5 (single top) |
| $`SU(9)_F`$ | 1.9 | 1.9 | 2.7 |
| $`SU(12)_F`$ | 10 (APV) | No reach | No reach |
Table I: The $`95\%`$ C.L. bounds (or sensitivity) on the models discussed. The numbers correspond to the mass of the gauge bosons in TeV if its coupling is critical. The first column comes from electroweak precision measurements and is taken from Ref. . The Run I bounds as well as the Run II sensitivities (for $`2fb^1`$) summarize our results. They come from the dijet analysis, with the exception of the Run II reach for the $`SU(3)`$ chiral quark model which comes from single top production.
xxxx
Acknowledgments
This work was supported in part by the Department of Energy under grants DE-FG02-91ER40676 and DE-FG02-95ER40896. N.E. is grateful for the support of a PPARC Advanced Research Fellowship. G.B. acknowledges the hospitality of the High Energy Physics Group at the University of Sao Paulo, where part of this work was completed.
## Appendix: Cross Sections
We present some standard tree-level expressions for cross sections.
$$\frac{d\sigma }{dt}=\frac{1}{16\pi }\frac{1}{s^2}||^2$$
(A.1)
To obtain the full cross section we must average over initial states and sum over final states. Summing over spins and splitting the matrix element into chiral components we have
$`\overline{L}L\overline{L}L:`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{spin}{}}||^2=`$ $`u^2|{\displaystyle \underset{i}{}}P_iQ_i|^2`$ (A.2)
$`\overline{L}R\overline{L}R:`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{spin}{}}||^2=`$ $`s^2|{\displaystyle \underset{i}{}}P_iQ_i|^2`$ (A.3)
$`\overline{L}L\overline{R}R:`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{spin}{}}||^2=`$ $`t^2|{\displaystyle \underset{i}{}}P_iQ_i|^2`$ (A.4)
where $`P_i`$ is the propagator factor associated with each diagram taking the form
$$P_i=\frac{i}{q_i^2M_F^2+i\mathrm{\Gamma }_FM_F}$$
(A.5)
and one must sum over all gauge bosons and $`q_i^2=s,t`$ channels. $`Q_i`$ are the group theory factors associated with each diagram. Application of the above construction kit and averaging over initial color states (1/9) and summing final color states gives the QCD backgrounds and flavor model contributions to dijet processes.
### A.1 QCD Backgrounds
$`{\displaystyle \frac{d\sigma }{dt}}(qqqq)`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s^2}{9s^2}}\left({\displaystyle \frac{u^2+s^2}{t^2}}+{\displaystyle \frac{t^2+s^2}{u^2}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{s^2}{ut}}\right)`$ (A.6)
$`{\displaystyle \frac{d\sigma }{dt}}(q\stackrel{~}{q}q\stackrel{~}{q})`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s^2}{9s^2}}\left({\displaystyle \frac{s^2+u^2}{t^2}}\right)`$ (A.7)
$`{\displaystyle \frac{d\sigma }{dt}}(q\overline{q}\stackrel{~}{q}\overline{\stackrel{~}{q}})`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s^2}{9s^4}}(t^2+u^2)`$ (A.8)
$`{\displaystyle \frac{d\sigma }{dt}}(q\overline{q}q\overline{q})`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s^2}{9s^2}}\left({\displaystyle \frac{s^2+u^2}{t^2}}+{\displaystyle \frac{t^2+u^2}{s^2}}{\displaystyle \frac{2}{3}}{\displaystyle \frac{u^2}{st}}\right)`$ (A.9)
$`{\displaystyle \frac{d\sigma }{dt}}(q\overline{\stackrel{~}{q}}q\overline{\stackrel{~}{q}})`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s^2}{9s^s}}\left({\displaystyle \frac{s^2+u^2}{t^2}}\right)`$ (A.10)
$`{\displaystyle \frac{d\sigma }{dt}}(ggq\overline{q})`$ $`=`$ $`{\displaystyle \frac{\pi \alpha _s^2}{6s^2}}\left({\displaystyle \frac{u}{t}}+{\displaystyle \frac{t}{u}}{\displaystyle \frac{9}{4}}{\displaystyle \frac{t^2+u^2}{s^2}}\right)`$ (A.11)
$`{\displaystyle \frac{d\sigma }{dt}}(q\overline{q}gg)`$ $`=`$ $`{\displaystyle \frac{32\pi \alpha _s^2}{27s^2}}\left({\displaystyle \frac{u}{t}}+{\displaystyle \frac{t}{u}}{\displaystyle \frac{9}{4}}{\displaystyle \frac{t^2+u^2}{s^2}}\right)`$ (A.12)
$`{\displaystyle \frac{d\sigma }{dt}}(qgqg)`$ $`=`$ $`{\displaystyle \frac{4\pi \alpha _s^2}{9s^2}}\left({\displaystyle \frac{u}{s}}{\displaystyle \frac{s}{u}}+{\displaystyle \frac{9}{4}}{\displaystyle \frac{s^2+u^2}{t^2}}\right)`$ (A.13)
$`{\displaystyle \frac{d\sigma }{dt}}(gggg)`$ $`=`$ $`{\displaystyle \frac{9\pi \alpha _s^2}{2s^2}}\left(3{\displaystyle \frac{tu}{s^2}}{\displaystyle \frac{su}{t^2}}{\displaystyle \frac{st}{u^2}}\right)`$ (A.14)
### A.2 Chiral Quark Family Symmetry: Matrix Elements into dijets.
$`\mathrm{\Delta }|(qqqq)|^2`$ $`=`$ $`(4\pi )^2\kappa ^2s^2\left|{\displaystyle \frac{1}{3}}P_t{\displaystyle \frac{1}{3}}P_u\right|^2`$ (A.15)
$`{\displaystyle \frac{(4\pi )^2\kappa \alpha _ss^2}{9}}Re\left({\displaystyle \frac{1}{t}}P_t+{\displaystyle \frac{1}{u}}P_u{\displaystyle \frac{1}{u}}P_t{\displaystyle \frac{1}{t}}P_u\right)`$
$`\mathrm{\Delta }|(udud)|^2`$ $`=`$ $`{\displaystyle \frac{(4\pi )^2\kappa ^2s^2}{9}}|P_t|^2+{\displaystyle \frac{(4\pi )^2\kappa \alpha _ss^2}{9}}Re\left({\displaystyle \frac{1}{t}}P_t\right)`$ (A.16)
$`\mathrm{\Delta }|(usus)|^2`$ $`=`$ $`{\displaystyle \frac{(4\pi )^2\kappa ^2s^2}{36}}|P_t|^2+{\displaystyle \frac{(4\pi )^2\kappa \alpha _ss^2}{18}}Re\left({\displaystyle \frac{1}{t}}P_t\right)`$ (A.17)
$`\mathrm{\Delta }|(dsds)|^2`$ $`=`$ $`(4\pi )^2\kappa ^2s^2\left|{\displaystyle \frac{1}{6}}P_t+{\displaystyle \frac{1}{2}}P_u\right|^2+{\displaystyle \frac{(4\pi )^2\kappa \alpha _ss^2}{3}}Re\left({\displaystyle \frac{1}{6t}}P_t+{\displaystyle \frac{1}{2t}}P_u\right)`$ (A.18)
$`\mathrm{\Delta }|(q\overline{q}q\overline{q})|^2`$ $`=`$ $`(4\pi )^2\kappa ^2u^2\left|{\displaystyle \frac{1}{3}}P_t{\displaystyle \frac{1}{3}}P_s\right|^2{\displaystyle \frac{(4\pi )^2\kappa \alpha _su^2}{9}}Re\left({\displaystyle \frac{1}{s}}P_t+{\displaystyle \frac{1}{t}}P_s\right)`$ (A.19)
$`\mathrm{\Delta }|(u\overline{u}d\overline{d})|^2`$ $`=`$ $`(4\pi )^2\kappa ^2u^2\left|{\displaystyle \frac{1}{3}}P_s\right|^2`$ (A.20)
$`\mathrm{\Delta }|(u\overline{u}s\overline{s})|^2`$ $`=`$ $`(4\pi )^2\kappa ^2u^2\left|{\displaystyle \frac{1}{6}}P_s\right|^2`$ (A.21)
$`\mathrm{\Delta }|(d\overline{d}s\overline{s})|^2`$ $`=`$ $`(4\pi )^2\kappa ^2u^2\left|{\displaystyle \frac{1}{6}}P_s+{\displaystyle \frac{1}{2}}P_t\right|^2{\displaystyle \frac{(4\pi )^2\kappa \alpha _su^2}{6}}Re\left({\displaystyle \frac{1}{s}}P_t\right)`$ (A.22)
$`\mathrm{\Delta }|(s\overline{d}s\overline{d})|^2`$ $`=`$ $`(4\pi )^2\kappa ^2u^2\left|{\displaystyle \frac{1}{2}}P_s+{\displaystyle \frac{1}{6}}P_t\right|^2{\displaystyle \frac{(4\pi )^2\kappa \alpha _su^2}{6}}Re\left({\displaystyle \frac{1}{t}}P_s\right)`$ (A.23)
$`\mathrm{\Delta }|(u\overline{d}u\overline{d})|^2`$ $`=`$ $`(4\pi )^2\kappa ^2u^2\left|{\displaystyle \frac{1}{3}}P_t\right|^2`$ (A.24)
$`\mathrm{\Delta }|(u\overline{s}u\overline{s})|^2`$ $`=`$ $`(4\pi )^2\kappa ^2u^2\left|{\displaystyle \frac{1}{6}}P_t\right|^2,`$ (A.25)
where $`P_s`$, $`P_t`$ and $`P_u`$ are defined by eqn.(A.5) and basically reflect the gauge boson propagator in the appropriate channel. Among the familon contributions we also include the interference with the gluon.
### A.3 SU(9) Chiral Flavor Symmetry. Matrix elements into dijets.
$`|(q_Lq_Lq_Lq_L)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left(\right|{\displaystyle \frac{\alpha _s}{t}}+{\displaystyle \frac{2\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}|^2`$
$`+\left|{\displaystyle \frac{\alpha _s}{u}}+{\displaystyle \frac{2\kappa }{3}}P_u^F+\alpha _s\mathrm{cot}^2\varphi P_u^F^{}\right|^2`$
$`{\displaystyle \frac{2}{3}}Re\left[({\displaystyle \frac{\alpha _s}{t}}+{\displaystyle \frac{2\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{})({\displaystyle \frac{\alpha _s}{u}}+{\displaystyle \frac{2\kappa }{3}}P_u^F+\alpha _s\mathrm{cot}^2\varphi P_u^F^{})\right])`$
$`|(q_Rq_Rq_Rq_R)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left(\right|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}|^2+|{\displaystyle \frac{\alpha _s}{u}}+\alpha _s\mathrm{tan}^2\varphi P_u^F^{}|^2`$ (A.28)
$`{\displaystyle \frac{2}{3}}Re\left[({\displaystyle \frac{\alpha _s}{u}}+\alpha _s\mathrm{tan}^2\varphi P_u^F^{})({\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{})\right])`$
$`|(q_Lq_Rq_Lq_R)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.29)
$`|(q_Lq_Rq_Rq_L)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2t^2}{9}}\left|{\displaystyle \frac{\alpha _s}{u}}\alpha _sP_u^F^{}\right|^2`$ (A.30)
$`|(u_Ld_Lu_Ld_L)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}+{\displaystyle \frac{2\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}\right|^2`$ (A.31)
$`|(u_Rd_Ru_Rd_R)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}\right|^2`$ (A.32)
$`|(u_Ld_Ru_Ld_R)|^2`$ $`=`$ $`|(u_Rd_Lu_Rd_L)|^2={\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.33)
$`|(u_Ls_Lu_Ls_L)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}{\displaystyle \frac{\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}\right|^2`$ (A.34)
$`|(u_Rs_Ru_Rs_R)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}\right|^2`$ (A.35)
$`|(u_Ls_Ru_Ls_R)|^2`$ $`=`$ $`|(u_Rs_Lu_Rs_L)|^2={\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.36)
$`|(d_Ls_Ld_Ls_L)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left(\right|{\displaystyle \frac{\alpha _s}{t}}{\displaystyle \frac{\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}|^2`$ (A.37)
$`+|\kappa P_s^F|^2{\displaystyle \frac{2}{3}}Re\left[\kappa P_s^F({\displaystyle \frac{\alpha _s}{t}}{\displaystyle \frac{\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{})\right])`$
$`|(d_Rs_Rd_Rs_R)|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}\right|^2`$ (A.38)
$`|(d_Ls_Rd_Ls_R)|^2`$ $`=`$ $`|(d_Rs_Ld_Rs_L)|^2={\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.39)
$`|(q_L\overline{q_L}q_L\overline{q_L})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left(\right|{\displaystyle \frac{\alpha _s}{s}}+{\displaystyle \frac{2\kappa }{3}}P_s^F+\alpha _s\mathrm{cot}^2\varphi P_s^F^{}|^2`$
$`+\left|{\displaystyle \frac{\alpha _s}{t}}+{\displaystyle \frac{2\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}\right|^2`$
$`{\displaystyle \frac{2}{3}}Re\left[({\displaystyle \frac{\alpha _s}{s}}+{\displaystyle \frac{2\kappa }{3}}P_s^F+\alpha _s\mathrm{cot}^2\varphi P_s^F^{})({\displaystyle \frac{\alpha _s}{t}}+{\displaystyle \frac{2\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{})\right])`$
$`|(q_R\overline{q_R}q_R\overline{q_R})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left(\right|{\displaystyle \frac{\alpha _s}{s}}+\alpha _s\mathrm{tan}^2\varphi P_s^F^{}|^2+|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}|^2`$ (A.42)
$`{\displaystyle \frac{2}{3}}Re\left[({\displaystyle \frac{\alpha _s}{s}}+\alpha _s\mathrm{tan}^2\varphi P_s^F^{})({\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{})\right])`$
$`|(q_L\overline{q_L}q_R\overline{q_R})|^2`$ $`=`$ $`|(q_R\overline{q_R}q_L\overline{q_L})|^2={\displaystyle \frac{2(4\pi )^2t^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}\alpha _sP_s^F^{}\right|^2`$ (A.43)
$`|(q_L\overline{q_R}q_L\overline{q_R})|^2`$ $`=`$ $`|(q_R\overline{q_L}q_R\overline{q_L})|^2={\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.44)
$`|(u_L\overline{u_L}d_L\overline{d_L})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}+{\displaystyle \frac{2\kappa }{3}}P_s^F+\alpha _s\mathrm{cot}^2\varphi P_s^F^{}\right|^2`$ (A.45)
$`|(u_R\overline{u_R}d_R\overline{d_R})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}+\alpha _s\mathrm{cot}^2\varphi P_s^F^{}\right|^2`$ (A.46)
$`|(u_L\overline{u_L}d_R\overline{d_R})|^2`$ $`=`$ $`|(u_R\overline{u_R}d_L\overline{d_L})|^2={\displaystyle \frac{2(4\pi )^2t^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}\alpha _sP_s^F^{}\right|^2`$ (A.47)
$`|(u_L\overline{u_L}s_L\overline{s_L})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}{\displaystyle \frac{\kappa }{3}}P_s^F+\alpha _s\mathrm{cot}^2\varphi P_s^F^{}\right|^2`$ (A.48)
$`|(u_R\overline{u_R}s_R\overline{s_R})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}+\alpha _s\mathrm{tan}^2\varphi P_s^F^{}\right|^2`$ (A.49)
$`|(u_L\overline{u_L}s_R\overline{s_R})|^2`$ $`=`$ $`|(u_R\overline{u_R}s_L\overline{s_L})|^2={\displaystyle \frac{2(4\pi )^2t^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}\alpha _sP_s^F^{}\right|^2`$ (A.50)
$`|(d_L\overline{d_L}s_L\overline{s_L})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left(\right|{\displaystyle \frac{\alpha _s}{s}}{\displaystyle \frac{\kappa }{3}}P_s^F+\alpha _s\mathrm{cot}^2\varphi P_s^F^{}|^2+{\displaystyle \frac{1}{2}}|\kappa P_t^F|^2`$ (A.52)
$`{\displaystyle \frac{1}{3}}Re\left[({\displaystyle \frac{\alpha _s}{s}}{\displaystyle \frac{\kappa }{3}}P_s^F+\alpha _s\mathrm{cot}^2\varphi P_s^F^{})\kappa P_t^F\right])`$
$`|(d_R\overline{d_R}s_R\overline{s_R})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}+\alpha _s\mathrm{tan}^2\varphi P_s^F^{}\right|^2`$ (A.53)
$`|(d_L\overline{d_L}s_R\overline{s_R})|^2`$ $`=`$ $`|(d_R\overline{d_R}s_L\overline{s_L})|^2={\displaystyle \frac{2(4\pi )^2t^2}{9}}\left|{\displaystyle \frac{\alpha _s}{s}}\alpha _sP_s^F^{}\right|^2`$ (A.54)
$`|(s_L\overline{d_L}s_L\overline{d_L})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left(\right|{\displaystyle \frac{\alpha _s}{t}}{\displaystyle \frac{\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}|^2+{\displaystyle \frac{1}{2}}|\kappa P_s^F|^2`$ (A.55)
$`{\displaystyle \frac{1}{3}}Re\left[\kappa P_s^F({\displaystyle \frac{\alpha _s}{t}}{\displaystyle \frac{\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{})\right])`$
$`|(s_R\overline{d_R}s_R\overline{d_R})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}\right|^2`$ (A.56)
$`|(s_L\overline{d_R}s_L\overline{d_R})|^2`$ $`=`$ $`|(s_R\overline{d_L}s_R\overline{d_L})|^2={\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.57)
$`|(u_L\overline{d_L}u_L\overline{d_L})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left(\left|{\displaystyle \frac{\alpha _s}{t}}+{\displaystyle \frac{2\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}\right|^2\right)`$ (A.58)
$`|(u_R\overline{d_R}u_R\overline{d_R})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}\right|^2`$ (A.59)
$`|(u_L\overline{d_R}u_L\overline{d_R})|^2`$ $`=`$ $`|(u_r\overline{d_L}u_R\overline{d_L})|^2={\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.60)
$`|(u_L\overline{s_L}u_L\overline{s_L})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left(\left|{\displaystyle \frac{\alpha _s}{t}}{\displaystyle \frac{\kappa }{3}}P_t^F+\alpha _s\mathrm{cot}^2\varphi P_t^F^{}\right|^2\right)`$ (A.61)
$`|(u_R\overline{s_R}u_R\overline{s_R})|^2`$ $`=`$ $`{\displaystyle \frac{2(4\pi )^2u^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}+\alpha _s\mathrm{tan}^2\varphi P_t^F^{}\right|^2`$ (A.62)
$`|(u_L\overline{s_R}u_L\overline{s_R})|^2`$ $`=`$ $`|(u_R\overline{s_L}u_R\overline{s_L})|^2={\displaystyle \frac{2(4\pi )^2s^2}{9}}\left|{\displaystyle \frac{\alpha _s}{t}}\alpha _sP_t^F^{}\right|^2`$ (A.63)
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# Charge, Orbital and Magnetic Order in 𝑁𝑑_0.5𝐶𝑎_0.5𝑀𝑛𝑂₃
## I Introduction
The manganites $`L_{1x}M_xMnO_3`$, where $`L`$ is a trivalent rare earth and $`M`$ a divalent alkaline earth element, present astonishing properties due to the interplay of magnetism, electric transport and crystallographic distortion. Much interest has been focused on the hole doped ferromagnetic phase $`(x0.3)`$ which exhibits colossal magnetoresistance. At fractional doping level, another mechanism is superimposed: charge ordering between $`Mn^{3+}`$ and $`Mn^{4+}`$ ions takes place. The crystallographic and magnetic environment is then strongly modified. For instance, $`Nd_{0.5}Sr_{0.5}MnO_3`$ is ferromagnetic and metallic below 245K. At 160K, the charge order transition occurs and the system becomes antiferromagnetic and insulating. Applying a strong magnetic field restores the ferromagnetic order together with the metallic conductivity.
However these properties depend crucially on the local lattice distortion and the mismatch between the rare earth and the alkaline element . In these perovskite structure, it is convenient to define the tolerance factor $`t=\left(r_{L,M}+r_O\right)/\sqrt{2}\left(r_{Mn}+r_O\right)`$ where $`r_{L,M}`$, $`r_{Mn}`$ and $`r_O`$ are the average ionic radii of the trivalent rare earth and divalent alkaline earth site, manganese site and oxygen site respectively. It describes the distortion of the perovskite structure. This distortion modifies the strength of the different magnetic interactions present in these systems. The double exchange mechanism couples the itinerant $`3de_g`$ electrons of $`Mn^{3+}`$ with the localized $`3dt_{2g}`$ spins of the $`Mn^{4+}`$ ions. It favors ferromagnetism together with metallic conductivity. It depends on the number of $`e_g`$ electrons as well as the bond angle Mn-O-Mn governed by the lattice distortion. On the other hand, the superexchange interaction between manganese ions favors antiferromagnetic order with localized spins with then insulating electric properties. Decreasing the ionic radius of the L and M site modifies the bond angle Mn-O-Mn . The ferromagnetic ordering interaction is then reduced . The bandwidth of the $`e_g`$ electron becomes narrower. Eventually, the system becomes insulating and lost its ferromagnetic order.
Additionally, for $`x=0.5`$, charge ordering will tend to localize the spins and destabilize the ferromagnetic order. This is what has been observed in $`Nd_{0.5}Sr_{0.5}MnO_3`$ and in $`Pr_{0.5}Sr_{0.5}MnO_3`$ where the tolerance factor is close to 0.95. The ferromagnetic, metallic order, established below $`T_C250K`$, is destroyed when charge ordering sets in at $`T_{CO}150K`$, and the system becomes antiferromagnetic and insulating ($`T_N=T_{CO})`$.
For smaller tolerance factor, the magnetic phase diagram is not so clear. In $`Pr_{0.5}Ca_{0.5}MnO_3`$ ,, with $`t=0.933`$,or in $`Sm_{0.5}Ca_{0.5}MnO_3`$ $`(t=0.924)`$, the charge order state occurs at higher temperature ($`T_{CO}250K`$) so that no ferromagnetic order is established in zero magnetic field. Antiferromagnetism is observed only at lower temperature ($`T_N150K<T_{CO})`$.
Surprisingly, the magnetic susceptibility presents a huge peak at $`T_{CO}`$, as observed for a long range antiferromagnetic order, and none at $`T_N`$. Such a behavior is very sensitive to stoichiometry. In $`Pr_{0.6}Ca_{0.4}MnO_3`$ ,, as well as in $`Nd_{0.6}Ca_{0.4}MnO_3`$, the anomaly at $`T_{CO}`$ persists while the signature of $`T_N`$ is quite visible. What is the exact nature of the phase when $`T_N<T<T_{CO}`$? What is the mechanism responsible for the susceptibility anomaly at $`T_{CO}`$? To answer these questions, we have focused on the compound $`Nd_{0.5}Ca_{0.5}MnO_3`$ for which $`t=0.930`$. We have undertaken an extensive study of its crystallographic, electric and magnetic properties using neutron diffraction, electric transport, magnetic susceptibility and high field magnetization measurements.
## II Experiments
The powdered sample was prepared by solid state reaction from the component oxides in appropriate ratios. The reagents were intimately mixed in an agate mortar and heated in air at $`900{}_{}{}^{}C`$ for 12 hours to allow decarbonation to take place. This powder was fired at $`1200{}_{}{}^{}C`$ for 12 hours. Final firing was performed at $`1500{}_{}{}^{}C`$ for five days with intermediate grinding to ensure homogeneity. Phase purity was checked by X ray diffraction using a Phillips diffractometer with Cu K<sub>α</sub> radiation.
The phase $`Nd_{0.5}Ca_{0.5}MnO_3`$ was subsequently registered on the high-resolution diffractometer D2b at the Institut Laue Langevin using a wavelength of $`\lambda =1.5938`$Å at room temperature and 1.5 K. The sample was contained in a vanadium can. The high-flux powder diffractometer CRG-CNRS D1b was used to investigate the thermal dependence of the magnetic structure in the temperature range $`1.5K293K`$. Thanks to the position sensitive detector covering $`80{}_{}{}^{}(2\theta )`$ and the high neutron flux available at the $`2.52`$Å wavelength, a neutron diffraction pattern has been measured every 6 Kelvin. The nuclear and magnetic refinement were performed using the profile fitting program Fullprof ($`b_{Mn}=\mathrm{0.373.10}^{12}cm`$, $`b_O=\mathrm{0.580.10}^{12}cm`$, $`b_{Nd}=\mathrm{0.769.10}^{12}cm`$, $`b_{Ca}=\mathrm{0.470.10}^{12}cm`$).
Magnetic susceptibility measurements were made with a SQUID magnetometer (MPMS Quantum Design) in the temperature range 4K-350K. These D.C. susceptibility measurements were carried out at low field (10 mT) with increasing temperature after the samples were either zero field cooled (ZFC) or field cooled (FC). High temperature susceptibility measurements were also taken up to 800K using a faraday balance. High field magnetization measurements were performed, up to 22T, with a conventional extraction set-up, in the temperature range 4K-300K.
The electric resistivity was measured between 5 K and 600 K, , using a standard four points method, on a bar with typical dimensions $`0.2\times 0.2\times 1cm^3`$. Magnetoresistance measurements were performed with magnetic fields up to 7T, the temperature ranging from 5 K to 400 K.
## III Results
### A Neutrons
The powder neutron diffraction pattern of $`Nd_{0.5}Ca_{0.5}MnO_3`$ , collected at room temperature in the paramagnetic domain, exhibits the orthorhombic $`GdFeO_3`$ structure ($`aa_p\sqrt{2}`$, $`b2a_p`$, $`ca_p\sqrt{2}`$). Refinements of the room temperature neutron diffraction patterns, carried out in the space group $`Pnma`$ attest the validity of the structure resolution. The data were first analyzed with a ’ whole pattern fitting ’ algorithm in order to determine accurately the profile shape function, background and cell parameters. This preliminary study provided a good estimate of the $`R_{wp}`$ and $`\chi ^2`$ that could be reached during the structure refinement. This whole pattern fitting led to an agreement factor $`R_{wp}=4.48\%`$ and $`\chi ^2=1.42`$. The refinement was undertaken with the room temperature neutron diffraction pattern, a temperature for which the magnetic measurements showed a paramagnetic behavior (see the magnetic properties section). 235 Bragg peaks were used to refine 7 positional parameters and 4 isotropic temperature factors. The refinement converged to give an agreement factor $`R_{wp}=4.84\%`$ and $`\chi ^2=1.71`$. Observed, calculated and difference diffraction profiles are shown in Figure 1. A Jahn-Teller (J-T) distortion of the $`Q_2`$ type (antiferrodistorsive), that displaces the basal plane oxygen atoms from their ideal positions, has been proposed to be responsible for the insulating state. This distortion does not change the lattice symmetry ($`Pnma`$) but modifies the cell deformation in such a way that $`b/\sqrt{2}ca`$ (the O’-type structure). Thus, the resulting Mn-O bond lengths of 1.943 Å and 1.948 Å (basal plane) and 1.936 Å (apical distance) at room temperature are at the origin of a small gap between the J-T split Mn $`e_g`$ bands ($`d_{x^2y^2}`$ and $`d_{z^2}`$).
On cooling the specimen $`Nd_{0.5}Ca_{0.5}MnO_3`$ below 250 K, very weak extra reflections appear on the electron microdiffraction patterns . The patterns can only be indexed by doubling the $`a`$ lattice parameter. However, the extra reflections are too weak to determine the point group from electron microdiffraction patterns and from convergent beam electron diffraction patterns. Therefore, subsequent refinements based on neutron diffraction data were carried out in the higher symmetry $`Pnma`$ space group, characteristic of the structure above the transition. It should be clear that this procedure will yield only a structure averaged over the symmetry operators of the $`Pnma`$ space group. Refinements of a $`Pnma`$ model lead to large displacement parameters of the oxygen atoms indicative of disorder and/or subtle distortions not taken into account by an orthorhombic model. The corresponding $`Nd_{0.5}Ca_{0.5}MnO_3`$ structural parameters, selected bond distances and angles at room temperature and 1.5 K are reported in Table 1. The most significant difference between the room temperature and the low temperature structures is in the Mn-O bond lengths. At room temperature, the octahedral coordination of manganese with oxygen is almost undistorted, with six approximately equal Mn-O distances. At low temperature, the two $`MnO_1`$ distances (along the b axis) become shorter than the four $`MnO_2`$ distances in the a-c plane, implying a J-T distortion of the ’apically compressed’ type as previously reported by Radaelli et al. in L$`a_{0.5}Ca_{0.5}MnO_3`$. The mean distortion in the $`MnO_6`$ octahedra $`\mathrm{\Delta }_d`$ is increased by a factor of 20 as the temperature decreases from room temperature to low temperature. These results clearly demonstrate that there is a close link between the lattice parameters and the presence of the J-T distorted $`Mn^{3+}O_6`$ octahedra with the $`d_{z^2}`$ orbital oriented in the a-c plane.
The high-flux powder diffractometer D1b was first used to investigate the thermal dependence of the orthorhombic cell parameters in the temperature range of 1.5 K-293 K (Figure 3). One observes that a and c parameters increase as the temperature is lowered from $`T_{CO}=250K`$ to $`T=160K`$ accompanied by a significant decrease in the b parameter. Our results demonstrate that the charge ordered state is progressively established from $`T_{CO}=250K`$ to $`T=160K`$.
The magnetic reflections in the low-temperature neutron powder diffraction pattern have been indexed on the basis of ($`aa_p\sqrt{2}`$, $`b2a_p`$, $`ca_p\sqrt{2}`$) unit cell, previously reported by Wollan and Koehler . The magnetic structure of $`Nd_{0.5}Ca_{0.5}MnO_3`$, known as CE-type (Figure 4), is quite complex : it entails a quadrupling of the volume of the original orthorhombic unit cell and consists of two magnetic sublattices with independent propagation vectors. This observation, associated with the fact that $`Nd_{0.5}Ca_{0.5}MnO_3`$ is an insulator (see electrical properties below), has confirmed the hypothesis, first formulated by Goodenough and developed by others , that the two sublattices result from charge ordering between $`Mn^{3+}`$ and $`Mn^{4+}`$. In Goodenough’s model, charge ordering is accompanied by orbital ordering forming ferromagnetic zig-zag chains in the a-c plane ($`d_{z^2}Mn^{3+}`$ orbitals associated with the long $`Mn^{3+}O`$ bonds in the J-T distorted$`Mn^{3+}O_6`$ octahedra). These ferromagnetic chains are connected antiferromagnetically each others in the a-c plane and along the b-axis. The magnetic structures of the two sublattices have different propagation vectors : $`\stackrel{}{k}_1=(0,0,1/2)`$ for $`Mn^{3+}`$ and $`\stackrel{}{k}_2=(1/2,0,1/2)`$ for $`Mn^{4+}`$. Rietveld refinements of the low-temperature neutron powder diffraction pattern clearly indicate that the magnetic moments are oriented along the a-axis (Figure 4). The refined values of the magnetic moments at 2K are reported in Table 2 while their temperature variations are plotted in Figure 4. Above $`T_N=160K`$, this magnetic moment vanishes. It is worth noticing that the long range antiferromagnetic order is established only well below $`T_{CO}.`$ It is important to note here that $`T_N`$ exactly coincides with the temperature below which the lattice parameters become constant (Figure 3).
### B Electric properties
The electric resistivity $`\rho \left(T\right)`$ is plotted in Figure 6. The compound remains insulating what ever the temperature. Below 50K, the resistivity is too high to be measured. Such an insulating state is in agreement with Goodenough predictions of the CE type antiferromagnetic state. The discontinuity around 250K corresponds to the charge order transition. Magnetic fields have little effects on the resistivity, up to 7 T at least. Above $`T_{CO}`$, a model of small polaron describes correctly the temperature dependence:
$`\rho \left(T\right)=\rho _0+AT^{1.5}e^{\left(\frac{E}{T}\right)}`$with $`\rho _0=1.76\times \mathrm{\hspace{0.17em}10}^3`$ $`\mathrm{\Omega }cm,A=1.40\times 10^8,E=947K`$. At lower temperature, Below $`T_N`$, a model of variable range hoping is more suitable: $`\rho \left(T\right)=Ae^{\left(E/T\right)^{0.25}}`$ with $`A=1.84\times 10^{16},E=3.4\times 10^8K`$. In the intermediate temperature range, there is a gradual change from the polaronic behavior towards the variable range hopping behavior.
### C Magnetic properties
The DC susceptibility $`\chi _{DC}`$ is plotted in Figure 7 as a function of temperature. A pronounced peak is visible at 250K which corresponds to the charge ordering temperature $`T_{CO}`$ established by neutron diffraction investigation. Below $`T_{CO}`$, the susceptibility is reduced suggesting the occurrence of antiferromagnetic correlations giving no long range order. On the other hand, no anomaly occurs at the antiferromagnetic ordering temperature $`T_N`$ determined by neutron diffraction analysis. The increase in $`\chi _{DC}`$ at low temperatures is attributed to $`Nd^{3+}`$ ions and will be discussed later.
We have also undertaken high temperature susceptibility measurements up to $`800K`$ using a Faraday balance. The susceptibility follows a Curie-Weiss law with a positive Curie Temperature $`\theta _1=216K`$ indicative of ferromagnetic interactions, (Figure 8). The Curie constant deduced from the linear part of the $`\chi _{DC}^1(T)`$ curve, (450 K $`<`$T $`<`$800 K), gives an effective moment $`\mu _{eff}=`$ $`5.30\mu _B`$. Assuming a rigid coupling of the moments of the $`Nd^{3+}`$, $`Mn^{4+}`$and $`Mn^{3+}`$ ions, one should measure, with 0.5 $`Nd^{3+}`$, 0.5 $`Mn^{4+}`$and 0.5 $`Mn^{3+}`$ions per formula unit, $`\mu _{eff}^{calc}=\sqrt{0.5.\mu _{eff}^2(Nd^{3+})+0.5.\mu _{eff}^2(Mn^{3+})+0.5.\mu _{eff}^2(Mn^{4+})}`$. The following assumptions have been made: the electronic levels of $`Nd^{3+}`$, at high temperature, are well described by $`g=\frac{8}{11}`$, $`J=\frac{9}{2}`$ , which leads to $`\mu _{eff}=gJ\left(J+1\right)\mu _B=3.62\mu _B`$; for $`Mn^{4+}`$and $`Mn^{3+}`$, the orbital momentum is quenched so that S is the appropriate quantum number: $`\mu _{eff}=gS\left(S+1\right)\mu _B`$, with $`g=2`$ and $`S=\frac{3}{2}`$ or $`2`$ , which leads to $`\mu _{eff}`$ $`=3.87\mu _B`$ for $`Mn^{4+}`$and $`\mu _{eff}`$ $`=4.90\mu _B`$ for $`Mn^{3+}`$. This gives $`\mu _{eff}^{calc}=`$ $`5.10\mu _B`$, not far from the experimental value.
DC magnetization measurements were performed up to 22 Tesla. Typical curves are plotted in Figure 9, one taken above $`T_{CO}`$, two others between $`T_{CO}`$ and $`T_N`$ and the last one below $`T_N`$. From these data, it seems that the compound presents a magnetic order occurring at the same time as the charge order. A field induced transition, analogous to a spin flop transition, is clearly visible below $`T_{CO}`$.This transition is first order, that is, strongly hysteretic.
We have further analyzed the magnetization curves in terms of Brillouin function, above the ”spin flop field”, by plotting the magnetization $`M`$ as a function of $`H_i/T`$ where $`H_i`$ is the internal field: $`H_i=H_a+H_m`$, $`H_a`$ is the applied field and $`H_m`$ the molecular field due to the magnetic interactions, proportional to the magnetization. All the curves in the temperature range 290K-170K collapse on a single curve at high field with the following fitting parameters: the Curie temperature is $`\theta _2=175K`$ and the effective moment is $`5.83\mu _B`$. It should be noted that this analysis is made in the temperature range where $`\chi _{DC}`$ slightly departs from the high temperature Curie-Weiss law, (Figure 8). This would indicate that the susceptibility could be described with two Curie-Weiss laws, one with $`\theta _2=175K=T_N`$, $`\mu _{eff}=5.83\mu _B`$ for the temperature range 170K$`<`$T$`<`$290K and for high magnetic fields and the other one with $`\theta _1=210KT_{CO},\mu _{eff}=5.30\mu _B`$ for T$`>`$400K.
As the temperature is lowered, the influence of the neodymium ions becomes visible in the magnetic response of the compound: the susceptibility increases at low temperature while the magnetization does not saturate. It is however quite difficult to estimate exactly the neodymium contribution at these temperature because of crystal field effects on the Nd electronic level. For instance, in $`NdGaO_3`$ which has a similar perovskite structure, the J=9/2 electronic level is split into 5 Kramer doublets separated by 137K, 133K, 363K and 183K so that only the fundamental level with effective spin 1/2 is occupied at low temperature. This will give a smaller effective moment and a strong deviation from a Curie law. This is exactly what is observed at low temperature in the magnetic susceptibility (see Figure 7).
## IV Discussion
From the neutron diffraction and magnetization data, we can extract a phase diagram (Figure 10). The paramagnetic phase (P), above $`T_{CO}`$, is a Jahn-Teller distorted phase. It is an insulator with a thermally activated behavior. At low temperature, below $`T_N`$, an antiferromagnetic phase (AF-O) with complete charge and orbital ordering is present. Ferromagnetic zig-zag chains of $`Mn^{3+}`$ and $`Mn^{4+}`$ are formed in the a-c plane. The Jahn-Teller distortion of the $`MnO_6`$ octahedra has considerably increased. At the same time, the electrical resistivity reflects some kind of disorder: a variable range hopping regime prevails. In the intermediate temperature range , $`Mn^{3+}`$ and $`Mn^{4+}`$ are ordered, giving rise to a doubling of the $`a`$ lattice parameter, while the orbital ordering is progressively established. It is clear that two different charge ordered phases are present: a charge ordered phase with no long range magnetic order in the temperature range $`T_N<T<T_{CO}`$, and a long range antiferromagnetic phase with complete orbital ordering for $`T<T_N`$. Both phases have similar response to magnetic field: at low field, they are characterized by a small magnetic susceptibility. At higher field, both phases present a spin flop transition, with a first order character; the charge ordered state is then destroyed by the magnetic field and a ferromagnetic state prevails.
In Figure 11, we have plotted the magnetization of this ferromagnetic phase as a function of temperature, for two different fields: $`15T`$ and $`20T`$ deduced from the $`M(H)`$ curves. The maximum at $`T_{CO}`$ has been completely removed which confirms that it is associated with the charge ordering transition. We can distinguish two different contributions to the magnetization: one arises from the Neodymium ions , the other one from the manganese ions.
The Neodymium ions are paramagnetic (no magnetic order has been detected by neutron diffraction in this temperature range). Their magnetic moment is no longer described by $`g=\frac{8}{11}`$, $`J=\frac{9}{2}`$ as has been already mentioned. Using the experimental data on $`NdGaO_3`$, we can estimate their contribution at 40K: $`M\left(15T\right)0.24\mu _B`$ and $`M\left(20T\right)0.32\mu _B`$ per formula unit. Experimentally, we get $`M\left(15T\right)3.6\mu _B`$ and $`M\left(20T\right)3.8\mu _B`$. The gap can therefore reasonably be attributed to the Neodymium ions, which leaves a value of $`3.5\mu _B`$ for the manganese ions : this is exactly the expected saturated magnetic moment for $`0.5`$ $`Mn^{4+}`$and 0.5 $`Mn^{3+}`$.
The ferromagnetic phase has then a conventional behavior: in the paramagnetic regime, a Curie Weiss law is observed with an effective moment $`5.30\mu _B`$ close to the theoretical one and with a Curie Temperature $`\theta _1=216K`$. The ordering temperature coincides with the Curie temperature $`T_C\theta _1`$. All these results suggests an isotropic ferromagnetic exchange interaction between each $`Mn^{4+}`$and $`Mn^{3+}`$ ions, in agreement with the double exchange model. At low temperature, all the manganese magnetic moment are aligned with a net magnetization of $`3.5\mu _B`$. Around the magnetic transition, the system is best described by $`\theta _2=175K`$ and $`\mu _{eff}=5.83\mu _B`$ , revealing the presence of antiferromagnetic correlations.
The low field, charge ordered phase is much more unconventional : whereas magnetization measurements suggest that $`T_N=T_{CO}`$, neutron diffraction measurements clearly demonstrate that the antiferromagnetic order is only established 100 K below $`T_{CO}`$. A huge peak in the susceptibility is observed at $`T_{CO}`$, as well as a jump in the magnetization as a function of field while no additional lines appear in the neutron diffraction spectra above $`T_N`$.
We can reconciliate all the experimental results as follows.
When lowering the temperature from room temperature, the susceptibility first increases due to the formation of a ferromagnetic phase. But this tendency to ferromagnetic order is destroyed by charge ordering which favors antiferromagnetism. The quenching of the double exchange interaction is responsible for the susceptibility peak at $`T_{CO}`$ as was proposed already in . In the temperature range $`T_N<T<T_{CO}`$ , orbital ordering sets in progressively; it is only achieved around 160K where long ranger magnetic order is established. This is suggested by the temperature variation of the lattice parameters: as the temperature is lowered from $`T_{CO}`$, they progressively decrease and remain constant only below $`T_N`$. In this picture, when $`T_N<T<T_{CO}`$, the manganese ions are charge ordered without complete orbital ordering. Several hypotheses remain.
The first one is that $`T_{CO}`$ corresponds to a magnetic transition towards some unconventional magnetic order, such as modulated or helicoidal phase or non commensurate phase, $`T_{CO}T_{N1}`$, and $`T_N`$ corresponds to a commensurate magnetic order. There are some experimental evidences, in electron microscopy measurements, that the charge ordered state has a non commensurable structure in the temperature range $`T_N<T<T_{CO}.`$
The second hypothesis is that it is not a long range magnetic order but rather fluctuations are present: ferromagnetic fluctuations above $`T_{CO}`$ and increasing antiferromagnetic fluctuations together with decreasing ferromagnetic fluctuations below $`T_{CO}`$, giving rise to the susceptibility peak at $`T_{CO}`$. A similar behavior has been observed in $`Bi_{0.18}Ca_{0.82}MnO_3`$ : ferromagnetic fluctuations have been observed above the charge order transition $`T_{CO}`$ but disappears at $`T_{CO}`$ while antiferromagnetic correlations set in. A long range antiferromagnetic order was found at lower temperature.
A third hypothesis is that there is a mixture of the ferromagnetic phase and the charge ordered phase. Similar mixing of phases has been observed in $`La_{0.5}Ca_{0.5}MnO_3`$ using NMR: antiferromagnetic and ferromagnetic domains are found to coexist at all temperatures below the first formation of charge ordered state.
The first hypothesis should be disregarded since no additional peak where observed in the high resolution neutron diffraction patterns. Besides, recent Electron Spin Resonance measurements on the same $`Nd_{0.5}Ca_{0.5}MnO_3`$ compound are in favor of the second hypothesis: the charge ordered phase is characterized by strong magnetic fluctuations with no long range magnetic order.
## V Conclusions
Combining neutron diffraction, electric transport and magnetization measurements, we have obtained a detailed description of $`Nd_{0.5}Ca_{0.5}MnO_3`$ phase diagram.
We have shown that, in the compound $`Nd_{0.5}Ca_{0.5}MnO_3`$ , the paramagnetic regime above $`T_{CO}=250K`$ corresponds to the progressive establishment of a ferromagnetic phase with $`T_C\theta _1T_{CO}`$.
From neutron diffraction experiments, a long range antiferromagnetic phase is observed only below $`T_N=160K`$. Then the manganese ions form ferromagnetic zig-zag chains coupled antiferromagnetically (CE type of ordering). The Jahn-Teller distortion is greatly enhanced and complete orbital ordering is established. The coumpound remains insulating.
A more complicated phase exists between $`T_N`$ and $`T_{CO}`$. In this phase, orbital ordering sets in progressively. antiferromagnetic and ferromagnetic interactions compete. At low magnetic field, a non magnetic, insulating state prevails whereas at high magnetic field the ferromagnetic state is favored. The transition towards the ferromagnetic phase, from the non-magnetic as well as the magnetic charge ordered phase is strongly hysteretic revealing a first order type transition.
## VI Acknowledgment
One of the authors (F.Millange) thanks O. Isnard and E. Suard for the neutron diffraction data collected on D1B diffractometer. We aknowledge also V. Caignaert for helpful discussion on the neutron diffraction results. The Grenoble High Magnetic Field Laboratory is ’laboratoire conventionné à l’Université Joseph Fourier’.
## VII Figure and table captions
Table 1: $`Nd_{0.5}Ca_{0.5}MnO_3`$ structural parameters at room temperature (RT) and 1.5K, as determined from Rietveld refinements based on neutron powder diffraction data.
Table 2: Refined magnetic moment components of the $`Mn^{3+}`$ and $`Mn^{4+}`$ sublattices. the site coordinates are expressed as fractions of the ($`2a\times b\times 2c`$) magnetic cell.
Figure 1: Rietveld refinement plot of the neutron diffraction data at room temperature. The ticks are given for the nuclear structure peaks.
Figure 2: Rietveld refinement plot of the neutron diffraction data at 1.5 K. The ticks are given for the nuclear structure peaks (upper), magnetic structure peaks ( $`Mn^{4+}`$sublattice, middle) and magnetic structure peaks ( $`Mn^{3+}`$ sublattice, lower). Points are the experimental data, line is the Rietveld fit and the lower line is the difference curve.
Figure 3: Lattices parameters of $`Nd_{0.5}Ca_{0.5}MnO_3`$ as a function of temperature. The charge order transition $`T_{CO}`$ is indicated by an arrow.
Figure 4: Representation of the CE-type magnetic structure. This structure is characterized by two Mn sublattices forming ferromagnetic zig-zag chains in the a-c plane.
Figure 5: Magnetic moment of the manganese ions as a function of temperature.
Figure 6: Electric resistivity as a function of temperature. The antiferromagnetic transition $`T_N`$ is indicated by an arrow.
Figure 7: DC susceptibility (10mT) as a function of temperature.
Figure 8: Inverse of the susceptibility at high temperature. The dotted line is a linear fit giving a Curie temperature $`\theta _1=216K`$ and an effective moment $`\mu _{eff}=`$ $`5.30\mu _B`$.
Figure 9: DC magnetization as a function of magnetic field:
(a) $`T=290K`$ (paramagnetic domain)
(b) $`T=211K`$ and (c) $`T=183K`$ (charge order domain)
(d) $`T=130K`$ (antiferromagnetic domain)
Figure 10: Magnetic phase diagram of $`Nd_{0.5}Ca_{0.5}MnO_3`$. $`T_{CO}`$ and $`T_N`$ are determined from neutron diffraction data. The different lines are guides to the eyes to separate the different phases : paramagnetic (P), charge ordered (CO), antiferromagnetic, orbital ordered (AF-O) and ferromagnetic (FM) phases. The dotted line corresponds to the average spin flop field.
Figure 11: Magnetization as a function of temperature for the high field ferromagnetic phase at 20T and 15T.
## VIII References
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# Untitled Document
OHSTPY-HEP-T-00-006 hep-th/0005072
$`R^4`$ terms in $`11`$ dimensions and conformal anomaly of (2,0) theory
A.A. Tseytlin e-mail address: tseytlin@mps.ohio-state.edu Also at Blackett Laboratory, Imperial College, London and Lebedev Physics Institute, Moscow.
Department of Physics
The Ohio State University
Columbus, OH 43210-1106, USA
Abstract
Using AdS<sub>7</sub>/CFT<sub>6</sub> correspondence we compute a subleading $`O(N)`$ term in the scale anomaly of (2,0) theory describing $`N`$ coincident M5 branes. While the leading $`O(N^3)`$ contribution to the anomaly is determined by the value of the supergravity action, the $`O(N)`$ contribution comes from a particular $`R^4`$ term (8-d Euler density invariant) in the 11-dimensional effective action. This $`R^4`$ term is argued to be part of the same superinvariant as the P-odd $`𝒞_3R^4`$ term known to produce $`O(N)`$ contribution to the R-symmetry anomaly of (2,0) theory. The known results for R-anomaly suggest that the total scale anomaly extrapolated to N=1 should be the same as the anomaly of a single free (2,0) tensor multiplet. A proposed explanation of this agreement is that the coefficient $`4N^3`$ in the anomaly (which was found previously to be also the ratio of the 2-point and 3-point graviton correlators in the (2,0) theory and in the free tensor multiplet theory) is shifted to $`4N^33N`$.
May 2000
1. Introduction
Two known maximally (2,0) supersymmetric conformal field theories in 6 dimensions are the free tensor multiplet theory describing low energy dynamics of a single M5 brane, and still largely mysterious interacting (2,0) conformal theory describing $`N`$ coincident M5 branes. A way to study the latter theory is provided by its conjectured duality to M-theory (or, for large $`N`$, 11-d supergravity corrected by higher derivative terms) on $`AdS_7\times S^4`$ background.
Comparison of the 2-point and 3-point correlators of the stress tensor of (2,0) theory as predicted by the $`AdS_7\times S^4`$ supergravity \[2,,3\] to those in the free tensor multiplet theory shows \[4,,5,,6\] that they differ only by the overall coefficient $`4N^3`$.<sup>1</sup> The same is true also for the correlators of R-symmetry currents . The remarkable coefficient $`4N^3`$ was originally found in \[5\] in the comparison of the M5 brane world volume theory and the $`D=11`$ supergravity expressions for the absorption cross-sections of longitudinally polarized gravitons by $`N`$ coincident M5 branes. The same coefficient $`4N^3`$ appears also as the ratio of the scale anomalies (or Weyl-invariant parts of conformal anomalies) of the interacting (2,0) theory and free theory of a single tensor multiplet .
The reason why the coefficient $`4N^3`$ was puzzling in was analogy with the $`d=4`$ case: a similar comparison of the gravitational and world-volume absorption cross-sections in the case of D3-branes \[10,,5\] led to the ratio $`N^2`$, which is equal to 1 for $`N=1`$. This agreement in the $`d=4`$ case was later understood as being a consequence of nonrenormalization of the conformal anomaly and thus of the 2-point stress tensor correlator in $`𝒩=4`$ SYM theory. The analogy between the $`d=4`$ and $`d=6`$ cases should not, of course, be taken too seriously, given that the (2,0) theory should have a different structure than SYM theory, being an interacting conformal fixed point without a free coupling parameter.
Still, one may expect that anomalies and 2- and 3-point correlators of currents of the (2,0) theory may have special “protected” form, with simple dependence on $`N`$, allowing one to interpolate between $`N1`$ and $`N=1`$ cases.
This was, in fact, observed for the R-symmetry anomaly of the (2,0) theory : the anomaly of the (2,0) theory obtained from the 11-d action containing the standard supergravity term plus a higher-derivative $`𝒞_3R^4`$ term is given by the sum of the leading supergravity $`O(N^3)`$ and subleading $`O(N)`$ terms, and for $`N=1`$ is equal to the R-symmetry anomaly corresponding to the single tensor multiplet \[13,,14\].
Since the conformal and R-symmetry anomalies of the (2,0) theory should belong to the same $`d=6`$ supermultiplet \[15,,11\], one should then expect to find a similar $`O(N)`$ correction to the $`O(N^3)`$ supergravity contribution to the (2,0) conformal anomaly. This $`O(N)`$ correction should originate from a higher-derivative $`R^4`$ term in the 11-d action which should be a part of the same superinvariant as $`𝒞_3R^4`$ term (just like the second-derivative supergravity terms $`R`$ and $`𝒞_3F_4F_4`$ are).
Our aim below is to discuss a mechanism of how this may happen. We shall argue that the 11-d action contains a particular $`R^4`$ term, which, upon compactification on $`S^4`$, leads to a special combination of $`R^3`$ terms in the effective 7-d action. These $`R^3`$ corrections produce extra $`O(N)`$ terms in the conformal anomaly of the boundary (2,0) conformal theory. As a result, the coefficient $`4N^3`$ in the ratio of the (2,0) theory and tensor multiplet scale anomalies may be shifted to $`4N^33N`$. Since the latter is equal to 1 for $`N=1`$, this would be a resolution of the “$`4N^3`$” puzzle.
Since this conclusion is sensitive to numerical values of coefficients in the 11-d low energy effective action we shall start with a critical review of what is known about the structure of $`R^4`$ terms in type IIA string theory in 10-d and their counterparts in M-theory. While the type IIB theory effective action contains the same $`J_0R^4`$ invariant at the tree and one-loop levels, the one-loop term in type IIA theory is a combination of two different $`R^4`$ structures. We shall argue that they should be organized into two different $`𝒩=2A`$ superinvariants – $`J_0`$ and $`_2`$ (containing P-odd $`B_2\mathrm{tr}R^4`$ term) in a way different than it was previously suggested (Section 2). The corresponding tow $`D=10`$ superinvariants “lifted” to $`D=11`$ represent the leading $`R^4`$ corrections to the 11-d supergravity action (Section 3).
These terms should be supplemented by proper $`F_4=d𝒞_3`$ dependent terms as required by supersymmetry and chosen in a specific “on-shell” scheme not to modify the $`AdS_7\times S^4`$ solution of the $`D=11`$ supergravity. Assuming that, in Section 4 we discuss higher derivative corrections to the 7-d action of $`S^4`$ compactified theory which follow from the presence of the $`R^4`$ terms in $`D=11`$ action. In Section 5 we compute the corresponding $`O(N)`$ contributions to the scale anomaly of the (2,0) theory using the method of , and draw analogy between the total $`O(N^3)+O(N)`$ result and the expression for the R-symmetry anomaly found in .
2. $`R^4`$ terms in 10 dimensions
Let us start with a review of the structure of the $`R^4`$ terms in the effective actions of type IIA superstring in 10 dimensions and the corresponding terms in M-theory effective action in 11 dimensions, paying special attention to explicit values of numerical coefficients.
The relevant terms in the tree + one loop type IIA string theory effective action can be written in the form
$$S=S_0+S_1,$$
$$S_0=\frac{1}{2\kappa _{10}^2}d^{10}x\sqrt{G}e^{2\varphi }\left(R\frac{1}{23!}H_3^2+\mathrm{}+b_0\alpha ^3J_0\right),$$
$$S_1=\frac{1}{2\pi \alpha ^{}}d^{10}x\sqrt{G}L_1,L_1=b_1𝒥_0+b_2K,$$
where $`H_{mnk}=3_{[m}B_{nk]}`$ and<sup>2</sup> We use Minkowski notation for the metric and $`ϵ`$ tensor, so that $`ϵ_{10}ϵ_{10}=10!`$, and upon reduction to 8 spatial dimensions $`ϵ_{mn\mathrm{}}ϵ_{mn\mathrm{}}2ϵ_8ϵ_8`$. For other notation see also .
$$J_0=𝒥_1+𝒥_2t_8t_8RRRR+\frac{1}{42!}ϵ_{10}ϵ_{10}RRRR,$$
$$𝒥_0=𝒥_1𝒥_2t_8t_8RRRR\frac{1}{42!}ϵ_{10}ϵ_{10}RRRR,$$
$$K=ϵ_{10}B_2[\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2].$$
In the notation we are using the numerical coefficients are
$$2\kappa _{10}^2=(2\pi )^7g_s^2\alpha ^4,$$
$$b_0=\frac{1}{32^{11}}\zeta (3),b_1=\frac{1}{(2\pi )^43^22^{13}},b_2=12b_1=\frac{1}{(2\pi )^432^{11}}.$$
The tree and one-loop coefficients of the well-known $`𝒥_1=t_8t_8RRRR`$ term<sup>3</sup> The more explicit form of this term is $`𝒥_1=24t_8[\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2],`$ where $`R=(R_{mn}^{ab})`$ and $`t_8\mathrm{tr}R^4\mathrm{tr}\left(16R_{mn}R_{rn}R_{ml}R_{rl}+8R_{mn}R_{rn}R_{rl}R_{ml}4R_{mn}R_{mn}R_{rl}R_{rl}2R_{mn}R_{rl}R_{mn}R_{rl}\right).`$ can be determined from the 4-graviton amplitude \[17,,18,,19\].<sup>4</sup> Note that the total coefficient of the $`t_8t_8RRRR`$ term in $`S`$ is thus $`\frac{1}{(2\pi )^732^{11}\alpha ^{}}(\frac{\zeta (3)}{g_s^2}+\frac{\pi ^2}{3})`$. The relative combination $`\frac{\zeta (3)}{g_s^2}+\frac{\pi ^2}{3}`$ is the same as in (where $`g^2=(2\kappa _{10})^2(2\alpha ^{})^4=16\pi ^7g_s^2`$) and in , but our overall normalization of this term is different (by factor $`2^5`$ compared to ).
The invariant $`𝒥_2=\frac{1}{42!}ϵ_{10}ϵ_{10}RRRR`$ which will play important role in what follows is the $`D=10`$ extension of the integrand of the Euler invariant in 8 dimensions
$$𝒥_2=\frac{1}{4}E_8,E_8=\frac{1}{(D8)!}ϵ_Dϵ_DR^4=\pm 8!\delta _{[m_1}^{n_1}\mathrm{}\delta _{m_8]}^{n_8}R_{n_1n_2}^{m_1m_2}\mathrm{}R_{n_7n_8}^{m_7m_8},$$
where $`\pm `$ correspond to the case of Euclidean or Minkowski signature.<sup>5</sup> The Euler number in 8 dimensions is $`\chi =\frac{1}{(4\pi )^432^7}d^8x\sqrt{g}E_8`$.
The expansion of $`E_8`$ near flat space ($`g_{mn}=\eta _{mn}+h_{mn}`$) starts with $`h^5`$ terms (see, e.g., ), so that its coefficient cannot be directly determined from the on-shell 4-graviton amplitude. The sigma-model approach implies \[22,,23\] that $`E_8`$ does appear in $`S_0`$, i.e. that (up to usual field redefinition ambiguities) the tree-level type II string $`R^4`$ term is indeed proportional to $`J_0`$ (2.1).
The structure of the kinematic factor $`(t_8+\frac{1}{2}ϵ_8)(t_8+\frac{1}{2}ϵ_8)`$ in the one-loop type IIA 4-point amplitude with transverse polarisations and momenta suggests \[24,,25,,26\] that the one-loop $`R^4`$ terms in $`D=10`$ type IIA theory should be proportional to the opposite-sign combination $`𝒥_0`$ (2.1) of the $`𝒥_1`$ and $`𝒥_2`$ terms, and this assumption passes some compactification tests \[25,,26\].
The presence of the P-odd one-loop term $`K`$ (2.1) can be established following similar calculations of anomaly-related terms in the heterotic string . Its coefficient $`b_2`$ can be fixed by considering compactification to 2 dimensions , and its value is in agreement with the coefficient required by 5-brane anomaly cancellation (see also below).
The low-energy effective string action should be supersymmetric.<sup>6</sup> The string S-matrix is invariant under on-shell supersymmetry, so the leading-order corrections to effective action evaluated on the supergravity equations of motion should be invariant under the standard supersymmetry transformations. Since the $`D=10`$ supersymmetry algebra does not close off shell, the full off-shell effective action should be invariant under “deformed” supersymmetry transformations (see, e.g., ). Remarkably, the coefficients in (2.1) are indeed consistent with what is known about the structure of possible $`R^4`$ super-invariants. First, the $`h^4`$ term in $`t_8t_8R^4`$ is the bosonic part of the on-shell linearized superspace invariant (i.e. $`d^{16}\theta \mathrm{\Phi }^4`$, $`\mathrm{\Phi }=\varphi +..+\theta ^4R+\mathrm{}`$ written in terms of $`𝒩=1`$ or $`𝒩=2B`$ \[31,,32\] on-shell superspace superfield $`\mathrm{\Phi }`$). If one first restricts consideration to $`𝒩=1,D=10`$ supersymmetry only, then one can use the classification of possible bosonic $`R^4`$ parts of on-shell non-linear $`𝒩=1`$ superinvariants given in . A basis of the three independent $`𝒩=1`$ invariants \[33,,16\] can be chosen as $`J_0,X_1,X_2`$
$$J_0=t_8t_8RRRR+\frac{1}{42!}ϵ_{10}ϵ_{10}RRRR,$$
$$X_1=t_8\mathrm{tr}R^4\frac{1}{4}ϵ_{10}B_2\mathrm{tr}R^4,X_2=t_8\mathrm{tr}R^2\mathrm{tr}R^2\frac{1}{4}ϵ_{10}B_2\mathrm{tr}R^2\mathrm{tr}R^2.$$
One may try to combine these $`𝒩=1`$ invariants to form potential $`𝒩=2A`$ superinvariants. Since $`t_8t_8R^4=24t_8[\mathrm{tr}R^4\frac{1}{4}\mathrm{tr}R^2\mathrm{tr}R^2]`$, one may consider two candidate invariants which contain combinations of $`𝒥_1`$ (2.1) or $`𝒥_2`$ (2.1) with $`\pm 6K`$ (2.1), i.e.
$$_1=24(X_1\frac{1}{4}X_2)=𝒥_16K$$
$$=t_8t_8RRRR6ϵ_{10}B_2[\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2],$$
or
$$_2=J_024(X_1\frac{1}{4}X_2)=𝒥_2+6K$$
$$=\frac{1}{42!}ϵ_{10}ϵ_{10}RRRR+6ϵ_{10}B_2[\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2],$$
$$_1+_2=J_0.$$
The 1-loop term $`L_1`$ (2.1) with $`b_2=12b_1`$ can thus be represented as a combination of two different $`R^4`$ superinvariants \[24,,25\], i.e. as
$$L_1=b_1𝒥_0+b_2K=b_1(𝒥_1𝒥_212K)=b_1(J_0+2_1),$$
or as
$$L_1=b_1(J_02_2).$$
The $`J_0`$-term should represent a separate $`𝒩=2`$ invariant.<sup>7</sup> In where non-linear extensions of $`𝒩=1`$ on-shell $`R^4`$ superinvariants were constructed the transformation of the dilaton prefactor was ignored. As a result, one was not able to make a distinction between $`J_0`$ terms appearing at the tree and 1-loop levels. It is natural to conjecture that $`f(\varphi )J_0`$ terms should combine into an $`𝒩=2A`$ superinvariant (invariant under deformed supersymmetry). For a discussion of supersymmetry of $`e^{2\varphi }R+f(\varphi )J_0`$ action in type IIB supergravity theory see . A non-trivial question is which of $`_1`$ and $`_2`$ can be actually extended to an invariant of $`𝒩=2A`$ supersymmetry.<sup>8</sup> Once the dilaton dependence of $`J_0`$ terms is taken into account, one will not be able to freely switch between $`_1`$ and $`_2`$ using (2.1).
We would like to argue that it is $`_2`$ and not $`_1`$ that is the true $`𝒩=2A`$ superinvariant. Namely, it is the Euler term $`𝒥_2=\frac{1}{4}E_8`$ and not $`𝒥_1=t_8t_8RRRR`$ that is the “superpartner” of the $`B_2`$-dependent term $`K`$ (2.1). The form of the 1-loop correction $`L_1`$ that admits a super-extension is then (2.1) and not (2.1). Then the tree + one-loop $`J_0`$ terms in the type IIA theory will be exactly the same as in the type IIB theory, $`\frac{1}{(2\pi )^432^{13}\alpha ^{}}(\frac{\zeta (3)}{g_s^2}+\frac{\pi ^2}{3})J_0`$, with the type IIA theory action containing in addition one extra one-loop contribution (2.1) proportional to the superinvariant $`_2`$.
Indeed, the weak-field expansions of both $`E_8`$ and $`K`$ start with 5-order terms, and the corresponding 5-point amplitudes should be related by global supersymmetry. At the same time, it is hard to imagine how the linearized on-shell “$`𝒲^4`$$`𝒩=2`$ superspace invariant corresponding to $`h^4`$ term in $`t_8t_8RRRR`$ may have a non-linear extension containing P-odd term $`K`$.
A more serious argument against $`t_8t_8RRRR`$ being a “superpartner” of $`ϵ_{10}B_2[\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2]`$ is the following. The $`D=10`$ type II supergravity is known to contain a one-loop quadratic $`\mathrm{\Lambda }^2`$ UV divergence proportional to $`t_8t_8RRRR`$ (this can be seen by taking the field theory limit, $`\alpha ^{}0,\mathrm{\Lambda }=`$fixed, in the one-loop 4-graviton amplitude, cf. (2.1)). At the same time, the Chern-Simons type terms like $`ϵ_{10}B_2R^4`$ can not appear in the UV divergent part of one-loop effective action.<sup>9</sup> Known examples of induced CS terms have finite coefficients and originate from IR effects (they appear from 1-loop contributions containing $`\frac{1}{^2}`$ massless poles, and thus can be re-written in a manifestly gauge invariant but nonlocal form). This can be proved directly by using the background field method: all one-loop UV divergent terms must be manifestly invariant under 2-form gauge transformations and as well as diffeomorphisms. Since, e.g., a proper time cutoff is expected to preserve supersymmetry at the level of one-loop UV divergences, one concludes that $`𝒥_1`$ and $`K`$ can not be parts of the same superinvariant.
Similar argument can be given in the context of $`D=11`$ theory. The $`t_8t_8RRRR`$ term appears \[36,,20,,37,,24\] as a cubic UV divergence (with a particular value of the UV cutoff being fixed by duality considerations ), but $`ϵ_{11}𝒞_3R^4`$ term can have only a finite coefficient (with a non-perturbative dependence on $`\kappa _{11}`$ on dimensional grounds). Thus (contrary to some previous suggestions in the literature, cf. \[20,,24,,25,,38\]) these terms can not be related by supersymmetry, and the superpartner of the $`ϵ_{11}𝒞_3R^4`$ term should be the $`D=11`$ analog of $`𝒥_2=\frac{1}{4}E_8`$ (see section 3).
Before turning to a detailed discussion of the $`D=11`$ terms, let us add few comments about the structure of the $`D=10`$ effective action (2.1),(2.1). In addition to the $`R^4`$ terms given explicitly in (2.1) and (2.1), it may contain also other Ricci tensor dependent terms as well as terms depending on other fields (cf. ), for example, terms involving two and more powers of $`H_3=dB_2`$ (which were not included in the discussion of super-invariants in ). The well-known field redefinition ambiguity \[18,,40\] allows one to change the coefficients of “on-shell” terms.<sup>10</sup> For example, ignoring other fields, one may use $`R_{mn}=0`$ to simplify the structure of $`R^4`$ invariants as the graviton legs in the string amplitudes they correspond to are on mass shell. In particular, the tree-level effective action (2.1) may contain other $`R_{mn}`$ dependent terms in addition to the full curvature contractions present in $`J_0`$ (see \[23,,41,,33\])
$$J_0=32^8(R^{hmnk}R_{pmnq}R_h^{rsp}R_{rsk}^q+\frac{1}{2}R^{hkmn}R_{pqmn}R_h^{rsp}R_{rsk}^q)+O(R_{mn}).$$
The field redefinition ambiguity allows one to choose the action in a specific “scheme” where only the Weyl tensor part of the curvature appears in $`J_0`$, i.e.
$$J_0\widehat{J}_0=32^8(C^{hmnk}C_{pmnq}C_h^{rsp}C_{rsk}^q+\frac{1}{2}C^{hkmn}C_{pqmn}C_h^{rsp}C_{rsk}^q).$$
That freedom of choice of a special scheme is crucial, in particular, in order to avoid corrections to certain highly symmetric leading-order solutions, both in 10 and in 11 dimensions (see section 3). For example, in type IIB theory the (scale of) $`AdS_5\times S^5`$ solution is not modified by the $`R^4`$ terms \[42\] only in the scheme where they have the form (2.1).
3. $`R^4`$ terms in 11 dimensions
Since the invariant $`_2`$ in (2.1) contains the P-odd CS type part $`K`$, its coefficient can not develop dilaton dependence without breaking $`B_2`$ gauge invariance, i.e. its value can not be renormalized from its coupling-independent one-loop value . Taking the limit $`g_s\mathrm{}`$ this term can then be lifted to a corresponding superinvariant in $`D=11`$ theory. Assuming that the coefficient of the $`J_0`$ invariant (2.1) does not receive higher than one loop perturbative string corrections, it can be also lifted \[20,,24,,25,,26\] to $`D=11`$ (with its tree-level part giving vanishing contribution). The resulting presence of the $`t_8t_8R^4`$ term in the M-theory effective action is indeed in agreement with what follows directly from the low-energy expansion of the 4-graviton amplitude in $`D=11`$ supergravity \[37,,24\].
In view of the above discussion, we conclude that the effective action of the $`D=11`$ theory should contain two distinct $`R^4`$ superinvariants: (i) $`J_0`$ with $`t_8t_8R^4`$ as its part, and (ii) $`_2`$ which is a sum of the $`E_8`$ and $`ϵ_{11}𝒞_3R^4`$ structures. With this separation, the coefficient in front of the $`J_0`$ term is then in agreement with the 4-graviton amplitude (with the M-theory cutoff ), and the coefficient of the $`_2`$ term (its $`𝒞_3R^4`$ part) is precisely the one implied by the M5 brane anomaly cancellation condition .
Explicitly, the $`D=11`$ action is then (cf. (2.1),(2.1))
$$𝒮=𝒮_0+𝒮_1,$$
$$𝒮_0=\frac{1}{2\kappa _{11}^2}d^{11}x\sqrt{g}\left[R\frac{1}{24!}F_4^2\frac{1}{63!(4!)^2}ϵ_{11}𝒞_3F_4F_4\right],$$
$$𝒮_1=b_1T_2d^{11}x\sqrt{g}(J_02_2).$$
Here $`F_{mnkl}=4_{[m}𝒞_{nkl]}`$ and the two $`R^4`$ super-invariants are (see (2.1),(2.1),(2.1))
$$J_0=t_8t_8RRRR+\frac{1}{4}E_8,E_8=\frac{1}{3!}ϵ_{11}ϵ_{11}RRRR,$$
$$_2=\frac{1}{4}E_8+2ϵ_{11}𝒞_3[\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2].$$
The constant $`b_1=\frac{1}{(2\pi )^43^22^{13}}`$ is the same as in (2.1) and the 10-d and 11-d parameters are related as follows ($`T_1`$ and $`T_2`$ are the string and the membrane tensions)<sup>11</sup> Note that $`B_2`$ and $`𝒞_3`$ are canonically normalized, so that the 10-d invariant $`T_1B_2\mathrm{tr}(R)^4`$ in (2.1) goes into the 11-d one $`T_2𝒞_3\mathrm{tr}(R)^4`$, where in the form notation $`B_2=\frac{1}{2}B_{mn}dx^mdx^n`$, $`𝒞_3=\frac{1}{3!}𝒞_{mnk}dx^mdx^ndx^k`$, $`R^{ab}=\frac{1}{2}R_{mn}^{ab}dx^mdx^n`$. Thus $`𝒮_1`$ (3.1) contains $`T_2𝒞_3\mathrm{tr}(R)^4`$ with the coefficient $`43!2^4b_1=\frac{1}{(2\pi )^432^6}`$ which is the same as in .
$$2\kappa _{11}^2=(2\pi )^5l_{11}^9,\kappa _{10}^2=\frac{\kappa _{11}^2}{2\pi R_{11}},l_{11}=(2\pi g_s)^{1/3}\sqrt{\alpha ^{}},R_{11}=g_s\sqrt{\alpha ^{}},$$
$$T_2=\frac{1}{2\pi l_{11}^3}=(2\pi )^{2/3}(2\kappa _{11}^2)^{1/3},T_1=\frac{1}{2\pi \alpha ^{}}=2\pi R_{11}T_2.$$
The subleading $`O(T_2)`$ term (3.1) in the effective action of 11-d theory may contain also other $`O(R_{mn})`$ and $`O(F_4)`$ terms. The invariant $`J_0`$ (supplemented with appropriate $`F_4`$ dependent terms) may be considered as a non-linear extension of the linearized “$`R^4`$” superinvariant in on-shell $`D=11`$ superspace . The P-even part of the second superinvariant starting with $`_2`$ (3.1) may also include extra $`O(F_4)`$ terms. Note that in the exterior form notation $`_2`$ may be written as
$$_2e^0e^1\mathrm{}e^{10}=\frac{2}{3}ϵ_{11}eeeRRRR$$
$$+2^53!𝒞_3\left[\mathrm{tr}(RRRR)\frac{1}{4}\mathrm{tr}(RR)\mathrm{tr}(RR)\right].$$
4. $`AdS_7\times S^4`$ solution and compactification on $`S^4`$
The $`D=11`$ supergravity admits the well known $`AdS_7\times S^4`$ solution with $`F_4`$ flux $`N`$ through $`S^4`$ . Compactifying on $`S^4`$, one may derive the corresponding $`d=7`$ supergravity action, which gives the $`O(N^3)`$ contribution to the conformal anomaly in the corresponding boundary conformal (2,0) theory.
Let us consider how the presence of the $`R^4`$ terms in the 11-d effective action $`𝒮_1`$ (3.1) may influence the existence of the $`AdS_7\times S^4`$ solution and expansion near it. Using the on-shell superspace description of 11-d supergravity and assuming that all local higher-order corrections to the equations of motion can be written again in terms of the basic on-shell supergravity superfield, it was argued in that these corrections cannot modify the maximally supersymmetric $`AdS_7\times S^4`$ solution. It should be possible to see explicitly that adding the $`J_0`$ term in (3.1) (supplemented with $`F_4`$ dependent terms as required by supersymmetry<sup>12</sup> In addition to $`F_4`$ dependent terms (which may contain up to 8 powers of $`F_4`$) there are also $`F_4`$ dependent terms which accompany $`t_8t_8R^4`$ part of $`J_0`$ in the 4-point S-matrix (as suggested by the analysis of tree-level 4-point scattering amplitudes in 11-d supergravity). These derivative terms vanish on $`AdS_7\times S^4`$ background. and chosen in a special “on-shell” scheme analogous but not equivalent<sup>13</sup> Note that in contrast to $`AdS_5\times S^5`$ space with equal radii the 11-d space $`AdS_7\times S^4`$ space with radii $`1`$ and $`\frac{1}{2}`$ is not conformally flat. to (2.1) in 10-d theory) does not change the leading-order $`AdS_7\times S^4`$ solution. One may view $`J_0`$ as originating from a restricted superspace integral of $`𝒲^4`$, where $`𝒲_{abcd}(x,\theta )`$ is the on-shell supergravity superfield , which has the structure $`𝒲=F_4+\mathrm{}+\theta \theta (\gamma \mathrm{}\gamma R+\gamma \mathrm{}\gamma F_4F_4+\gamma \mathrm{}\gamma DF_4)+\mathrm{}`$ ($`\gamma \mathrm{}\gamma `$ stand for products of gamma matrices). Then $`J_0(R+F_4F_4)^4`$ and its first, second and third variation over the metric evaluated on $`AdS_7\times S^4`$ \+ $`F_4`$-flux background ($`R_{mn}(F_4^2)_{mn},F_4=0`$) will vanish, essentially as in the case of $`AdS_5\times S^5`$ solution of type IIB theory corrected by $`J_0`$ term (taken in the form (2.1)).<sup>14</sup> The vanishing of the first variation is equivalent to the vanishing of the first correction to the 11-d supergravity equations of motion $`\gamma ^{abc}D𝒲_{abcd}=0`$ due to the supercovariant constancy of $`𝒲`$ . The argument of should certainly apply to the first subleading correction to the 11-d supergravity equations of motion coming from $`R^4`$ terms in the action.
The fact that the $`AdS_7\times S^4`$ solution (and, in particular, the radii of its factors) is not modified by the $`J_0`$ correction can be also represented as a consequence of the fact that upon compactification of the 11-d theory on $`S^4`$ with $`F_4`$ flux the $`J_0`$ term (taken in the special “on-shell” scheme) reduces to the Weyl tensor dependent $`C^4`$ term (2.1), now defined in 7 dimensions.<sup>15</sup> The tree + one-loop $`J_0`$ term in type IIB theory leads to the same $`C^4`$ term (2.1) in the 5-d effective action obtained by compactifying the type IIB theory on $`S^5`$ with $`F_5`$ flux. This term produces an $`O(N)`$ correction to the leading $`N^3`$ term in the entropy of (2,0) theory describing multiple M5 branes. As in the $`AdS_5\times S^5`$ case in type IIB theory, this $`C^4`$ term does not, however, modify the expression for the conformal anomaly of the boundary conformal theory.<sup>16</sup> It is important to stress for what follows that in the above discussion we treated $`J_0`$ (3.1) as a whole, without splitting it into $`t_8t_8R^4`$ and $`E_8`$ parts. It is only that particular combination of $`R^4`$ terms that takes the “irreducible” form (2.1) (cf. ), and thus should lead only to $`C^4`$ terms upon compactification to $`d=7`$. At the same time, $`E_8`$ contains “reducible” curvature contractions like $`((R_{mnkl})^2)^2+R(R_{mnkl})^3+\mathrm{}`$ and thus may, in principle, lead to $`O(R^n),n<4,`$ terms upon compactification to $`d=7`$.
Let us now discuss the second invariant $`_2`$ (3.1) in (3.1). It is easy to see that its P-odd part $`ϵ_{11}𝒞_3[\mathrm{tr}R^4\frac{1}{4}(\mathrm{tr}R^2)^2]`$ does not modify the $`AdS_7\times S^4`$ solution. Upon reduction on $`S^4`$ it leads to $`O(N)`$ CS terms in $`d=7`$ action . As for the $`E_8`$ part of $`_2`$, we shall assume that, as in the case of $`J_0`$, there exists an “on-shell” scheme in which this term, supplemented with proper $`F_4`$-dependent terms, also does not modify the leading-order $`AdS_7\times S^4`$ solution.
The main point is that upon compactification on $`S^4`$ the $`E_8`$ term in (3.1) should produce additional $`R^3`$ higher-derivative terms in the 7-d effective action which, while not changing the vacuum solution, will give subleading $`O(N)`$ corrections to the conformal anomaly of the boundary CFT.<sup>17</sup> These terms will give also another $`O(N)`$ correction to the entropy of (2,0) theory, in addition to the one coming from the $`J_0`$ term (2.1) found in .
It is known that the $`𝒞_3R^4`$ part of $`_2`$ (3.1) gives a subleading $`O(N)`$ correction to the R-symmetry anomaly of the (2,0) theory \[12,,11\]. Since the R-symmetry and conformal anomalies should belong to the same 6-d supermultiplet, it is natural to expect that the “superpartner” of the $`𝒞_3R^4`$ term, i.e. the $`E_8`$ term in $`_2`$, should lead to an $`O(N)`$ correction to the conformal anomaly of the boundary 6-d theory. This is what we are going to suggest below.
Since we do not know the $`F_4`$ (and $`R_{mn}`$) dependent terms which supplement $`E_8`$ to a superinvariant, to determine the terms in the 7-d action that originate from the $`E_8`$ part of the invariant $`_2`$ in (3.1) we shall use the following heuristic strategy. We shall start with $`E_8`$ and compute it in the case when the 11-d space is a direct product, $`M^{11}=M^7\times M^4`$. It is easy to see that
$$E_8(M^7\times M^4)=4E_2(M^4)E_6(M^7)+12E_4(M^4)E_4(M^7),$$
where, as in (2.1),
$$E_{2n}(M^d)\frac{1}{(d2n)!}ϵ_dϵ_dR^n,d2n,$$
and $`E_{2n}(M^d)=0`$ for $`d<2n`$. In the case when $`M^4`$ is a 4-sphere of radius $`L`$ ($`R_{S^4}=\frac{12}{L^2}`$) and $`M^7`$ has curvature $`R`$ we get
$$E_8(M^7\times S^4)=\frac{32^5}{L^2}E_6(M^7)+\frac{3^22^7}{L^4}E_4(M^7)$$
$$=\frac{32^5}{L^2}ϵ_7ϵ_7RRR+\frac{32^6}{L^4}ϵ_7ϵ_7RR.$$
A remarkable property of the $`E_8`$ invariant is that it does not produce a correction to the cosmological or Einstein term in the 7-d action.
Next, we shall assume that when the same reduction is repeated for the analog of $`E_8`$ term in a special “on-shell” scheme (i.e. for $`E_8`$ supplemented by $`F_4`$ and $`R_{mn}`$ dependent terms so that it does not produce a modification of the leading-order $`AdS_7\times S^4`$ solution) then the resulting terms in the 7-d action will be the same as in (4.1) but with the curvature tensor $`R`$ of $`M^7`$ replaced by its Weyl tensor $`C`$ part.
In what follows we shall consider only on the $`E_6(M^7)C^3+\mathrm{}`$ term in (4.1) coming from $`E_8`$. The reason is that we shall compute the corresponding contribution to the scale anomaly of the boundary theory only modulo $`R_{mn}`$-dependent terms, but it is easy to see that a potential $`C^2`$ term in the 7-d action (coming from $`E_4`$ in (4.1)) can lead only to terms in the conformal anomaly which vanish when the 6-d boundary space is Ricci flat.
Choosing the normalization in which the radii of $`AdS_7`$ and $`S^4`$ are $`1`$ and $`L=\frac{1}{2}`$ so that Vol$`(S^4)=\frac{8\pi ^2}{3}L^4=\frac{\pi ^2}{6}`$, and assuming that the value of the quantized $`F_4`$ flux is $`N`$, we get (see (3.1),(3.1) and \[48,,43\])
$$\frac{1}{2\kappa _{11}^2}=\frac{N^3}{2^8\pi ^5L^9}=\frac{2N^3}{\pi ^5},\frac{1}{2\kappa _7^2}=\frac{\mathrm{Vol}(S^4)}{2\kappa _{11}^2}=\frac{N^3}{3\pi ^3},T_2=\frac{2N}{\pi }.$$
The relevant $`[N^3(R2\lambda )+NC^3]`$ terms in the 7-d action<sup>18</sup> Here we consider the Euclidean signature and change overall sign of the action, i.e. $`RR`$. are then
$$S^{(7)}=\frac{N^3}{3\pi ^3}d^7x\sqrt{g}(R+30)+\frac{\gamma N}{3^22^{11}\pi ^3}d^7x\sqrt{g}\widehat{E}_6+\mathrm{},$$
where the explicit form of the $`\widehat{E}_6C^3`$ correction term is (cf. (2.1))
$$\widehat{E}_6=(E_6)_{_{R_{mn}=0}}=ϵ_7ϵ_7CCC$$
$$=6!\delta _{[m_1}^{n_1}\mathrm{}\delta _{m_6]}^{n_6}C_{n_1n_2}^{m_1m_2}C_{n_3n_4}^{m_3m_4}C_{n_5n_6}^{m_5m_6}=32(2I_1+I_2),$$
where $`I_1`$ and $`I_2`$ are defined as
$$I_1=C_{amnb}C^{mpqn}C_p{}_{}{}^{ab}{}_{q}{}^{},I_2=C_{ab}{}_{}{}^{mn}C_{mn}^{}{}_{}{}^{pq}C_{pq}^{}{}_{}{}^{ab}.$$
As follows from (4.1) the numerical coefficient $`\gamma `$ is
$$\gamma =1,$$
but we shall keep it arbitrary, given the uncertainties in the above derivation of the correction term in (4.1) (for example, the presence of $`(F_4)^2(R_{mnkl})^3`$ terms in $`_2`$ would shift the value of $`\gamma `$).
5. Conformal anomaly of (2,0) theory
Let us now determine the contribution of the $`C^3`$ correction term in the 7-d action (4.1) which originated from the $`E_8`$ part of the $`_2`$ superinvariant in the 11-d action (3.1) to the conformal anomaly of the $`d=6`$ boundary conformal theory. We shall follow the same method as used in in computing the leading $`N^3`$ term in the anomaly.<sup>19</sup> Similar computation of subleading corrections to conformal anomaly of 4-d boundary conformal field theories (with $`𝒩<4`$ supersymmetry) coming from $`R^2`$ curvature terms in 5-d effective action were discussed in \[49,,50,,51,,38\]. We shall compute only the $`O(N)`$ contribution to the scale anomaly (which is the same as integrated conformal anomaly, assuming topology of 6-space is trivial) and ignore terms which depend on $`R_{mn}`$, i.e. concentrate only on the Weyl-invariant non total derivative $`C^3`$ terms (“type B” part) in the 6-d conformal anomaly.
To obtain the conformal anomaly one is to solve the 7-d equations for the metric (as in (4.1) we set the radius of $`AdS_7`$ to be equal to 1)
$$ds^2=\frac{1}{4}\rho ^2d\rho ^2+\rho ^1g_{ij}(x,\rho )dx^idx^j,$$
evaluate the action on the solution $`g=g_0(x)+\rho g_2(x)+\mathrm{}`$, and compute its variation under the Weyl rescaling of the 6-d boundary metric. The anomaly is essentially determined by the coefficient of the logarithmic divergence produced by the integral over $`\rho `$ . In the present case of (4.1) we find (using (4.1),(4.1) and $`R_{AdS_7}=42`$)
$$S^{(7)}=d^6x\left[\frac{N^3}{3\pi ^3}6_ϵ\frac{d\rho }{\rho ^4}\sqrt{g(x,\rho )}\frac{\gamma N}{3^22^6\pi ^3}\frac{1}{2}_ϵ\frac{d\rho }{\rho }\sqrt{g}(2I_1+I_2)+\mathrm{}\right].$$
Since
$$6_ϵ\frac{d\rho }{\rho ^4}\sqrt{g(x,\rho )}=\sqrt{g_0}\left[a_0(x)ϵ^3+\mathrm{}a_6(x)\mathrm{ln}ϵ\right]+\mathrm{},$$
the anomaly is given by the sum of the $`O(N^3)`$ and $`O(N)`$ terms<sup>20</sup> To obtain the $`O(N)`$ contribution we evaluate the $`C^3`$ term in the 7-d action on the leading-order solution for the metric (5.1) (see for a similar computation in the case of the $`R_{mnkl}^2`$ action in $`d=5`$), separate the $`C^3`$ part depending on the 6-d metric $`g_0`$, and omit other parts that depend on the Ricci tensor of $`g_0`$.
$$𝒜_{(2,0)}=𝒜_{(2,0)}^{N^3}+𝒜_{(2,0)}^N=\frac{N^3}{3\pi ^3}2a_6+\frac{\gamma N}{3^22^6\pi ^3}(2I_1+I_2+\mathrm{}).$$
Here $`a_6`$ and $`I_1,I_2`$ are evaluated for the boundary metric $`g_0`$, and dots stand for $`O(N)`$ $`R_{mn}`$-dependent and total derivative terms we are ignoring.
The result of for the leading-order contribution $`𝒜_{(2,0)}^{N^3}`$ written as a sum of the type A (Euler), type B (Weyl invariant) and scheme-dependent (covariant total derivative) terms \[52,,53\] is
$$𝒜_{(2,0)}^{N^3}=\frac{4N^3}{(4\pi )^33^22^5}\left[E_6+8(12I_1+3I_2I_3)+O(_iJ^i)\right],$$
where $`E_6=ϵ_6ϵ_6RRR`$. The invariants $`I_1,I_2`$ (4.1) and $`I_3`$
$$I_3=C_{mnbc}^2C^{mnbc}+O(R_{mn})+O(_iJ^i),$$
which form the basis of 3 Weyl invariants are the same as used in . They are related to the invariants used in \[52,,8\] as follows: $`E_{(6)},I_1,I_2`$ and $`I_3`$ in are equal to $`\frac{1}{3^32^{11}}E_6,I_1,I_2`$ and $`5I_3^{}`$, $`I_3^{}=I_3\frac{8}{3}(2I_1+I_2)\frac{1}{12}E_6+O(_iJ^i),`$ in terms of the invariants $`E_6,I_1,I_2`$ and $`I_3`$ used in and here.<sup>21</sup> Our curvature tensor $`R_{bmn}^a=_m\mathrm{\Gamma }_{bn}^a\mathrm{}`$ has the opposite sign to that of . Note also that was assuming Euclidean signature where $`E_6`$ is defined as $`ϵ_6ϵ_6RRR`$.
We use this opportunity to point out that the curvature invariant $`I_3=5I_3^{}`$ as defined in \[52,,8\] is not, in fact, covariant under Weyl transformations, contrary to what was assumed in (this can be easily checked by computing it for the metric of a sphere $`S^6`$: one finds that while $`I_1(S^6)=I_2(S^6)=0`$, $`I_3^{}(S^6)0`$). The proper third Weyl invariant of type $`C^2C`$ (5.1) was given in and is equivalent to the Weyl invariant $`I_3`$ used in and here. Since $`I_3`$ of or $`I_3^{}`$ is a mixture of the true Weyl invariants $`I_1,I_2,I_3`$ with $`E_6`$, the separation of the leading $`N^3`$ Weyl anomaly of the (2,0) theory into type A and type B parts was not presented correctly in . The correct separation was given in and is used here.<sup>22</sup> Note that when $`R_{mn}=0`$ the two invariants – $`I_3^{}`$ and $`I_3`$ – coincide, up to a covariant total derivative term. In fact, a separation of the conformal anomaly into type A and type B parts becomes ambiguous on a Ricci flat background.
Note that modulo terms that vanish for $`R_{mn}=0`$ and total derivative terms, one has the following relations (cf. (4.1))
$$E_6=32(2I_1+I_2)+O(R_{mn}),I_3=4I_1I_2+O(R_{mn})+O(_iJ^i),$$
so that $`𝒜_{(2,0)}^{N^3}`$ vanishes for $`R_{mn}=0`$, as it should .
Eq. (5.1) is to be compared with the expression for the conformal anomaly for the free (2,0) tensor multiplet found in :
$$𝒜_{tens.}=\frac{1}{(4\pi )^33^22^5}\left[\frac{7}{4}E_6+8(12I_1+3I_2I_3)+O(_iJ^i)\right].$$
As was concluded in , the Weyl-invariant (type B) parts of the leading (2,0) theory anomaly (5.1) and the tensor multiplet anomaly (5.1) have exactly the same form, up to the overall factor $`4N^3`$ in (5.1).
Since we have found the $`O(N)`$ correction to the anomaly of the (2,0) theory in (5.1) only modulo $`R_{mn}`$-dependent and total derivative terms, we are able to compare only type B anomalies, or scale anomalies (assuming that the $`d=6`$ space has trivial topology, so that we can ignore the integral of the Euler term $`E_6`$)
$$𝐀_{(2,0)}=d^6x\sqrt{g_0}𝒜_{(2,0)},𝐀_{tens.}=d^6x\sqrt{g_0}𝒜_{tens.}.$$
Using (5.1) to express $`I_3`$ in terms of $`I_1`$ and $`I_2`$, we find from (5.1),(5.1) and (5.1)
$$𝐀_{(2,0)}^{N^3}=\frac{4N^3}{(4\pi )^33^2}d^6x\sqrt{g_0}(2I_1+I_2),$$
$$𝐀_{(2,0)}^N=\frac{\gamma N}{(4\pi )^33^2}d^6x\sqrt{g_0}(2I_1+I_2),$$
$$𝐀_{tens.}=\frac{1}{(4\pi )^33^2}d^6x\sqrt{g_0}(2I_1+I_2).$$
The total scale anomaly of the (2,0) theory following from (4.1),(5.1) is then
$$𝐀_{(2,0)}=𝐀_{(2,0)}^{N^3}+𝐀_{(2,0)}^N=\frac{4N^3\gamma N}{(4\pi )^33^2}d^6x\sqrt{g_0}(2I_1+I_2).$$
Equivalently,
$$𝐀_{(2,0)}=\frac{4(N^3N)}{(4\pi )^33^3}d^6x\sqrt{g_0}(2I_1+I_2)+(4\gamma )N𝐀_{tens.}.$$
Thus if the true value of $`\gamma `$ is 3 instead of the naive value $`1`$ (4.1) which follows directly from reduction of $`E_8`$ (4.1), ignoring possible $`F_4`$-dependent ($`F_4^2R^3`$) terms in the 11-d super-invariant $`_2`$, then $`𝐀_{(2,0)}`$ reproduces the scale anomaly (5.1) of a single (2,0) tensor multiplet. This $`N=1`$ relation should be expected, given that a similar correspondence is true for the R-symmetry anomalies (see below). Though we are unable to show that $`\gamma =3`$ does follow from the $`d=7`$ reduction of the 11-d super-invariant $`_2`$ containing P-odd $`𝒞_3R^4`$ term, we find it remarkable that the required value of $`\gamma `$ differs from the naive value 1 simply by factor of 3.<sup>23</sup> In the original version of the present paper we mistakenly used the basis of type B invariants including $`I_3`$ of instead of the correct invariant of and as a result got the $`O(N)`$ term with extra coefficient 3, concluding that $`\gamma =1`$ gives already the desired coefficient $`4N^33N`$ in (5.1). <sup>24</sup> Note that if we were comparing the full local conformal anomalies evaluated for $`R_{mn}=0`$ then, since the $`N^3`$ contribution (5.1) vanishes in this case, we would need $`\gamma =\frac{3}{4}`$ in order to reproduce the non-zero $`R_{mn}=0`$ value of the tensor multiplet anomaly (5.1) by the $`N=1`$ limit of the $`O(N)`$ term in (5.1).
Making a natural conjecture that the same relation $`𝒜_{tens.}=(𝒜_{(2,0)})_{N=1}`$ should be true between the full expressions for the conformal anomalies of the (2,0) theory and tensor multiplet, one can make a prediction about the complete structure of the $`O(N)`$ term in the (2,0) theory anomaly $`𝒜_{(2,0)}`$ (5.1) (cf. (5.1),(5.1))<sup>25</sup> The shift of the coefficient of the $`E_6`$ term in the conformal anomaly seems to imply a contradiction between our assumption that the $`R^4`$ terms in the 11-d action (3.1) do not change the scale of $`AdS_7\times S^4`$ solution (i.e. that the value of the 7-d action (4.1) evaluated on the $`AdS_7`$ solution is not changed), and the claim of that the coefficient of the type A (Euler) term in the anomaly of a generic effective theory is determined only by the value of the action on the $`AdS`$ solution.
$$𝒜_{(2,0)}=\frac{1}{(4\pi )^33^22^5}\left[(4N^3\frac{9}{4}N)E_6+(4N^33N)8(12I_1+3I_2I_3)+O(_iJ^i)\right],$$
or, equivalently,
$$𝒜_{(2,0)}=\frac{N^3N}{(4\pi )^33^22^3}\left[E_6+8(12I_1+3I_2I_3)+O(_iJ^i)\right]+N𝒜_{tens.}.$$
Using (5.1), we can rewrite (5.1) also as
$$𝒜_{(2,0)}=\frac{N}{(4\pi )^332^7}\left[E_6+O(R_{mn})+O(_iJ^i)\right],$$
in agreement with the fact that for $`R_{mn}=0`$ the conformal anomaly of the tensor multiplet becomes \[9,,36$`𝒜_{tens.}=\frac{1}{(4\pi )^332^7}\left[E_6+O(_iJ^i)\right]`$.
It is useful to compare the above expressions (5.1),(5.1) with the previously known results for the R-symmetry anomalies of the interacting (2,0) theory and free tensor multiplet theory. The 1-loop effective action $`\mathrm{\Gamma }`$ for a free 6-d tensor multiplet in a background of 6-d Lorentz curvature $`R`$ and $`SO(5)`$ R-symmetry gauge field $`F`$ has local $`SO(6)`$ and $`SO(5)`$ anomalies. They satisfy the descent relations $`d(\delta \mathrm{\Gamma })=\delta I_7,I_8=dI_7`$, with the 8-form anomaly polynomial $`I_8`$ being \[13,,14\]
$$I_8^{tens.}(F,R)=\frac{1}{32^4}\left[p_2(F)p_2(R)+\frac{1}{4}[p_1(F)p_1(R)]^2\right],$$
with (here $`F^2FF`$, etc.)
$$p_1(F)=\frac{1}{2}\mathrm{tr}\overline{F}^2,p_2(F)=\frac{1}{4}\left(\mathrm{tr}\overline{F}^4\frac{1}{2}\mathrm{tr}\overline{F}^2\mathrm{tr}\overline{F}^2\right),\overline{F}=\frac{i}{2\pi }F.$$
The corresponding anomalies of the interacting (2,0) theory describing multiple M5 branes derived (by assuming that the total M5-brane anomaly \+ inflow anomaly should cancel) from the 11-d supergravity action (3.1) with the $`R^4`$ correction term (3.1) is \[11\]
$$I_8^{(2,0)}(F,R)=\frac{1}{32^4}\left[(2N^3N)p_2(F)Np_2(R)+\frac{1}{4}N[p_1(F)p_1(R)]^2\right].$$
Here the $`O(N^3)`$ term comes from the CS term in supergravity action (3.1) and the $`O(N)`$ term \[12,,14\]– from the P-odd $`𝒞_3R^4`$ part of the superinvariant $`_2`$ (3.1),(3.1). Equivalently,
$$I_8^{(2,0)}(F,R)=\frac{1}{32^3}(N^3N)p_2(F)+NI_8^{tens.}(F,R).$$
Thus for $`N=1`$ the anomaly of the (2,0) theory is the same as the anomaly of a single tensor multiplet. This is the same type of a relation we have established above (cf. (5.1)) for the scale anomalies, with the crucial $`O(N)`$ contribution coming from the P-even $`E_8`$ part of the superinvariant $`_2`$ (3.1). This is obviously consistent with the fact that R-symmetry and conformal anomalies should be parts of the same 6-d supermultiplet.
Acknowledgements We are grateful to S. Frolov for a collaboration at an initial stage and many useful discussions. We would like also to acknowledge J. Harvey, P. Howe, K. Intriligator, R. Metsaev, Yu. Obukhov, H. Osborn, T. Petkou and M. Shifman for helpful discussions and comments. This work was supported in part by the DOE grant DOE/ER/01545, EC TMR grant ERBFMRX-CT96-0045, INTAS grant No. 99-0590, NATO grant PST.CLG 974965 and PPARC SPG grant PPA/G/S/1998/00613.
References
relax J. Maldacena, “The large N limit of superconformal field theories and supergravity,” Adv. Theor. Math. Phys. 2, 231 (1998) \[hep-th/9711200\]. relax H. Liu and A.A. Tseytlin, “D = 4 super Yang-Mills, D = 5 gauged supergravity, and D = 4 conformal supergravity,” Nucl. Phys. B533, 88 (1998), hep-th/9804083. relax G. Arutyunov and S. Frolov, “Three-point Green function of the stress-energy tensor in the AdS/CFT correspondence,” Phys. Rev. D60, 026004 (1999), hep-th/9901121. relax F. Bastianelli, S. Frolov and A.A. Tseytlin, “Three-point correlators of stress tensors in maximally-supersymmetric conformal theories in $`d=3`$ and $`d=6`$,” hep-th/9911135. relax S.S. Gubser, I.R. Klebanov and A.A. Tseytlin, “String theory and classical absorption by three-branes,” Nucl. Phys. B499, 217 (1997), hep-th/9703040. relax S.S. Gubser and I.R. Klebanov, “Absorption by branes and Schwinger terms in the world volume theory,” Phys. Lett. B413, 41 (1997) \[hep-th/9708005\]. relax R. Manvelian and A.C. Petkou, “A note on R-currents and trace anomalies in the (2,0) tensor multiplet in d = 6 AdS/CFT correspondence,” hep-th/0003017. relax M. Henningson and K. Skenderis, “The holographic Weyl anomaly,” JHEP 07, 023 (1998), hep-th/9806087. relax F. Bastianelli, S. Frolov and A.A. Tseytlin, “Conformal anomaly of (2,0) tensor multiplet in six dimensions and AdS/CFT correspondence,” JHEP 02, 013 (2000), hep-th/0001041. relax I.R. Klebanov, “World-volume approach to absorption by non-dilatonic branes,” Nucl. Phys. B496, 231 (1997) \[hep-th/9702076\]. relax J.A. Harvey, R. Minasian and G. Moore, “Non-abelian tensor-multiplet anomalies,” JHEP 09, 004 (1998), hep-th/9808060. relax M.J. Duff, J.T. Liu and R. Minasian, “Eleven-dimensional origin of string / string duality: A one-loop test,” Nucl. Phys. B452, 261 (1995) \[hep-th/9506126\]. relax L. Alvarez-Gaume and E. Witten, “Gravitational Anomalies,” Nucl. Phys. B234, 269 (1984). relax E. Witten, “Five-brane effective action in M-theory,” J. Geom. Phys. 22, 103 (1997), hep-th/9610234; relax P.S. Howe, G. Sierra and P.K. Townsend, “Supersymmetry In Six-Dimensions,” Nucl. Phys. B221, 331 (1983). relax A.A. Tseytlin, “Heterotic - type I superstring duality and low-energy effective actions,” Nucl. Phys. B467, 383 (1996) \[hep-th/9512081\]. relax M.B. Green and J.H. Schwarz, “Supersymmetrical Dual String Theory. 2. Vertices And Trees,” Nucl. Phys. B198, 252 (1982). relax D.J. Gross and E. Witten, “Superstring Modifications Of Einstein’s Equations,” Nucl. Phys. B277, 1 (1986). relax N. Sakai and Y. Tanii, “One Loop Amplitudes And Effective Action In Superstring Theories,” Nucl. Phys. B287, 457 (1987). relax M.B. Green and P. Vanhove, “D-instantons, strings and M-theory,” Phys. Lett. B408, 122 (1997) \[hep-th/9704145\]. relax B. Zumino, “Gravity Theories In More Than Four-Dimensions,” Phys. Rept. 137, 109 (1986). relax M.T. Grisaru, A.E. van de Ven and D. Zanon, “Four Loop Beta Function For The N=1 And N=2 Supersymmetric Nonlinear Sigma Model In Two-Dimensions,” Phys. Lett. B173, 423 (1986); “Four Loop Divergences For The N=1 Supersymmetric Nonlinear Sigma Model In Two-Dimensions,” Nucl. Phys. B277, 409 (1986). relax M.T. Grisaru and D. Zanon, “Sigma Model Superstring Corrections To The Einstein Hilbert Action,” Phys. Lett. B177, 347 (1986); M.D. Freeman, C.N. Pope, M.F. Sohnius and K.S. Stelle, “Higher Order Sigma Model Counterterms And The Effective Action For Superstrings,” Phys. Lett. B178, 199 (1986); Q. Park and D. Zanon, “More On Sigma Model Beta Functions And Low-Energy Effective Actions,” Phys. Rev. D35, 4038 (1987). relax J.G. Russo and A.A. Tseytlin, “One-loop four graviton amplitude in eleven dimensional supergravity,” Nucl. Phys. B578, 139 (2000), hep-th/9707134. relax E. Kiritsis and B. Pioline, “On $`R^4`$ threshold corrections in type IIB string theory and (p,q) string instantons,” Nucl. Phys. B508, 509 (1997) \[hep-th/9707018\]. relax I. Antoniadis, S. Ferrara, R. Minasian and K.S. Narain, “$`R^4`$ couplings in M- and type II theories on Calabi-Yau spaces,” Nucl. Phys. B507, 571 (1997) \[hep-th/9707013\]. relax C. Vafa and E. Witten, “A One loop test of string duality,” Nucl. Phys. B447, 261 (1995) \[hep-th/9505053\]. relax W. Lerche, B.E. Nilsson and A.N. Schellekens, “Heterotic String Loop Calculation Of The Anomaly Cancelling Term,” Nucl. Phys. B289, 609 (1987); W. Lerche, B.E. Nilsson, A.N. Schellekens and N.P. Warner, “Anomaly Cancelling Terms From The Elliptic Genus,” Nucl. Phys. B299, 91 (1988). relax E.A. Bergshoeff and M. de Roo, “The Quartic Effective Action Of The Heterotic String And Supersymmetry,” Nucl. Phys. B328, 439 (1989). relax R. Kallosh, “Strings And Superspace,” Phys. Scripta T15, 118 (1987). B. E. Nilsson and A. K. Tollsten, “Supersymmetrization Of $`\zeta (3)R^4`$ In Superstring Theories,” Phys. Lett. B181, 63 (1986); relax P.S. Howe and P.C. West, “The Complete N=2, D = 10 Supergravity,” Nucl. Phys. B238, 181 (1984). relax M.B. Green, “Interconnections between type II superstrings, M theory and N = 4 Yang-Mills,” hep-th/9903124. relax M. de Roo, H. Suelmann and A. Wiedemann, “Supersymmetric $`R^4`$ actions in ten-dimensions,” Phys. Lett. B280, 39 (1992); “The Supersymmetric effective action of the heterotic string in ten-dimensions,” Nucl. Phys. B405, 326 (1993) \[hep-th/9210099\]; J. H. Suelmann, “Supersymmetry and string effective actions,” Ph.D. thesis, Groningen, 1994. relax M.B. Green and S. Sethi, “Supersymmetry constraints on type IIB supergravity,” Phys. Rev. D59, 046006 (1999) \[hep-th/9808061\]. relax R.R. Metsaev and A.A. Tseytlin, “On Loop Corrections To String Theory Effective Actions,” Nucl. Phys. B298, 109 (1988). relax E.S. Fradkin and A.A. Tseytlin, “Quantum Properties Of Higher Dimensional And Dimensionally Reduced Supersymmetric Theories,” Nucl. Phys. B227, 252 (1983). “Present State Of Quantum Supergravity,” In Proc. of Third Seminar on Quantum Gravity, Oct. 23-25 1984, Moscow, M.A. Markov et al eds., p. 303 (World Scientific, 1985) relax M.B. Green, M. Gutperle and P. Vanhove, “One loop in eleven dimensions,” Phys. Lett. B409, 177 (1997) \[hep-th/9706175\]. relax D. Anselmi and A. Kehagias, “Subleading corrections and central charges in the AdS/CFT correspondence,” Phys. Lett. B455, 155 (1999) \[hep-th/9812092\]. relax A. Kehagias and H. Partouche, “The exact quartic effective action for the type IIB superstring,” Phys. Lett. B422, 109 (1998) \[hep-th/9710023\]. relax A.A. Tseytlin, “Ambiguity In The Effective Action In String Theories,” Phys. Lett. B176, 92 (1986). relax R. Myers, “Superstring Gravity And Black Holes,” Nucl. Phys. B289, 701 (1987). relax T. Banks and M.B. Green, “Non-perturbative effects in $`AdS_5\times S^5`$ string theory and d = 4 SUSY Yang-Mills,” JHEP 9805, 002 (1998) \[hep-th/9804170\]. relax S.S. Gubser, I.R. Klebanov and A.A. Tseytlin, “Coupling constant dependence in the thermodynamics of N = 4 supersymmetric Yang-Mills theory,” Nucl. Phys. B534, 202 (1998), hep-th/9805156. relax E. Cremmer and S. Ferrara, “Formulation Of Eleven-Dimensional Supergravity In Superspace,” Phys. Lett. B91, 61 (1980); L. Brink and P. Howe, “Eleven-Dimensional Supergravity On The Mass - Shell In Superspace,” Phys. Lett. B91, 384 (1980). relax K. Pilch, P. van Nieuwenhuizen and P.K. Townsend, “Compactification Of D = 11 Supergravity On $`S^4`$ (Or 11 = 7 + 4, Too),” Nucl. Phys. B242, 377 (1984). relax R. Kallosh and A. Rajaraman, “Vacua of M-theory and string theory,” Phys. Rev. D58, 125003 (1998) \[hep-th/9805041\]. relax S. Deser and D. Seminara, “Tree amplitudes and two-loop counterterms in D = 11 supergravity,” hep-th/0002241. relax I.R. Klebanov and A.A. Tseytlin, “Entropy of Near-Extremal Black p-branes,” Nucl. Phys. B475, 164 (1996), hep-th/9604089. relax M. Blau, K. S. Narain and E. Gava, “On subleading contributions to the AdS/CFT trace anomaly,” JHEP 9909, 018 (1999) \[hep-th/9904179\]. relax C. Imbimbo, A. Schwimmer, S. Theisen and S. Yankelowicz, “Diffeomorphisms and Holographic Anomalies,” Class. Quant. Grav. 17, 1129 (2000), hep-th/9910267. relax S. Nojiri and S. D. Odintsov, “Weyl anomaly from Weyl gravity,” hep-th/9910113; “On the conformal anomaly from higher derivative gravity in AdS/CFT correspondence,” hep-th/9903033. relax L. Bonora, P. Pasti and M. Bregola, “Weyl Cocycles,” Class. Quant. Grav. 3, 635 (1986). relax S. Deser and A. Schwimmer, “Geometric classification of conformal anomalies in arbitrary dimensions, Phys. Lett. B309, 279 (1993), hep-th/9302047. relax C. Fefferman and C. Graham, “Conformal invariants”, Asterisque, hors serie, 95 (1985); T. Parker and S. Rosenberg, “Invariants of conformal Laplacians”, J. Diff. Geom. 25, 199 (1987); J. Erdmenger, “Conformally covariant differential operators: Properties and applications,” Class. Quant. Grav. 14, 2061 (1997), hep-th/9704108. relax D. Freed, J.A. Harvey, R. Minasian and G. Moore, “Gravitational anomaly cancellation for M-theory fivebranes,” Adv. Theor. Math. Phys. 2, 601 (1998) \[hep-th/9803205\].
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# Lines on Non-degenerate surfaces
## 1. Introduction
Let $`(X,O)`$ be a germ of analytic varieties embedded in $`(^n,O)`$ with a singularity at $`O`$. By abuse of language, we say that $`L`$ is a line in $`(X,O)`$ if $`(L,O)`$ is a smooth curve germ in $`(X,O)`$ and $`L\{0\}`$ is contained in the regular part of $`X`$.
In , lines on hypersurfaces with simple singularities are classified by using the classification machinery. All the hypersurfaces of dimension 2 and 3 with simple or simple elliptic singularities passing through $`x`$-axis are equivalent to (under the coordinate transformation preserving the $`x`$-axis) some surfaces defined by explicit equations. It turns out that the $`A,D,E`$ singularities split in this classification. This says that different smooth curves on the same surface might have different properties.
Let $`\pi :\stackrel{~}{X}(X,O)`$ be a resolution of a surface $`(X,O)`$ with an isolated singularity at the origin $`O`$ and let $`\{E_1,\mathrm{},E_r\}`$ be the exceptional divisors of $`\pi `$. For an exceptional divisor $`E_i`$, let $`_{E_i}`$ denote the set of lines on $`(X,0)`$ whose strict transform intersect $`E_i`$ transversally. It is known that $`_{E_i}`$ is non-empty if and only if there exist a function germ $`h`$ in the maximal ideal $`𝔪`$ such that the multiplicity of $`\pi ^{}h`$ along $`E_i`$ is one and conversely any line in $`X`$ is contained in some $`_{E_i}`$ (). We call $`E_i`$ a normally smooth divisor if $`_{E_i}\mathrm{}`$. Geometrically this implies that $`d\pi (v)0`$ for any tangent vector $`vT_P\stackrel{~}{X}`$ as long as $`PE_i_{ji}E_j`$ and $`v`$ is not tangent to $`E_i`$. If $`E_i`$ is normally smooth, any germ of a curve intersecting $`E_i_{ji}E_j`$ transversely defines a line in $`X`$. Any two lines in the same $`_{E_i}`$ can be connected by an analytic family of lines in $`(X,O)`$.
For a given resolution $`\pi :\stackrel{~}{X}X`$, we consider the integer $`\rho (\pi ):=\mathrm{}\{E_i;_{E_i}\mathrm{}\}`$. This number depends on the resolution. Put $`\rho (X,O)`$ to be the minimal value of $`\rho (\pi )`$. Obviously $`\rho (\pi )=\rho (X,O)`$ if $`\pi :\stackrel{~}{X}X`$ is a minimal resolution. We call $`\rho (\pi )`$ the line index of the resolution $`\pi :\stackrel{~}{X}X`$ and we call $`\rho (X,O)`$ the line index of $`(X,O)`$.
In this paper, we study $`\rho (\pi )`$ where $`\pi `$ is a toric resolution of a non-degenerate surface singularity. Let $`(X,0)(^3,0)`$ be a surface defined by $`f(z_1,z_2,z_3)=0`$ with isolated singularity at the origin. We assume that $`f`$ is non-degenerate in the sense of the Newton boundary (). Let $`\mathrm{\Sigma }^{}`$ be a regular simplicial cone subdivision of the dual Newton diagram $`\mathrm{\Gamma }^{}(f)`$ and let $`\pi :X_\mathrm{\Sigma }^{}(X,0)`$ be the associated toric resolution. We denote $`\rho (\pi )`$ by $`\rho (\mathrm{\Sigma }^{})`$ for simplicity. To each vertex $`P={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)`$ of $`\mathrm{\Sigma }^{}`$, there corresponds an exceptional divisor $`E(P)`$ of $`\pi `$, which may have several components. The multiplicity of $`\pi ^{}z_i`$ along $`E(P)`$ is equal to $`p_i`$ (). Thus by the result of Gonzalez-Sprinberg and Lejeune-Jalabert (), $`E(P)`$ is normally smooth if and only if $`\mathrm{min}(p_1,p_2,p_3)=1`$. We observe that $`\rho (\mathrm{\Sigma }^{})`$ is independent of the choice of $`\mathrm{\Sigma }^{}`$ under certain conditions (see Proposition 6). This allows us to use the canonical toric resolution to determine $`\rho (\mathrm{\Sigma }^{})`$. Note that a toric resolution is not necessarily minimal. So, in general, $`\rho (\mathrm{\Sigma }^{})`$ may be bigger than $`\rho (X,O)`$ (see Example 28). However to have the equality $`\rho (\mathrm{\Sigma }^{})=\rho (X,O)`$, it is enough that $`\pi :X_\mathrm{\Sigma }^{}X`$ is line-equivalent to the minimal resolution (see § 2 for the definition). The purpose of this paper is to give a method to compute $`\rho (\mathrm{\Sigma }^{})`$.
## 2. Line-admissible blowing-ups
Let $`(X,O)`$ be a germ of a surface with an isolated singularity at $`O`$. Suppose that we have a good resolution $`\pi _1:X_1X`$ and let $`E_1,\mathrm{},E_r`$ be the exceptional divisors of $`\pi _1`$. Take a divisor $`E_{i_0}`$ and a point $`Q`$ on $`E_{i_0}`$ and let $`\pi _Q:\stackrel{~}{X}_1X_1`$ be the blowing-up at $`Q`$ and let $`E_Q`$ be the exceptional divisor of $`\pi _Q`$. The following statements are obvious.
###### Proposition 1.
Take a function $`h𝔪`$ and let $`m_i`$ be the multiplicity of $`\pi _1^{}h`$ along $`E_i`$. Then the multiplicity $`m_Q`$ of the pull-back $`\pi _Q^{}(\pi _1^{}h)`$ along $`E_Q`$ is the sum of $`m_i`$ for all $`i`$ such that $`QE_i`$. In particular, $`m_Q1`$, and $`m_Q=1`$ if and only if $`m_{i_0}=1`$ and $`QE_{i_0}_{ii_0}E_i`$.
###### Corollary 2.
Under the situation of Proposition 1, $`E_Q`$ is a normally smooth divisor of the composition $`\pi _1\pi _Q:\stackrel{~}{X}_1X`$ if and only if $`E_{i_0}`$ is a normally smooth divisor of $`\pi _1:X_1X`$ and $`Q`$ is contained in $`E_{i_0}_{ji_o}E_j`$.
We call $`\pi _Q:\stackrel{~}{X}_1X_1`$ a line-admissible blowing-up if either the center $`Q`$ is at the intersection of two exceptional divisor or the supporting divisor is not normally smooth. Suppose that we have another good resolution $`\pi _2:X_2X`$. We say that $`\pi _2:X_2X`$ is line-equivalent to $`\pi _1:X_1X`$ if there exist a finite chain of resolutions $`\pi _i^{}:Y_iX,i=1,\mathrm{},s`$ such that (1) $`Y_1=X_1`$ and $`\pi _1^{}=\pi _1`$ and $`Y_s=X_2`$ and $`\pi _s^{}=\pi _2`$ and (2) any consecutive resolutions factor by either $`\sigma _i:Y_iY_{i+1}`$ or $`\sigma _i^{}:Y_{i+1}Y_i`$, where $`\sigma _i`$ and $`\sigma _i^{}`$ are line-admissible blowing-ups.
An immediate consequence of the definition and Corollary 2 is:
###### Corollary 3.
Assume that $`\pi _i:X_1X,i=1,2`$ are line-equivalent. Then $`\rho (\pi _1)=\rho (\pi _2)`$.
## 3. Toric resolution and the computation of $`\rho (\mathrm{\Sigma }^{})`$
### 3.1. Non-degenerate surfaces
We begin with recalling the toric resolutions of surface singularities since this also helps us to fix some notations. We use the notations of . Let $`(X,O)`$ be the germ of a surface in $`(^3,O)`$ defined by a function $`f:(^3,O)(,O)`$. Hereafter we always assume that $`X`$ has an isolated singularity at $`O`$. Let $`\underset{\nu }{}a_\nu z^\nu `$ be the Taylor expansion of $`f`$. The Newton polyhedron $`\mathrm{\Gamma }_+(f)`$ is by definition the convex hull of $`_{\{\nu ;a_\nu 0\}}\{\nu +^3\}.`$ The Newton boundary $`\mathrm{\Gamma }(f)`$ is by definition the union of the compact faces of $`\mathrm{\Gamma }_+(f)`$.
Let $`N:=\mathrm{Hom}_{}(^3,)`$ be the set of covectors. We identify $`N`$ with $`^3`$ and we denote the elements of $`N`$ by column vectors. Let $`N_+`$ be the set of covectors $`P={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)N`$ with $`p_i0,i=1,2,3`$. Put $`E_1:={}_{}{}^{\mathrm{t}}(1,0,0),E_2:={}_{}{}^{\mathrm{t}}(0,1,0)`$ and $`E_3:={}_{}{}^{\mathrm{t}}(0,0,1)`$. $`P`$ is called strictly positive covector if $`p_j>0`$ for all $`j`$. We denote the minimal value of the linear function $`P`$ on $`\mathrm{\Gamma }_+(f)`$ by $`d(P;f)`$. Put $`\mathrm{\Delta }(P;f)=\{z\mathrm{\Gamma }_+(f)P(z)=d(P;f)\}`$. The face function of $`f`$ with respect to $`P`$ is by definition $`f_P(z)=f_{\mathrm{\Delta }(P;f)}:=_{\nu \mathrm{\Delta }(P;f)}a_\nu z^\nu `$. Two covectors $`P,P^{}N_+`$ are equivalent if and only if $`\mathrm{\Delta }(P;f)=\mathrm{\Delta }(P^{};f)`$. The dual Newton diagram $`\mathrm{\Gamma }^{}(f)`$ of $`X`$ is a conical polyhedral subdivision of $`N_+`$ given by the above equivalent classes.
A surface $`X`$ is called non-degenerate (with respect to the local coordinate $`z`$) if for any strictly positive covector $`PN_+`$, $`X^{}(P):=\{z_{}^{}{}_{}{}^{3}f_P(z)=0\}`$ is a reduced non-singular surface in the complex torus $`_{}^{}{}_{}{}^{3}.`$ The notion of non-degeneracy can be extended to complete intersection varieties (cf. ).
### 3.2. Canonical subdivisions
We assume that $`X`$ is defined by $`f(z_1,z_2,z_3)=0`$ and $`f`$ is non-degenerate. Let $`\mathrm{\Gamma }^{}(f)_2^+`$ be the union of the two-dimensional cones $`\text{Cone}(P,Q)`$ of $`\mathrm{\Gamma }^{}(f)`$ such that the interior points are strictly positive. Let $`\mathrm{\Sigma }^{}`$ be a regular simplicial subdivision of the dual Newton diagram $`\mathrm{\Gamma }^{}(f)`$ and let $`\pi :X_\mathrm{\Sigma }^{}X`$ be the associated toric modification. Let $`𝒱(\mathrm{\Sigma }^{})`$ be the set of strictly positive vertices $`P`$’s of $`\mathrm{\Sigma }^{}`$ such that $`dim\mathrm{\Delta }(P;f)1`$. The exceptional divisors correspond bijectively to $`𝒱(\mathrm{\Sigma }^{})`$ and for each $`P𝒱(\mathrm{\Sigma }^{})`$ we denote the corresponding divisor by $`E(P)`$. Note that $`E(P)`$ need not to be irreducible but it is a disjoint union of rational spheres if $`dim\mathrm{\Delta }(P;f)=1`$. The number of connected components is given by $`r(P)+1`$, where $`r(P)`$ is the number of integral points on the interior of $`\mathrm{\Delta }(P;f)`$ (\[9, III§6\]). The structure of this resolution $`\pi :X_\mathrm{\Sigma }^{}X`$ depends only on the restriction of $`\mathrm{\Sigma }^{}`$ to $`\mathrm{\Gamma }^{}(f)_2^+`$. This follows from the following observation:
###### Proposition 4.
Assume that $`\mathrm{\Sigma }_1^{}`$ is a regular subdivision of $`\mathrm{\Sigma }^{}`$ such that $`𝒱(\mathrm{\Sigma }_1^{})=𝒱(\mathrm{\Sigma }^{})`$. Then the canonical morphism $`\psi :X_{\mathrm{\Sigma }_1^{}}X_\mathrm{\Sigma }^{}`$, which is induced by the morphism of the ambient toric varieties, is an isomorphism.
For any two dimensional cone $`\sigma =\mathrm{Cone}(P,Q)\mathrm{\Gamma }^{}(f)`$, there exists a canonical regular subdivision of $`\sigma `$ which is described as follows. Denote by $`d:=det(P,Q)`$ the greatest common divisor of the absolute values of the $`2\times 2`$ minors of the matrix $`(P,Q)`$. If $`d>1`$, there exists a unique integer $`d_1,1d_1<d`$ such that $`Q_1:=(P+d_1Q)/d`$ is an integral covector. If $`d_1>1`$, repeat the process for $`\mathrm{Cone}(P,Q_1)`$, until a regular subdivision of $`\mathrm{Cone}(P,Q)`$ is obtained. Let $`Q_1,\mathrm{},Q_k`$ be the covectors obtained in this way. Let $`d/d_1=[m_1,\mathrm{},m_{\mathrm{}}]`$ be the continuous fraction expansion. Then $`\mathrm{}=k`$ and the self-intersection number of each component of $`E(Q_i)`$ is $`m_i`$ (cf. \[9, III\]). Note that $`\mathrm{\Delta }(Q_i;f)=\mathrm{\Delta }(P;f)\mathrm{\Delta }(Q;f)`$. This implies $`r(Q_i)`$ is independent of $`i=1,\mathrm{},k`$ and we denote this number by $`r(P,Q)`$. Recall that the continuous fraction is defined inductively by $`[m_1]=m_1`$ and $`[m_1,m_2,\mathrm{},m_k]=m_11/[m_2,\mathrm{},m_k]`$.
A regular simplicial cone subdivision of $`\mathrm{\Gamma }^{}(f)`$ is called a canonical regular subdivision if its restriction to each cone $`\sigma `$ in $`\mathrm{\Gamma }^{}(f)_2^+`$ is canonical in the above sense, and we denote it by $`\mathrm{\Sigma }_{\mathrm{can}}^{}`$. The associated toric resolution is called the canonical toric resolution of $`X`$.
Let $`Q={}_{}{}^{\mathrm{t}}(q_1,q_2,q_3)`$ and $`P={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)`$. Put $`Q_0=Q`$ and $`Q_{k+1}=P`$ and let $`Q_j:={}_{}{}^{\mathrm{t}}(q_{1,j},q_{2,j},q_{3,j}),j=0,\mathrm{},k+1`$. The canonical subdivision enjoys the following property:
###### Lemma 5.
Assume that $`\mathrm{Cone}(P,Q)\mathrm{\Gamma }^{}(f)_2^+`$. Fix an $`\mathrm{}=1,2,3`$.
* If $`q_{\mathrm{}}1`$ , then $`\{q_{\mathrm{},j}\}_{j=0}^{k+1}`$ is monotone increasing in $`j`$ i.e. $`q_{\mathrm{},j+1}q_{\mathrm{},j}`$ for $`0jk`$.
* If $`q_{\mathrm{}}2`$, then either $`\{q_{\mathrm{},j}\}`$ is monotone increasing or monotone decreasing in $`j`$ or there exists a $`j_0`$ ($`1j_0k`$) such that $`q_{\mathrm{},j_0}1`$ and
$$p_{\mathrm{}}=q_{\mathrm{},k+1}\mathrm{}q_{\mathrm{},j_0+1}q_{\mathrm{},j_0}q_{\mathrm{},j_01}\mathrm{}q_{\mathrm{},0}=q_{\mathrm{}}.$$
###### Proof.
We prove the assertion 2). If the assertion does not hold, there exists an index $`j,1jk`$ such that $`q_{\mathrm{},j1}q_{\mathrm{},j}>q_{\mathrm{},j+1}`$. This implies that the self intersection number of each component of $`E(Q_j)`$ is $`(q_{\mathrm{},j1}+q_{\mathrm{},j+1})/q_{\mathrm{},j}>2`$, which is a contradiction (cf. \[9, II(2.3) and III(6.3)\]). The assertion 1) follows from 2) as $`Q_j,j=1,\mathrm{},k`$ are strictly positive. $`\mathrm{}`$
Let $`\mathrm{\Sigma }^{}`$ be any regular simplicial cone subdivision of $`\mathrm{\Gamma }^{}(f)`$ and let $`\pi :\stackrel{~}{X}X`$ be the corresponding toric modification. We denote the line index of $`\pi `$ by $`\rho (\mathrm{\Sigma }^{})`$. Take a two dimensional cone $`\sigma =\text{Cone}(P,Q)\mathrm{\Gamma }^{}(f)_2^+`$. Let $`Q_0:=Q,Q_1,\mathrm{},Q_k,Q_{k+1}:=P`$ be the canonical subdivision of $`\sigma `$ and let $`S_0:=Q,S_1,\mathrm{},S_\eta ,S_{\eta +1}:=P`$ be the vertices of $`\mathrm{\Sigma }^{}`$ on this cone. By \[9, II(2.3)\], $`\{Q_0,\mathrm{},Q_{k+1}\}\{S_0,\mathrm{},S_{\eta +1}\}`$. We consider the condition:
($`\mathrm{}`$): $`\mathrm{\Sigma }^{}`$ has no vertex in the interior of $`\text{Cone}(Q,Q_1)`$.
We say that $`\mathrm{\Sigma }^{}`$ satisfies the ($`\mathrm{}`$)-condition if it satisfies ($`\mathrm{}`$)-condition for any $`\mathrm{Cone}(P,Q)`$ in $`\mathrm{\Gamma }^{}(f)_2^+`$ such that $`Q`$ is not strictly positive. The inclusion $`𝒱(\mathrm{\Sigma }_{\mathrm{can}}^{})𝒱(\mathrm{\Sigma }^{})`$ implies that the following statements.
###### Theorem 6.
There exists a canonical morphism $`\varphi :X_\mathrm{\Sigma }^{}X_{\mathrm{\Sigma }_{\mathrm{can}}^{}}`$. Furthermore $`\varphi `$ is a composition of line-admissible blowing-ups if $`\mathrm{\Sigma }^{}`$ satisfies ($`\mathrm{}`$)-condition. In particular, the line index $`\rho (\mathrm{\Sigma }^{})`$ does not depend on the choice of a toric resolution associated with any regular simplicial subdivision satisfying ($`\mathrm{}`$)-condition and $`\rho (\mathrm{\Sigma }^{})=\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})`$.
###### Proof.
Take a two dimensional cone $`\sigma =\text{Cone}(P,Q)\mathrm{\Gamma }^{}(f)_2^+`$ and assume that $`P`$ is strictly positive. Let $`Q_0:=Q,Q_1,\mathrm{},Q_k,Q_{k+1}:=P`$ be the canonical subdivision of $`\sigma `$ and let $`S_0:=Q,S_1,\mathrm{},S_\eta ,S_{\eta +1}:=P`$ be the vertices of $`\mathrm{\Sigma }^{}`$ on this cone. Write $`S_i={}_{}{}^{\mathrm{t}}(s_{1,j},s_{2,j},s_{3,j})`$. Assume that $`Q_{i_0}=S_\nu `$ and $`Q_{i_0+1}=S_\mu `$ and $`\mu \nu >1`$. Take $`S_j`$ with $`\nu <j<\mu `$ and put $`\alpha _j=det(Q_{i_0},S_j)`$ and $`\beta _j=det(S_j,Q_{i_0+1})`$. Then $`\alpha _j`$ and $`\beta _j`$ are positive integers and $`S_j=\alpha _jQ_{i_0+1}+\beta _jQ_{i_0}`$. This implies that $`s_{1,j}>s_{1,\nu }+s_{1,\mu }`$. Suppose that $`s_1^{\mathrm{max}}=\mathrm{max}\{s_{1,j};\nu <j<\mu \}`$ and put $`\gamma =\mathrm{min}\{\gamma ;s_{1,\gamma }=s_1^{\mathrm{max}}\}`$. Then by \[9, II(2.3)\] the intersection number of (each component of) $`E(S_\gamma )`$ is $`(s_{1,\gamma 1}+s_{1,\gamma +1})/s_{1,\gamma }>2`$. Then the negativity of the intersection number implies that $`s_{1,\gamma 1}+s_{1,\gamma +1}=s_{1,\gamma }`$. Thus each component of $`E(S_\gamma )`$ is a rational sphere of the first kind. This implies also that $`S_\gamma =S_{\gamma 1}+S_{\gamma +1}`$ and $`det(S_{\gamma 1},S_{\gamma +1})=1`$. Put $`𝒱^{}=𝒱(\mathrm{\Sigma }^{})\{S_\gamma \}`$. Then we can extend $`𝒱^{}`$ to get a regular simplicial subdivision $`\mathrm{\Sigma }_{}^{}{}_{}{}^{}`$ such that its restriction to $`\mathrm{\Gamma }^{}(f)_2^+`$ is defined by the vertices $`𝒱^{}`$. Thus we get a toric resolution $`\pi ^{}:X_\mathrm{\Sigma }_{}^{}{}_{}{}^{}X`$. Changing $`\mathrm{\Sigma }^{}`$ outside of $`\mathrm{\Gamma }^{}(f)_2^+`$ if necessary, we may assume by Proposition 4 that $`\mathrm{\Sigma }^{}`$ is a subdivision of $`\mathrm{\Sigma }_{}^{}{}_{}{}^{}`$. Thus we get a canonical morphism $`\psi :X_\mathrm{\Sigma }^{}X_\mathrm{\Sigma }_{}^{}{}_{}{}^{}`$ which factors $`\pi `$ by $`\pi ^{}`$. By the definition, $`\psi `$ is the composition of blowing-up at $`r(S_\gamma )+1`$ intersection points of respective components of $`E(S_{\gamma 1})`$ and $`E(S_{\gamma +1})`$ in $`X_\mathrm{\Sigma }_{}^{}{}_{}{}^{}`$. Note that $`\psi `$ is line-admissible unless $`Q`$ is not strictly positive and $`S_\nu =Q_0`$ and $`S_\mu =Q_1`$. This is the situation where $`\psi `$ is the blowing up at the intersection of $`E(Q_1)`$ and $`E(Q)`$. This does not occur if $`\mathrm{\Sigma }^{}`$ satisfies $`(\mathrm{})`$-condition. Now the assertion follows by the induction on the cardinality of $`𝒱(\mathrm{\Sigma }^{})𝒱(\mathrm{\Sigma }_{\mathrm{can}}^{})`$. $`\mathrm{}`$
### 3.3. Computation of $`\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})`$
Let $`\pi :X_\mathrm{\Sigma }^{}X`$ be a toric resolution. We assume that $`\mathrm{\Sigma }^{}`$ satisfies the ($`\mathrm{}`$)-condition. We define $`𝒱_{\mathrm{ns}}(\mathrm{\Sigma }^{}):=\{P𝒱(\mathrm{\Sigma }^{})P\text{ has 1 as a coordinate }\}`$. We know that $`E(P)`$ is a normally smooth divisor if and only if $`P𝒱_{\mathrm{ns}}(\mathrm{\Sigma }^{})`$. Thus for each $`\text{Cone}(P,Q)\mathrm{\Gamma }^{}(f)_2^+`$, we define $`\rho _{PQ}:=\mathrm{\#}𝒱_{\mathrm{ns}}(\mathrm{\Sigma }^{})\mathrm{Cone}(P,Q)^{}`$, where $`\mathrm{Cone}(P,Q)^{}`$ is the interior of $`\mathrm{Cone}(P,Q)`$. This number is independent of $`\mathrm{\Sigma }^{}`$ by Theorem 6. Recall that $`r(P,Q)`$ is the number of integral points in the interior of $`\mathrm{\Delta }(P;f)\mathrm{\Delta }(Q;f)`$. By the definition we have
(1) $`\rho (\mathrm{\Sigma }^{})=\mathrm{}\{P𝒱_{\mathrm{ns}}(\mathrm{\Sigma }^{});dim\mathrm{\Delta }(P;f)=2\}+{\displaystyle \underset{\mathrm{Cone}(P,Q)\mathrm{\Gamma }^{}(f)_2^+}{}}(r(P,Q)+1)\rho _{PQ}`$
Thus we need only to compute $`\rho _{PQ}`$ for the calculation of $`\rho (\mathrm{\Sigma }^{})`$. Take a cone $`\sigma =\mathrm{Cone}(P,Q)`$ in $`\mathrm{\Gamma }^{}(f)_2^+`$. The following gives a practical method to compute $`\rho _{PQ}`$.
###### Theorem 7.
Let $`P={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)`$ be strictly positive and let $`Q={}_{}{}^{\mathrm{t}}(q_1,q_2,q_3)`$ and assume that $`d:=det(P,Q)>1`$. Let $`Q_i={}_{}{}^{\mathrm{t}}(q_{1,i},q_{2,i},q_{3,i}),i=0,\mathrm{},k+1`$ be the vertices defining the canonical subdivision from $`Q`$ with $`Q_0=Q`$ and $`Q_{k+1}=P`$. Fix an $`\mathrm{}\{1,2,3\}`$. Then
1. For each $`1ik`$, there exists positive integers $`0<\alpha _i,\beta _i<d`$ such that $`Q_i=(\beta _iP+\alpha _iQ)/d`$. Putting $`\alpha _0=\beta _{k+1}=d,\alpha _{k+1}=\beta _0=0`$, they satisfy the inequality:
$$\alpha _i>\alpha _{i+1},\beta _i<\beta _{i+1},i=0,\mathrm{},k$$
2. Let $`𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q)`$ be the set of integral covectors $`R`$ expressed as $`R=(\beta P+\alpha Q)/d`$ where $`\alpha ,\beta `$ are positive integers satisfying
(2) $`\{\begin{array}{cc}& \alpha q_{\mathrm{}}+\beta p_{\mathrm{}}=d,0<\alpha ,\beta <d\hfill \\ & \alpha q_k+\beta p_k0modd(k\mathrm{})\hfill \end{array}`$
and let $`𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q;\mathrm{\Sigma }_{\mathrm{can}}^{})`$ be the set of covectors $`Q_i`$, $`1ik`$ such that $`q_{\mathrm{},i}=1`$. Then $`𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q)=𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q;\mathrm{\Sigma }_{\mathrm{can}}^{})`$. Note that the inequality $`\alpha ,\beta <d`$ follows automatically from the positivity if both $`p_{\mathrm{}}`$ and $`q_{\mathrm{}}`$ are positive.
###### Proof.
The first assertion follows by an inductive argument. Write $`Q_i=(\beta _iP+\alpha _iQ)/d`$ with positive rational numbers $`\alpha _i,\beta _i`$. As $`det(P,Q_i)=\alpha _i`$ and $`det(Q_i,Q)=\beta _i`$, $`\alpha _i,\beta _i`$ are positive integers. By the definition of $`Q_1`$, we can write $`Q_1=(P+\alpha _1Q)/d`$ for some $`0<\alpha _1<d`$. The assertion for $`Q_1`$ holds and $`det(P,Q_1)=\alpha _1`$. Assume that $`Q_j=(\beta _jP+\alpha _jQ)/d`$ with $`0<\alpha _j<d`$. As $`det(P,Q_j)=\alpha _j`$ and $`\{Q_j,\mathrm{},Q_{k+1}\}`$ is the vertices of the canonical subdivision of $`\mathrm{Cone}(P,Q_j)`$, there exists $`\alpha ^{}`$, $`0<\alpha ^{}<\alpha _i`$, such that
$$Q_{j+1}=\frac{1}{\alpha _j}P+\frac{\alpha ^{}}{\alpha _j}Q_j=\frac{1}{\alpha _j}P+\frac{\alpha ^{}}{\alpha _j}\frac{(\beta _jP+\alpha _jQ)}{d}=(\frac{1}{\alpha _j}+\frac{\alpha ^{}\beta _j}{\alpha _jd})P+\frac{\alpha ^{}}{d}Q$$
Thus $`\alpha _{j+1}=\alpha ^{}<\alpha _j<d`$. The inequality $`\beta _{j+1}>\beta _j`$ can be proved similarly by using the fact that $`\{P,Q_k,\mathrm{},Q_1,Q\}`$ is the vertices of the canonical subdivision of the cone $`\mathrm{Cone}(P,Q)`$ from $`P`$ (cf. \[9, II(2.3)\]). Now we show the second assertion. The inclusion $`𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q;\mathrm{\Sigma }_{\mathrm{can}}^{})𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q)`$ is obvious. Suppose that $`R=(\beta P+\alpha Q)/d𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q)`$ is not contained in $`𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q;\mathrm{\Sigma }_{\mathrm{can}}^{})`$. Suppose that $`R\mathrm{Cone}(Q_i,Q_{i+1})^{}`$. Then we can write $`R=mQ_i+nQ_{i+1}`$ for some positive integers $`m,n`$. If $`i1`$, this gives a contradiction by comparing the $`\mathrm{}`$-th coefficient: $`1=mq_{\mathrm{},i}+nq_{\mathrm{},i+1}m+n`$. Suppose that $`i=0`$. Write $`Q_1=(P+\alpha _1Q)/d`$ as above. Then $`R=mQ+(P+\alpha _1Q)n/d=nP/d+(md+n\alpha _1)Q/d`$. Thus we get $`\alpha =md+\alpha _1nd`$ which contradicts to the assumption. $`\mathrm{}`$
###### Remark 8.
The computation of $`𝒱_{\mathrm{ns}}(P,Q)`$ is most difficult for the case $`p_{\mathrm{}},q_{\mathrm{}}>1`$. Assume that $`p_{\mathrm{}},q_{\mathrm{}}>0`$. If we have a solution $`(\alpha _0,\beta _0)`$, the other solutions are reduce to the following equation. Put $`\alpha =\alpha _0+\alpha ^{},\beta =\beta _0+\beta ^{}`$. Then
(3) $`\{\begin{array}{cc}& \alpha ^{}q_{\mathrm{}}+\beta ^{}p_{\mathrm{}}=0\hfill \\ & \alpha ^{}q_k+\beta ^{}p_k0modd(k\mathrm{})\hfill \end{array}`$
Let $`\mathrm{\Delta }:=\mathrm{\Delta }(P;f)\mathrm{\Delta }(Q;f)`$. Let $`T={}_{}{}^{\mathrm{t}}(t_1,t_2,t_3)`$ be a covector in $`𝒱_{\mathrm{ns}}^{(\mathrm{})}(P,Q)`$ (thus $`t_{\mathrm{}}=1`$). Geometrically this implies that $`\mathrm{\Delta }(T;f)=\mathrm{\Delta }`$. In particular, $`\mathrm{\Gamma }_+(f)\{(\nu _1,\nu _2,\nu _3);t_1\nu _1+t_2\nu _2+t_3\nu _3d(T;f)\}`$. This gives a practical way to find such a $`T`$.
The case $`q_{\mathrm{}}=0`$ or $`1`$, the computation is much easier. See Corollary 11.
The canonical subdivision of $`\mathrm{Cone}(P,Q)`$ takes sometimes a lot of computations (see Example 9). Theorem 7 gives us a criterion on the existence or non-existence of normally smooth divisors, without computing the whole subdivision $`Q_i,i=1,\mathrm{},k`$.
###### Example 9.
For simplicity, we write $`x=z_1,y=z_2,z=z_3`$. Let us consider $`f(x,y,z)=x^m+y^n+x^ry^r+z^2`$. We assume that $`m,n>2r`$. Put $`n=n_1r+n_0,m=m_1r+m_0`$ with $`0m_0,n_0r1`$. Then $`\mathrm{\Gamma }(f)`$ has two compact faces whose covectors are $`P={}_{}{}^{\mathrm{t}}(2(nr),2r,nr)/\delta _1`$ and $`Q={}_{}{}^{\mathrm{t}}(2r,2(mr),mr)/\delta _2`$ where $`\delta _1=\mathrm{gcd}(2(nr),2r,nr)`$ and $`\delta _2=\mathrm{gcd}(2r,2(mr),mr)`$ and the corresponding dual Newton diagram is as in Figure 1. Note that $`d:=det(P,Q)`$ is given by $`d=2(mnmrnr)/(\delta _1\delta _2)`$. We consider $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)`$. First we consider the covector $`T_0={}_{}{}^{\mathrm{t}}(1,1,r)`$, which is a weight vector of $`x^ry^r+z^2`$. As $`m,n>2r`$, $`T_0`$ must be on $`\mathrm{Cone}(P,Q)`$. To proceed the further computation, let us assume that $`n,m,r`$ are odd and $`\mathrm{gcd}(m,r)=\mathrm{gcd}(n,r)=1`$. This implies $`\delta _1=\delta _2=1`$. By Theorem 7, we have
$$\{\begin{array}{cc}& 2\beta (nr)+2\alpha r=d\hfill \\ & 2\beta r+2\alpha (mr)0modd\hfill \\ & \beta nr+\alpha mr0modd\hfill \end{array}$$
First we have a canonical solution $`(\alpha _0,\beta _0)=(n2r,m2r)`$ which corresponds to the covector $`T_0={}_{}{}^{\mathrm{t}}(1,1,r)`$. Thus putting $`\alpha =\alpha _0+a`$ and $`\beta =\beta _0+b`$, we can reduce the equation as
$$\{\begin{array}{cc}& 2b(nr)+2ar=0\hfill \\ & 2br+2a(mr)0modd\hfill \\ & bnr+amr0modd\hfill \end{array}$$
Taking the positivity of $`\alpha ,\beta `$ into account, we have the solution
$$\{(\alpha ,\beta )\}=\left\{((n2r)+2j(nr),(m2r)2jr);0j\left[\frac{m_12}{2}\right]\right\}$$
For example, consider the easiest case $`m=n`$. This has a unique solution $`(\alpha ,\beta )=(n2r,n2r)`$ and $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)=\{B\}`$ where $`B={}_{}{}^{\mathrm{t}}(1,1,r)`$. By symmetry, we have $`𝒱_{\mathrm{ns}}^{(2)}=\{B\}`$. Note $`r(P,Q)=1`$. By writing down the equation described by Theorem 7, we can show $`𝒱_{\mathrm{ns}}^{(3)}(P,Q)=\mathrm{}`$.
Now we look at $`\mathrm{Cone}(P,E_1)`$ and $`\mathrm{Cone}(P,E_3)`$. Note that $`det(P,E_1)=r`$ and $`det(P,E_3)=2`$. It is easy to see that there are no normally smooth divisor on these cones. Observe that the computation of canonical subdivision of $`\mathrm{Cone}(P,Q)`$ is not so easy. For example, if $`r=15,n=37`$, then $`B={}_{}{}^{\mathrm{t}}(1,1,15)`$ and first covector $`B_1`$ (from $`Q`$) is given by $`(P+223Q)/518={}_{}{}^{\mathrm{t}}(13,19,240)`$ and $`518/223=[3,2,2,12,2,2,3]`$ and it takes some computation to complete the subdivision.
The following lemma describes the covectors corresponding to the non-compact faces.
###### Lemma 10.
Assume that $`X=\{f(z_1,z_2,z_3)=0\}`$ and assume that $`f`$ is non-degenerate and $`\mathrm{\Gamma }(f)`$ has at least one compact two dimensional face for simplicity. Suppose that $`z_2=z_3=0`$ is a line in $`X`$. (So $`f`$ is not convenient.) Then there is a unique covector $`Q={}_{}{}^{\mathrm{t}}(q_1,q_2,q_3)\mathrm{Vertex}(\mathrm{\Gamma }^{}(\mathrm{f}))`$ such that $`q_1=0`$. Furthermore $`Q`$ takes the form $`{}_{}{}^{\mathrm{t}}(0,1,q_3)`$ or $`{}_{}{}^{\mathrm{t}}(0,q_2,1)`$.
There exists a unique covector $`P={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)`$ which corresponds to a compact divisor and adjacent to $`Q`$ in $`\mathrm{\Gamma }^{}(f)_2^+`$. Then we have $`det(P,Q)=p_1`$.
###### Proof.
As $`X`$ has an isolated singularity, $`f`$ must contain a monomial of type $`z_1^az_2`$ or $`z_1^az_3`$. Suppose that $`B:=(a,1,0)\mathrm{\Gamma }(f)`$. Let $`C=(b,0,c)`$ be the vertex of $`\mathrm{\Gamma }(f)\{z_2=0\}`$ adjacent to $`B`$ by an edge. It is clear that the non-compact face $`\mathrm{\Xi }`$ which has $`\overline{BC}`$ as a face and is unbounded to the direction of the $`z_1`$-axis has covector $`Q={}_{}{}^{\mathrm{t}}(0,c,1)`$. One can see that there exists no other non-compact face which is unbounded to the $`z_1`$-axis direction and bounded to $`z_2,z_3`$-direction. Let $`\mathrm{\Delta }`$ be the compact face which has $`\overline{BC}`$ as a boundary and let $`P={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)`$ be the corresponding covector. As $`\mathrm{\Delta }(P;f)`$ contains $`B,C`$, we need to have $`p_1a+p_2=bp_1+cp_3`$. Now the last assertion follows from $`det(P,Q)=\mathrm{gcd}(p_1,p_2cp_3)=\mathrm{gcd}(p_1,p_1(ba))=p_1`$. $`\mathrm{}`$
The following corollary describes explicitly $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)`$ in the case $`q_1=0`$ or $`1`$.
###### Corollary 11.
With the assumptions of Theorem 7, we have the following.
* Assume $`q_1=0`$. Then $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)\mathrm{}`$ if and only if $`d:=det(P,Q)>1`$ and $`d=p_1`$. In this cases, $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)=\{Q_1\}`$. If $`QE_2,E_3`$, then $`\{y=z=0\}X`$ and $`d=det(P,Q)=p_1`$.
* Assume $`q_1=1`$. Then $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)\mathrm{}`$ if and only if $`d>p_1`$. In this case, we have $`Q_i=(iP+(dip_1)Q)/d`$ for $`i=1,\mathrm{},[d/p_1]`$ and $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)=\{Q_i;i=1,\mathrm{},[d/p_1]\}`$.
###### Proof.
Assume that $`Q^{}=(\beta P+\alpha Q)/d𝒱_{\mathrm{ns}}^{(1)}(P,Q)`$ with $`0<\alpha ,\beta <d`$.
1) If $`q_1=0`$, we have $`\mathrm{gcd}(q_2,q_3)=1`$. As $`d=\mathrm{gcd}(p_1q_2,p_1q_3,p_2q_3p_3q_2)=\mathrm{gcd}(p_1,p_2q_3p_3q_2)`$, $`d`$ divides $`p_1`$. Thus $`Q^{}𝒱_{\mathrm{ns}}^{(1)}(P,Q)`$ if and only if $`d=p_1`$ and $`\beta =1`$. In this case, $`Q^{}=Q_1`$ and $`𝒱_{\mathrm{ns}}^{(1}(P,Q)=\{Q_1\}`$. Assume that $`QE_2,E_3`$. By the definition of $`\mathrm{\Gamma }^{}(f)_2^+`$, $`\mathrm{\Delta }(Q;f)`$ is a non-compact face with dimension 2. In particular, $`\{y=z=0\}X`$. By Lemma 10, we have $`d=p_1`$.
2) Suppose that $`q_1=1`$. Then $`\beta p_1+\alpha =d`$. This implies $`d>p_1`$. Put $`d=rp_1+d^{}`$ with $`0d^{}<p_1`$ and $`r=[d/p_1]`$. Then by the above equality, we have $`(\alpha ,\beta )=(djp_1,j),j=1,\mathrm{},[d/p_1]`$. Put $`Q_j^{}:=(jP+(djp_1)Q)/d`$. By the definition, $`d`$ divides the minors of $`(P,Q)`$ which are $`p_1q_2p_2,p_1q_3p_3,p_2q_3p_3q_2`$. Thus $`\beta p_j+\alpha q_j=\beta p_j+(d\beta p_1)q_j\beta (p_jp_1q_j)0modd`$ for $`j=2,3`$. Thus $`Q_j^{}`$ is an integral covector for $`\beta =1,\mathrm{},r`$. It is clear that $`Q_1^{}=Q_1`$. Assume that $`Q_r^{}=Q_\iota `$ for some $`\iota `$. By the monotonity of the coefficients (Lemma 5), we have $`Q_j𝒱_{\mathrm{ns}}^{(1)}(P,Q)`$ for $`j\iota `$. Thus $`\iota =r`$ and $`Q_j^{}=Q_j`$ for $`jr`$. $`\mathrm{}`$
###### Remark 12.
In the case of non-convenient surface with $`q_1=0`$, the divisor $`E(Q_1)`$ corresponds to the deformations of the line $`z_2=z_3=0`$. In fact, $`E(Q)`$ is a non-compact divisor which is the strict transform of $`z_1`$-axis and $`E(Q)`$ intersects transversely with $`E(Q_1)`$.
For $`R𝒱_{\mathrm{ns}}^{(\mathrm{})}`$, write $`R=(\beta P+\alpha Q)/d`$. We call $`\beta /d`$ the P-coefficient of $`R`$.
###### Corollary 13.
With the assumptions of Theorem 7, suppose that $`q_1>1`$. Let $`\overline{Q}=(\overline{\beta }P+\overline{\alpha }Q)/d𝒱_{\mathrm{ns}}^{(\mathrm{})}`$ and $`\underset{¯}{Q}=(\underset{¯}{\beta }P+\underset{¯}{\alpha }Q)/d𝒱_{\mathrm{ns}}^{(\mathrm{})}`$ be the covectors with maximal and minimal $`P`$-coefficients in $`𝒱_{\mathrm{ns}}^{(\mathrm{})}`$. Then
(4)
$$\rho _{PQ}^{(\mathrm{})}=1+|det(\overline{Q},\underset{¯}{Q})|=1+\frac{|\overline{\beta }\underset{¯}{\alpha }\overline{\alpha }\underset{¯}{\beta }|}{d}$$
###### Proof.
Denote by $`d^{}:=|det(\overline{Q},\underset{¯}{Q})|`$. Suppose that $`\underset{¯}{Q}=Q_i`$ and $`\overline{Q}=Q_{i+j}`$. Then $`𝒱_{\mathrm{ns}}^{(\mathrm{})}=\{Q_i,\mathrm{},Q_{i+j}\}`$ by Lemma 5 and $`\rho _{PQ}^{(\mathrm{})}=j+1`$. By the assumption, we have $`Q_{i+1}=(Q_{i+j}+(d^{}1)Q_i)/d^{}`$. As the continuous fraction $`d^{}/(d^{}1)`$ is given by $`[2,\mathrm{},2]`$ ($`(d^{}1)`$ copies of 2), we get $`j1=d^{}1`$ and the assertion follows immediately. $`\mathrm{}`$
## 4. Applications
### 4.1. Weighted homogeneous surfaces
In this section we study lines on weighted homogeneous surface singularities, which are classified as follows ( ):
$`\begin{array}{cc}X_\mathrm{I}:\hfill & h_\mathrm{I}=x^a+y^b+z^c=0,\hfill \\ X_{\mathrm{II}}:\hfill & h_{\mathrm{II}}=x^ay+y^b+z^c=0,\hfill \\ X_{\mathrm{III}}:\hfill & h_{\mathrm{III}}=x^ay+xy^b+z^c=0,\hfill \\ X_{\mathrm{IV}}:\hfill & h_{\mathrm{IV}}=x^ay+y^bz+z^c=0,\hfill \\ X_\mathrm{V}:\hfill & h_\mathrm{V}=x^ay+y^bz+z^cx=0,\hfill \\ X_{\mathrm{VI}}:\hfill & h_{\mathrm{VI}}=xy+z^c=0,\hfill \\ X_{\mathrm{VII}}:\hfill & h_{\mathrm{VII}}=x^az+y^bz+z^c+tx^{c_1}y^{c_2}=0,t0\hfill \\ X_{\mathrm{VIII}:}\hfill & h_{\mathrm{VIII}}=x^ay+xy^b+xz^c+ty^{c_1}z^{c_2}=0,t0.\hfill \end{array}`$
The surface $`X_\mathrm{I}`$ is called a Pham-Brieskorn surface. This type of surfaces have been studied in the previous paper . The surface $`X_{\mathrm{VI}}`$ is an $`A_{c1}`$ type singularity. There are exact $`c1`$ families of lines on this surface (see ). On surface $`X_{\mathrm{VII}}`$ and $`X_{\mathrm{VIII}}`$, the term $`y^{c_1}z^{c_2}`$ must be on the supporting plane of the previous three monomials. Thus $`a,b,c`$ are not arbitrary. The Newton boundaries of the surfaces other than $`X_{\mathrm{VI}},X_{\mathrm{VII}}`$ and $`X_{\mathrm{VIII}}`$ are triangles. Note that for a weighted homogeneous surface, the Newton boundary has only one compact 2-dimensional face. Let $`P={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)`$ be the corresponding covector. The formula (1) in §2 reduces to
(5) $`\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})=\epsilon +{\displaystyle \underset{\mathrm{Cone}(P,Q)\mathrm{\Gamma }^{}(f)_2^+}{}}(r(P,Q)+1)\rho _{PQ}(\mathrm{\Sigma }_{\mathrm{can}}^{}).`$
where $`\epsilon =1`$ if $`P𝒱_{\mathrm{ns}}(\mathrm{\Sigma }_{\mathrm{can}}^{})`$ and $`\epsilon =0`$ otherwise.
For each type of surfaces, one can calculate $`\rho _{PQ}(\mathrm{\Sigma }_{\mathrm{can}}^{})`$ for each $`\mathrm{Cone}(P,Q)`$ in the dual Newton diagram by using the method described in the previous sections.
###### Lemma 14.
Assume that $`\mathrm{Cone}(P,E_i)`$ be a cone in $`\mathrm{\Gamma }^{}(f)_2^+`$. Then $`det(P,E_i)`$ is given by $`\delta _i:=\mathrm{gcd}(p_j,p_k)`$ with $`\{i,j,k\}=\{1,2,3\}`$. Assume that $`\delta _i>1`$.
* $`𝒱_{\mathrm{ns}}^{(i)}(P,E_i)\mathrm{}`$ if and only if $`\delta _i>p_i`$ and $`\rho _{PE_i}^{(i)}=\left[\frac{\delta _i}{p_i}\right]`$.
* $`𝒱_{\mathrm{ns}}^{(j)}(P,E_i)\mathrm{}`$ if and only if $`p_j|p_k`$. In this case, $`\rho _{PE_i}^{(j)}=1`$.
* $$\rho _{PE_i}=\{\begin{array}{cc}& 0,\text{if}\left[\frac{\delta _i}{p_i}\right]=0\text{and}\delta _i<\mathrm{min}\{p_j,p_k\}\hfill \\ & \mathrm{max}\{1,\left[\frac{\delta _i}{p_i}\right]\},\text{otherwise}\hfill \end{array}$$
###### Proof.
This follows from Corollary 11. $`\mathrm{}`$
###### Lemma 15.
Let $`\mathrm{Cone}(P,Q)`$ be a cone in $`\mathrm{\Gamma }^{}(f)_2^+`$ with $`Q={}_{}{}^{\mathrm{t}}(0,c,1)`$. Suppose that $`det(P,Q)=p_1>1`$. Then
$$\rho _{PQ}=\{\begin{array}{cc}& \mathrm{max}\{1,\left[\frac{p_1}{p_2}\right],\left[\frac{p_1}{p_3}\right]\},c=1\hfill \\ & \rho _{PQ}^{(2)}+\mathrm{max}\{1,\left[\frac{p_1}{p_3}\right]\}\epsilon ,c>1\hfill \end{array}$$
where $`\epsilon =1`$ if either $`Q_1𝒱_{\mathrm{ns}}^{(2)}(P,Q)`$ or $`Q_{j_1}𝒱_{\mathrm{ns}}^{(2)}(P,Q)`$ with $`j_1:=\left[\frac{p_1}{p_3}\right]1`$ and $`\epsilon =0`$ otherwise.
###### Proof.
Let $`Q_1,\mathrm{},Q_k`$ be the primitive covectors in $`\mathrm{Cone}(P,Q)`$ inserted by the canonical subdivision from $`Q`$. If $`c=1`$, the assertion is immediate from Corollary 11, as $`q_{1,1}=1`$. We assume that $`c>1`$. If $`\left[p_1/p_3\right]=0`$, the assertion is obvious. Assume that $`\left[p_1/p_3\right]1`$. By Corollary 11, $`Q_j`$ is given by $`(jP+(p_1jp_3)Q)/p_1`$ for $`1jj_1`$. Thus $`q_{2,j}=cj(cp_3p_2)/p_1`$. If $`cp_3p_2<0`$, $`q_{2,j}`$ is monotone increasing by Lemma 5 and we see that $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)=\mathrm{}`$ and the assertion follows immediately. Assume that $`cp_3p_20`$. Then $`q_{2,j}`$ is monotone decreasing for $`0jj_1`$. Thus $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)𝒱_{\mathrm{ns}}^{(3)}(P,Q)\mathrm{}`$ if and only if $`q_{2,j_1}=1`$. If this is the case, $`Q_{j_1}`$ is the unique covector in common. Thus the assertion follows from these observations. $`\mathrm{}`$
### 4.2. Normally smooth divisors on $`X_{\mathrm{II}}`$
By using Lemmas 14 and 15, we can compute the number $`\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})`$. We show this by considering the surface $`X_{\mathrm{II}}`$. One can do the same consideration for the other types of surfaces. Let $`X_{\mathrm{II}}:h_{\mathrm{II}}(x,y,z)=x^ay+y^b+z^c=0`$. Put $`\widehat{a}:=\mathrm{gcd}(a,b1),e:=\mathrm{gcd}(b,c)`$ and $`d:=\mathrm{gcd}(c(b1),ac,ab)=e\mathrm{gcd}(a,c(b1)/e)`$. The dual Newton diagram $`\mathrm{\Gamma }^{}(h_{\mathrm{II}})_2^+`$ consists of three cones: $`\mathrm{Cone}(P,Q),\mathrm{Cone}(P,E_1)`$ and $`\mathrm{Cone}(P,E_3)`$ where $`P:={}_{}{}^{\mathrm{t}}(c(b1)/d,ac/d,ab/d)`$ and $`Q:={}_{}{}^{\mathrm{t}}(0,c,1)`$.
The following three propositions are special cases of Lemmas 14 and 15.
###### Proposition 16.
$`\mathrm{Cone}(P,E_1)`$ is regular if and only if $`a`$ divides $`c(b1)/e`$. Assume that $`a(c(b1)/e)`$. Then
* $`𝒱_{\mathrm{ns}}^{(1)}(P,E_1)\mathrm{}`$ if and only if $`ae>(b1)c`$. And in this case $`\rho _{PE_1}^{(1)}=\left[\frac{ae}{(b1)c}\right]`$.
* $`𝒱_{\mathrm{ns}}^{(2)}(P,E_1)\mathrm{}`$ if and only if $`c|b`$.
* $`𝒱_{\mathrm{ns}}^{(3)}(P,E_1)\mathrm{}`$ if and only if $`b|c`$.
* $`\rho _{PE_1}=\mathrm{max}\{\rho _{PE_1}^{(2)},\rho _{PE_1}^{(3)},\left[\frac{ae}{(b1)c}\right]\}`$. $`\mathrm{}`$
###### Proposition 17.
As $`det(P,E_3)=c\widehat{a}/d`$, $`\mathrm{Cone}(P,E_3)`$ is regular if and only if $`d=c\widehat{a}`$. Assume that $`c\widehat{a}>d`$. Then
* $`𝒱_{\mathrm{ns}}^{(1)}(P,E_3)\mathrm{}`$ if and only if $`(b1)|a`$.
* $`𝒱_{\mathrm{ns}}^{(2)}(P,E_3)\mathrm{}`$ if and only if $`a|(b1)`$.
* $`𝒱_{\mathrm{ns}}^{(3)}(P,E_3)\mathrm{}`$ if and only if $`c\widehat{a}>ab`$ and $`\rho _{PE_3}^{(3)}=\left[\frac{c\widehat{a}}{ab}\right]`$.
* $`\rho _{PE_3}=\mathrm{max}\{\rho _{PE_3}^{(1)},\rho _{PE_3}^{(2)},\left[\frac{c\widehat{a}}{ab}\right]\}`$.
Recall that $`\rho _{P,E_i}^{(j)}1`$ for $`i=1,3`$ and $`ji`$ by Lemma 5.
###### Proposition 18.
$`\mathrm{Cone}(P,Q)`$ is regular if and only if $`(b1)c`$ divides $`ae`$, or equivalently $`(b1)|a`$ and $`c|b\frac{a}{b1}`$. Assume that $`\mathrm{Cone}(P,Q)`$ is not regular. Then we have
* $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)=\{Q_1\}`$.
* $`𝒱_{\mathrm{ns}}^{(3)}(P,Q)\mathrm{}`$ if and only if $`c(b1)>ab`$. And in this case $`\rho _{PQ}^{(3)}=\left[\frac{c(b1)}{ab}\right]`$.
* $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)\mathrm{}`$ if and only if there exist positive integers $`\alpha `$ and $`\beta `$ such that
(6) $`a\beta +d\alpha =b1,`$
(7) $`ab\beta +d\alpha 0modc(b1).`$
The second condition can be replaced by $`a\beta +10\mathrm{modulo}\mathrm{c}`$.
###### Proof.
The last assertion follows from by (6) as $`ab\beta +d\alpha =(b1)(a\beta +1)`$. $`\mathrm{}`$
The non-trivial computation is required only for $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)`$ which we will explain more in detail. Write $`b=eb_1`$ and $`c=ec_1`$.
###### Corollary 19.
I. For $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)\mathrm{}`$, it is necessary that
(8) $`\mathrm{gcd}(a,c)=1,b>a,c`$
In this case, we have $`d=e\widehat{a}`$ and $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)`$ is the set of covectors $`T=(\alpha Q+\beta P)/d`$ which satisfies
(9) $`a\beta +e\widehat{a}\alpha =b1`$
(10) $`0<\alpha ,\beta `$
(11) $`be\widehat{a}\alpha 0\mathrm{modulo}\mathrm{c}`$
II. Furthermore $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)`$ is non-empty if $`[b/c]a+\widehat{a}`$.
###### Proof.
From the congruence $`a\beta +10`$ modulo $`c`$, it is clear that $`\mathrm{gcd}(a,c)=1`$. Thus $`d=e\mathrm{gcd}(a,c_1(b1))=e\widehat{a}`$. The equality (11) results from
$$a\beta +1=bd\alpha =e(b_1\widehat{a}\alpha )0\mathrm{modulo}\mathrm{c}$$
Thus $`b>a\beta a`$ and $`b>c`$. The last congruence equation is equivalent to $`b_1\widehat{a}\alpha 0`$ $`\mathrm{modulo}\mathrm{c}_1`$.
Assume that $`[b/c]a\widehat{a}0`$. As $`\mathrm{gcd}(\widehat{a},b_1)=1`$, there exists positive integer $`\alpha _0`$, $`0<\alpha _0<c_1`$, such that $`b_1\widehat{a}\alpha _00`$ modulo $`c_1`$. Put $`b_1\alpha _0\widehat{a}=j_0c_1`$. We see that $`j_0=b_1/c_1\alpha _0\widehat{a}/c_1>[b/c]\widehat{a}`$. Take $`\alpha `$ which satisfies the congruence $`a\beta +10`$ modulo $`c`$. Then $`\alpha `$ takes the form $`\alpha =\alpha _0+jc_1`$ with $`j𝐍`$ and thus $`b_1\widehat{a}\alpha =(j_0j\widehat{a})c_1`$. For the positivity of $`\beta `$, we need to have $`0j<j_0/\widehat{a}`$. The integrity of $`T`$ implies
$$e(b_1\widehat{a}\alpha )1=ec_1(j_0j\widehat{a})10\mathrm{modulo}\mathrm{a}$$
As $`j`$ can move $`0j<j_0/\widehat{a}`$ and $`j_0>[b/c]\widehat{a}a`$ or $`j_0/\widehat{a}>a/\widehat{a}`$, this congruence equation has a positive solution $`j_1,0j_1j_0/\widehat{a}`$. Then put $`\beta =(ec_1(j_0j_1\widehat{a})1)/a`$ for such a solution $`j_1`$. This gives a covector $`T=(\alpha Q+\beta P)𝒱_{\mathrm{ns}}^{(2)}(P,Q)`$. $`\mathrm{}`$
###### Example 20.
Consider $`X_{\mathrm{II}}:x^9y+y^b+z^8=0`$ with $`b=22+36k`$. Then $`e=2,\widehat{a}=3`$ and the equation is
$`9\beta +6\alpha =21+36k,9\beta +10\mathrm{modulo}8`$
In this case, $`[b/c]a\widehat{a}=(22+36k)/8120`$ if $`k37/18`$. For $`k3`$ (in fact, for $`k2`$), we have a solution $`(\alpha ,\beta )=(6k7,7)`$. In this case, $`P={}_{}{}^{\mathrm{t}}(28+48k,12,33+54k)`$ and $`Q={}_{}{}^{\mathrm{t}}(0,8,1)`$ and $`T:=(\alpha Q+\beta Q)/(28+48k)={}_{}{}^{\mathrm{t}}(7,1,8)`$. We leave the computation of the other covectors in $`𝒱_{\mathrm{ns}}^{(2)}(P,Q)`$ to the reader.
### 4.3. The minimality of the canonical toric resolutions
We study when the canonical toric resolution of a weighted homogeneous surface is minimal. Though the canonical toric resolution is not always minimal (see Example 28), we can expect that the minimality hold for almost all classes of non-degenerate surfaces. By \[9, III(6.3)\], for each weighted homogeneous surface the resolution graph associated with the canonical toric resolution is star-shaped. Hence, when the resolution graph has at least three arms, the canonical resolution is minimal.
We have the following general statement which is very helpful to see if a given toric modification is minimal.
###### Lemma 21.
Let $`X:=f^1(0)`$ be a non-degenerate surface. Suppose that $`P\mathrm{\Gamma }^{}(f)`$ is the strictly positive covector corresponding to a compact face $`\mathrm{\Delta }`$ of the Newton boundary $`\mathrm{\Gamma }(f).`$
* Let $`\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_{\mathrm{}}`$ be the boundary edges of $`\mathrm{\Delta }`$. The exceptional divisor $`E(P)`$ is rational if and only if
$$\frac{6\mathrm{V}\mathrm{o}\mathrm{l}(\mathrm{Cone}\mathrm{\Delta })}{d(P;f)}+\underset{i=1}{\overset{l}{}}(r(\mathrm{\Delta }_i)+1)=2$$
where $`\mathrm{Cone}\mathrm{\Delta }`$ is the cone over $`\mathrm{\Delta }`$ with vertex $`O`$ and $`r(\mathrm{\Delta }_i)`$ is the number of integral points in the interior of $`\mathrm{\Delta }_i`$.
* The canonical toric resolution $`\pi :\stackrel{~}{X}(X,0)`$ is not minimal if and only if there exists a compact face $`\mathrm{\Delta }`$ of $`\mathrm{\Gamma }(f)`$ such that $`E(P)`$ is rational, $`E(P)^2=1`$ and $`E(P)`$ intersects at most two other exceptional divisors where $`P`$ is the covector corresponding to $`\mathrm{\Delta }`$.
###### Proof.
The first statement is a conclusion of \[9, III(6.4)\]. The assertion 2) follows from the Castelnuovo-Enriques criterion and \[9, III §4(A) and §6\]. $`\mathrm{}`$
###### Theorem 22.
Let $`X`$ be one of the surfaces of type $`X_{\mathrm{II}}`$, $`X_{\mathrm{III}}`$, $`X_{\mathrm{IV}}`$, $`X_\mathrm{V}`$, $`X_{\mathrm{VII}}`$ or $`X_{\mathrm{VIII}}`$. We assume that $`a,b,c>1`$ in 4.1. Then the canonical toric resolution of $`X`$ is minimal. In particular, $`\rho (X,0)=\rho (\mathrm{\Sigma }_{\mathrm{can}}^{}).`$
###### Proof.
We first check when the central exceptional divisor $`E(P)`$ is rational by using Lemma 21 (see also \[9, III(6.9)\]). If this is the case, we compute the number of arms from $`E(P)`$. If this number is less than 3, we show that $`E(P)^22`$. Recall that the number of arms in the resolution graph is the sum of $`r(P,Q)+1`$ for non-regular cones $`\mathrm{Cone}(P,Q)\mathrm{\Gamma }^{}(f)_2^+`$.
(II). Let $`X=X_{\mathrm{II}}:x^ay+y^b+z^c=0`$. Put $`e=\mathrm{gcd}(b,c),\widehat{a}=\mathrm{gcd}(a,b1)`$. Then $`P={}_{}{}^{\mathrm{t}}(c(b1),ac,ab)/d`$ with $`d=e\mathrm{gcd}(a,c(b1)/e)`$. Note that $`r(P,Q)+1=1`$, $`r(P,E_1)+1=e`$ and $`r(P,E_3)+1=\widehat{a}`$. By loc. cit. $`E(P)`$ is rational if and only if 1) $`e=\mathrm{gcd}(c,a/\widehat{a})=1`$ or 2) $`\widehat{a}=\mathrm{gcd}(a,c/e)=1`$. If 1) holds, then $`d=\widehat{a}`$. We have $`det(P,Q)=c(b1)/\widehat{a}>1`$, $`det(P,E_3)=c>1`$ and $`det(P,E_1)=a/\widehat{a}`$. If $`\widehat{a}=a`$, $`\mathrm{Cone}(P,E_3)`$ gives $`\widehat{a}=a`$ arms. Hence, in any case the resolution graph of $`X_{\mathrm{II}}`$ has at least three arms centered at $`E(P)`$.
In case 2), we have $`det(P,Q)=c(b1)/e>1`$, $`det(P,E_1)=a>1`$ and $`det(P,E_3)=c/e`$. If $`e<c`$, we have at least three arms in the resolution graph. Suppose that $`e=c`$. Then the number of arms at $`E(P)`$ is $`e+13`$, unless $`b=2`$ and $`e=c=2`$. In this case, the resolution graph has two similar arms and $`E(P)`$ is normally smooth with $`E(P)^22`$.
Outline of other cases:
(III) Let $`X_{\mathrm{III}}:x^ay+xy^b+z^c=0`$. Then $`P={}_{}{}^{\mathrm{t}}(c(b1),c(a1),ab1)/d`$ with $`d=e\mathrm{gcd}(c,(ab1)/e)`$ and $`e=\mathrm{gcd}(a1,b1)`$. The dual Newton diagram $`\mathrm{\Gamma }^{}(f)_2^+`$ has 3 arms $`\mathrm{Cone}(P,E_3)`$, $`\mathrm{Cone}(P,Q)`$, $`\mathrm{Cone}(P,R)`$ where $`Q={}_{}{}^{\mathrm{t}}(0,c,1)`$ and $`R={}_{}{}^{\mathrm{t}}(c,0,1)`$. The central divisor $`E(P)`$ is rational if and only if $`\mathrm{gcd}(c,(ab1)/e)=1`$. If $`E(P)`$ is rational, then $`d=e`$ and $`det(P,Q)=c(b1)/e>1,det(P,R)=c(a1)/e>1,`$ and $`det(P,E_3)=c>1`$. Hence, the resolution graph has at least three arms.
(IV) Let $`X_{\mathrm{IV}}:x^ay+y^bz+z^c=0`$. Then $`P:={}_{}{}^{\mathrm{t}}(bcc+1,a(c1),ab)/d`$ with $`d=e\mathrm{gcd}(a,(bcc+1)/e)`$ and $`e:=\mathrm{gcd}(b,c1)`$. The dual Newton diagram $`\mathrm{\Gamma }^{}(f)_2^+`$ has 3 arms $`\mathrm{Cone}(P,E_1)`$, $`\mathrm{Cone}(P,Q)`$, $`\mathrm{Cone}(P,S)`$ where $`Q={}_{}{}^{\mathrm{t}}(0,c,1)`$ and $`S={}_{}{}^{\mathrm{t}}(1,0,a)`$. The divisor $`E(P)`$ is rational if and only if $`\mathrm{gcd}(a,(bcc+1)/e)=1`$ which is equivalent to $`d=e`$. We have $`det(P,E_1)=a>1`$, $`det(P,S)=a(c1)/e>1`$ and $`det(P,Q)=(bcc+1)/e`$. As $`\mathrm{Cone}(P,E_1)`$ has $`e`$-copies of arms, $`E(P)`$ has at least three arms.
(V) Let $`X_\mathrm{V}:x^ay+y^bz+z^cx=0`$. Then $`P:={}_{}{}^{\mathrm{t}}(bcc+1,caa+1,abb+1)/d`$ with $`d=\mathrm{gcd}(bcc+1,caa+1,abb+1)`$. The dual Newton diagram $`\mathrm{\Gamma }^{}(f)_2^+`$ has 3 arms $`\mathrm{Cone}(P,Q)`$, $`\mathrm{Cone}(P,S)`$, $`\mathrm{Cone}(P,T)`$ where $`Q={}_{}{}^{\mathrm{t}}(0,c,1)`$, $`S={}_{}{}^{\mathrm{t}}(1,0,a)`$ and $`T:={}_{}{}^{\mathrm{t}}(b,1,0)`$. The divisor $`E(P)`$ is rational if and only if $`d=1`$. In this case, we have $`det(P,Q)=bcc+1>1`$, $`det(P,S)=caa+1>1`$ and $`det(P,T)=abb+1>1`$. Thus $`E(P)`$ has three arms.
(VII) Let $`X_{\mathrm{VII}}:x^az+y^bz+z^c+tx^{c_1}y^{c_2}=0`$. Then $`P={}_{}{}^{\mathrm{t}}(b(c1),a(c1),ab)/\delta `$ with $`\delta =\mathrm{gcd}(b(c1),a(c1),ab)`$. The dual Newton diagram $`\mathrm{\Gamma }^{}(f)_2^+`$ has 4 arms $`\mathrm{Cone}(P,Q)`$, $`\mathrm{Cone}(P,S)`$, $`\mathrm{Cone}(P,E_1)`$, $`\mathrm{Cone}(P,E_2)`$ where $`Q={}_{}{}^{\mathrm{t}}(0,1,c_2)`$ and $`S={}_{}{}^{\mathrm{t}}(1,0,c_1)`$. By the weighted homogenuity, we have the equality $`b(c1)c_1+a(c1)c_2=abc`$ which implies that $`(c1)|ab`$. Hence $`\delta =(c1)\mathrm{gcd}(a,b,ab/(c1))`$. By loc. cit., $`E(P)`$ is rational if and only if either (i) $`\mathrm{gcd}(a,b)=\mathrm{gcd}(a,c1)=1`$, or (ii) $`\mathrm{gcd}(a,b)=\mathrm{gcd}(b,c1)=1`$. By symmetry, we may assume that the first case (i). Then $`\delta =c1`$, $`det(P,Q)=b>1`$, $`det(P,S)=a>1`$, $`det(P,E_1)=a>1`$. Thus the resolution graph has at least three arms.
(VIII) Let $`X_{\mathrm{VIII}}:x^ay+xy^b+xz^c+ty^{c_1}z^{c_2}=0`$. Then $`P={}_{}{}^{\mathrm{t}}(c(b1),c(a1),b(a1))/\delta `$ with $`\delta =\mathrm{gcd}(c(b1),c(a1),b(a1)).`$ By the weighted homogenuity, we must have $`c(a1)c_1+b(a1)c_2=c(ab1)`$ which implies that $`(a1)|c(ab1)`$ and $`cc_1+bc_2=bc+c(b1)/a1`$. Thus $`\delta =(a1)\mathrm{gcd}(b,c,c(b1)/(a1))`$. The dual Newton diagram $`\mathrm{\Gamma }^{}(f)_2^+`$ has 4 arms $`\mathrm{Cone}(P,E_3)`$, $`\mathrm{Cone}(P,Q)`$, $`\mathrm{Cone}(P,S)`$ and $`\mathrm{Cone}(P,T)`$ where $`Q={}_{}{}^{\mathrm{t}}(0,c,1)`$, $`S={}_{}{}^{\mathrm{t}}(c_2,0,1)`$ and $`T={}_{}{}^{\mathrm{t}}(c_1,1,0)`$. The divisor $`E(P)`$ is rational if and only if $`(b1)=k(a1)`$ for some $`k𝐍`$ and $`\mathrm{gcd}(b,c)=1`$. Then $`d=a1`$ and $`det(P,Q)=ck>1`$, $`det(P,S)=c>1`$, $`det(P,T)=b>1`$ and $`det(P,E_3)=c`$. Thus the $`E(P)`$ has at least 3 arms. $`\mathrm{}`$
### 4.4. Normally smooth divisors on $`T_{p,q,r}`$-surfaces
Let $`T_{p,q,r}:x^p+y^q+z^r+xyz=0`$ with $`1/p+1/q+1/r<1`$.
(1) Suppose that $`p,q,r`$ are pairwisely coprime and $`p<q<r`$. The diagram $`\mathrm{\Gamma }^{}(f)_2^+`$ has three strictly positive vertices $`P:={}_{}{}^{\mathrm{t}}(rqrq,r,q),Q:={}_{}{}^{\mathrm{t}}(r,prpr,p),`$ and $`R:={}_{}{}^{\mathrm{t}}(q,p,pqqp)`$. The cones $`\mathrm{Cone}(P,E_1),\mathrm{Cone}(Q,E_2)`$ and $`\mathrm{Cone}(R,E_3)`$ are regular. Put $`\delta :=pqrprqrpq`$. Then $`det(P,Q)=det(Q,R)=det(P,R)=\delta `$.
###### Proposition 23.
Under the above assumption, we have
$$\rho (X_{p,q,r},O)=\rho _{QR}^{(1)}+\rho _{QR}^{(2)}+\rho _{QR}^{(3)}+\rho _{PR}^{(2)}+\rho _{PR}^{(3)}+\rho _{PQ}^{(3)}2ϵ,$$
where $`\epsilon =1`$ if $`p=3`$, and $`ϵ=0`$ if $`p3`$.
###### Proof.
This is a summary of the following three lemmas. $`\mathrm{}`$
###### Lemma 24.
* $`𝒱_{\mathrm{ns}}^{(1)}(Q,R)=\{P_k={}_{}{}^{\mathrm{t}}(1,k,pk1)p/q<k<(rprp)/r\}.`$
* $`𝒱_{\mathrm{ns}}^{(2)}(Q,R)=\{P_k^{}={}_{}{}^{\mathrm{t}}(k,1,pkk1)r/(prpr)<k<q/p\}.`$
* $`𝒱_{\mathrm{ns}}^{(3)}(Q,R)=\{P_k^{\prime \prime }={}_{}{}^{\mathrm{t}}(k,pkk1,1)q/(pqpq)<k<r/p\}.`$
* $`𝒱_{\mathrm{ns}}^{(1)}(Q,R)𝒱_{\mathrm{ns}}^{(2)}(Q,R)𝒱_{\mathrm{ns}}^{(3)}(Q,R)\mathrm{}`$ if and only if $`p=3`$.
* $`\rho _{QR}=\rho _{QR}^{(1)}+\rho _{QR}^{(2)}+\rho _{QR}^{(3)}1ϵ`$, where $`ϵ=1`$ if $`p=3`$, and $`ϵ=0`$ if $`p3`$.
###### Proof.
We mainly use Theorem 7. Let $`P^{}:=(\beta Q+\alpha R)/\delta ={}_{}{}^{\mathrm{t}}(p_1,p_2,p_3)`$. The equation is
$$\{\begin{array}{cc}\beta r+\alpha q=p_1\delta \hfill & \\ \beta (prpr)+\alpha p=p_2\delta \hfill & \\ \beta p+\alpha (pqpq)=p_3\delta \hfill & \end{array}\text{ this implies }\{\begin{array}{cc}\alpha =(prpr)p_1rp_2\hfill & \\ \beta =qp_2pp_1\hfill & \\ p_2+p_3=(p1)p_1\hfill & \end{array}$$
Hence, we have the following conclusions.
1) $`p_1=1`$ if and only if there exists an integer $`p_2>0`$ such that $`\alpha >0`$ and $`\beta >0`$. This is equivalent to $`p/q<p_2<(prpr)/r`$. And in this case $`P^{}=(1,p_2,p1p_2)`$.
2) $`p_2=1`$ if and only if there exists an integer $`p_1>0`$ such that $`r/(prpr)<p_1<q/p`$. And in this case $`P^{}=(p_1,1,(p1)p_11)`$.
3) $`p_3=1`$ if and only if there exists an integer $`p_1>0`$ such that $`q/(pqpq)<p_1<r/p`$. And in this case $`P^{}={}_{}{}^{\mathrm{t}}(p_1,pp_1p_11,1)`$.
4) is obvious now.
5) One can see this by comparing the three sets $`𝒱_{\mathrm{ns}}^{(i)}(Q,R)`$. In case $`p=2`$, we have $`𝒱_{\mathrm{ns}}^{(1)}(Q,R)=\mathrm{}`$ and $`𝒱_{\mathrm{ns}}^{(2)}(Q,R)𝒱_{\mathrm{ns}}^{(3)}(Q,R)=\{{}_{}{}^{\mathrm{t}}(2,1,1)\}`$. Hence, $`\rho _{QR}=\rho _{QR}^{(2)}+\rho _{QR}^{(3)}1`$.
In case $`p=3`$, we have $`𝒱_{\mathrm{ns}}^{(i)}(Q,R)𝒱_{\mathrm{ns}}^{(j)}(Q,R)=𝒱_{\mathrm{ns}}^{(1)}(Q,R)𝒱_{\mathrm{ns}}^{(2)}(Q,R)𝒱_{\mathrm{ns}}^{(3)}(Q,R)=\{{}_{}{}^{\mathrm{t}}(1,1,1)\}`$ for $`ij`$. Hence, $`\rho _{QP}=\rho _{QR}^{(1)}+\rho _{QR}^{(2)}+\rho _{QR}^{(3)}2`$.
In case $`p>3`$, we have $`𝒱_{\mathrm{ns}}^{(1)}(Q,R)𝒱_{\mathrm{ns}}^{(2)}(Q,R)=\{{}_{}{}^{\mathrm{t}}(1,1,p2)\}`$ and $`𝒱_{\mathrm{ns}}^{(1)}(Q,R)𝒱_{\mathrm{ns}}^{(3)}(Q,R)=𝒱_{\mathrm{ns}}^{(2)}(Q,R)𝒱_{\mathrm{ns}}^{(3)}(Q,R)=\mathrm{}`$. Hence, $`\rho _{QP}=\rho _{QR}^{(1)}+\rho _{QR}^{(2)}+\rho _{QR}^{(3)}1`$. $`\mathrm{}`$
Similarly, one can prove the following two lemmas.
###### Lemma 25.
* $`𝒱_{\mathrm{ns}}^{(1)}(P,R)=\mathrm{}`$.
* $`𝒱_{\mathrm{ns}}^{(2)}(P,R)=\{Q_{\mathrm{}}^{}={}_{}{}^{\mathrm{t}}(q\mathrm{}1,1,\mathrm{})q/r<\mathrm{}<(pqpq)/p\}.`$
* $`𝒱_{\mathrm{ns}}^{(3)}(P,R)=\{Q_{\mathrm{}}^{\prime \prime }={}_{}{}^{\mathrm{t}}(q\mathrm{}\mathrm{}1,\mathrm{},1)p/(pqpq)<\mathrm{}<r/q\}.`$
* Let $`Q^{}={}_{}{}^{\mathrm{t}}(q_1,q_2,q_3)=(\beta P+\alpha R)/\delta `$. Then $`(q1)q_2=q_1+q_3`$.
* $`\rho _{PR}=\rho _{PR}^{(2)}+\rho _{PR}^{(3)}1.`$ $`\mathrm{}`$
###### Lemma 26.
* $`𝒱_{\mathrm{ns}}^{(1)}(P,Q)=𝒱_{\mathrm{ns}}^{(2)}(P,Q)=\mathrm{}`$.
* $`𝒱_{\mathrm{ns}}^{(3)}(P,Q)=\{R_{\mathrm{}}^{}={}_{}{}^{\mathrm{t}}(r\mathrm{}1,\mathrm{},1)r/q<\mathrm{}<(prpr)/p\}`$ and $`\rho _{PQ}=\rho _{PQ}^{(3)}`$. $`\mathrm{}`$
###### Example 27.
(1) Let $`p=2,q=3`$ and $`r7`$. By the canonical subdivisions of the three cones, we see that $`\rho _{QR}=\left[\frac{r6}{2}\right]1`$, $`\rho _{PR}=\left[\frac{r6}{3}\right]1`$, and $`\rho _{PQ}=\left[\frac{r3}{6}\right]`$.
(2) Let $`p=3,q=4`$ and $`r>4`$. By the canonical subdivisions of the three cones, we see that $`\rho _{QR}=\left[\frac{r}{3}\right]1`$, $`\rho _{PR}=\left[\frac{r}{4}\right]1`$ and $`\rho _{PQ}=[\frac{2r}{3}][\frac{r}{4}]1`$.
(2) Another case. Let $`f(x,y,z)=x^n+y^n+z^n+xyz`$ ($`n4`$). The dual Newton diagram has three covectors $`P_i,i=1,2,3`$ corresponding to the three compact faces. They are given by $`{}_{}{}^{\mathrm{t}}(n2,1,1),{}_{}{}^{\mathrm{t}}(1,n2,1),{}_{}{}^{\mathrm{t}}(1,1,n2)`$. And for $`ij`$, $`det(P_i,P_j)=n3`$. Let $`B_1,\mathrm{},B_k`$ be the vertices of the canonical subdivision of $`\mathrm{Cone}(P_1,P_2)`$ from $`P_1`$. Then $`B_1=(P_2+(n4)P_1)/(n3)={}_{}{}^{\mathrm{t}}(n3,2,1)`$. Thus $`(n3)/(n4)=[2,\mathrm{},2]`$ with $`(n4)`$-copies of 2. This implies $`k=n4`$ and $`B_j={}_{}{}^{\mathrm{t}}(n2j,1+j,1),j=1,\mathrm{},n4`$. In fact, by Lemma 5 the third coordinate of $`B_j`$ is always 1 as both of $`P_1,P_2`$ have 1 as the third coordinate. Hence $`\rho _{P_1P_2}=n4`$. The branch $`\mathrm{Cone}(P_i,E_i)`$ is regular. Thus $`\rho (V,O)=\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})=3n9`$ and every exceptional divisor is normally smooth.
## 5. Remarks
### 5.1. Example of the inequality $`\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})>\rho (X,O)`$
Let us consider $`A_{2c1}`$-singularity, $`X=\{x^2+y^2+z^{2c}=0\}`$. The resolution graph has two arms and the central divisor $`E(P)`$ is a rational curve with $`E(P)^2=1`$. Thus we have to blow-down the central divisor once ( Example (6.7.1) in \[9, III\] ). However in this example, the central exceptional divisor is not normally smooth, i.e., the extra blowing-up is line-admissible. So $`\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})=\rho (X,O)`$. The following gives an example of $`\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})>\rho (X,O)`$.
###### Example 28.
Let $`X`$ be defined by $`h=xy+y^{bc}+z^c`$ with $`b,c2`$. This is an $`A_{c1}`$-singularity and a special case of $`X_{\mathrm{II}}`$ with $`P:={}_{}{}^{\mathrm{t}}(bc1,1,b)`$ and $`Q:={}_{}{}^{\mathrm{t}}(0,c,1).`$
Since $`det(P,E_1)=det(P,E_3)=1`$ and $`det(P,Q)=bc1`$, we make the canonical subdivision of $`\mathrm{Cone}(P,Q)`$. The first covector $`T_1`$ from $`P`$ is given by
$$T_1=(Q+(bcc1)P)/(bc1)={}_{}{}^{\mathrm{t}}(bcc1,1,b1)$$
We have the continuous fraction expansion $`(bc1)/(bcc1)=[2,\mathrm{},2,3,2,\mathrm{},2]`$ where the number of $`2`$ in the first 2-series (respectively in the second $`2`$-series) is $`(b2)`$ (resp. $`c2`$). Thus we have $`c+b3`$ covectors $`T_1,\mathrm{},T_{b+c3}`$. The exceptional divisor $`E(P)`$ is rational with $`E(P)^2=1`$ and $`E(T_j)`$ with self intersection number $`E(T_j)^2=2`$ for $`jb1`$ and $`3`$ for $`j=b1`$ (see Theorem (6.3), Chapter III, ). In fact first $`b2`$ covectors are given by
$`Q_j={}_{}{}^{\mathrm{t}}(cbjc1,1,bj),j=1,\mathrm{},b1`$
$`Q_{b1+j}={}_{}{}^{\mathrm{t}}(cj1,j+1,1),j=1,\mathrm{},c2`$
and we see that they are normally minimal. To get a minimal reslution, we need to blow down $`b1`$ divisors $`E(P),E(T_1),\mathrm{},E(T_{b2})`$ in this order. Then the self-intersection number of $`E(T_{b1})`$ changes to $`2`$ and we get $`A_{c1}`$ graph. In this example, we have $`\rho (X,O)=c1`$ and $`\rho (\mathrm{\Sigma }_{\mathrm{can}}^{})=b+c2`$.
### 5.2. Parametrization of lines
The normally smooth divisors on a surface $`X`$ correspond to the lines on $`X`$. By using a toric resolution, one can give the exact parameterizations of the lines on $`X`$. This was done already for the Pham-Brieskorn surfaces in .
###### Proposition 29.
Suppose that we have a line $`L`$ in a non-degenerate surface $`X:f(x,y,z)=0`$ and assume that $`L`$ is parametrized as
$$x(t)=\alpha t^a+\alpha _1t^{a+1}\mathrm{},y(t)=\beta t^b+\beta _1t^{b+1}+\mathrm{},z(t)=\gamma t^c+\gamma _1t^{c+1}+$$
with $`\alpha ,\beta ,\gamma 0`$ and $`\mathrm{min}(a,b,c)=1`$. Let $`P={}_{}{}^{\mathrm{t}}(a,b,c)`$. Then the pull back of $`L`$ intersects $`E(P)`$ transversally and $`f_P(\alpha ,\beta ,\gamma )=0`$. Conversely any curve in $`_{E(P)}`$ has such a parametrization.
###### Example 30.
(1) Let $`X`$ be defined by $`h=x^ay+y^bz^b=0`$ with $`a=a_1(b1)`$ and $`a_1>1`$. This is a special case of $`X_{\mathrm{II}}`$. We use the notations in §4.2. Note that $`P={}_{}{}^{\mathrm{t}}(1,a_1,a_1)`$, $`Q={}_{}{}^{\mathrm{t}}(0,b,1)`$, $`det(P,Q)=det(P,E_3)=1`$ and $`det(P,E_1)=a_1`$. By canonical subdivision of $`\mathrm{Cone}(P,E_1)`$ we have $`R_i:={}_{}{}^{\mathrm{t}}(1,i,i)`$ with $`i=0,1,\mathrm{},i_1=a_1`$, where $`R_0:=E_1`$ and $`R_{i_1}:=P`$. Hence $`\rho _{PE_1}=a_11`$. Since $`r(P,E_1)+1=b`$, each $`E(R_i)`$ has $`b`$ components. By \[9, III(6.3)\], $`E(P)^2=b<1`$. Hence $`\pi `$ is minimal and $`\rho (X,0)=b(a_11)+1`$. The restriction of $`\pi `$ on the toric chart associated with $`\sigma _i:=\mathrm{Cone}(R_i,R_{i1},E_2)`$ is given by
$$\pi _{\sigma _i}:x=uv,y=u^iv^{i1}w,z=u^iv^{i1}.$$
and the pull-back of $`h`$ is given by
$$h\pi _{\sigma _i}=u^{ib}v^{(i1)b}\left(u^{(a_1i)(b1)}v^{(a_1i+1)(b1)}w+w^b1\right)$$
The divisor $`E(R_i)`$ is defined by $`u=0`$ and $`w^b1=0`$, hence $`E(R_i)`$ has $`b`$ components. On this toric chart, the resolution $`\stackrel{~}{X}`$ of $`X`$ is defined by
$$\stackrel{~}{h}_i(u,v,w):=u^{(a_1i)(b1)}v^{(a_1i+1)(b1)}w+w^b1=0$$
and in a neighborhood of $`qE(R_i)`$ we take $`u,v`$ to be the local coordinates of $`\stackrel{~}{X}`$. Let $`q=(0,s)`$ in this coordinates. We consider the lines $`C_s`$ defined by $`t(t,s)`$. The image of $`C_s`$ by $`\pi _{\sigma _i}`$ is given by
$$\pi _{\sigma _i}(C_s):x=st,y=s^{i1}w_k(t,s)t,z=s^{i1}t^i,$$
where $`w_k(t,s)`$ is the solution of $`\stackrel{~}{h}_i(t,s,w)=0`$ with $`w_k(0)=\mathrm{exp}(2\pi ki/b)`$. As a special case, take $`i=1`$. Then $`C_s`$ is a normal line on $`E(Q_1)`$. When we moves $`s0`$, this line approaches to $`E(E_1)`$ and $`w_k(t)\mathrm{exp}(2k\pi i/b)`$ and the image is the obvious line $`t(x,y,z)=(0,w_kt,t)`$.
(2) Let $`X=T_{2,3,7}:x^2+y^3+z^7+xyz=0`$. We have three covectors
$$P={}_{}{}^{\mathrm{t}}(11,7,3),Q={}_{}{}^{\mathrm{t}}(7,5,2),R={}_{}{}^{\mathrm{t}}(3,2,1)$$
and we do not need any other covector. Consider the toric chart $`\sigma :=(Q,R,E_3)`$ with coordinates $`(u,v,w)`$. Then the line $`u=1,v=t`$ produces a line parametrized as $`t(t^3,t^2,2t+128t^2+\mathrm{})`$.
### 5.3. Obvious lines on surfaces
We consider a surface $`X=\{f(x,y,z)=0\}`$ where $`f`$ has a non-degenerate Newton boundary. There are surfaces having obvious lines which can be read off from the polynomial defining the surface.
(1) Assume that $`f(x,y,z)`$ is not convenient and assume for example $`\{y=z=0\}X`$. Then as we have seen in Lemma 10, there is a unique non-compact face, different from the coordinate planes, which has the covector of the type $`Q={}_{}{}^{\mathrm{t}}(0,c,1)`$ or $`{}_{}{}^{\mathrm{t}}(0,1,c)`$ and a unique covector $`P`$ such that $`\mathrm{Cone}(P,Q)`$ is in $`\mathrm{\Gamma }^{}(f)_2^+`$ and $`P`$ corresponds to a compact face. Let $`Q_1,\mathrm{},Q_k`$ be the covectors defining the canonical regular subdivision from $`Q`$. Then $`Q_1`$ is a normally smooth divisor and $`_{\mathrm{Q}_1}`$ contains the canonical line $`\{y=z=0\}`$.
(2) Assume that $`h(x,y):=f(x,y,0)`$ (the section of $`f`$ with $`z=0`$) is a non-monomial homogeneous polynomial of degree $`d`$. Then we can factor $`h(x,y)=cx^ay^b_{i=1}^k(y\alpha _ix)`$. Thus $`X`$ has the lines $`z=0,y=\alpha _ix`$ for $`i=1,\mathrm{},k`$. Combinatorially this says the following. There exists a compact face $`\mathrm{\Delta }`$ such that $`\mathrm{\Delta }\mathrm{\Delta }(h)`$. The corresponding covector takes the form $`P={}_{}{}^{\mathrm{t}}(p,p,r)`$ with $`\mathrm{gcd}(p,r)=1`$. Then the first covector $`Q_1`$ from $`E_3`$ in the canonical regular subdivision of $`\mathrm{Cone}(P,E_3)`$ takes the form $`Q_1={}_{}{}^{\mathrm{t}}(1,1,s)`$ with $`s=1+[r/p]`$. So we can see that $`Q_1𝒱_{\mathrm{ns}}(P,E_3)`$. A typical example is $`T_{n,n,n}:x^n+y^n+z^nxyz=0`$. Another example is (1) of Example 30.
(3) Assume that the monomial $`x^A`$ in $`f`$ such that $`(A,0,0)\mathrm{\Gamma }(f)`$. We say that $`x^A`$ is negligibly truncatable if $`f_t(x,y,z)=(f(x,y,z)f(x,0,0))+tf(x,0,0)`$ defines a $`\mu `$-constant family for $`0t1`$ (cf. ). Assume for example, the monomials $`x^ay`$ and $`x^bz^c`$ are on the non-compact face of $`\mathrm{\Gamma }(f_0)`$. Let $`Q^{}:={}_{}{}^{\mathrm{t}}(c/d,c(Aa)/d,(Ab)/d)`$ with $`d=\mathrm{gcd}(c,Ab)`$. The covector $`Q^{}`$ corresponds to the negligible compact face of $`f_1`$ containing $`(a,1,0),(b,0,c)`$ and $`(A,0,0)`$. Then there is a normally smooth divisor on $`\mathrm{Cone}(Q,E_3)`$. In fact, $`det(Q^{},E_3)=c/d`$. If $`c=d`$, $`Q`$ gives normally smooth divisor. If $`c>d`$, the first covector $`Q_1^{}`$ of the canonical regular subdivision of $`\mathrm{Cone}(Q^{},E_3)`$ is normally smooth. An example is given by $`f(x,y,z)=x^2y+y^2+z^5+x^5`$. Then $`x^5`$ is negligibly truncatable.
(4) Assume that $`\mathrm{\Gamma }(f)`$ has a compact face whose covector $`P`$ has 1 in its coefficients. Then $`E(P)`$ is a normally smooth divisor. This is the case, for example, if $`P={}_{}{}^{\mathrm{t}}(1,1,1)`$ and $`f_P(x,y,z)`$ has a two-dimensional support. We can see easily that $`E(P)`$ is isomorphic to the projective curve $`f_P(x,y,z)=0`$ in $`^2`$. The tangent cone of $`X`$ at $`O`$ is given by the cone of $`f_P=0`$.
### 5.4. Normally smooth divisors on complete intersections
In this paper we mainly considered normally smooth divisors on two dimensional hypersurface singularities. However every assertion can be generalized to non-degenerate complete intersections. We give an example. Consider the surface given by $`X=\{f_1(x,y,z,w)=f_2(x,y,z,w)=0\}`$ where $`f_1`$ and $`f_2`$ has the same Newton boundary. Assume that $`f_1,f_2`$ are Pham-Brieskorn polynomials of the same type, with generic coefficients:
$$f_i=a_ix^{p_1}+b_iy^{p_2}+c_iz^{p_3}+d_iw^{p_4},i=1,2$$
We assume that $`p_1,\mathrm{},p_42`$ and mutually coprime. Then the dual Newton diagram $`\mathrm{\Gamma }^{}(f_1,f_2)`$ is the same with $`\mathrm{\Gamma }^{}(f_i)`$ and $`\mathrm{\Gamma }^{}(f)_2^+`$ is star-shaped with the center
$`P={}_{}{}^{\mathrm{t}}(p_2p_3p_4,p_1p_3p_4,p_1p_2p_4,,p_1p_2p_3)`$ and four arms $`\mathrm{Cone}(P,E_i),i=1,\mathrm{},4`$. We consider the $`\mathrm{Cone}(P,E_1)`$. First $`det(P,E_1)=p_1`$. By Lemma 11, $`𝒱_{ns}^{(i)}(P,E_1)=\mathrm{}`$ for $`2i4.`$ As for $`𝒱_{ns}^{(1)}(P,E_1)\mathrm{}`$ if and only if $`p_2p_3p_4<p_1`$ and putting $`r=[p_1/p_2p_3p_4]`$, $`𝒱_{ns}^{(1)}(P,E_1)=\{Q_j=(jP+(p_1jp_2p_3p_4)E_1)/p_1;j=1,\mathrm{},r\}`$.
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# Radio Galaxy Clustering at 𝑧∼0.3
## 1 Introduction
Powerful radio sources are almost exclusively associated with giant elliptical galaxies, and appear to be in richer than average environments (e.g. Hill & Lilly 1991). This suggests that they should be more biased tracers of the mass distribution than normal galaxies. A study of the clustering of local ($`z<0.1`$) radio galaxies by Peacock & Nicholson (1991) showed that this was indeed the case, with radio galaxies having a cross correlation function of the usual form (assumed throughout this paper), $`\xi _{gg}=(r/r_0)^{1.8}`$, with a correlation length of $`r_0=11h^1`$ Mpc. <sup>1</sup><sup>1</sup>1We assume a cosmology with $`\mathrm{\Omega }_\mathrm{M}=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`H_0=100h^1\mathrm{kms}^1\mathrm{Mpc}^1`$ throughout This can be compared to $`r_0=5.7h^1`$Mpc for normal galaxies (Loveday et al. 1992); $`r_0=4.5h^1`$ for IRAS-selected galaxies (Fisher et al. 1994), and $`r_0=14.3h^1`$Mpc for rich clusters of galaxies (Dalton et al. 1994). Radio galaxies thus cluster with a strength intermediate between normal galaxies and clusters (Bahcall & Chokshi 1992). The associated cosmological bias factor is lower than for clusters, but about twice as high as that for IRAS-selected galaxies (Peacock & Dodds 1994).
Studies of the angular correlation function of radio sources by Magliocchetti et al. (1999) and Cress & Kamionkowski (1998) show that the data from radio surveys such as FIRST are consistent with little evolution in the clustering amplitude in the range $`0<z\stackrel{<}{_{}}1`$. However, these conclusions are based on extrapolating the luminosity functions of Dunlop & Peacock (1990; hereafter DP), to the faint flux densities near the limit of the FIRST survey, $`S_{1.4}3`$ mJy. At these levels the DP luminosity functions are constrained by source count data only, as direct redshift surveys were only available at $`S_{1.4}\stackrel{>}{_{}}200`$mJy. In particular, at faint flux levels the radio source population is a mix of nearby star-forming galaxies and AGN-powered radio sources with a range in redshifts from 0 to $`>`$4 (Condon et al. 1998), which have quite different clustering properties. It is therefore important to test the results of the angular correlation function studies with direct measurements of clustering from radio galaxy redshift surveys.
Studies of the clustering of radio-quiet AGN seem to show a generally similar correlation length, but there is a wide range in estimates of $`r_0`$ from different samples. This can probably be explained if the correlation function depends both on redshift and AGN luminosity \[e.g. Sabbey et al. (2000), La Franca, Andreani & Christiani (1998)\].
Magliocchetti et al. (1999) discuss theoretical predictions for the evolution of the two-point correlation function of radio sources. Perhaps the most appropriate case to take is that where radio galaxies form at high redshift ($`z1`$). We can trace the evolution of the host population out to $`z3`$, and find that the hosts vary little with redshift, apart from some passive evolution. The host magnitudes are also only weakly dependent on radio luminosity (Lacy, Bunker & Ridgway 2000). Hence uncertainties in the evolution in the bias factor are unlikely to be as important an issue for radio galaxies as they are for normal galaxies or radio-quiet quasars. Fry (1996) shows that in this “galaxy conservation” scenario, the bias factor increases with redshift according to $`b(z)=1+(b_01)(1+z)`$, where $`b_0`$ is the bias factor at the present epoch. This is because fluctuations in the galaxy density field are fixed at the epoch of formation, but the fluctuations in the matter density field grow with time. The decrease in the bias factor with time is mostly compensated for by the clustering of matter under gravity, for which the growth factor $`D(z)=(1+z)^1`$ for an $`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ cosmology. The two-point correlation function, $`D^2(z)b^2(z)`$, should therefore show little evolution.
To test this model, and to examine the nature of intermediate redshift superclusters, we therefore decided to begin a survey of large-scale structure at moderate redshifts ($`z0.20.65`$). This paper describes the initial result from this survey and also the prospects for future surveys.
## 2 Survey strategy and observations
We tried to optimize the survey to detect supercluster-scale objects at $`z0.4`$. We therefore picked a point on the radio luminosity function where the space density of objects close to the flux limit would be $`10^5h^3`$Mpc<sup>-3</sup>, thus obtaining several objects in structures of linear sizes $`100h^1`$ Mpc. This corresponds to a radio luminosity of $`2\times 10^{23}h^2`$ WHz<sup>-1</sup>sr<sup>-1</sup>, or a flux limit of $`20`$mJy at 1.4 GHz. This is comfortably above the completeness limits of the FIRST and NVSS radio surveys (Becker et al. 1995; Condon et al. 1998).
Initial selection was made from the NVSS catalogue, each NVSS source was examined in FIRST to check for confusion, and to estimate the position of the identification. This technique combines the sensitivity to extended flux of the NVSS survey with the positional accuracy of FIRST. The sample discussed in this paper consists of 322 objects within right ascension and declination ranges of $`01^h<\mathrm{R}.\mathrm{A}.<01^h48^m`$ and $`02^{}<\mathrm{Dec}.<01^{}20^{^{}}`$ (an area of $`40\mathrm{deg}^2`$).
Identifications were made on the UK Schmidt Telescope plates using the Cambridge APM. The plates were approximately calibrated onto $`R`$-band using CCD images of star fields close to each of the plate centres. Initially, only the 34 objects classed as non-stellar in $`R`$ and with $`17.0R20.2`$ were considered for spectroscopy (one of these turned out to be a misclassified high redshift quasar). Six stellar objects were also selected later as checks on the APM classifier. For compact sources, or those with identifiable central components the probability of a chance misidentification is low. In the survey region the density of objects with $`17.0R20.2`$ is $`0.6`$ arcmin<sup>-2</sup>. The error on a FIRST position is $`\stackrel{<}{_{}}2^{^{\prime \prime }}`$, so there is only an $`0.2`$% chance of a misidentification. For double sources with no central component, (18 out of 36 objects) the probability of a misidentification is larger. For these objects we followed the prescription of Lacy et al. (1993) by searching an ellipse with a major axis equal to the radio source size $`d`$ and minor axis of $`d/2`$. The number of objects falling into the search region by chance, $`N`$ is given in Table 1, and is $`<1`$ for all our objects, and $`1`$ for most of them. There are very unlikely to be any $`z<0.7`$ quasars in the sample – the fraction of quasars in complete samples drops rapidly at radio luminosities below $`L_{1.4}10^{25}h^2`$WHz<sup>-1</sup>sr<sup>-1</sup> (Willott et al. 2000).
Spectra were obtained on the Shane 3-m Telescope at Lick Observatory on 1999 October 13-14, 1999 November 12-14 and 1999 December 11-12 and at the Nordic Optical Telescope (NOT) on 2000 January 5 (all dates UT). The plate calibration observations were made on the Shane on 1999 September 15. The Kast spectrograph was used for all observations on the Shane, and the Andalucia Faint Object Spectrograph (ALFOSC) for the NOT observations. Redshifts were determined from emission lines where present, otherwise absorption features were matched both by eye and by a cross-correlation of the spectrum with that of a nearby elliptical; the dominant error sources were the accuracy of the wavelength calibration, and centroiding of spectral features, both $`0.3`$nm, resulting in a typical redshift error of $`\pm 0.001`$. Full details and the spectra will be presented in a future paper.
## 3 Analysis
To measure the two-point correlation function simulated surveys were constructed with the same selection function as the original survey, but with randomly assigned coordinates. The redshift selection function was estimated by using the scatter in the observed $`Rz`$ relation to estimate the probability of an object with redshift $`z`$ having a magnitude within the selected range (Fig. 1). A least deviation fit to the $`Rz`$ relation for $`z<0.7`$ radio galaxies from the 8C-NEC redshift survey (Lacy et al. 1999) was used to establish a mean $`Rz`$ relation:
$$R=21.35.9\mathrm{lg}z.$$
This is close to that derived for the 100-times radio-brighter 3C sample (Eales 1985). The scatter about this relation was measured to be 0.57 mag, using the actual $`R`$ and $`z`$ values obtained for sources in the survey. A selection function was then constructed by taking the luminosity-density evolution model of DP and integrating to obtain a redshift distribution. This was then multiplied by the effect of the photometric selection, which was modelled as a pair of oppositely-tailed error functions with a width in log redshift corresponding to the scatter in the $`Rz`$ relation, and half-power points corresponding to the predicted redshifts for objects at the bright and faint magnitude limits of the survey.
In practice an upper cutoff of $`z=0.45`$ was placed on the sample. Up to this redshift the agreement with the predicted redshift distribution is good, with 28.7 objects predicted in the range $`0.19<z<0.45`$ compared to 29 observed. Above $`z0.45`$, however, the predicted selection function and the observed redshift distribution diverge rapidly (Fig. 1), and of 47.5 objects predicted to be present in the redshift range $`0.19<z<0.65`$, only 36 are found. This could be due either to the mistaking of high redshift galaxies for quasars by the APM classification program, a problem with the assumed luminosity function or a genuine underdense region.
We next estimate the possible effect of incompleteness due to misclassification at $`z<0.45`$. There are 12 objects within the survey selection criteria but with stellar identifications on the UKST plates for which spectra have not yet been obtained. With one exception all have radio morphologies more consistent with being quasars than galaxies (i.e. unresolved, triple or core-halo structures). Of the 13 $`R`$ plate stellar identifications in the range $`17.0R20.2`$ for which we have spectra in this region of sky, from both this work and the FIRST spectroscopic database, three are in fact galaxies, but only one is at $`z<0.45`$. We therefore think it very unlikely that more than 2-3 objects below $`z0.45`$ are missing from the sample other than those accounted for in the selection function.
In Fig. 2 we show a three dimensional representation of the survey volume. Although this figure needs to be interpreted with caution as the space density of objects is dropping with redshift, there is a clump at $`z0.28`$ and R.A. $`01^h05^m`$, which seems to be comprised mostly of FRI sources. There is also a looser association at $`z0.35`$, and a group of three at $`z0.52`$. The sizes of the associations seem to be $`30100h^1\mathrm{Mpc}`$, comparable to low redshift superclusters.
We have estimated the statistical significance of our detection of clustering in the complete $`0.19<z<0.45`$ sample (median redshift $`0.3`$) by binning the distribution of pair separations in 10$`h^1`$ Mpc bins, and comparing with the mean distribution obtained from 10000 Monte-Carlo simulations of the survey using a $`\chi ^2`$ test. This gives the probability of our distribution arising by chance as $`4\times 10^6`$ ($`\chi ^2=112`$ with 53 degrees of freedom). However, a large part of the $`\chi ^2`$ is contributed by the first bin which contains two pairs with separations $`<5h^1`$Mpc. This is close enough that the members of each pair could be in the same galaxy cluster. We therefore removed these two from the first bin and recalculated the $`\chi ^2`$ statistic, this gave a probability of 0.006 ($`\chi ^2=82`$ with 53 degrees of freedom).
To estimate the two-point correlation function, the numbers of data-data (DD), data-random (DR) and random-random (RR) pairs were measured from the data and the simulations and the two-point correlation function calculated according to the formula espoused by Landy & Szalay (1993):
$$\xi (r)=\frac{DD2DR+RR}{RR}.$$
The results for the complete $`0.19<z<0.45`$ sample are shown in Fig. 3. The correlation function was fit over the range 0-100 $`h^1`$Mpc, the best value for $`r_0`$ was $`17h^1`$Mpc, with a range of 5–24$`h^1`$Mpc over which the probability of obtaining the $`\chi ^2`$ was $`>`$5 % (assuming Poisson errors).
As a check on the effect of possible incompleteness we added three galaxies (the largest number we expect based on the discussion above) to the sample with redshifts drawn at random from the selection function, and with random sky positions within the survey region. This reduced $`r_0`$ to 14$`h^1`$Mpc. Incompleteness is thus unlikely to affect our estimate by more than the random error.
## 4 Conclusions and future surveys
We have succeeded in developing an effective method for studying the clustering of moderate redshift radio galaxies directly, and have detected clustering of radio galaxies at $`z0.3`$. The amplitude of the cross-correlation function we measure is consistent with that for radio galaxies locally. This is as expected in the simple model discussed in the introduction, in which radio source hosts evolve little with redshift, and is higher than that for normal galaxies at $`z0.3`$, for which Small et al. (1999) measure $`r_0=3.7h^1\mathrm{Mpc}`$. At present, however, the small size of our survey prevents us ruling out all but the most extreme evolution in the correlation function. Expansion of this survey to $`\stackrel{>}{_{}}100`$ redshifts will allow a measurement of the two-point correlation function to be made which has comparable accuracy to that for normal galaxies at these redshifts. The relatively large volume probed by a larger survey will allow us to define a sample of superclusters and to examine their structures and the evolution of those structures to the present day.
The discovery of redshift clustering in deep pencil beam galaxy surveys, e.g. that of the HDF (Cohen et al. 2000) has raised the possibility that large-scale structures continue to be present in the Universe at least to $`z1`$. Therefore the evolution of these structures should place interesting constraints on cosmology. Crucially, however, because we can trace the evolution of radio galaxy hosts to $`z3`$, we can, in principle, predict how the bias should evolve with redshift, removing an important uncertainty from the interpretation of the results of correlation function studies.
We thank the staff at the Lick Observatory and NOT for their help with the observations, in particular the night assistants on the Shane, Andy, Wayne and Keith, and our support astronomers Elinor Gates and Hugo Schwartz. We also thank Dan Stern, Susan Ridgway, Sally Laurent-Meulheisen and Margrethe Wold for help with various aspects of the project. We thank Bob Becker, Mike Brotherton, Michael Gregg, and Sally Laurent-Meuleisen for spectra obtained from the FIRST spectroscopic database. We also thank Mike Irwin for the APM service, and Devinder Sivia for the pgxtal software. This work was performed under the auspices of the U.S. Department of Energy by University of California Lawrence Livermore National Laboratory under contract No. W-7405-Eng-48, with support from NSF grants AST-98-02791 and AST-98-02732. The NOT is operated on the Island of La Palma jointly by Denmark, Finland, Iceland, Norway and Sweden, in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias.
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# References
CERN-TH/2000-138
LAPTH-794/2000
HARRY LEHMANN AND THE ANALYTICITY UNITARITY PROGRAMME<sup>1</sup><sup>1</sup>1To appear in a volume of Communications in Mathematical Physics, dedicated to the memory of Harry Lehmann.
André MARTIN
Theoretical Physics Division, CERN
CH - 1211 Geneva 23
and
LAPP<sup>2</sup><sup>2</sup>2URA 1436 du CNRS, associée à l’Université de Savoie.
F - 74941 Annecy le Vieux Cedex
ABSTRACT
I try to describe the extremely fruitful interaction I had with Harry Lehmann and the results which came out of the analyticity unitarity programme, especially the proof of the Froissart bound, which, with recent and future measurements of total cross-sections and real parts, remains topical.
Dedication
I dedicate this paper to Marie-Noëlle Fontaine,
the last of the many papers she typed so skillfully
for me, wishing her a happy retirement.
CERN-TH/2000-138
LAPTH-794/2000
May 2000
My first meeting with Harry Lehmann was not with his person but with the famous paper of the trio Lehmann-Symanzik-Zimmermann, LSZ , the importance of which everybody in the Theory group of Maurice Lévy at Ecole Normale realized immediately. In spite of the fact that I did not know German (I still don’t) I read it, Nuovo Cimento in one hand, dictionary in the other hand (I am a “corrected” left hander). Later Harry visited the Ecole Normale in person and I was immediately impressed. That was the time where there was a wave of interest into what is an unstable particle and Lehmann and Lévy were some of the people involved. I remember also quite vividly our meeting at the La Jolla Conference in 1961 which I attended, coming from CERN. It was, as I realized a posteriori, a very important conference, for physicists and for people (some of the people I met there became my very best friends). I remember that Marcel Froissart gave a talk on his famous Froissart bound on the total cross-section, $`\sigma _t<c\mathrm{log}s)^2`$, $`s`$ square of the centre-of-mass energy, and Harry with his very meticulous mind found out that some of the estimates of Froissart were not quite correct, though this did not affect the result (some year later, I published a sum rule on pion-nucleon scattering and Harry discovered a very well hidden mistake. I was very impressed). Anyway we were both admirative of the achievement of Froissart and for me it was a decisive turning point, since I left almost completely for many years potentials and the Schrödinger equation for the study of high-energy scattering and high-energy bounds. The Froissart bound was derived from a combination of the Mandelstam representation where the scattering amplitude is the boundary value of an analytic function of two variables, which implies automatically dispersion relations proved from field theory in one variable as well as the Lehmann ellipse which is probably the most celebrated result of Harry, a fundamental result presented in 10 small pages of Nuovo Cimento (compare with the incredibly lengthy papers on what I would call “rigorous atomic physics” which appeared during the last 15 years!).
The trouble with the Mandelstam representation is that nobody was ever able to prove it even in perturbation theory (through some wrong proofs were published!). Both Harry and I were anxious to obtain high-energy bounds with minimal assumptions. A step in this direction was made by Greenberg and Low who used the Lehmann Ellipse to derive a bound on the total cross-section where $`(\mathrm{log}s)^2`$ was replaced by $`s(\mathrm{log}s)^2`$. Myself, I realized that the whole Mandelstam representation was not needed to get the Froissart bound and that it was sufficient to replace the Lehmann ellipse by a larger one . Later, in Princeton, Y.S. Jin (a former student of Harry) and I found a way to control the growth of the scattering amplitude for unphysical momentum transfer using positivity but at the time we made no progress on the derivation of the Froissart bound. In the autumn of 1965 I was visiting IHES (Institut des Hautes Etudes Scientifiques) and Harry was there. He attracted my attention on a paper by Nakanishi which contained the claim that the Lehmann Ellipse could be enlarged by using results from perturbative field theory, leading to the obtention of the Froissart bound. As I shall explain later, we tried to make sense of the paper of Nakanishi but in the end could not. Nevertheless it started again my interest in the subject, and after a visit to Cambridge where I learnt that the Nakanishi perturbative domain of analyticity had been obtained independently and in a simpler way by T.T. Wu , I came back to CERN and finally succeeded, using positivity properties not terribly different from those I had used with Jin, to enlarge the Ellipse without using perturbation and prove the Froissart bound from first principles. Something rather rare happened: Harry sent a postcard to congratulate me, but while moving from one apartment to another one or maybe from one office to another, I lost it!
I had many occasions to meet Harry later, but the last one was in the spring of 1998 at CERN where he came to work with T.T. Wu after an operation which seemed successful. At the Ringberg Castle meeting, in the honor of Wolfhart Zimmermann, he was supposed to be the first speaker and could not come because he was ill. I became the first speaker. Then I knew I would never meet him again.
Now I believe that it is necessary to give some technical details.
In 3+1 dimensions (3 space, 1 time) the scattering amplitude depends on two variables energy and angle. For a reaction $`A+BA+B`$
$$E_{c.m.}=\sqrt{M_A^2+k^2}+\sqrt{M_B^2+k^2},$$
(1)
$`k`$ being the centre-of-mass momentum. The angle is designated by $`\theta `$. There are alternative variables:
$$s=(E_{CM})^2,t=2k^2(\mathrm{cos}\theta 1)$$
(2)
(Notice that physical $`t`$ is NEGATIVE).
We shall need later an auxiliary variable $`u`$, defined by
$$s+t+u=2M_A^2+2M_B^2$$
(3)
The Scattering amplitude (scalar case) can be written as a partial wave expansion, the convergence of which will be justified in a moment:
$$F(s,\mathrm{cos}\theta )=\frac{\sqrt{s}}{k}(2\mathrm{}+1)f_{\mathrm{}}(s)P_{\mathrm{}}(\mathrm{cos}\theta )$$
(4)
$`f_{\mathrm{}}(s)`$ is a partial wave amplitude.
The Absorptive part, which coincides for $`\mathrm{cos}\theta `$ real (i.e., physical) with the imaginary part of $`F`$, is defined as
$$A_s(s,\mathrm{cos}\theta )=\frac{\sqrt{s}}{k}(2\mathrm{}+1)\mathrm{Im}f_{\mathrm{}}(s)(\mathrm{cos}\theta )$$
(5)
The Unitarity condition, implies, with the normalization we have chosen
$$\mathrm{Im}f_{\mathrm{}}(s)|f_{\mathrm{}}(s)|^2$$
(6)
which has, as a consequence
$$\mathrm{Im}f_{\mathrm{}}(s)>0,|f_{\mathrm{}}|<1.$$
(7)
The differential cross-section is given by
$$\frac{d\sigma }{d\mathrm{\Omega }}=\frac{1}{s}|F|^2,$$
and the total cross-section is given by the “optical theorem”
$$\sigma _{total}=\frac{4\pi }{k\sqrt{s}}A_s(s,\mathrm{cos}\theta =1).$$
(8)
With these definitions, a dispersion relation can be written as:
$$F(s,t,u)=\frac{1}{\pi }\frac{A_s(s^{},t)ds^{}}{s^{}s}+\frac{1}{\pi }\frac{A_u(u^{},t)du^{}}{u^{}u}$$
(9)
with possible subractions, i.e., for instance the replacement of $`1/(s^{}s)`$ by $`s^N/s^N(s^{}s)`$ and the addition of a polynomial in $`s`$, with coefficients depending on $`t`$.
The scattering amplitude in the $`s`$ channel $`A+BA+B`$ is the boundary value of $`F`$ for $`s+iϵ`$, $`ϵ>00`$, $`s>(M_A+M_B)^2`$. In the same way the amplitude for $`A+\overline{B}A+\overline{B}`$, $`\overline{B}`$ being the antiparticle of $`B`$ is given by the boundary value of $`F`$ for $`u+iϵ,ϵ0u>(M_A+M_B)^2`$. Here we understand the need for the auxiliary variable $`u`$.
The dispersion relation implies that, for fixed $`t`$ the scattering amplitude can be continued in the $`s`$ complex plane with two cuts. The scattering amplitude possesses the reality property, i.e., for $`t`$ real it is real between the cuts and takes complex conjugate values above and below the cuts.
In the most favourable cases, like $`\pi \pi \pi \pi `$ or $`\pi N\pi N`$ scattering dispersion relations have been established for $`T<t0`$, $`T>0`$ .
In the general case, even if dispersion relations are not proved, the crossing property of Bros, Epstein and Glaser states that the scattering amplitude is analytic in a twice cut plane, minus a finite region, for any negative $`t`$ . So it is possible to continue the amplitude directly from $`A+BA+B`$ to the complex conjugate of $`A+\overline{B}A+\overline{B}`$. By a more subtle argument, using a path with fixed $`u`$ and fixed $`s`$ it is possible to continue directly from $`A+BA+B`$ to $`A+\overline{B}A+\overline{B}`$
At this point, we see already that one cannot dissociate analyticity, i.e., dispersion relations, and unitarity, since the discontinuity in the dispersion relations is given by the absorptive part. In the simple case of $`t=0`$, the absorptive part is given by the total cross-section and the forward amplitude is given, as we said already for the case of Compton Scattering, by an integral over physical quantities.
It was recognized very early that the combination of analyticity and unitarity might lead to very interesting consequences and might give some hope to fulfill at least partially the $`S`$ matrix Heisenberg program. This was very clearly stated already in 1956 by Murray Gell-Mann at the Rochester conference. Later this idea was taken over by many people, in particular by Geff Chew. To make this program as successful as possible it seemed necessary to have an analyticity domain as large as possible. Dispersion relations are fixed $`t`$ analyticity properties, in the other variable $`s`$, or $`u`$ as one likes.
Another property derived from local field theory was the existence of the Lehmann ellipse , which states that for fixed $`s`$, physical, the scattering amplitude is analytic in $`\mathrm{cos}\theta `$ in an ellipse with foci at $`\mathrm{cos}\theta =\pm 1`$. $`\mathrm{cos}\theta =1`$ corresponds to $`t=0`$ the ellipse therefore contains a circle
$$|t|<T_1(s).$$
(10)
$`T_1(s)`$ is given by
$$\begin{array}{cc}x_0& =1+\frac{T_1(s)}{2k^2}\\ & \\ x_0& =\left[1+\frac{(M_1^2M_A^2)(M_2^2M_B^2)}{k^2(s(M_1M_2)^2)}\right]\end{array}^{1/2}$$
where $`M_A`$ and $`M_B`$ are the masses of the particles, $`M_1`$ and $`M_2`$ are the lowest intermediate states in the currents associated to the fields of the incoming particles.
Hence $`T_1(s)0`$ for $`s(M_A+M_B)^2`$ and $`s\mathrm{}`$.
The absorptive part is analytic in the larger ellipse, the “large” Lehmann ellipse, containing the circle
$$|t|<T_2(s),$$
(11)
$`T_2`$ is given by
$$2x_0^21=1+\frac{T_2(s)}{2k^2}$$
So $`T_2(s)c>0`$ for $`s(M_A+M_B)^2`$, $`T_2(s)0`$ for $`s\mathrm{}`$.
It was thought by Mandelstam that these two analyticity properties, dispersion relations and Lehmann ellipses, were insufficient to carry very far the analyticity-unitarity program. he proposed the Mandelstam representation which can be written schematically as
$`F=`$ $`{\displaystyle \frac{1}{\pi ^2}}{\displaystyle \frac{\rho (s^{},t^{})ds^{}dt^{}}{(s^{}s)(t^{}t)}}`$ (12)
$`+\mathrm{circular}\mathrm{permutations}\mathrm{in}s,t,u`$
$`+\mathrm{one}\mathrm{dimensional}\mathrm{dispersion}\mathrm{integrals}`$
$`+\mathrm{subtractions}`$
This representation is nice. It gives back the ordinary dispersion relations and the Lehmann ellipse when one variable is fixed, but it was never proved nor disproved for all mass cases, even in perturbation theory,. One contributor, Jean Lascoux, refused to co-sign a “proof”, which, in the end, turned out to be imperfect.
One very impressive consequence of Mandelstam representation was the proof, by Marcel Froissart, that the total cross-section cannot increase faster than (log $`s)^2`$, the so-called “Froissart Bound” .
My own way to obtain the Froissart bound was to use the fact that the Mandelstam representation implies the existence of an ellipse of analyticity in $`\mathrm{cos}\theta `$ qualitatively larger than the Lehmann ellipse, i.e., such that it contains a circle $`|t|<R`$, $`R`$ fixed, independent of the energy. This has a consequence that $`\mathrm{Im}f_{\mathrm{}}(s)`$ decreases with $`\mathrm{}`$ at a certain exponential rate because of the convergence of the Legendre polynomial expansion and of the polynomial boundedness, but on the other hand the $`\mathrm{Im}f_{\mathrm{}}(s)`$’s are bounded by unity because of unitarity \[Eq. (7)\]. Taking the best bound for each $`\mathrm{}`$ gives the Froissart bound.
Let me now try to recall the exchange Harry Lehmann and I had in the Autumn of 1965 in Bures sur Yvette. We had in common the same desire to find a proof of the Froissart bound without using the Mandelstam representation and to find a way to enlarge the Lehmann ellipse. Harry pointed out to me a paper published by N. Nakanishi a few months earlier where he claimed that he had a proof of the Froissart bound.
Let me remind you that the largest possible ellipse of convergence of the Legendre Polynomial series for the absorptive part has necessarily a singularity at its right extremity. This is the analogue of a classical theorem on power series with positive coefficients. This means that if you succeed (take the $`\pi \pi `$ case, $`m_\pi =1`$) in proving that the absorptive part is analytic in the neighbourhood of the segment
$$t=0t=4,$$
then it is automatically analytic in the ellipse with foci
$$t=4st=0,$$
and right extremity $`t=4`$, and a fortiori it is analytic in the circle
$$|t|<4,$$
entirely contained in the ellipse.
Nakanishi had obtained a representation valid for any Feynman diagram
$$T_N(s,t)=\frac{d^n\alpha }{\left[f(\alpha )+sg(\alpha )+th(\alpha )\right]^p}$$
Later on I learnt from P. Landshoff that this representation had also been obtained, independently and in a simpler way by T.T. Wu . A minimal analytic domain for $`T_N(s,t)`$ is obtained when the denominator in the integral representaiton does not vanish.
This domain, for the $`\pi \pi `$ case for fixed complex $`s`$ is a kind of strip containing the straight line going through $`t=0`$ and $`t=4s`$, (which corresponds to $`\mathrm{cos}\theta _s`$ real), and the segments $`4<t<+4`$ and $`s<t<8s`$. When $`s`$ tends to a real value the domain shrinks to zero for $`s>4`$ and for $`s<t`$ (for $`t`$ real).
This means that for $`t`$ fixed, real $`4<t<+4`$, dispersion relations hold.
The Nakanishi-Wu representation also implies the validity of partial wave dispersion relations but this is irrelevant for our problem.
However, there is nothing like a small or a large Lehmann ellipse in this domain. The absorptive part in perturbation theory, which is defined only in the limit $`ss_R+iϵ`$, $`ϵ0`$ has a priori no analyticity in $`t`$. A priori, it is just a distribution.
In fact in perturbation theory, unitarity connects amplitudes of different orders and positivity properties of the absorptive part are completely hidden. In three-space dimensions, nobody knows if the perturbation series can be resummed (probably not!) and it is not “legal” to combine the results of axiomatic field theory and perturbation theory. Of course one can always try it as a game, which is what Harry and I tried to do, but we went nowhere.
It is only in December 1965, after a visit to Cambridge, that I found a way to enlarge the Lehmann ellipse in the framework of axiomatic field theory , without using at all the results of perturbation theory. I was maybe a bit unfair not to quote the Nakanishi-Wu representation because the “wrong” paper of Nakanishi was undoubtedly a source of stimulation but, on the other hand, I did not use it at all.
Our method was the following.
The positivity of $`\mathrm{Im}f_{\mathrm{}}`$ implies, by using expansion (5),
$$\left|\left(\frac{d}{dt}\right)^nA_S(s,t)\right|_{4k^2t0}\left|\left(\frac{d}{dt}\right)^nA_S(s,t)\right|_{t=0}.$$
(13)
To calculate
$$F(s,t)=\frac{1}{\pi }_{s_0}\frac{A_s(s^{}t)ds^{}}{s^{}s}$$
(forget the left-hand cut and subtractions!), for $`s`$ real $`<s_0`$ one can expand $`F(s,t)`$ around $`t=0`$. From the property (13) one can prove that the successive derivatives can be obtained by differentiating under the integral. When one resums the series one discovers that this can be done not only for $`s`$ real $`<s_0`$, but for any $`s`$ and that the expansion has a domain of convergence in $`t`$ independent of $`s`$. This means that the large Lehmann ellipse must contain a circle $`|t|<R`$. This is exactly what is needed to get the Froissart bound. In fact, in favourable cases, $`R=4m_\pi ^2`$, $`m_\pi `$ being the pion mass. A recipe to get a lower bound for $`R`$ was found by Sommer
$$R\mathrm{sup}_{s_0<s<\mathrm{}}T_1(s)$$
(14)
It was already known that for $`|t|<4m_\pi ^2`$ the number of subtractions in the dispersion relations was at most two , and it lead to the more accurate bound
$$\sigma _T<\frac{\pi }{m_\pi ^2}(\mathrm{log}s)^2$$
(15)
Notice that this is only a bound, not an asymptotic estimate.
In spite of many efforts the Froissart bound was never qualitatively improved, and it was shown by Kupsch that if one uses only $`\mathrm{Im}f_{\mathrm{}}|f_{\mathrm{}}|^2`$ and full crossing symmetry one cannot do better than Froissart.
On the more theoretical side one might wonder if using crossing symmetry and analytic completion one could not prove Mandelstam representation at least for the pion-pion case using only axiomatic results. This is not the case, as I showed it in 1967 at a meeting organized by Bob Marshak in Rochester where Harry was present . One can write a representation of the scattering amplitude
$$\begin{array}{cc}F_\nu =& _0^1𝑑x_{p_0q_0}^{\mathrm{}}\frac{dpdqw(x,p,q)}{x(ps)^\nu +(1x)(qt)^\nu }\\ & \\ & +\mathrm{circular}\mathrm{permutations}\end{array}$$
For $`\nu =1/2`$ this is just a funny way to write the Mandelstam representation.
For $`\nu =1`$, you get back to the Nakanishi-Wu representation.
For $`\nu =2/3`$ you get a natural domain bigger than all you can get from axiomatic field theory and positivity.
Before 1972, rising cross-sections were a pure curiosity. Almost everybody believed that the proton-proton cross-section was approaching 40 millibarns at infinite energy. Yet, Khuri and Kinoshita took seriously very early the possibility that cross-sections rise and proved, in particular, that if the scattering amplitude is dominantly crossing even, and if $`\sigma _t(\mathrm{log}s)^2`$ then
$$\rho =\frac{ReF}{ImF}\frac{\pi }{\mathrm{log}s},$$
where $`ReF`$ and $`ImF`$ are the real and imaginary part of the forward scattering amplitude.
As early as 1970, Cheng and Wu proposed a model in which cross-sections were rising and eventually saturating the Froissart bound. However, at that time there was no experimental indication of this. It is only in 1972 that it was discovered at the ISR, at CERN, that the $`pp`$ cross-section was rising by 3 millibarns from 30 GeV c.m. energy to 60 GeV c.m. energy . I suggested to the experimentalists that they should measure $`\rho `$ and test the Khuri-Kinoshita predictions. They did it and this kind of combined measurements of $`\sigma _T`$ and $`ReF`$ are still going on. In $`\sigma _T`$ we have now more than a 50 % increase with respect to low energy values. For an up to date review I refer to the article of Matthiae . it is my strong conviction that this activity should be continued with the future LHC. A breakdown of dispersion relation might be a sign of new physics due to the presence of extra compact dimensions of space according to N.N. Khuri . Future experiments, especially for $`\rho `$, will be difficult because of the necessity to go to very small angles, but not impossible .
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# Spin dynamics of Mn12-acetate in the thermally-activated tunneling regime: ac-susceptibility and magnetization relaxation
## I Introduction
In recent years, numerous experimental results on macroscopic samples of molecular magnets, especially Mn<sub>12</sub>-acetate and Fe<sub>8</sub>-triazacyclononane, have drawn attention to the peculiar resonant structure observed in the hysteresis loops and relaxation time measurements, as well as in the dynamic susceptibility. In this paper, we concentrate on Mn<sub>12</sub> (shorthand for Mn<sub>12</sub>-acetate). At low temperature, the observed relaxation times $`\tau `$ are long, up to several months and more, and display a series of resonances with faster relaxation as a function of an external magnetic field directed along the easy ($`z`$) axis of the sample. These are considered as signs of macroscopic quantum tunneling (MQT) of magnetization.
Typical experimental samples consist of single crystals or ensembles of aligned crystallites of identical Mn<sub>12</sub>-molecules. Each molecule has eight Mn<sup>3+</sup> and four Mn<sup>4+</sup> ions which, in their ferrimagnetic ground state, have a total spin $`S=10`$, see Fig. 1. Due to strong anisotropy along one of the crystalline axes ($`z`$-direction), there is a high potential barrier
$`U(S_z)=AS_z^2BS_z^4,`$ (1)
with $`A/k_\mathrm{B}0.54`$K and $`B/k_\mathrm{B}0.0011`$K, between the opposite orientations of the spin ($`S_z=\pm 10`$); the easy axis is the same for all the molecules. The dipolar interaction between the molecular spins, a possible relaxation mechanism, has been found to be weak in Mn<sub>12</sub>. Instead, the observed resonant phenomena are attributed to quantum tunneling of single spins – the response being magnified by the large number of them – interacting with the phonons in the lattice. The role of the hyperfine interactions is still under some controversy and is only briefly touched upon in the following.
The main features of the experimental findings can be understood in terms of two competing relaxation mechanisms: quantum mechanical tunneling through and thermal activation over the anisotropy barrier. At high temperatures ($`T>3\mathrm{K}`$ or $`T>6\mathrm{K}`$, depending on the experiment ), the spins relax predominantly via thermal activation due to the phonons in the lattice. In this regime, the relaxation time follows the Arrhenius law $`\tau =\tau _0\mathrm{exp}(U/k_\mathrm{B}T)`$, where $`U/k_\mathrm{B}60\mathrm{K}`$ denotes the barrier height and $`\tau _010^8\mathrm{s}`$ the inverse attempt frequency. When temperature is lowered to $`2\mathrm{K}<T<3\mathrm{K}`$, the time required by the over-barrier relaxation increases exponentially but several of the excited states still remain thermally populated. In this regime, pairs of states on the opposite sides of the barrier can be brought to degeneracy by tuning the external magnetic field. This enhances the probability to tunnel across the barrier; the tunneling arises due to crystalline anisotropy and possible transverse magnetic field at the site of the spin. The tunneling amplitudes are the larger the closer to the top of the barrier the states are and, consequently, the thermal population of the higher states plays a key role in relaxation, cf. e.g. Refs. and . At still lower temperatures, tunneling and the relaxation becomes sensitive to fluctuations in the dipolar and hyperfine fields. In this paper, we concentrate in the regime of thermally-activated tunneling.
Several authors have investigated the spin dynamics theoretically with the emphasis ranging from “minimal” models, assuming as simple a spin Hamiltonian $`_S`$ and a model of the surroundings as possible (in order to explain experiments, that is), to more specific models for investigating the role of the dipolar and/or hyperfine interactions, and combinations of these . The thermally activated relaxation has typically been studied using a master equation approach to describe the time evolution of the spin density matrix. The susceptibility, on the other hand, has only been treated within a phenomenological model.
The existing theories have been successfull in explaining the general features seen in experiments. However, several points call for further attention. 1) A microscopic calculation of the dynamic susceptibility is missing altogether. 2) The effect of a strong transverse magnetic field has not been thoroughly studied and, in particular, not in the context of the susceptibility. 3) Several authors including ourselves have found a series of side resonances to arise in their calculations \- the fact that these peaks are not in general (see Ref. for exceptions) observed in experiments is not quite clear. In this work, we aim at elucidating these points and present calculations for the relaxation rates and susceptibility in a unified language that can be conveniently extended to systems with stronger couplings. We work with a Hamiltonian similar to, e.g., Ref. , cf. Eqs. (2), (3), and (4), but introduce an alternative framework to work with the density matrix. It is well known that all the resonances can be enhanced by a strong transverse magnetic field but we show that the resonances can also be reduced and even suppressed: Both the relaxation rate, see also Ref. , and susceptibility are found to display significant dependence on the direction of the transverse field suggesting that interference effects of a geometrical phase could be observed also in Mn<sub>12</sub> and, what is more, do so in the regime of thermally-activated tunneling.
The paper is organized as follows. Part II introduces the microscopic model used for Mn<sub>12</sub> and a discussion on the different interaction mechanisms in the system. In part III, we develop a time-dependent description of the system in terms of a real-time diagrammatic technique. This approach is then used to derive and evaluate the kinetic equation and the resulting master equation governing the spin dynamics. In part III B, we solve the kinetic equation for the field-dependent relaxation times $`\tau (H)`$ and the static susceptibility $`\chi _0(H)`$. The dynamic susceptibility $`\chi (\omega ;H)`$ is calculated in part III C and a Kubo-type formula is found. Part IV displays the numerical results for both $`\tau (H)`$ and $`\chi (\omega ;H)`$ accompanied with a discussion on the results and their relevance to experiments. In part V we sum up the work.
## II System
The spin Hamiltonian of a single Mn<sub>12</sub> molecule can be written in the form $`_S=_z+_\mathrm{T}`$. The first term,
$`_z=AS_z^2BS_z^4g\mu _\mathrm{B}H_zS_z,`$ (2)
with $`S_z`$ being the spin component along the easy axis (here the $`z`$-direction), describes the part that commutes with $`S_z`$. It consists of the anisotropy terms of Eq. (1) and a Zeeman term which enables external biasing of the energies. The anisotropy constants have been experimentally estimated as $`A/k_\mathrm{B}=0.520.56\mathrm{K}`$ and $`B/k_\mathrm{B}=0.00110.0013\mathrm{K}`$; the g-factor is 1.9. The resulting energy levels $`E_m`$ for the eigenstates of $`S_z|m=m|m`$ together with the potential barrier are shown schematically in Fig. 2.
The second term in the Hamiltonian,
$`_\mathrm{T}={\displaystyle \frac{1}{2}}B_4\left(S_+^4+S_{}^4\right)g\mu _\mathrm{B}(H_xS_x+H_yS_y)`$ (3)
does not commute with $`S_z`$ and gives rise to tunneling. The $`B_4`$-term arises from crystalline anisotropy, $`B_4=(4.314.4)10^5\mathrm{K}`$ (below we use $`B_4=8.610^5`$K, but the particular choice is unimportant for the results obtained), while the second term is the Zeeman term corresponding to a transverse magnetic field $`H_{}=H\mathrm{sin}\theta `$ (in spherical coordinates, $`\theta `$ is the polar angle away from the $`z`$-axis; the azimuth angle is denoted $`\varphi `$: $`H_x=H_{}\mathrm{cos}\varphi `$ and $`H_y=H_{}\mathrm{sin}\varphi `$).
Figure 3 shows the eigenenergies of $`_S`$ as a function of the longitudinal magnetic field $`H_z`$. Away from the resonances, the eigenstates $`|d`$ resemble the states $`|m`$ and are also localized onto the different sides of the barrier – the linear field dependence of the eigenenergies stems from the Zeeman term in Eq. (2). Close to the resonances, $`_\mathrm{T}`$ couples the $`|m`$ states across the barrier and gives rise to avoided crossings in the energy diagram for $`E_d`$. The magnitude of these splittings directly gives the tunneling strengths. Depending on the states in question as well as on the magnitude of $`B_4`$ ($`H_{}=0`$ for the moment), the splittings are found to vary enormously: from $`10^{10}`$K for the choice $`B_4/k_\mathrm{B}=4.310^5`$K and the states $`m=\pm 10`$ up to almost 2K for $`B_4/k_\mathrm{B}=14.410^5`$K and the resonance $`m=\pm 2`$. This upper limit is already of the same order of magnitude as the level spacing and in fact even exceeds it rendering perturbative calculations of the tunneling couplings/strengths somewhat questionable. Therefore, even though we first formulate everything in a general form, independent of the choice of basis for the spins, we calculate the actual results working in the eigenbasis of $`_S`$. This choice of basis provides the additional advantage that it allows us to consider arbitrarily strong magnetic fields.
In absence of $`H_{}`$, there is a selection rule to $`_\mathrm{T}`$: only states $`m`$ and $`m^{}`$ that are a multiple of four apart are coupled. Experimental data do not lend support to such a rule, however, but rather suggest that all transitions are allowed. It turns out that already a tiny transverse field in Eq. (3) is sufficient in achieving this – such a field may arise due to dipolar and/or hyperfine interactions within the sample, see below, as well as due to uncertainity in the precise angle between the external field and the easy axis of the sample. We assume throughout the paper a small constant misalignment angle $`\theta =1^{}`$ and $`H_{}=|𝐇|\mathrm{sin}\theta `$. In places, we wish to investigate the effects of a significantly stronger transverse field and state so explicitly. The transverse field increases the tunnel splittings if $`\varphi `$ is close to $`n\pi /2`$ ($`n`$ is an integer), i.e., along the $`x`$ and $`y`$-axes, or leads to oscillations in the splittings if $`\varphi `$ is close to one of the directions $`\pi (2n+1)/4`$. Below we denote these special directions as the hard axes of the molecules. The tunnel splittings are illustrated in Fig. 4 as a function of $`H_{}`$ and for different angles $`\varphi `$, see Ref. .
The usage of $`_S`$ of an isolated spin is based, first of all, on the assumption that the molecules indeed reside in their $`S=10`$ ground state. This is well justified for the experimental temperatures below $`36\mathrm{K}`$ – the energy required to excite the system to the lowest excited state with $`S=9`$ is around 30K. The second assumption is the absence of interaction. In reality, the spins interact with each other via dipolar interaction, with the nuclear spins via hyperfine interaction, and with the phonons of the surrounding lattice. Experimental evidence shows that the dipolar interactions are small for $`T>22.4\mathrm{K}`$, while the hyperfine interactions produce an intrinsic broadening of the order of 10mT to the spin states. We neglect these for the moment and return to them in Sec. IV C.
The spin-phonon interaction is mediated by variations in the local magnetic field induced by lattice vibrations and distortions. For tetragonal symmetry, this can be generally formulated as, cf. Ref. ,
$`_{\mathrm{sp}}`$ $`=`$ $`g_1(ϵ_{xx}ϵ_{yy})\left(S_x^2S_y^2\right)+{\displaystyle \frac{1}{2}}g_2ϵ_{xy}\{S_x,S_y\}`$ (4)
$`+`$ $`{\displaystyle \frac{1}{2}}g_3\left(ϵ_{xz}\{S_x,S_z\}+ϵ_{yz}\{S_y,S_z\}\right)`$ (5)
$`+`$ $`{\displaystyle \frac{1}{2}}g_4\left(\omega _{xz}\{S_x,S_z\}+\omega _{yz}\{S_y,S_z\}\right)`$ (6)
where $`ϵ_{\alpha \beta }`$ and $`\omega _{\alpha \beta }`$ are the symmetric and antisymmetric strain tensors, respectively, and $`g_i`$ ($`i`$=1,2,3,4) are the spin-phonon coupling constants. For these, we adopt the values from Ref. : $`g_1=g_4/2=A`$ and $`|g_2|g_1`$ and $`|g_3|g_4`$. In leading order in $`g_i`$’s, $`_{\mathrm{sp}}`$ produces transitions such that for, $`g_{1,2}`$-terms, $`\mathrm{\Delta }m=\pm 2`$ and, for $`g_{3,4}`$-terms, $`\mathrm{\Delta }m=\pm 1`$. The phonons themselves are described by
$`_{\mathrm{ph}}={\displaystyle \underset{\stackrel{}{k}\sigma }{}}\omega _{\stackrel{}{k}\sigma }b_{\stackrel{}{k}\sigma }^{}b_{\stackrel{}{k}\sigma }`$ (7)
as a bath of noninteracting bosons.
To be more quantitative, the phonons are assumed to be plane waves with a linear spectrum and three modes – two transverse and one longitudinal – denoted by $`\sigma `$. The elements of the strain tensor, $`ϵ_{\alpha \beta }(_\alpha u_\beta +_\beta u_\alpha )/2`$ and $`\omega _{\alpha \beta }(_\alpha u_\beta _\beta u_\alpha )/2`$, are defined in terms of the local displacement vector
$`u_\alpha (\stackrel{}{r})`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{k}\sigma }{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2NM\omega _{\stackrel{}{k}\sigma }}}}e_\alpha ^{(\sigma )}[b_{\stackrel{}{k}\sigma }^{}+b_{\stackrel{}{k}\sigma }]e^{i\stackrel{}{k}\stackrel{}{r}}.`$ (8)
Here $`b_{\stackrel{}{k}\sigma }^{()}`$ are the bosonic operators for phonons with wave vector $`\stackrel{}{k}`$, $`\omega _{\stackrel{}{k}\sigma }`$ is the correponding frequency, and $`e_\alpha ^{(\sigma )}`$ the $`\alpha `$th element (of $`x`$, $`y`$, and $`z`$) of the polarization vector; $`N`$ is the number of unit cells and $`M`$ the mass per unit cell.
Above we considered the energy scales inherent to $`_S`$, and the spin-phonon rates, which are found below (Apps. A and B) to be typically of the order of $`10^510^4`$K, fall in between the extremes of the tunnel splittings. This value is very small compared to the stronger tunneling couplings and seems to suit well for a perturbative treatment. On the other hand, for the low-lying and weakly coupled states the spin-phonon rates may be several orders of magnitude larger than the tunnel splittings and one would expect the tunneling to be suppressed. However, the intermediate regime, where the tunneling and spin-phonon rates are of the same order of magnitude, requires some extra care, see below.
## III Dynamics
The magnetization measured in experiments is the molecular magnetization $`M(t)`$ magnified by the large number of molecules in the samples. Let us define $`M(t)`$ in terms of the reduced density matrix $`\rho (t)=\mathrm{Tr}_{\mathrm{ph}}[\rho ^{\mathrm{tot}}(t)]`$ ($`\rho ^{\mathrm{tot}}(t)`$ is the full density matrix)
$`M(t)`$ $`=`$ $`g\mu _\mathrm{B}S_z(t)`$ (9)
$``$ $`g\mu _\mathrm{B}\mathrm{Tr}_S[S_z\rho (t)]=g\mu _\mathrm{B}{\displaystyle \underset{m}{}}m\rho _{m,m}(t)`$ (10)
The reduced density matrix $`\rho (t)`$ describes the spin degrees of freedom in the presence of the phonon reservoir – its diagonal elements are just the probabilities for the spin to be in the states $`|m`$. Our strategy is to start with the general formulation
$`S_z(t)=\mathrm{Tr}[S_z(t)\rho _0^{\mathrm{tot}}],`$ (11)
where $`\rho _0^{\mathrm{tot}}=\rho ^S\rho ^{\mathrm{ph}}`$ denotes the initial density matrix of the whole system encompassing the spin and phonon degrees of freedom. However, this does not contain the interaction between the two: we assume that the coupling is turned on adiabatically and only enters the time evolution of $`S_z(t)`$ (in Heisenberg picture and with the convention $`\mathrm{}=1`$)
$`S_z(t)=\mathrm{Tr}[e^{+i_{t_0}^t𝑑t^{}(t^{})}S_ze^{i_{t_0}^t𝑑t^{}(t^{})}\rho _0^{\mathrm{tot}}].`$ (12)
In order to evaluate Eq. (12) and to describe the dynamics of the system more quantitatively, we introduce a real-time diagrammatic language that is applicable to any mesoscopic system comprising a part with a finite number of states linearly coupled to an external heat (or particle) reservoir.
### A Diagrammatic language
Equation (12) can be written as a diagram depicted in Fig. 5. In this equation, the spin-phonon interaction terms can be separated from the noninteracting part of the Hamiltonian by shifting to the interaction picture with respect to $`_0=_S+_{\mathrm{ph}}`$ (indicated with the subscript I)
$`S_z(t)`$ $`=`$ $`\mathrm{Tr}[\rho _0^{\mathrm{tot}}S_z(t)]`$ (13)
$`=`$ $`\mathrm{Tr}[\rho _0^{\mathrm{tot}}\stackrel{~}{T}e^{+i_{t_0}^t𝑑t^{}_{\mathrm{sp}}(t^{})_\mathrm{I}}S_z(t)_\mathrm{I}Te^{i_{t_0}^t𝑑t^{}_{\mathrm{sp}}(t^{})_\mathrm{I}}]`$ (14)
where ($`\stackrel{~}{T}`$) $`T`$ is the (anti-)time-ordering operator and $`\rho _0^{\mathrm{tot}}`$ has been moved to the front (this is allowed by the invariance of trace under a cyclical permutation of its arguments).
We proceed by eliminating the phonon degrees of freedom with the aim to obtain an effective theory for the dynamics of the spin system. We first separate the trace and introduce time-ordering $`T_\mathrm{K}`$ along the closed Keldysh-contour in Fig. 5
$`\mathrm{Tr}[\rho _0^{\mathrm{tot}}`$ $``$ $`S_z(t)]=`$ (15)
$`=`$ $`\mathrm{Tr}_S\{\rho _0^S\mathrm{Tr}_{\mathrm{ph}}\rho _0^{\mathrm{ph}}T_\mathrm{K}e^{i_\mathrm{K}𝑑t^{}_{\mathrm{sp}}(t^{})_\mathrm{I}}S_z(t)_\mathrm{I}\}.`$ (16)
The exponential, expanded in powers of $`_{\mathrm{sp}}`$, reads
$`T_\mathrm{K}e^{i_\mathrm{K}𝑑t^{}_{\mathrm{sp}}(t^{})_\mathrm{I}}={\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}(i)^m`$ (17)
$`{\displaystyle _\mathrm{K}}dt_1{\displaystyle _\mathrm{K}}dt_2\mathrm{}{\displaystyle _\mathrm{K}}dt_mT_\mathrm{K}\{_{\mathrm{sp}}(t_1)_\mathrm{I}_{\mathrm{sp}}(t_2)_\mathrm{I}\mathrm{}_{\mathrm{sp}}(t_m)_\mathrm{I}\}`$ (18)
with $`t_1>t_2>\mathrm{}>t_m`$. Each $`_{\mathrm{sp}}(t_i)_\mathrm{I}`$ is represented as a vertex on the Keldysh contour. In the next step, this expansion together with the explicit expression (4) for $`_{\mathrm{sp}}`$ is inserted into Eq. (16) and the trace over the phonons is performed by using Wick’s theorem. As a consequence, the phonon operators are pairwise contracted: $`\mathrm{Tr}_{\mathrm{ph}}\rho _0^{\mathrm{ph}}b^{}(t)b(t^{})=b^{}(t)b(t^{})`$ and $`\mathrm{Tr}_{\mathrm{ph}}\rho _0^{\mathrm{ph}}b(t)b^{}(t^{})=b(t)b^{}(t^{})`$. Here, the phonons are considered as a reservoir that is not perturbed by the single spin – the bath remains in equilibrium and the contractions are given by the Bose-Einstein distribution $`b_{\stackrel{}{k}\sigma }^{}b_{\stackrel{}{k}\sigma }=[\mathrm{exp}(\beta \omega _{\stackrel{}{k}\sigma })1]^1`$. In terms of the diagrams this means that the vertices get coupled (in all possible ways) by interaction or phonon lines which correspond to just these contractions, cf. Fig. 6. The general rules for evaluating the diagrams are given in App. A and the contributions from the interaction lines are calculated in App. B.
All the diagrams, e.g., those in Fig. 6, are composed of two kinds of elements that can be distinguished by drawing a vertical line through the diagram and checking whether it cuts phonon lines or not. If it does not, the corresponding part of the diagram represents free evolution of the spin; if it does, the part of the diagram between successive periods of free evolution corresponds to spin-phonon interaction and, more particularly, to a transition rate between different spin states. The sum of all the possible diagrams with the interaction lines is denoted as $`\mathrm{\Sigma }`$ and, when evaluated, its terms are $`O(g^{2n})`$ where $`g`$ is the spin-phonon coupling constant and $`n`$ the number of interaction lines in the diagram, cf. Appendices A and B, and below. The full time evolution of $`S_z(t)`$ can be expressed as in Fig. 7a in terms of the two kinds of contributions, just discussed. In lowest order in $`g`$, $`\mathrm{\Sigma }`$ is given by the diagrams shown in part (b) of the figure.
### B Kinetic equation
The reduced density matrix defined as
$`\rho _{m,m^{}}(t)(|m^{}m|)(t)`$ (19)
accounts for the full time evolution of the spin in the presence of the spin-phonon interaction. Its diagrammatic expansion is analog to Fig. 7a with $`S_z`$ being replaced by $`|m^{}m|`$. Setting $`t_0=0`$ and differentiating with respect to time, we obtain the kinetic equation
$`\dot{\rho }(t)+i[_S,\rho (t)]={\displaystyle _0^t}𝑑t^{}\mathrm{\Sigma }(t,t^{})\rho (t^{})`$ (20)
The commutator on the l.h.s. of Eq. (20) corresponds to the free (Hamiltonian) time evolution of the spin, while the integral on the r.h.s. describes a dissipative interaction which is nonlocal in time and contains the spin-phonon coupling terms.
The time dependence of the kernel $`\mathrm{\Sigma }(t,t^{})=\mathrm{\Sigma }(tt^{})`$ is evaluated explicitly in Apps. A and B up to order $`O(g^2)`$, compare Fig. 7b. It depends only on the time difference since the Hamiltonian is time-translationally invariant. According to Eq. (A8), we find that $`\mathrm{\Sigma }(tt^{})`$ is a fast decaying function of time, and we make the simplifying Markov assumption that $`\rho (t)`$ remains essentially constant over the time period $`\mathrm{\Delta }t`$ within which $`\mathrm{\Sigma }(tt^{})`$ decays to zero and take $`\rho (t)`$ in front of the integral. This is justified at least for the longest relaxation time $`\tau _1`$ (see below). For convenience, we also take the upper integration limit to infinity.
After taking $`\rho (t)`$ out of the integral in Eq. (20) the integration over time (up to infinity) can be performed and we are left with a constant $`\mathrm{\Sigma }`$. This is evaluated in App. B and corresponds to the diagrams in Figs. 7b and 21a and b, see also below. The kinetic equation (20) becomes
$`\dot{\rho }(t)=iL_0\rho (t)+\mathrm{\Sigma }\rho (t)W\rho (t)`$ (21)
with $`L_0[_S,.]`$. This is similar to the master equation formulations in the literature and may be solved for the eigensolutions of $`W`$
$`\dot{\rho }^{(i)}(t)`$ $`=`$ $`W\rho ^{(i)}(t)=\lambda _i\rho ^{(i)}(t)`$ (22)
$`\rho ^{(i)}(t)`$ $`=`$ $`\rho ^{(i)}(0)e^{\lambda _it}.`$ (23)
The real parts of the eigenvalues correspond to relaxation rates $`Re(\lambda _i)=\tau _i^1`$.
The eigensolution of $`W`$ with $`\lambda _0=0`$ corresponds to the stationary state, defined as $`\dot{\rho }^{(0)}(t)=0`$ with $`\rho ^{(0)}(t)=\rho ^{(0)}`$. The diagonal elements of $`\rho `$ are the thermal probabilities to find the spin in the respective states. In the $`m`$-basis, the static magnetization and susceptibility are readily expressed in terms of the diagonal components $`\rho _m^{(0)}=\rho _{m,m}^{(0)}`$ yielding
$`M_0=g\mu _\mathrm{B}{\displaystyle \underset{m}{}}m\rho _m^{(0)}`$ (24)
$`\chi _0={\displaystyle \frac{M_0}{H_z}}`$ $`=`$ $`g\mu _\mathrm{B}{\displaystyle \underset{m}{}}m{\displaystyle \frac{\rho _m^{(0)}}{H_z}}.`$ (25)
Since the stationary probabilities are given by Boltzmann factors we obtain in the absence of tunneling
$$\chi _0=\frac{(g\mu _\mathrm{B})^2}{k_\mathrm{B}T}\underset{m,m^{}}{}m(m^{}m)\rho _m^{(0)}\rho _m^{}^{(0)}.$$
(26)
In the precence of tunneling, this form remains essentially the same up to transverse fields of the order of $`1\mathrm{T}`$ (in this range most of the tunneling splittings are too small compared to the level spacing in order to change the result considerably).
For all the other solutions $`Re(\lambda _i)<0`$ and the vectors $`\rho ^{(i)}(t)`$ correspond to deviations from $`\rho ^{(0)}`$. The longest relaxation time $`\tau _1`$ is several orders of magnitude larger than $`\tau _2`$ and the respective solution is interpreted to correspond to the over-barrier relaxation.
Above we have considered the full (reduced) density matrix $`\rho (t)`$. In appendix C, we discuss the choice of the basis and argue that the $`d`$-basis has certain advantages and, e.g., in the case of strong tunneling compard to the spin-phonon rates, it allows the restriction to the diagonal elements of $`\rho (t)_{d,d^{}}`$ only. The nondiagonal elements become important when the rates for the tunneling and the spin-phonon coupling are of the same order of magnitude.
### C Ac-susceptibility
In this section we derive an expression for the dynamic susceptibility
$`\chi (\omega )`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑\tau e^{i\omega \tau }\chi (\tau )`$ (27)
$`\chi (tt^{})`$ $`=`$ $`{\displaystyle \frac{\stackrel{´}{}M(t)}{\stackrel{´}{}H_z(t^{})}}=g\mu _\mathrm{B}{\displaystyle \frac{\stackrel{´}{}S_z(t)}{\stackrel{´}{}H_z(t^{})}}`$ (28)
starting from the Hamiltonian $``$ and formulating the expectation value in Eq. (28) in terms of diagrams. For convenience, we assume in the following that $`H_{}`$ can be tuned to any (static) value independent of $`H_z`$ and that the actual measurement is done by applying a tiny ac-excitation field $`h_z(\omega )`$ on top of the static $`H_z`$. The more general calculation of $`\stackrel{´}{}M_\alpha (t)/\stackrel{´}{}H_\beta (t^{})`$, $`\alpha ,\beta =x,y,z`$, can be carried out along the same lines.
The derivation with respect to $`H_z(t^{})`$ acts on the terms $`\mathrm{exp}[idt^{}_z(t^{})]\mathrm{exp}[i(g\mu _\mathrm{B})dt^{}S_z(t^{})H_z(t^{}))]`$ in Eq. (12) and may occur on either the forward or backward propagator of the diagrams (minus and plus signs, respectively). The derivation takes down a factor $`\pm ig\mu _\mathrm{B}S_z(t^{})`$ and we obtain
$`\chi (tt^{})=i(g\mu _\mathrm{B})^2[S_z(t),S_z(t^{})]`$ (29)
which, when inserted to Eq. (27), is just the Kubo formula for the linear response to an external magnetic field.
Susceptibility can be expressed diagrammatically starting from either of the equations (28) or (29). In either case, $`S_z(t)`$ is written at the latest point to the right. For Eq. (28), $`\stackrel{´}{}/\stackrel{´}{}H_z(t^{})`$ may act on either of the propagators, while, for Eq. (29), the terms of the commutator are ordered on the contour reading them from right (earlier) to left (later times). It is the most convenient to work in the eigenbasis where it is sufficient to consider the diagonal elements only, see App. C. In this case, if $`S_z(t^{})`$ falls in between $`\mathrm{\Sigma }`$’s in the diagram in Fig. 7a, the commutator equals zero – we only need to be concerned with the $`\mathrm{\Sigma }`$ parts.
Let us set $`t_0=0`$ and define $`\stackrel{~}{\mathrm{\Sigma }}(t^{\prime \prime },0)`$ as in Fig. 7b but with $`S_z(t^{})`$ ($`0<t^{}<t^{\prime \prime }`$) inserted onto one of the propagators. Furthermore, let us assign the part $`e^{i\omega (t^{\prime \prime }t^{})}`$ of the Laplace transform to the definition of $`\stackrel{~}{\mathrm{\Sigma }}(t^{\prime \prime },0)`$. This is then the part of the integrand in Eq. (27) appearing between times $`t=0`$ and $`t=t^{\prime \prime }`$. On the forward propagator, this yields a factor $`ig\mu _\mathrm{B}m_{d,d^{}}ig\mu _\mathrm{B}d^{}|S_z|d`$ when acting on the state $`d`$ on the propagator and changing this to the state $`d^{}`$ ($`d`$-basis is not the eigenbasis of $`S_z`$). On the backward propagator, there is a further minus sign. As to the whole diagram, see Fig. 8, the system is in the stationary state before $`t_0=0`$, described by $`\rho ^{(0)}`$ – this corresponds to the requirement that the expectation value $`S_z`$ is taken with respect to the stationary state. \[One could extend the definition of the susceptibility to account for experiments with an additional time-dependent (e.g. constantly sweeped) magnetic field $`H_z(t)`$ and consider initial states before $`t^{}`$ that are combinations of the stationary and the relaxing modes.\] After $`\stackrel{~}{\mathrm{\Sigma }}(t^{\prime \prime },0)`$, there may be any number of $`\mathrm{\Sigma }`$’s between this and the final time $`t`$ – the sum of such series is denoted $`\mathrm{\Pi }(t,t^{\prime \prime })`$. Finally the resulting diagram/expression is Laplace transformed – $`e^{i\omega (tt^{})}`$ being divided into two parts, one denoted by the dashed line running through the diagram from time $`t^{\prime \prime }`$ to $`t`$ and the other one extending from $`t^{}`$ to $`t^{\prime \prime }`$ and already belonging to the above definition of $`\stackrel{~}{\mathrm{\Sigma }}`$.
In evaluating the diagram in Fig. 8, we vary $`t^{}`$ within $`\stackrel{~}{\mathrm{\Sigma }}(t^{\prime \prime },0)`$, and extend the integrations of $`t`$ and $`t^{\prime \prime }`$ to infinity. This yields
$`{\displaystyle _t^{}^{\mathrm{}}}𝑑te^{i\omega (tt^{})}[S_z(t),S_z(t^{})]={\displaystyle \underset{d,d^{},d^{\prime \prime }}{}}m_d`$ (30)
$`{\displaystyle _0^{\mathrm{}}}𝑑t^{\prime \prime }{\displaystyle _{t^{\prime \prime }}^{\mathrm{}}}𝑑t{\displaystyle _0^{t^{\prime \prime }}}𝑑t^{}\mathrm{\Pi }(t,t^{\prime \prime })_{d,d^{\prime \prime }}\stackrel{~}{\mathrm{\Sigma }}(t^{\prime \prime },0)_{d^{\prime \prime },d^{}}\rho _d^{}^{(0)}e^{i\omega (tt^{\prime \prime })},`$ (31)
where, for brevity, we denote the diagonal states by just one index $`d`$, e.g., $`m_d=m_{d,d}`$, see above. The integrations can be carried out one by one in the order $`t^{}`$, $`t`$ (which yields $`\mathrm{\Pi }(z=\omega )`$), and $`t^{\prime \prime }`$ to yield
$`\chi (\omega )=(g\mu _\mathrm{B})^2`$ (32)
$`{\displaystyle \underset{d,d^{},d^{\prime \prime }}{}}m_d\left[{\displaystyle \frac{1}{i\omega +iL_0\mathrm{\Sigma }(\omega )}}\right]_{d,d^{\prime \prime }}\stackrel{~}{\mathrm{\Sigma }}(\omega )_{d^{\prime \prime },d^{}}\rho _d^{}^{(0)},`$ (33)
where $`\mathrm{\Sigma }(\omega )`$ is defined as the Laplace transform of $`\mathrm{\Sigma }(t,0)`$. The explicit expressions for $`\stackrel{~}{\mathrm{\Sigma }}(\omega )_{d,d^{\prime \prime }}`$ turn out lengthy and are omitted here.
The resolvent in Eq. (32) is treated in terms of the eigensolutions of $`W\rho (t)`$ and the appropriate projections of the terms $`\stackrel{~}{\mathrm{\Sigma }}(\omega )\rho ^{(0)}=_ic_i\rho ^{(i)}`$ are used. In particular, for the component along the relaxing mode $`\rho ^{(1)}(t)`$, the resolvent in Eq. (32) reads
$`{\displaystyle \frac{1}{i\omega +iL_0\mathrm{\Sigma }(\omega )}}={\displaystyle \frac{1}{i\omega W(\omega )}}{\displaystyle \frac{1}{i\omega +\tau _1^1}}.`$ (34)
The approximate sign is due to the Markov approximation where $`\mathrm{\Sigma }(\omega )\mathrm{\Sigma }(0)`$, cf. App. D. Apart from one $`\tau _1`$ this is the well-known factor in expressions for susceptibility, see e.g. Refs. or . The rest of Eq. (32) reduces to a prefactor roughly proportional to $`\tau _1^1`$ but with weak $`H_z`$ and $`\omega `$-dependences. For the low frequency limit, we recover the static susceptibility, i.e., $`lim_{\omega 0}\chi (\omega )=\chi _0`$ where $`\chi `$ is given by Eq. (26).
## IV Results
In this section, we concentrate on results obtained in the eigenbasis of $`_S`$. This choice is advocated by three points. First, it allows us to consider much stronger transverse fields (see the discussion on energy scales in Part II) than the common approach to calculate tunnel splittings in the leading-order perturbation theory; second, the origin of the Lorentzian shape of the resonances is seen to arise naturally from the tunnel splittings. The third point is that, in this basis, all the relevant properties of the system are captured by the diagonal elements $`\rho (t)_{d,d}`$ only. This feature also improves the performance of the numerics. We also present results obtained with the full $`\rho (t)`$ and discuss the differences as examples of phonon-induced decoherence. The energies are expressed in kelvin, magnetic fields in teslas, and rates in $`s^1`$.
### A Relaxation rates
The relaxation rate as a function of the longitudinal magnetic field $`H_z`$ is shown in Fig. 9. The overall behaviour is of the Arrhenius form, i.e., $`\tau \tau _0\mathrm{exp}[\beta U(H_z)]`$ where $`\tau _0`$ only gives a constant offset to the semilogarithmic figure and $`U`$ is the effective height of the barrier, see Fig. 3. On top of this exponential field dependence, there are series of resonances located at
$`H_z^{m,m^{}}={\displaystyle \frac{m+m^{}}{g\mu _\mathrm{B}}}\left[A+B(m^2+m_{}^{}{}_{}{}^{2})\right]`$ (35)
corresponding to the values of the external field that brings the states $`m`$ and $`m^{}`$ on-resonance; the resonances form groups close to the fields $`H_znH_1`$ where $`n=m+m^{}`$ and $`H_1=\frac{A}{g\mu _\mathrm{B}}`$, see Fig. 9. The broadest of the resonances resemble those seen in usual experiments, while most of the resonances are too narrow to be seen and, as is discussed below, already the phonon coupling is sufficient to suppress them. One of the main topics of this paper is the possibility to detect some of the satellite peaks by the application of a relatively strong transverse field $`H_{}`$ as well as to suppress the already visible peaks by pointing $`H_{}`$ along one of the hard axes.
The whole curve in Fig. 9 can be understood in terms of the expression
$`\tau ={\displaystyle \underset{n}{}}\tau _0^{(n)}D_n(H_z)e^{\beta U_n(H_z)}`$ (36)
where the index $`n`$ enumerates the resonances (pairs of states $`m_n`$ and $`m_n^{}`$). As discussed in App. E, there is a certain $`\tau _0^{(n)}`$ for each resonance; the function $`D_n(H_z)`$ is a tunneling-induced Lorentzian that yields the peak shapes, and $`U_n(H_z)`$ is the effective barrier height for the $`n`$th resonance, i.e., $`U_n=E_{m_n}E_{10}`$ ($`E_{10}`$ is the energy of the metastable minimum). The widths of the Lorentzian peaks are given by
$`\delta H_z={\displaystyle \frac{4|\mathrm{\Delta }_n|}{g\mu _\mathrm{B}|m_nm_n^{}|}},`$ (37)
where $`2|\mathrm{\Delta }_n|`$ is the tunnel splitting between the resonant states, cf. App. E. As is seen in the figure and also suggested by the inset in Fig. 3, the background is not given by the over-barrier relaxation but by the “leakage” due to the direct coupling of states with $`mm^{}=\pm 4`$, e.g., $`m=\pm 2,m^{}=2`$ close to $`H_z=0`$ and $`m=\pm 1,m^{}=3`$ close to $`H_z=2H_1`$. This effectively lowers the barrier height to $`UE_2E_{10}`$. The rates that determine $`\tau _0^{(n)}`$, depend only weakly on temperature, the resonant states, and $`H_z`$, cf. App. E. This is in line with the experimental observation that the estimated prefactor of the Arrhenius law, $`\tau _0`$, varies within an order of magnitude for different samples and resonances.
The sharpness of the narrowest peaks in Fig. 9 is an artefact of the reduced model used so far, i.e., neglect of the nondiagonal matrix elements, and let us next consider two things affecting these peaks: decoherence due to spin-phonon coupling and the broadening effect of a transversal magnetic field $`H_{}`$. For other contributions such as the hyperfine and dipolar interactions, see Sec. IV C.
The decay of the nondiagonal elements of the density matrix is a classical definition of decoherence and this is also what is seen here when the above calculation is performed using the full $`\rho (t)`$. Inclusion of the off-diagonal elements still produces Lorentzian peaks but now the narrow resonances – with widths of the same order of magnitude as the spin-phonon coupling – are broadened and reduced in height. The narrowest peaks may even be completely suppressed, see the solid-line curves in Fig. 10. Despite its intuitive appeal, the suppression of the narrow peaks should be taken only qualitatively because of the limitations of the Born approximation used in treating the spin-phonon coupling. For this reason, the following considerations are restricted to regimes where the effect of the nondiagonal elements of $`\rho (t)`$ is negligible.
The suppressing effect of the spin-phonon coupling being essentially a constant, the sharper resonances can be made observable by broadening them with a transverse field $`H_{}`$; even ground state tunneling has been observed in high enough fields. The various tunnel splittings are shown in Fig. 4 as a function of $`H_{}`$ pointed to three different directions $`\varphi `$. It can be seen that for small values of the field, the splittings are essentially independent of the chosen angle, see also Fig. 11. Furthermore, tunnel splittings between states that can be coupled with sole $`B_4`$ terms of $`_\mathrm{T}`$ are quite insensitive to the transverse field below 0.1-0.2T. Figure 10 illustrates the effect of $`H_{}=H_x`$ onto the two first series of peaks increasing $`H_x`$ monotonously increases the tunnel splittings. The resonances that initially ($`H_x=0`$) were broader than $`\tau _{0}^{(n)}{}_{}{}^{1}`$ are seen to retain their height; the narrower peaks, such as the two shown in the insets, are strongly broadened with increasing $`H_x`$ and also their heights increase once their widths become larger than the spin-phonon rates.
A transverse field cannot only increase the tunnel splittings but it can also reduce them, see Fig. 4. In the ideal case, the splittings can even be suppressed by varying the angle $`\varphi `$ where the transverse field is pointed. For the zeroth peak, i.e., for $`H_z0T`$, all ten resonances occur at the same position and the suppression of any one of them is obscured in the relaxation rate curves by the simultaneous broadening of other resonances. For this reason, let us first focus on the resonances occuring at finite $`H_z`$ and return to the $`H_z0`$ case in the next section.
Figure 11 shows the tunnel splittings $`\mathrm{\Delta }_{4,5}`$ and $`\mathrm{\Delta }_{3,5}`$, respectively, as a function of $`H_{}`$ and for four different values of $`\varphi `$. The former corresponds to the resonance shown in the left inset of Fig. 10, while the latter corresponds to the broad resonance in the right panel of Fig. 10 which is also the one seen in experiments at $`H_z2H_1`$. Figure 11 also exemplifies some general properties of the tunnel splittings. As already mentioned above, the $`\mathrm{\Delta }`$’s depend only weakly on $`\varphi `$ for low $`H_{}`$, while for larger fields and $`\varphi `$ close to the hard axes, the tunneling amplitudes are successively increased and reduced as $`H_{}`$ is increased. On the other hand, with $`H_{}`$ held fixed to one of the suppression points, the $`\mathrm{\Delta }`$’s oscillate as a function of $`\varphi `$. Similar to recent work on Fe<sub>8</sub>, see e.g. Ref. , these phenomena are attributed to alternating constructive and destructive interference of the geometrical phase in the tunneling amplitudes. The number of minima for a given $`\mathrm{\Delta }_{m,m^{}}`$ appears to be, as was also pointed out in Ref. , given by the number of $`B_4`$-terms involved in coupling the states $`m`$ and $`m^{}`$. This ranges from zero at the top of the barrier to five for the gound states $`m=\pm 10`$, cf. Fig. 4.
Figure 12 shows the relaxation rates corresponding to the resonances of Fig. 11 with $`H_{}`$’s chosen to match the first points of suppression $`H_{}^{m,m^{}}`$. The resonance width is found to be very sensitive to $`\varphi `$ quite as expected and, for $`\varphi =45^{}`$, the left-most peak becomes suppressed. Note, that the suppression is also quite sensitive to the value of $`H_{}`$ as well and the higher peaks in Fig. 12 are hardly affected at all by changes in $`\varphi `$.
### B Ac-susceptibility
In explaining experimental results, the susceptibility $`\chi (\omega )`$ is often described by the formula
$`\chi (H,T,\varphi ,\theta ,\omega )={\displaystyle \frac{\chi _0(H,T,\theta )}{1i\omega \tau (H,T,\varphi ,\theta )}}`$ (38)
(39)
where $`\tau (H,T,\varphi ,\theta )`$ is the relaxation time, $`\omega `$ the frequency of the small excitation field, and $`\chi _0(H,T,\theta )`$ the static susceptibility, cf. Ref. . We emphasize here the $`\varphi `$-dependence as we wish to investigate the interference effects and resulting oscillations in $`\tau _1`$ (below we do not explicitly write the parametric dependences of $`\tau `$ but they are implicitly assumed). Equation (39) relates to the results of Sec. III C as follows. Equation (32) formally accounts for all the modes $`i`$ and for each mode divides into two contributions, one being $`1/(1i\omega \tau _i)`$ and the other the respective prefactor. It turns out that, for most parameter values, the over-barrier relaxation, i.e., the mode $`i=1`$ strongly dominates over the others and one can use Eq. (39) to a good approximation. Figure 13 illustrates the two contributions to Eq. (39); they form the basis for understanding the following results. In all figures, the susceptibility is expressed for a single spin and in units of $`\mathrm{KT}^2`$.
In the present model, the main correction to Eq. (39) corresponds to an intra-valley mode describing transitions between the two lowest states on the same side of the anisotropy barrier. This correction increases with increasing $`H_z`$ and temperature and, for the parameters and the curves of $`\chi _0`$ in Fig. 13, becomes of the same order of magnitude as the actual relaxing mode once $`H_z0.91`$T. However, such a mode corresponds to very high frequencies and is neglected below. There could be also other sources of high-frequency contributions such as dipole-dipole flip-flops, nuclear spin dynamics, or moving impurities but these are beyond the scope of this work.
As to the actual results, it is in principle sufficient to combine the curves for $`\tau `$ of the previous section with those shown in Fig. 13. The shapes of the real and imaginary parts of $`1/(1i\omega \tau )`$ readily imply how the resulting susceptibility behaves if one fixes $`\omega `$ – it turns out that all the structure found in $`\tau ^1`$ in the previous section can also be found in $`\chi (\omega )`$. First, the response is the most sensitive to changes in $`\tau `$ when $`\omega \tau ^1`$ or $`\omega \tau 1`$, cf. the inset of Fig. 13. Second, $`\mathrm{Re}(\chi (H_z;\omega ))`$ replicates the shape of $`\tau (H_z)`$ and exhibits peaks and valleys as $`\tau ^1`$ increases and decreases. This is the case with the imaginary part, $`\mathrm{Im}(\chi (H_z;\omega ))`$, as well, but only as long as $`\omega \tau 1`$, i.e., as long as we remain on the right hand side of the peak in $`\mathrm{Im}(\chi (H_z;\omega ))`$. For lower $`\omega `$, $`\omega \tau `$ can cross the maximum point and the picture with the peaks and valleys turns upside down. These points are illustrated in Fig. 14 for a case corresponding to one of the curves in 12.
Let us next consider a different scheme, keeping $`H_z`$ and $`\varphi `$ fixed and varying $`H_{}`$ starting from zero field. The case $`H_z0`$T is of particular interest because it allows comparison to recent experiments by Wernsdorfer done on the related material Fe<sub>8</sub>, cf. Ref. – the experimental data showed clear oscillations of the relaxation rates as a function of $`H_{}`$ in the thermally activated regime. By choosing parameters in the feasible range of these experiments, we find somewhat similar oscillations also for Mn<sub>12</sub>. Figure 15 shows first the relaxation rates for two different combinations of temperature and the longitudinal field: $`T=3.0`$K and $`H_z=0.2`$mT to show the behaviour close to (or at) the peak maximum and at a lower temperature, and $`T=5.0`$K and $`H_z=10`$mT as an example of the behaviour further away from the maximum and at a higher temperature. Also the $`\varphi `$ dependence is shown. We would like to stress the novelty of this result: such oscillations have neither been observed nor predicted before.
For $`\varphi =0^{}`$, the gentle oscillations resembling a strongly-smeared staircase stem from the fact that the lower (in terms of energy; higher in terms of rates) resonances are broadened and one by one start to dominate the relaxation, see the dot-dashed curves in Fig. 16 for the tunnel splittings and Fig. 17 for the general idea. All the additional structure, seen for $`\varphi =45^{}`$, is due to the oscillations in the tunnel splittings and the suppression of some of the resonances. The positions of the notches can be compared with the structure of $`\mathrm{\Delta }_{m,m}`$ in Fig. 16 and one finds that for smaller $`H_z`$ the relaxation takes place via the lower-lying resonances such as $`m=\pm 6`$ and $`m=\pm 7`$; further away from the maximum, the broader resonances between states $`m=\pm 4`$ and $`m=\pm 5`$ act as the dominant relaxation paths.
Figures 18 and 19 show the real part of the susceptibility $`\chi ^{}`$ corresponding to the two cases of Fig. 15. The purpose of the different frequencies is to show that also here one can choose the structure of interest and study it by tuning the frequency to fulfill $`\omega \tau 1`$. In both figures, there are regimes where the susceptibility can be varied by a factor of five; in Ref. the oscillations were quite clear already with the amplitude being a mere 20% of the signal.
### C Discussion – relevance to experiments
So far we have considered the simple model comprising a single spin coupled to a phonon bath but, as was pointed out in the introduction, in real samples there are also other kinds of interactions. In this section, we consider the additional features arising from the hyperfine and/or dipolar interactions and aim to point out the experimentally relevant aspects of the results obtained above.
Let us first recall some experimental facts concerning relaxation measurements and results – these underlie also the understanding and appreciation of the susceptibility measurements. Relaxation rates are typically measured by first magnetizing the sample to saturation, and then reversing the direction of the field and measuring the resulting magnetization as a function of time. The initial relaxation is observed to be nonexponential – this is attributed to dipolar interactions, see below – while, at later times, it becomes exponential. Several authors have proposed an extended exponential $`M(t)=M(0)\mathrm{exp}[(t/\tau )^\beta ]`$ to account for both of these regimes with just one additional fitting parameter $`\beta `$. In Ref. it was found that $`\beta `$ varies from $`\beta 0.5`$ below 2.0K to $`\beta 1`$ – usual exponential relaxation – roughly above 2.4K. The thus obtained relaxation rates exhibit a series of broad Lorentzian-shaped resonances; their height and location correspond to tunneling-assisted relaxation 3-4 levels below the top of the barrier.
This shows two clear differences as compared to the present work: in experiment, the relaxation may be nonexponential even though the single-spin model always yields exponential behaviour, and no satellite peaks are observed (see Ref. for exceptions). In order to understand these discepancies, let us first consider the effect of the nuclei via the hyperfine interactions and then the intermolecular dipolar interactions.
Hyperfine interactions. In Mn<sub>12</sub>, all the manganese nuclei have magnetic momenta and the hyperfine interaction between the nuclei and the molecular spin state is relatively large, of the order of 10mT. Recently, several authors have investigated how this affects tunneling and the relaxation in Mn<sub>12</sub>. For the present purposes, the relevant effect of the hyperfine interactions is to induce an intrinsic Gaussian broadening, of the width $`\sigma _{\mathrm{hyp}}6`$mT, to all levels, i.e., the nuclear spins are importantly dynamic and their influence on the molecular spin cannot be reduced to a rigid but spatially varying background field. Simultaneously with the broadening, the hyperfine interactions reduce the tunneling amplitudes of the resonances for which $`\mathrm{\Delta }<\sigma _{\mathrm{hyp}}`$ – this should lead to reduced peak heights in $`\tau ^1(H_z)`$.
With this in mind, let us consider the different regimes in terms of the relative magnitudes of the tunnel splitting and the hyperfine broadening. First, for resonances with $`\mathrm{\Delta }\sigma _{\mathrm{hyp}}`$ the shapes of the resonances are expected to be Lorentzian with the widths determined by the tunnel splittings $`\mathrm{\Delta }`$. This is the regime, where all our results apply. In the other extreme, $`\mathrm{\Delta }\sigma _{\mathrm{hyp}}`$, the resonances should be essentially suppressed providing one possible explanation why the sharp satellite peaks are not observed in experiments. Note that the minuscule phonon-induced broadening or dephasing is hidden under the hyperfine broadening and cannot be seen. In this regime, one could in principle try and extend the present theory by adding by hand a strong dephasing term to the nondiagonal density matrix elements. The intermediate regime where $`\mathrm{\Delta }\sigma _{\mathrm{hyp}}`$ is the most interesting of the three cases. In this regime, the peak shape should be a combination of Lorentzian and Gaussian curves and, depending on which one of $`\mathrm{\Delta }`$ or $`\sigma _{\mathrm{hyp}}`$ is larger, one of the shapes should dominate. In the tail region, i.e., away from the peak maxima, the Lorentzian tails dominate and it has been suggested that this together with experimental error bars may obscure the resolution between the two types of curves, cf. e.g. Ref. .
The immediate conclusion from these considerations is that, if the application of $`H_{}`$ broadens some of the tunnel splittings $`\mathrm{\Delta }_{m,m^{}}`$ to exceed $`\sigma _{\mathrm{hyp}}`$, the corresponding resonance should become observable. If, on the other hand, the transverse field is applied along one of the hard axes, a given resonance becomes suppressed for certain special values of $`H_{}`$; in the presence of the hyperfine interactions this should happen already when $`\mathrm{\Delta }_{m,m^{}}`$ becomes smaller than $`\sigma _{\mathrm{hyp}}`$. The intermediate regime can be intentionally achieved by tuning the tunnel splitting from being well below $`\sigma _{\mathrm{hyp}}`$ to above it. This may provide means to probe the Gaussian broadening, see also the subsection below focusing on the advantages of susceptibility measurements.
Dipolar interactions. The intermolecular spin-spin interactions are of dipolar form and they are weaker in Mn<sub>12</sub> than, e.g., in Fe<sub>8</sub>. Due to their short range, the dipolar fields can vary in space changing the local field at the position of the individual molecules. In our view, the essential difference between the hyperfine and dipolar interactions can be stated as follows: even if one could measure the response of a single molecule, this would always be dressed by the level broadening and reduction in tunneling amplitudes due to hyperfine interactions intrinsic to each molecule; the dipolar fields, on the other hand, just change the molecule’s local electromagnetic environment.
In experiment, the relaxation of the magnetization $`M(t)`$ leads to time-dependent dipolar fields and, in order to describe the relaxation correctly, it would be necessary to solve for $`M(t)`$ self-consistantly, for simulations see e.g. Refs. and . However, it is this time-dependent field that provides an explanation to the initial nonexponential relaxation. For example in Ref. , it is therefore concluded that the deviation from a single-exponential relaxation, i.e., $`\beta 1`$, demonstrates the important role of the dipolar interactions and dynamics of the spin distribution. It should be kept in mind, though, that also a static distribution of local fields (be it dipolar or not) – and hence relaxation rates $`\tau ^1(H_z^{\mathrm{local}})`$ – leads to a superposition of exponential rates, which looks nonexponential. Such a distribution of fields also hides all features in $`\tau ^1(H_z)`$ that are sharper than this distribution.
The time-dependence of the dipolar fields owes to the fact that the sample is first magnetized and, as the field direction is abruptly reversed, the dipolar distribution finds itself far from equilibrium and quickly starts to relax. The reason for such experiments is the strong response from almost all the spins.
In anticipation of the discussion on susceptibility, let us consider the dipolar distributions at equilibrium. By distribution we mean spatial variations in the dipolar field at the locations of the individual molecules. First, for $`H_z0`$T, the annealed (not quenched) distribution is random but due to the low temperatures, $`k_\mathrm{B}TE_{\pm 9}E_{\pm 10}`$, almost all the spins are aligned with the easy $`z`$-axis – randomly pointing to the positive and negative directions – and only contribute to the local longitudinal field. On the other hand, the equilibrium magnetization is close to saturation already for $`H_zH_1`$, thus drastically narrowing the dipolar distribution – e.g. for $`T=3.0`$K (5.0K) and $`H_z=H_1`$, more than 95% (85%) of all the spins are aligned parallel to $`H_z`$ and all the molecules feel essentially the same field. Such distributions have been experimentally verified in Fe<sub>8</sub>, cf. Ref. .
Susceptibility. The influence of the dipolar dynamics on relaxation can be avoided almost completely by measuring linear response to a small ac-field, i.e., the ac-susceptibility, instead of $`M(t)`$. This has the advantage that the system is probed in its equilibrium state and ideally by a small enough field that in itself does not perturb the equilibrium. Therefore we propose that the susceptibility measurements provide a gentle or noninvasive means to probe the relaxation dynamics in absence of the time-dependent dipolar distribution. Furthermore, while the hyperfine interactions cannot be tuned, the static distribution can be made markedly narrower by a finite $`H_z`$, see above. As on the other hand $`\chi _0`$ decreases with increasing $`H_z`$, the first group of resonances, close to $`H_1`$, is especially attractive for investigating the oscillations in the relaxation rates as well as the hyperfine fields themselves. All of the resonances around $`H_1`$ depend strongly on $`H_{}`$ and can be broadened such that $`\mathrm{\Delta }>\sigma _{\mathrm{hyp}}`$ making them observable; the peaks can also be selectively suppressed if $`\varphi 45^{}`$. The hyperfine fields may even simplify the observation of the suppressions as very narrow peaks are strongly reduced in height. For a sharp dipolar distribution, we expect that also the crossover between Lorentzian and Gaussian shapes of $`\tau ^1(H_z)`$ should be observable when changing the tunnel splittings with the transverse field.
## V Conclusions
To conclude, we present a diagrammatic description of the spin dynamics of the molecular magnet Mn<sub>12</sub>. The work focuses on the regime of thermally-activated tunneling, i.e., $`T>2.0\mathrm{K}`$, and emphasizes the phenomena that could be observed for strong transverse magnetic fields. In the calculations, we study the dynamics of a single spin $`S=10`$ coupled to a phonon bath. The role of the phonons is in the thermal activation of the spins to states with higher energies and larger tunneling amplitudes.
As the first main result, we calculate the dynamic susceptibility $`\chi (\omega )`$ starting from the same microscopic Hamiltonian as is used for the relaxation rates. Susceptibility is found to reflect the rich structure found in $`\tau ^1(H_z)`$ and we argue that susceptibility measurements are in fact more sensitive and better controlled in terms of time scales and also the dipolar interactions than the relaxation experiments.
All the results obtained are calculated using the eigenbasis of the spin Hamiltonian, which naturally accounts for strong transverse magnetic fields. A strong transverse magnetic field enhances tunneling through the anisotropy barrier and enables relaxation via eigenstates further away from the top of the barrier. In relaxation or susceptibility measurements, this would lead to shifted and higher resonances. The tunnel splittings are found to be very sensitive to the azimuth angle $`\varphi `$ of the transverse field $`H_{}`$. It is found that, in the directions $`\varphi =\pi (2n+1)/4`$, the tunnel splittings exhibit alternating minima and maxima and become totally suppressed at certain values of $`H_{}`$. This phenomenon attributed to the interference of the geometrical or Berry phase of alternative tunneling paths, with a destructive interference leading to the suppressions. As the second major result, we predict that these oscillations in the tunnel splittings should be observable both in the relaxation rates and the susceptibility.
## Acknowledgements
We are grateful to Myriam Sarachik and Yicheng Zhong as well as Wolfgang Wernsdorfer for the discussions on their experiments and M.S. and Y.Z. for providing us with their unpublished data. We would also like to thank Michael Leuenberger and Daniel Loss for the fruitful exchange of theoretical ideas. The calculations were in part carried out in the Center for Scientific Computing (CSC) in Finland. This work has been supported by the Finnish Academy of Science and Letters, the Finnish Cultural Foundation, the EU TMR network “Dynamics of Nanostructures”, the Swiss National Foundation, and DFG through SFB 195.
## A Diagrammatic rules
In this Appendix, we list the rules for evaluating the diagrammatic expressions for the kernels $`\mathrm{\Sigma }=𝑑t\mathrm{\Sigma }(t)`$ arising in the evaluation of $`S_z(t)`$ and the kinetic equation. We assume that the Markov approximation is valid and check the results for self-consistency.
### 1 Real-time representation
In the time domain, the rules for $`\mathrm{\Sigma }(t)`$ are
1. Propagators
Assign a factor $`\mathrm{exp}[i_S(t_{i+1}t_i)]`$ to each piece of a forward propagator between two vertices at times $`t_i`$ and $`t_{i+1}`$ $`(t_i<t_{i+1})`$. For a particular diagram with specific states at the ends of the propagators, cf. Fig. 20a, take the element
$`\left[e^{i_S(t_{i+1}t_i)}\right]_{m_{i+1}^{},m_i}.`$ (A1)
For a backward propagator, change the sign in the exponent and invert the order of the states $`m_i`$ and $`m_i^{}`$. In the $`m`$-basis, the tunneling contained in $`_S`$ changes the states along the propagator but in the $`d`$-basis the above exponential reads
$`e^{iE_{d_i}(t_{i+1}t_i)}.`$ (A2)
A compact way to account for the tunneling in Eq. (A1) is to rewrite it with the help of Eq. (A2) as
$`{\displaystyle \underset{d_i}{}}m_{i+1}^{}|d_ie^{iE_{d_i}(t_{i+1}t_i)}d_i|m_i.`$ (A3)
In so doing, we come to incorporate the tunneling elements exactly to the corresponding diagram.
2. Vertices
Assign a prefactor $`i`$ (+i) to each vertex lying on the forward (backward) propagator. For each pair of vertices at times $`t_i`$ and $`t_{i+1}`$ and coupled by a phonon line, cf. Fig. 20b, assign a factor
$`G_{\overline{m}m^{},m_1\overline{m}_1}=\sqrt{S_{\overline{m},m^{}}S_{m_1,\overline{m}_1}}C_{\xi ,\xi ^{}},`$ (A4)
where $`\xi =\overline{m}m^{}`$ and $`\xi =m_1\overline{m}_1`$, see App. B.
3. Spin-phonon interaction lines
A phonon line connecting a pair of vertices corresponds to a phonon correlation function – this is worked out in detail in App. B. It is more convenient to first calculate this in the energy representation
$`\mathrm{\Gamma }(\omega )={\displaystyle \frac{A^2}{12\rho c^5\mathrm{}^4}}{\displaystyle \frac{\omega ^3}{e^{\beta \omega }1}}`$ (A5)
and then Fourier transform it
$`\mathrm{\Gamma }(tt^{})`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega \mathrm{\Gamma }(\omega )e^{i\omega (tt^{})}`$ (A6)
$``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega {\displaystyle \frac{A^2}{12\rho c^5\mathrm{}^4}}{\displaystyle \frac{\omega ^3}{e^{\beta \omega }1}}e^{i\omega (tt^{})}.`$ (A7)
Here $`A`$ is the same as in $`_z`$, cf. Ref. , $`\rho =1.8310^3\mathrm{kg}/\mathrm{m}^3`$ is the density of Mn<sub>12</sub>, $`c`$ is the sound velocity (the only “fitting” parameter of the theory), and $`\omega `$ is the phonon energy, $`\omega >0`$ meaning absorption and $`\omega <0`$ emission. Note that the prefactor here differs from that found in the literature for the spin-phonon rates. This is just because for convenience we assign part of it to the vertices, cf. Eqs. (B18)-(B17). Fourier transforming the correlator into the time domain gives
$`\mathrm{\Gamma }(\sigma t)={\displaystyle \frac{A^2}{12\rho c^5\mathrm{}^4}}`$ (A8)
$`\left\{2\left({\displaystyle \frac{\pi }{\beta }}\right)^4\left[{\displaystyle \frac{1+2\mathrm{ch}^2\left(\frac{\pi \sigma t}{\beta }\right)}{\mathrm{sh}^4\left(\frac{\pi \sigma t}{\beta }\right)}}\right]+i\pi \delta ^{\prime \prime \prime }(\sigma t)\right\}`$ (A9)
where $`\sigma =+1`$ (-1) if the vertex with the earlier time lies on the forward (backward) propagator; the function $`\delta ^{\prime \prime \prime }(t)`$ is the third time derivative of the delta function. This expression diverges at $`t=0`$ but this is just an artefact that is removed by including a high-energy cutoff in Eq. (A7). For instance, an exponential cutoff $`e^{\omega /D}`$ would lead to the replacement $`tti/D`$ in Eq. (A8) and thus avoiding the divergence. We are not interested in the limit $`t0`$ but rather the longer-time behaviour of the correlations: the correlator decays exponentially on the time scale of $`1/k_\mathrm{B}T`$ which is 5-10 orders of magnitude faster than typical relaxation times $`\tau `$ of the reduced density matrix. This observation is used to justify the Markov approximation in the text.
4. Summations and integrations
Sum over the states internal to the diagram. For example, in evaluating the element $`\mathrm{\Sigma }_{mm_1,m^{}m_1^{}}`$ in Fig. 20b, sum over the states $`\overline{m}`$ and $`\overline{m}_1`$, and integrate over all the times (time differences) in the diagram.
We might add a fifth rule to capture a general property of combinations of diagrams: a series of irreducible diagrams is evaluated iteratively such that all the information needed from earlier processes/diagrams is incorporated into the reduced density matrix $`\rho (t)`$. As an example, the diagram in Fig. 20b equals
$`ii{\displaystyle \underset{\overline{m},\overline{m}_1}{}}G_{\overline{m}m^{},m_1\overline{m}_1}`$ (A10)
$`{\displaystyle \underset{d,d_1}{}}m_1^{}|d_1d_1|m_1m|dd|\overline{m}`$ (A11)
$`{\displaystyle _0^t}dt^{}e^{i(E_dE_{d_1})(tt^{})}\mathrm{\Gamma }(t^{}t)\rho (t^{})_{m^{}m_1^{}}.`$ (A12)
### 2 Energy representation
Once the rules for calculating $`\mathrm{\Sigma }`$ have been laid down in time domain, it is more convenient to shift into the energy representation with respect to the phonon energies and to do this in the $`d`$-basis. The rule 2. remains as it is, Eq. (A5) in the rule 3. is already in the energy representation, and we only need to reformulate the first and fourth rules.
1. Propagators
Assign a resolvent
$`{\displaystyle \frac{i}{\sigma \omega E_d^{}+E_d+i\eta }}`$ (A13)
to each piece of a diagram temporally between two vertices, i.e., irrespective of the propagator they lie on. Here $`\omega `$ is the phonon energy, and $`E_d^{}`$ and $`E_d`$ are the energies of the states on the forward and backward propagators, respectively; $`\sigma =+1`$ $`(1)`$ if the vertex with the earlier time lies on the forward (backward) propagator. The factor $`\eta >0`$ arises from the adiabatic turning on of the interaction and is taken to zero at the end.
1. Summations and integrations
Sum over all internal indices and integrate over $`\omega `$. Due to the factor $`i\eta `$, the integrations are to be understood as combinations of Cauchy’s principal value integrals and delta functions.
In diagrams such as those in Fig. 21a the prefactors due to the vertices are equal. Therefore these can be pairwise combined and their integral parts written together as
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega \mathrm{\Gamma }(\omega )`$ (A14)
$`\left[{\displaystyle \frac{i}{\omega +(E_d^{}E_{d_1})+i\eta }}+{\displaystyle \frac{i}{\omega (E_{d_1^{}}E_d)+i\eta }}\right].`$ (A15)
If $`d^{}=d_1^{}`$ and $`d=d_1`$, i.e., if the density matrix is diagonal before and after the transition, this expression simplifies into
$`\pi {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega \mathrm{\Gamma }(\omega )\delta (\omega (E_d^{}E_d))`$ (A16)
$`=\pi \mathrm{\Gamma }(E_d^{}E_d)={\displaystyle \frac{\pi A^2}{12\rho c^5\mathrm{}^4}}{\displaystyle \frac{\mathrm{\Delta }E^3}{e^{\beta \mathrm{\Delta }E}1}}.`$ (A17)
In the last step we inserted the definition of $`\mathrm{\Gamma }(\omega )`$ from the equation (B13) and identified $`\mathrm{\Delta }E=E_d^{}E_d`$ as the energy required for the transition from the state $`d`$ to $`d^{}`$. This expression, Eq. (A16), together with the prefactor $`G_{d^{}d,d_1d_1^{}}`$ is the phonon induced transition rate $`\mathrm{\Sigma }_{d,d^{}}`$ used in the literature and in Parts III and IV of this paper. From Eqs. (A14) and (A16) we see that the energy is only conserved (the delta function, cf. Fermi’s golden rule) under some special circumstances and, in general, the rate need not be the simple real function used in the literature.
## B The spin-phonon rates
In this Appendix, we outline the calculation of the phonon-induced transition rates $`\mathrm{\Sigma }`$ between different spin states. In order to calculate $`\mathrm{\Sigma }`$ in lowest order in the spin-phonon coupling constants $`g_i`$, we need to evaluate the contractions $`_{\mathrm{sp}}(t)_{\mathrm{sp}}(t^{})_{\mathrm{ph}}`$, where the expectation value is taken with respect to the phonon degrees of freedom. Let us first insert Eq. (8) into Eq. (4) and explicitly write down all the resulting terms
$`_{\mathrm{sp}}(t)=i{\displaystyle \underset{\stackrel{}{k}\sigma }{}}\sqrt{{\displaystyle \frac{\mathrm{}}{2MN\omega _{\stackrel{}{k}\sigma }}}}[b_{\stackrel{}{k}\sigma }^{}e^{i\omega _{\stackrel{}{k}\sigma }t}+b_{\stackrel{}{k}\sigma }e^{i\omega _{\stackrel{}{k}\sigma }t}]e^{i\stackrel{}{k}\stackrel{}{r}}`$ (B1)
$`\{g_1[S_x^2(t)S_y^2(t)][e_x^{(\sigma )}k_xe_y^{(\sigma )}k_y]`$ (B2)
$`+{\displaystyle \frac{g_2}{4}}\{S_x(t),S_y(t)\}[e_x^{(\sigma )}k_y+e_y^{(\sigma )}k_x]`$ (B3)
$`+{\displaystyle \frac{g_3}{4}}(\{S_x(t),S_z(t)\}[e_x^{(\sigma )}k_z+e_z^{(\sigma )}k_x]`$ (B4)
$`+\{S_y(t),S_z(t)\}[e_y^{(\sigma )}k_z+e_z^{(\sigma )}k_y])`$ (B5)
$`+{\displaystyle \frac{g_4}{4}}(\{S_x(t),S_z(t)\}[e_x^{(\sigma )}k_ze_z^{(\sigma )}k_x]`$ (B6)
$`+\{S_y(t),S_z(t)\}[e_y^{(\sigma )}k_ze_z^{(\sigma )}k_y])\}.`$ (B7)
In evaluating the contraction, we need to consider all the possible states before and after the action of $`_{\mathrm{sp}}(t^{()})`$ as well as all the possible orderings of the times $`t`$ and $`t^{}`$ along the contour. As a characteristic example, let us consider the diagram in Fig. 20b and hold to the same usage of indices. The time dependences of the spin operators are accounted for by the propagators and they can be put aside for the moment. The remaining part with the spin operators is fully determined by the spin states before and after the vertices yielding
$`\overline{m}_1|S^\xi ^{}|m_1\overline{m}|S^\xi |m^{}=\sqrt{S_{\overline{m},m^{}}S_{m_1,\overline{m}_1}}.`$ (B8)
Here $`\xi =\overline{m}m^{},\xi ^{}=m_1\overline{m}_1`$ can take values $`\pm 1`$ or $`\pm 2`$ due to the structure of $`_{\mathrm{sp}}`$, and, for brevity, $`S^{\pm |\xi |}S_\pm ^{|\xi |}`$. The quantities on the right hand side of the equation are
$`S_{\overline{m},m^{}}=(2m^{}+\xi )\sqrt{S(S+1)m^{}(m^{}+\xi )}`$ (B9)
for $`\xi =\pm 1`$ or
$`S_{\overline{m},m^{}}`$ $`=`$ $`\{[S(S+1)m^{}(m^{}+\nu )]`$ (B11)
$`[S(S+1)(m^{}+\nu )(m^{}+2\nu )]\}^{1/2}`$
for $`\xi =\pm 2`$, $`\nu =sign(\xi )`$.
The part of the contraction depending on the phonon degrees of freedom is evaluated assuming acoustic phonons with a linear dispersion relation and three modes indexed by $`\sigma =1,2,3`$ (one longitudinal and two transverse modes). The polarization vectors $`e_\alpha ^{(\sigma )}`$ are assumed unit vectors. The resulting expression consists of two parts. The first one of them contains all the information concerning the phonon spectrum, energies, and temperature,
$`\mathrm{\Gamma }(tt^{})`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega \mathrm{\Gamma }(\omega )e^{i\omega (tt^{})}`$ (B12)
$``$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega {\displaystyle \frac{A^2}{12\rho c^5\mathrm{}^4}}{\displaystyle \frac{\omega ^3}{e^{\beta \omega }1}}e^{i\omega (tt^{})}.`$ (B13)
This part is defined such that it is independent of the spin states and only contains the term $`A^2`$ from the coupling constants that turns out to be constant for all the rates. In the diagrams, this corresponds to the spin-phonon interaction line. The second part,
$`C_{\xi ,\xi ^{}}=\{\begin{array}{cccc}0,& \mathrm{for}|\xi ||\xi ^{}|\hfill & & \\ 1,& \xi =\xi ^{}=\pm 1\hfill & & \\ \frac{15}{16}+\frac{1}{8}\delta _{\xi ,\xi ^{}},& |\xi |=|\xi ^{}|=2\hfill & & \end{array},`$ (B17)
however, does depend on the spin states and implies additional selection rules. It should be noted that the third clause allows for $`\xi \xi ^{}`$ if $`|\xi |=|\xi ^{}|=2`$. For convenience, we combine the term $`C_{\xi ,\xi ^{}}`$ with the spin operators and obtain the term
$`G_{\overline{m}m^{},m_1\overline{m}_1}=\sqrt{S_{\overline{m},m^{}}S_{m_1,\overline{m}_1}}C_{\xi ,\xi ^{}},`$ (B18)
which in the diagrammatic language corresponds to a pair of vertices.
In the $`d`$-basis, the time dependence of the spin operators is of a simple exponential form, cf. Eq. (A3), and the diagram in Fig. 20b can be evaluated to yield
$`\mathrm{\Sigma }(tt^{})_{mm_1,m^{}m_1^{}}=`$ (B19)
$`{\displaystyle \underset{\overline{m},\overline{m}_1}{}}G_{\overline{m}m^{},m_1\overline{m}_1}{\displaystyle \underset{d,d_1}{}}m_1^{}|d_1d_1|m_1m|dd|\overline{m}`$ (B20)
$`e^{i(E_dE_d^{})(tt^{})}\mathrm{\Gamma }(t^{}t).`$ (B21)
The actual transition rates are obtained by integrating Eq. (B19) over the time difference $`\tau =tt^{}`$. This takes us to the energy representation, combining the exponential factor in Eq. (B19) with the $`e^{i\omega (tt^{})}`$ of the Fourier transform in Eq. (B13). The intergration over $`\tau `$ yields
$`\mathrm{\Sigma }_{mm_1,m^{}m_1^{}}=`$ (B22)
$`{\displaystyle \underset{\overline{m},\overline{m}_1}{}}G_{\overline{m}m^{},m_1\overline{m}_1}{\displaystyle \underset{d,d_1}{}}m_1^{}|d_1d_1|m_1m|dd|\overline{m}`$ (B23)
$`i{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega {\displaystyle \frac{\mathrm{\Gamma }(\omega )}{\omega E_d+E_d^{}+i\eta }}.`$ (B24)
As the final step in calculating $`\mathrm{\Sigma }`$, let us evaluate the integral in Eq. (B22) or
$`\stackrel{~}{I}_\sigma (\mathrm{\Delta }E,\beta )`$ $`={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega {\displaystyle \frac{\omega ^3}{e^{\beta \omega }1}}{\displaystyle \frac{1}{\sigma (\omega \mathrm{\Delta }E)+i\eta }}`$ (B25)
$`=`$ $`P{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega {\displaystyle \frac{\omega ^3}{e^{\beta \omega }1}}{\displaystyle \frac{1}{\sigma (\omega \mathrm{\Delta }E)}}i\pi {\displaystyle \frac{\mathrm{\Delta }E^3}{e^{\beta \mathrm{\Delta }E}1}}`$ (B26)
for the general case. Here $`\sigma =\pm 1`$. The imaginary part of the integral (the last term) is, up to a prefactor, just $`i\pi \mathrm{\Gamma }(\mathrm{\Delta }E)`$ and is independent of $`\sigma `$. The real part in turn can be evaluated using the calculus of residues. By combining the integral along the real axis with an infinite semicircle in the upper half plane to form a closed contour we obtain $`\mathrm{Re}\{\stackrel{~}{I}_\sigma (\mathrm{\Delta }E,\beta )\}=2\pi \mathrm{Im}\{enclosed\mathrm{residues}\}`$.
The real part of the integral $`\stackrel{~}{I}_\sigma (\mathrm{\Delta }E,\beta )`$ is divergent due to the $`\omega ^3`$-term and some kind of a cutoff procedure is needed here. We have chosen a functional cutoff using instead of $`\stackrel{~}{I}(\mathrm{\Delta }E,\beta )`$ the integral
$`I_\sigma (\mathrm{\Delta }E,\beta ,D)=`$ (B27)
$`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega {\displaystyle \frac{\omega ^3}{e^{\beta \omega }1}}{\displaystyle \frac{1}{\sigma (\omega \mathrm{\Delta }E)+i\eta }}\left({\displaystyle \frac{D^2}{D^2+\omega ^2}}\right)^2,`$ (B28)
i.e., the cutoff function is a Lorentzian squared and the additional parameter $`D`$ is the cutoff parameter of the Lorentzian. The integrand in (B27) has single poles at $`\omega =\mathrm{\Delta }Ei\sigma \eta `$, and at $`\omega =2m\pi i/\beta `$, where $`m`$ is a positive integer; there are also second-order poles at $`\omega =\pm iD`$. The residues $`C^{(1)}`$ from the first pole are real and can be neglected. Collecting the other poles in the upper half plane and evaluating the residues yields
$`\mathrm{Re}\left\{I_\sigma (\mathrm{\Delta }E,\beta ,D)\right\}=2\pi \left[{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}C_m^{(2)}+C^{(3)}\right],`$ (B29)
where
(B30)
$`C_m^{(2)}=\sigma {\displaystyle \frac{\beta ^2\mathrm{\Delta }ED^4}{(2\pi )^3}}{\displaystyle \frac{m^3}{\left(\frac{\mathrm{\Delta }E\beta }{2\pi }\right)^2+m^2}}{\displaystyle \frac{1}{\left[\left(\frac{D\beta }{2\pi }\right)^2m^2\right]^2}}`$ (B31)
$`C^{(3)}=\sigma {\displaystyle \frac{D^5}{4(\mathrm{\Delta }E^2+D^2)}}[{\displaystyle \frac{\beta \mathrm{\Delta }E}{4\mathrm{sin}^2\left(\frac{\beta D}{2}\right)}}`$ (B32)
$`{\displaystyle \frac{1}{\mathrm{\Delta }E^2+D^2}}({\displaystyle \frac{3}{2}}\mathrm{\Delta }E^2+{\displaystyle \frac{1}{2}}D^2+{\displaystyle \frac{\mathrm{\Delta }E^3}{D}}\mathrm{cot}\left({\displaystyle \frac{\beta D}{2}}\right))].`$ (B33)
If $`D`$ is accidentally chosen to equal $`2\pi m^{}/\beta `$ for some $`m^{}`$, the residues $`C_m^{}^{(2)}+C^{(3)}`$ should be replaced by
$`C^{}=\sigma {\displaystyle \frac{D^3}{16\beta }}{\displaystyle \frac{3\mathrm{\Delta }E^514\mathrm{\Delta }E^3D^2\mathrm{\Delta }ED^4}{(\mathrm{\Delta }E^2+D^2)^3}}.`$ (B34)
In the $`d`$-basis and using only the diagonal elements of the reduced density matrix, the integrals $`I_\sigma (\mathrm{\Delta }E,\beta ,D)`$ can always be combined such that their real parts cancel each other. When using the nondiagonal elements of $`\rho (t)_{d,d^{}}`$, as well, this is no longer true and the real parts of $`I_\sigma (\mathrm{\Delta }E,\beta ,D)`$ give rise to $`D`$-dependent shifts in the energies $`E_d`$. In this work, we have used $`D`$’s of the order of the anisotropy barrier, i.e., 50-100K, and found that this yields tiny shifts also in the relaxation rate curves but leaves the qualitative picture unchanged. In the figures in the section IV, $`D`$ is chosen as low as 25K in order to simplify the comparison between the calculations with and without the nondiagonal elements.
With all the contributions to $`\mathrm{\Sigma }`$ written down, we can find an estimate for the order of magnitude of the elements (the individual rates) $`\mathrm{\Sigma }_{mm_1,m^{}m_1^{}}`$. The most interesting piece of information for each state is the largest rate coupling that state to other states – this rate plays a key role in justifying the neglect of the nondiagonal states in $`\rho (t)`$, see App. C, as well as in the suppression of the narrow resonances found in the text.
The prefactor in $`\mathrm{\Gamma }(\omega )`$, cf. Eq. (B13), amounts to $`7.010^5c^5\mathrm{s}^5\mathrm{m}^5\mathrm{K}^2`$, where the sound velocity $`c`$ is expressed in meters per second. The units are chosen such that, when $`\omega `$ is expressed in kelvin, also $`\mathrm{\Gamma }(\omega )`$ is given in kelvin. For $`\omega >0`$ (and also $`\omega >k_\mathrm{B}T`$), i.e., for transitions related with phonon absorption, the energy-dependent part of $`\mathrm{\Gamma }(\omega )`$ strongly decreases for increasing $`\omega `$; for $`\omega <0`$, corresponding to phonon emission, $`\mathrm{\Gamma }(\omega )`$ approaches the temperature-independent power-law dependence $`\omega ^3`$. The largest $`\mathrm{\Gamma }(\omega )`$’s are attained for these latter processes in connection with low-energy spin states. The contribution from the spin operators, cf. Eq.(B8), on the other hand, is larger for spin states closest to the top of the barrier and tends to balance the changes in $`\mathrm{\Gamma }(\omega )`$ and reduce the variations in $`\mathrm{\Sigma }_{mm_1,m^{}m_1^{}}`$ for different states and for varying $`H_z`$. The typical energy scale arising from the spin-phonon rates is found to be $`10^510^4`$K.
## C Choice of basis
In sections III A and III B, we decidedly formulated the more general equations independent of the chosen basis for $`_S`$. In this appendix, we consider the eigenbases of $`_S`$ or the $`d`$-basis, in more detail.
The (strong) tunneling poses a problem for the diagrammatic formulation in the $`m`$-basis, but this can be easily solved by first diagonalizing $`_S`$ and then expressing all the equations in its eigenbasis. In this $`d`$-basis, the kinetic equation for the diagonal and off-diagonal density matrix elements reads
$`\dot{\rho }(t)_{d,d}`$ $`=`$ $`{\displaystyle \underset{d_1,d_1^{}}{}}\mathrm{\Sigma }_{dd,d_1d_1^{}}\rho (t)_{d_1,d_1^{}}`$ (C1)
and
$`\dot{\rho }(t)_{d,d^{}}`$ $`=`$ $`i(E_dE_d^{})\rho (t)_{d,d^{}}+{\displaystyle \underset{d_1,d_1^{}}{}}\mathrm{\Sigma }_{dd^{},d_1d_1^{}}\rho (t)_{d_1,d_1^{}},`$ (C2)
respectively. From the knowledge of (the full) $`\rho (t)`$ we can again obtain, e.g., the magnetization
$`M(t)`$ $`=`$ $`g\mu _\mathrm{B}{\displaystyle \underset{d,d^{}}{}}{\displaystyle \underset{m}{}}d^{}|mmm|d\rho _{d,d^{}}(t)`$ (C3)
$``$ $`g\mu _\mathrm{B}{\displaystyle \underset{d,d^{}}{}}m_{d,d^{}}\rho _{d,d^{}}(t)`$ (C4)
where $`m_{d,d^{}}`$ is defined in this way as the matrix element of $`S_z`$ in the $`d`$-basis.
When the tunneling rates dominate the kinetic equations, the main features of the eigenstates can be understood in even simpler terms as follows. When a given state is off-resonant, there is essentially a one-to-one correspondence between each of the $`m`$\- and $`d`$-states, i.e., also the $`d`$-states are localized on one or the other side of the barrier. Close to a resonance, two states $`m_l`$ and $`m_r`$ on different sides of the barrier, see Fig. 22, get coupled and form an approximate two-state system described by
$`_2=\left(\begin{array}{cccc}E_{m_l}& \mathrm{\Delta }\hfill & & \\ \mathrm{\Delta }^{}& E_{m_r}\hfill & & \end{array}\right).`$ (C7)
The nondiagonal elements denote the tunnel splitting as obtained from the diagonalization of the full spin Hamiltonian; the subscripts stand for the left and right sides of the barrier. The eigensolutions to this are the symmetric and antisymmetric combinations of the respective $`m`$-states
$`|d_s`$ $`=`$ $`\alpha |m_l+\beta |m_r`$ (C8)
$`|d_a`$ $`=`$ $`\beta |m_l\alpha |m_r`$ (C9)
that extend through the barrier, see Fig. 22. The factors $`\alpha `$ and $`\beta `$ are the normalized constants
$`\alpha `$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }}{\sqrt{\widehat{\epsilon }^2+|\mathrm{\Delta }|^2}}}`$ (C10)
$`\beta `$ $`=`$ $`{\displaystyle \frac{\widehat{\epsilon }}{\sqrt{\widehat{\epsilon }^2+|\mathrm{\Delta }|^2}}}`$ (C11)
with $`\widehat{\epsilon }=\frac{1}{2}[(E_lE_r)\sqrt{(E_lE_r)^2+4|\mathrm{\Delta }|^2}]`$.
The biggest simplification is attained when we argue that, for the most values of $`H_z`$, we can restrict our considerations to the diagonal elements of the density matrix. A naive justification for this concerns the stationary values of the density matrix elements \[obtained by requiring $`\dot{\rho }(t)_{d,d^{}}=0`$\]. This leads to the immediate conclusion that all the off-diagonal elements between nonresonant states are negligibly small. Furthermore, the nondiagonal elements are also very small for any pair of resonant states as long as the tunnel splitting of that particular resonance is larger than the spin-phonon rates coupling these states to others, see the end of App. B.
We also investigated the temporal behaviour of the off-diagonal elements in terms of the reduced model shown in Fig. 22 and the results lend support to the above conclusions. The idea of this simulation was to prepare the system into the state $`d_i`$ at the initial time $`t_0`$, let the system then evolve in time according to the kinetic equation, and see how the off-diagonal elements $`\rho (t)_{d_s,d_a}`$ and $`\rho (t)_{d_a,d_s}`$ behave. The resonant pair of states in the figure is similar to the one in Eqs. (C8) and (C9) and it is coupled to two lower, nonresonant states $`d_i`$ and $`d_f`$. The rates depicted in the figure are $`\mathrm{\Sigma }_u=\mathrm{\Sigma }_{ll,ii}`$ and $`\mathrm{\Sigma }_d(\mathrm{\Sigma }_{ii,ll}+\mathrm{\Sigma }_{ff,rr})/2`$. The magnitudes of these rates – as compared to the tunnel splitting $`|2\mathrm{\Delta }|`$ – determine two regimes. If $`2|\mathrm{\Delta }|\mathrm{\Sigma }_\mathrm{d}`$, the amplitudes of the nondiagonal elements are found to quickly reach their maxima $`\mathrm{\Sigma }_\mathrm{u}/|2\mathrm{\Delta }|`$ and their values orbit around and “decay” towards the respective complex stationary values. On the other hand, according to the detailed-balance relation, the stationary values of the diagonal elements are proportional to $`\mathrm{\Sigma }_\mathrm{u}/\mathrm{\Sigma }_\mathrm{d}`$. Hence $`\rho (t)_{d,d^{}}/\rho (t)_{d^{()}}\mathrm{\Sigma }_\mathrm{d}/|2\mathrm{\Delta }|`$ and we can neglect the nondiagonal elements, if $`\mathrm{\Sigma }_\mathrm{d}2|\mathrm{\Delta }|`$. In this case the kinetic equation, Eqs. (C1) and (C2), becomes very simple: Eq. (C2) can be neglected and the rate $`\mathrm{\Sigma }`$ acquires the form, cf. App. B,
$`\mathrm{\Sigma }_{d^{}d^{},dd}=\pm G_{dd^{},d^{}d}{\displaystyle \frac{\pi A^2}{12\rho c^5\mathrm{}^4}}{\displaystyle \frac{\mathrm{\Delta }E^3}{e^{\beta \mathrm{\Delta }E}1}},`$ (C12)
where
$`G_{dd^{},d^{}d}`$ $`=`$ $`{\displaystyle \underset{m_1,m_2}{}}{\displaystyle \underset{m_3,m_4}{}}G_{m_4m_3,m_2m_1}`$ (C14)
$`d^{}|m_4m_3|dd|m_2m_1|d^{}`$
corresponds to the vertices and contains the spin-phonon coupling constants, see App. A. In the opposite case, $`2|\mathrm{\Delta }|\mathrm{\Sigma }_\mathrm{d}`$, the nondiagonal elements do not perform orbiting motion in the complex plane but increase motonously to roughly one half of $`\rho _{d^{()}}^{(0)}`$. In this case, the off-diagonal elements clearly cannot be neglected.
For Mn<sub>12</sub> we can attain the whole range of cases: for the most strongly coupled level(s) $`2|\mathrm{\Delta }|\mathrm{\Sigma }_\mathrm{d}`$, while for the lower levels, $`2|\mathrm{\Delta }|\mathrm{\Sigma }_\mathrm{d}`$. In the former case, the diagonal elements $`\rho _d(t)`$ are sufficient in describing the system whereas, in the latter case, we either have to include the nondiagonal states or restrict our considerations to magnetic fields for which $`\mathrm{\Delta }E\mathrm{\Sigma }_\mathrm{d}`$ for all the levels, cf. Ref. . In the text, we neglect the nondiagonal elements in the calculation of $`\chi (\omega )`$ and in some of the analytical considerations but compare the two cases in section IV.
## D Laplace transformation
In this appendix, we consider the Laplace tranformation
$`f(z){\displaystyle _0^{\mathrm{}}}𝑑te^{izt}f(t)`$ (D1)
of the kinetic equation, Eq. (20). We also give another proof of the applicability of the Markov approximation in calculating the relaxation rates.
The kinetic equation is readily transformed into
$`iz\rho (z)\rho (t`$ $`=`$ $`0)=iL_0\rho (z)+\mathrm{\Sigma }(z)\rho (z)`$ (D2)
$`\rho (z)`$ $`=`$ $`{\displaystyle \frac{\rho (t=0)}{iz+iL_0\mathrm{\Sigma }(z)}}.`$ (D3)
The poles of Eq. (D3), i.e., the solutions of $`iz_i+iL_0\mathrm{\Sigma }(z_i)=0`$, yield the exact eigenvalues to the kinetic equation: $`z_i=\omega _i+i/\tau _i`$. For the slowest mode of the time evolution, one can consider the expansion
$`\mathrm{\Sigma }(z_1)\mathrm{\Sigma }(0)+z_1{\displaystyle \frac{\mathrm{\Sigma }(z)}{z}}|_{z=0}+\mathrm{}`$ (D4)
The prefactor of the $`z_1`$ may be evaluated to be proportional to $`[\tau _1\mathrm{min}\{\mathrm{\Delta }E,k_\mathrm{B}T,D\}]^1`$, i.e., to the maximal ratio between the over-barrier relaxation rate $`1/\tau _1`$ and the other characteristic energy scales in the problem: level spacing and/or splitting $`\mathrm{\Delta }E`$, temperature $`k_\mathrm{B}T`$, and the cutoff of the phonon spectrum $`D`$ (the cutoff $`D`$ is introduced in order to assure that the real part of Eq. (A14) is convergent also when we consider the nondiagonal density matrix elements). It turns out that the actual relaxation rates are several orders of magnitude smaller than any other energy scale and it becomes safe to approximate
$`\rho (z){\displaystyle \frac{\rho (t=0)}{iz+iL_0\mathrm{\Sigma }(0)}}={\displaystyle \frac{\rho (t=0)}{izW}}`$ (D5)
where $`\mathrm{\Sigma }(0)`$ has been identified as the constant $`\mathrm{\Sigma }`$ of the Markov approximation above and $`W`$ is defined accordingly, cf. Eq. (21). This approximation is valid for the relaxation mode and time $`\tau _1`$
$`\rho ^{(1)}(z)={\displaystyle \frac{\rho ^{(1)}(t=0)}{iz1/\tau _1}}`$ (D6)
but the Markov approximation may give erroneous results for the faster eigenmodes for which Eq. (D5) no longer holds true.
## E Lorentzian peak shapes
The series of peaks found in the relaxation rates/times, cf. Fig. 9, may be understood in terms of different relaxation paths, each path with a possible tunneling channel giving rise to a peak – see for nice illustrations of the paths. In this appendix, we sketch a derivation that aims to show that the Lorentzian peak shapes are actually something to be expected.
When the spin system is somehow disturbed away from equilibrium and then let relax, it quickly acquires a metastable state, a thermal equilibrium separately on each side of the barrier. This initial thermalization into the metastable state is driven by the spin-phonon interaction that can change the spin states $`m`$ by $`\pm 1`$ or $`\pm 2`$. At a much longer time scale, the system relaxes over the barrier towards the real equilibrium configuration. For the relaxation to take place, the crucial step is the final transition that transfers the spin onto the other side of the barrier. We can distinguish two regimes in terms of how this critical transition takes place. In absence of tunneling, e.g., in off-resonance conditions, the relaxation is only possible over the top of the barrier, while for relatively strong tunnel splitting and for resonant conditions, the dominant path is via tunneling across the barrier well below its top. When the tunneling is weak compared with the spin-phonon interaction, the tunneling rate is the bottle neck for the relaxation to take place. In Mn<sub>12</sub>, this is the case for tunneling between the low-lying states with $`|m|>4`$ (for $`H_z0`$T). However, for the experimentally relevant resonances, the tunneling is strong and takes place between the higher states. In this case, the spin actually oscillates back and forth through the barrier until it relaxes to some lower state on either side of the barrier. This is the case of interest here.
In the strong-tunneling regime, the system is best described in terms of the $`d`$-basis where it suffices to consider the diagonal elements of the density matrix. In order to get a more intuitive picture of the relaxation, let us consider a situation where the system has been prepared onto one side of the barrier and has reached the metastable thermal equilibrium there. This initial condition is convenient for two purposes: first, the relaxation only proceeds into one direction and, second, the phonon-induced transitions on this one side of the barrier are accounted for by the thermal probabilities $`\stackrel{~}{\rho }_d`$ (tilde denotes the metastable state and we write just one index for the diagonal matrix elements). Let us further consider relaxation via a single tunneling resonance and take into account the states and transition processes illustrated in Fig. 23. The system starts in the initial state $`d_i`$ (localized onto the left side of the barrier, $`m_i`$), is then activated onto resonance into either the symmetric or antisymmetric state, denoted by $`d_s`$ and $`d_a`$, respectively, and at some point is transfered down to the final state $`d_f`$ (localized onto the right side of the barrier, $`m_f`$). The subsequent intravalley relaxation is so fast that after the transition to $`d_f`$ the relaxation can be considered complete. The states $`d_s`$ and $`d_a`$ extend through the barrier and this is the key point of the present discussion: the spin is transferred through the barrier in a single step, the rate being determined by thermal activation but also by the magnetic-field dependent amplitudes $`\alpha `$ and $`\beta `$, from Eqs. (C8) and (C9), for the two extended states to be on either side of the barrier. These amplitudes determine the relative probabilities for the activation process to couple to the resonant states.
The above discussion can be formulated in the language of a master equation:
$`\dot{\stackrel{~}{\rho }}_{d_s}(t)=0`$ $``$ $`\mathrm{\Sigma }_{d_s,d_i}\stackrel{~}{\rho }_{d_i}(t)(\mathrm{\Sigma }_{d_f,d_s}+\mathrm{\Sigma }_{d_i,d_s})\stackrel{~}{\rho }_{d_s}(t)`$ (E5)
$`|d_s|m_l|^2\mathrm{\Sigma }_{m_l,m_i}\stackrel{~}{\rho }_{d_i}(t)`$
$`(|d_s|m_r|^2\mathrm{\Sigma }_{m_f,m_r}+|d_s|m_l|^2\mathrm{\Sigma }_{m_i,m_l})\stackrel{~}{\rho }_{d_s}(t)`$
$`=|\alpha |^2\mathrm{\Sigma }_{m_l,m_i}\stackrel{~}{\rho }_{d_i}(t)`$
$`(|\beta |^2\mathrm{\Sigma }_{m_f,m_r}+|\alpha |^2\mathrm{\Sigma }_{m_i,m_l})\stackrel{~}{\rho }_{d_s}(t)`$
$`\dot{\stackrel{~}{\rho }}_{d_a}(t)=0`$ $``$ $`\mathrm{\Sigma }_{d_a,d_i}\stackrel{~}{\rho }_{d_i}(t)(\mathrm{\Sigma }_{d_f,d_a}+\mathrm{\Sigma }_{d_i,d_a})\stackrel{~}{\rho }_{d_a}(t)`$ (E10)
$`|d_a|m_l|^2\mathrm{\Sigma }_{m_l,m_i}\stackrel{~}{\rho }_{d_i}(t)`$
$`(|d_a|m_r|^2\mathrm{\Sigma }_{m_f,m_r}+|d_a|m_l|^2\mathrm{\Sigma }_{m_i,m_l})\stackrel{~}{\rho }_{d_a}(t)`$
$`=|\beta |^2\mathrm{\Sigma }_{m_l,m_i}\stackrel{~}{\rho }_{d_i}(t)`$
$`(|\alpha |^2\mathrm{\Sigma }_{m_f,m_r}+|\beta |^2\mathrm{\Sigma }_{m_i,m_l})\stackrel{~}{\rho }_{d_a}(t).`$
The approximate equalities are just a reminder that we have neglected, e.g., the contributions from states above the resonance as well as the return possibility from state $`d_f`$. These equations can be simplified by the assumption $`\mathrm{\Sigma }_{m_i,m_l}\mathrm{\Sigma }_{m_f,m_r}`$ which is reasonable for a pair of resonant states – as a result, the $`\mathrm{\Sigma }`$’s can be taken out of the brackets in the last forms of the above formulas. By further noting that due to normalization $`|\alpha |^2+|\beta |^2=1`$ and that the resulting probabilities are time independent, we obtain
$`\stackrel{~}{\rho }_{d_s}`$ $``$ $`|\alpha |^2{\displaystyle \frac{\mathrm{\Sigma }_{m_l,m_i}}{\mathrm{\Sigma }_{m_i,m_l}}}\stackrel{~}{\rho }_{d_i}`$ (E11)
$`\stackrel{~}{\rho }_{d_a}`$ $``$ $`|\beta |^2{\displaystyle \frac{\mathrm{\Sigma }_{m_l,m_i}}{\mathrm{\Sigma }_{m_i,m_l}}}\stackrel{~}{\rho }_{d_i}.`$ (E12)
The ratio of the $`\mathrm{\Sigma }`$’s is just the thermal factor $`\mathrm{exp}[\beta (E_lE_i)]`$, cf. detailed balance, and $`\stackrel{~}{\rho }_{d_i}`$ is the thermal probability to be in a state with energy $`E_i`$ over the lowest energy $`E_{10}`$ on left hand side (for $`H_z>0`$). Together these yield a factor $`c\mathrm{exp}[\beta (E_lE_{10})]`$, where $`c`$ is a normalization constant equal to $`\stackrel{~}{\rho }_{10}`$ which is close to unity for the temperatures of interest.
In the next and final step, the relaxation rate is obtained from the knowledge of these probabilities and the rates to be dragged down on the right hand side of the barrier, i.e.,
$`\tau ^1`$ $``$ $`\mathrm{\Sigma }_{d_f,d_s}\stackrel{~}{\rho }_{d_s}+\mathrm{\Sigma }_{d_f,d_a}\stackrel{~}{\rho }_{d_a}`$ (E13)
$``$ $`|\beta |^2\mathrm{\Sigma }_{m_f,m_r}\stackrel{~}{\rho }_{d_s}+|\alpha |^2\mathrm{\Sigma }_{m_f,m_r}\stackrel{~}{\rho }_{d_a}`$ (E14)
$``$ $`2|\alpha |^2|\beta |^2\mathrm{\Sigma }_{m_f,m_r}ce^{\beta (E_iE_{10})}.`$ (E15)
The exponential factor is just the effective Arrhenius factor seen in experiments, $`\frac{1}{2}c\mathrm{\Sigma }_{m_f,m_r}=\tau _0^1`$, and
$`4|\alpha |^2|\beta |^2={\displaystyle \frac{(2|\mathrm{\Delta }|)^2}{(E_{m_l}E_{m_r})^2+(2|\mathrm{\Delta }|)^2}}.`$ (E16)
Here $`2|\mathrm{\Delta }|`$ is the tunnel splitting. It is more or less independent of $`H_z`$ but $`\xi E_{m_l}E_{m_r}`$ can be tuned with the magnetic field. In terms of the field, the width of the resonant peak at its half maximum is
$`\delta H_z={\displaystyle \frac{4|\mathrm{\Delta }|}{g\mu _\mathrm{B}|m_lm_r|}}.`$ (E17)
This sketch of a derivation introduces all the factors seen in experiments: the Arrhenius law with a reasonable prefactor $`\tau _0^1`$, that depends weakly on temperature and the particular resonance, see the two last paragraphs of App. B, and peaks of accelerated relaxation superimposed on it. The peak heights or the relaxation rates on resonance are found to correspond to the Boltzmann or Arrhenius factor with the energy corresponding to the effective barrier height. The peak shape is Lorentzian as observed in experiment with widths given by precisely the tunnel splittings.
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# 1 Introduction
## 1 Introduction
Supersymmetry (SUSY) has been considered as an attractive candidate for physics beyond the standard model. Many mechanisms have been proposed for SUSY breaking and its transmission to our sector. Among them, the gauge-mediated SUSY breaking (GMSB) models are fascinating since they beautifully solve the problem of flavor changing processes inherent in the SUSY standard model. Moreover, GMSB models are determined by a few parameters, and thus have high predictability. In general, they predict that the gravitino is the lightest SUSY particle (LSP) and the next-to-lightest SUSY particle (NLSP) is either the lightest neutralino (almost purely bino $`\stackrel{~}{B}`$) or the lighter stau $`\stackrel{~}{\tau }_1`$ . This mass spectrum gives striking features to the low energy phenomenology. Here we focus on the cosmological aspects of GMSB models.
The light stable gravitino<sup>1</sup><sup>1</sup>1 We assume here that the $`R`$-parity is exact and hence the LSP gravitino is stable.gives rise to a serious cosmological problem, known as the “gravitino problem” . The energy density of the gravitinos which are produced in the early universe may exceed the present critical density if the gravitino mass is $`m_{3/2}^{}{}_{}{}^{>}`$ 1 keV . Even if a primordial inflation is assumed, this problem cannot be solved completely since gravitinos are reproduced after the inflation ends. At the reheating epoch, gravitinos are produced by scatterings in the thermal bath. It has been pointed out in Ref. that gravitinos might also be produced non-thermally at the preheating epoch. Although this non-thermal mechanism might dominate the thermal production, we will not consider it, since it highly depends on inflation models and we would like to have a general discussion.
Stable gravitinos which are produced thermally do not overclose the universe if the reheating temperature of the inflation $`T_R`$ is low enough. The upper bound on $`T_R`$ is given by
$$T_{R}^{}{}_{}{}^{<}\{\begin{array}{cc}100\mathrm{GeV}\mathrm{to}\mathrm{\hspace{0.33em}1}\mathrm{TeV}\hfill & \mathrm{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}\mathrm{keV}_{}^<m_{3/2}^{}{}_{}{}^{<}100\mathrm{keV}\hfill \\ 10^8\mathrm{GeV}\times \left(\frac{m_{3/2}}{1\mathrm{GeV}}\right)\left(\frac{m_{\stackrel{~}{B}}}{100\mathrm{GeV}}\right)^2\hfill & \mathrm{for}m_{3/2}^{}{}_{}{}^{>}100\mathrm{keV}\hfill \end{array},$$
(1)
where $`m_{\stackrel{~}{B}}`$ denotes the bino mass. Since this upper bound on $`T_R`$ is so severe for a lighter gravitino mass region<sup>2</sup><sup>2</sup>2 Of course, there is no gravitino problem for $`m_{3/2}`$ $`_{}^<`$ 1 keV ., inflation models which predict high reheating temperatures are allowed only for a larger $`m_{3/2}`$ region.
However, when one considers a heavy gravitino (say, $`m_{3/2}`$ $`_{}^>`$ 100 MeV) in GMSB, there is another cosmological difficulty associated with the NLSP decay. The NLSP decays into the LSP gravitino only through gravitational interaction and its lifetime will be comparable to the big-bang nucleosynthesis (BBN) era ($`t`$1–$`10^2`$ sec). The decay of the NLSP during or after the BBN is dangerous, since the decay products might alter the abundances of the light elements, which spoils the success of the BBN . To avoid this difficulty, the lifetime and the energy density of the NLSP are severely restricted .
In this letter, we investigate cosmological difficulties associated with the stable gravitino and derive the most conservative upper bound on the reheating temperature. Especially, we consider in detail the cosmological consequences of the NLSP decay into the gravitino around the BBN epoch. Similar analysis had been made in Ref. , in which the BBN constraint is obtained by considering mainly the effects of hadrons produced in the NLSP decay. However, the BBN with high energy hadron injection, examined in Refs. , has not completely been settled yet, since there exist some uncertainties<sup>3</sup><sup>3</sup>3 For example, there exist uncertainties of the experimental data of the hadron scattering processes, and also of the statistical treatment of the errors. . Here, since we would like to obtain the most conservative constraint on the reheating temperature, we will focus on the BBN constraint coming from the photo-dissociation of the light elements by the NLSP decay. Furthermore, we consider the case when the lighter stau $`\stackrel{~}{\tau }_1`$ is the NLSP. This is because, compared to the $`\stackrel{~}{B}`$ NLSP, the $`\stackrel{~}{\tau }_1`$ NLSP has larger annihilation cross sections and so its abundance when it decays is smaller, which results in a weaker constraint on the reheating temperature.
## 2 Estimation of $`\stackrel{~}{\tau }_1`$ NLSP abundance
First of all, we discuss cosmological evolution of the $`\stackrel{~}{\tau }_1`$ NLSP and estimate its abundance when it decays, since the abundance is crucial for obtaining BBN constraints. Let us start from a brief thermal history of $`\stackrel{~}{\tau }_1`$. By the time when the cosmic temperature is comparable to the $`\stackrel{~}{\tau }_1`$ mass ($`Tm_{\stackrel{~}{\tau }_1}`$), only the $`\stackrel{~}{\tau }_1`$ NLSP among the SUSY particles is in thermal equilibrium with the standard model particles. When $`\stackrel{~}{\tau }_1`$’s become non-relativistic for $`T`$ $`_{}^<`$ $`m_{\stackrel{~}{\tau }_1}`$, they pair-annihilate into standard model particles and the abundance decreases exponentially. At $`TT_f`$ ($`T_f`$: the freeze-out temperature of $`\stackrel{~}{\tau }_1`$<sup>4</sup><sup>4</sup>4 In our case, we find $`T_f/m_{\stackrel{~}{\tau }_1}1/28`$$`1/25`$ by numerical calculation.), $`\stackrel{~}{\tau }_1`$ decouples from the thermal bath and its abundance freezes out. Finally, when the Hubble parameter $`H`$ becomes comparable to the decay rate, $`\stackrel{~}{\tau }_1`$ decays into the gravitino. The evolution of $`\stackrel{~}{\tau }_1`$ described above can be traced by solving the Boltzmann equation.
Before writing down the equation, we must take care of the following two facts. The first one is that the masses of charged sleptons $`\stackrel{~}{\tau }_1`$, $`\stackrel{~}{\mu }_1`$, and $`\stackrel{~}{e}_1`$ (the subscript 1 denotes the lighter mass eigenstates) are almost degenerate in GMSB models. This is because they are mostly right-handed, and hence receive at the messenger scale the same soft masses which are determined by the gauge quantum numbers. The mass differences between them at the weak scale are induced from the renormalization group effects due to the leptonic Yukawa couplings $`y_i`$ ($`i=e,\mu ,\tau `$) and also from the left-right mixings through $`y_i`$. Considering the observed masses of charged leptons, we safely neglect $`y_e`$ and $`y_\mu `$, and hence the masses for $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$, $`m_{\stackrel{~}{\mu }_1}`$ and $`m_{\stackrel{~}{e}_1}`$, are the same. On the other hand, the tau Yukawa coupling $`y_\tau `$ gives negative contribution to the stau mass when it is evolved from the messenger scale to the weak scale, and the left-right mixing between the staus also decreases $`m_{\stackrel{~}{\tau }_1}`$. Both effects lead to the fact that $`m_{\stackrel{~}{\tau }_1}`$ is always smaller than $`m_{\stackrel{~}{\mu }_1(\stackrel{~}{e}_1)}`$. This mass splitting $`\mathrm{\Delta }m`$ $``$ $`m_{\stackrel{~}{\mu }_1(\stackrel{~}{e}_1)}m_{\stackrel{~}{\tau }_1}`$ plays a crucial role in calculating the $`\stackrel{~}{\tau }_1`$ abundance. If $`\mathrm{\Delta }m`$ is very small, the relic abundance of $`\stackrel{~}{\tau }_1`$ is determined by not only its annihilation processes, but also by the annihilation with $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$, i.e., we should include the coannihilation effects . We have verified numerically that these coannihilation effects become significant for $`\mathrm{\Delta }m`$ $`_{}^<`$ $`T_f`$. In GMSB models, both $`\mathrm{\Delta }m`$ $`_{}^<`$ $`T_f`$ and $`\mathrm{\Delta }m`$ $`_{}^>`$ $`T_f`$ are allowed<sup>5</sup><sup>5</sup>5 The mass difference $`\mathrm{\Delta }m`$ becomes larger as $`\mathrm{tan}\beta `$ increases, where $`\mathrm{tan}\beta v_u/v_d`$ ($`v_u`$ and $`v_d`$ denote the vacuum expectation values of Higgs fields which couple to up-type and down-type quarks, respectively). , so we respect both cases.
Secondly, if the bino mass $`m_{\stackrel{~}{B}}`$ is close to $`m_{\stackrel{~}{\tau }_1}`$, we should also take into account the coannihilation of $`\stackrel{~}{\tau }_1`$ with $`\stackrel{~}{B}`$, as well as with $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$. However, this effect increases the abundance of $`\stackrel{~}{\tau }_1`$, which results in more stringent upper bound on $`T_R`$ (see the following discussions). Thus, we forbid the coannihilation of $`\stackrel{~}{\tau }_1`$ with $`\stackrel{~}{B}`$ and restrict ourselves in the parameter region
$$m_{\stackrel{~}{B}}^{}{}_{}{}^{>}m_{\stackrel{~}{\tau }_1}+T_f.$$
(2)
In the actual calculation below, we take the lower bound on $`m_{\stackrel{~}{B}}`$ as $`m_{\stackrel{~}{B}}1.1m_{\stackrel{~}{\tau }_1}`$ to be conservative.
Now we are at the point to present the Boltzmann equation including the coannihilation effects. This equation describes the evolution of the total number density of charged sleptons $`n=n_{\stackrel{~}{\tau }_1}+n_{\stackrel{~}{\tau }_1^{}}+n_{\stackrel{~}{\mu }_1}+n_{\stackrel{~}{\mu }_1^{}}+n_{\stackrel{~}{e}_1}+n_{\stackrel{~}{e}_1^{}}`$, where $`n_i`$ stands for the number density of one slepton species $`i`$<sup>6</sup><sup>6</sup>6 We distinguish particles from their anti-particles. ,
$$\frac{dn}{dt}=3Hn\sigma _{\mathrm{eff}}v\left[n^2(n^{\mathrm{eq}})^2\right],$$
(3)
where
$$\sigma _{\mathrm{eff}}v=\underset{i,j}{}\sigma _{ij}v\frac{n_i^{\mathrm{eq}}}{n^{\mathrm{eq}}}\frac{n_j^{\mathrm{eq}}}{n^{\mathrm{eq}}}.$$
(4)
Here $`n^{\mathrm{eq}}`$ ($`n_i^{\mathrm{eq}}`$) is the equilibrium value of $`n`$ ($`n_i`$) and $`v`$ is the relative velocity of particles $`i`$ and $`j`$. The bracket denotes the thermal average and $`\sigma _{ij}`$ is the total annihilation cross section of $`i+jX+X^{}`$:
$$\sigma _{ij}=\underset{X,X^{}}{}\sigma (i+jX+X^{}),$$
(5)
where $`X`$ and $`X^{}`$ represent the possible standard model particles. Note that the Boltzmann equation (3) also holds for the case where the coannihilations of $`\stackrel{~}{\tau }_1`$ with $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$ do not occur.
Next we discuss the thermal-averaged cross sections $`\sigma _{\mathrm{eff}}v`$ in Eq. (3). The thermal-averaged cross sections $`\sigma _{ij}v`$ in Eq. (4) can be expanded in terms of $`T/m_{\stackrel{~}{\tau }_1}`$ \[see Eqs. (25) – (27) below\]. Since the final abundance is determined by $`\sigma _{\mathrm{eff}}v`$ at $`TT_f`$ and the freeze out temperature is typically $`T_fm_{\stackrel{~}{\tau }_1}/25`$, it is sufficient only to consider the leading term, i.e., the $`s`$-wave cross sections of the sleptons. Note that the $`s`$-wave component of the thermal-averaged cross section $`\sigma _{ij}v`$ is equal to that of $`\sigma _{ij}`$. The relevant (co)annihilation channels are
$`(\mathrm{I})`$ $`\{\begin{array}{c}\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1^{}\hfill \\ \stackrel{~}{\mu }_1\stackrel{~}{\mu }_1^{}\hfill \\ \stackrel{~}{e}_1\stackrel{~}{e}_1^{}\hfill \end{array}\gamma \gamma ,Z\gamma ,ZZ,W^+W^{},f\overline{f},h^0h^0`$ (9)
$`(\mathrm{II})`$ $`\{\begin{array}{c}\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1\tau \tau \hfill \\ \stackrel{~}{\tau }_1^{}\stackrel{~}{\tau }_1^{}\overline{\tau }\overline{\tau }\hfill \\ \stackrel{~}{\mu }_1\stackrel{~}{\mu }_1\mu \mu \hfill \\ \stackrel{~}{\mu }_1^{}\stackrel{~}{\mu }_1^{}\overline{\mu }\overline{\mu }\hfill \\ \stackrel{~}{e}_1\stackrel{~}{e}_1ee\hfill \\ \stackrel{~}{e}_1^{}\stackrel{~}{e}_1^{}\overline{e}\overline{e}\hfill \end{array}`$ (16)
$`(\mathrm{III})`$ $`\{\begin{array}{c}\stackrel{~}{\tau }_1\stackrel{~}{\mu }_1\tau \mu \hfill \\ \stackrel{~}{\tau }_1^{}\stackrel{~}{\mu }_1^{}\overline{\tau }\overline{\mu }\hfill \\ \stackrel{~}{\tau }_1\stackrel{~}{e}_1\tau e\hfill \\ \stackrel{~}{\tau }_1^{}\stackrel{~}{e}_1^{}\overline{\tau }\overline{e}\hfill \\ \stackrel{~}{\mu }_1\stackrel{~}{e}_1\mu e\hfill \\ \stackrel{~}{\mu }_1^{}\stackrel{~}{e}_1^{}\overline{\mu }\overline{e}\hfill \end{array}`$
$`(\mathrm{IV})`$ $`\{\begin{array}{c}\stackrel{~}{\tau }_1\stackrel{~}{\mu }_1^{}\tau \overline{\mu },\nu _\tau \overline{\nu }_\mu \hfill \\ \stackrel{~}{\tau }_1^{}\stackrel{~}{\mu }_1\overline{\tau }\mu ,\overline{\nu }_\tau \nu _\mu \hfill \\ \stackrel{~}{\tau }_1\stackrel{~}{e}_1^{}\tau \overline{e},\nu _\tau \overline{\nu }_e\hfill \\ \stackrel{~}{\tau }_1^{}\stackrel{~}{e}_1\overline{\tau }e,\overline{\nu }_\tau \nu _e\hfill \\ \stackrel{~}{\mu }_1\stackrel{~}{e}_1^{}\mu \overline{e},\nu _\mu \overline{\nu }_e\hfill \\ \stackrel{~}{\mu }_1^{}\stackrel{~}{e}_1\overline{\mu }e,\overline{\nu }_\mu \nu _e\hfill \end{array}`$
Here $`f`$ denotes the ordinary quarks and leptons.
We calculate cross sections for all the processes listed in (I)-(IV). The processes (I) and (II) represent the annihilation of each slepton species, while (III) and (IV) describe coannihilation processes. In the following we briefly explain their features. First of all, let us estimate the cross sections for processes (I). Although they have many final states, it turns out that most of them give very small contributions to the total cross section. Processes which are induced by the tau Yukawa coupling and/or the left-right mixing of sleptons are suppressed due to their smallness<sup>7</sup><sup>7</sup>7The mixing angles $`\alpha _i`$ of sleptons are in general small in GMSB models ($`\mathrm{sin}\alpha _i1`$). However, in the extremely large $`\mathrm{tan}\beta `$ region, the stau mixing angle $`\alpha _{\stackrel{~}{\tau }}`$ can be large as $`\mathrm{sin}\alpha _{\stackrel{~}{\tau }}𝒪`$(0.1). Although $`\sigma _{\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1^{}}v/\sigma (\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1^{}\gamma \gamma )`$ might become of $`𝒪`$(10) as opposed to Eq. (25) in this case, we find that the final upper bound on $`T_R`$ only increases by less than a factor of 2.. Furthermore, in calculating $`s`$-wave cross sections, derivative couplings of the sleptons vanish. From these considerations, we find that $`\gamma \gamma `$ and $`Z\gamma `$ are the dominant channels<sup>8</sup><sup>8</sup>8 The ratios of the cross sections for the dominant processes $`\gamma \gamma `$$`Z\gamma `$, and $`ZZ`$ are 1.0:0.6:(0.1–0.2). The contributions of the other channels are less than a few percents., and $`\sigma _{\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1^{}}v`$ is approximately given by<sup>9</sup><sup>9</sup>9 Actually, if $`m_{\stackrel{~}{\tau }_1}<m_Z/2`$, the $`Z\gamma `$ channel is forbidden kinematically. Furthermore, if $`m_{\stackrel{~}{\tau }_1}m_{h^0}/2`$ ($`m_{h^0}`$: mass of the lightest Higgs boson $`h^0`$), the annihilation cross section is enhanced by the pole contribution of $`h^0`$. These effects may change the value of $`Y`$ in Fig. 1. However, we simply assume Eq. (25), since the mass region of $`m_{\stackrel{~}{\tau }_1}<`$ 90 GeV is excluded from the current experimental limit.:
$$\sigma _{\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1^{}}v2\sigma (\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1^{}\gamma \gamma )=\frac{4\pi \alpha _{\mathrm{em}}^2}{m_{\stackrel{~}{\tau }_1}^2}+𝒪\left(\frac{T}{m_{\stackrel{~}{\tau }_1}}\right).$$
(25)
Next, we turn to the processes (II), which have only $`t`$-channel and $`u`$-channel diagrams. Diagrams induced by the exchange of Higgsino and neutral Wino can be safely neglected due to the smallness of $`y_\tau `$ and left-right mixing, respectively. Then, there is only $`\stackrel{~}{B}`$ contribution and its cross section is calculated as
$`\sigma _{\stackrel{~}{\tau }_1\stackrel{~}{\tau }_1}v`$ $``$ $`{\displaystyle \frac{16\pi \alpha _{\mathrm{em}}^2m_{\stackrel{~}{B}}^2}{\mathrm{cos}^4\theta _W(m_{\stackrel{~}{\tau }_1}^2+m_{\stackrel{~}{B}}^2)^2}}+𝒪\left({\displaystyle \frac{T}{m_{\stackrel{~}{\tau }_1}}}\right).`$ (26)
The annihilation channels listed in (III) can be calculated in the same way as
$`\sigma _{\stackrel{~}{\tau }_1\stackrel{~}{\mu }_1}v`$ $``$ $`{\displaystyle \frac{8\pi \alpha _{\mathrm{em}}^2m_{\stackrel{~}{B}}^2}{\mathrm{cos}^4\theta _W(m_{\stackrel{~}{\tau }_1}^2+m_{\stackrel{~}{B}}^2)^2}}+𝒪\left({\displaystyle \frac{T}{m_{\stackrel{~}{\tau }_1}}}\right),`$ (27)
for $`m_{\stackrel{~}{\mu }_1}=m_{\stackrel{~}{\tau }_1}`$. Cross sections for the processes (IV) are also dominated by the $`\stackrel{~}{B}`$ exchange. However, the left-right mixings of sleptons are necessary for these processes, and so they give only tiny contributions to the total thermal-averaged cross section $`\sigma _{\mathrm{eff}}v`$. Cross sections for other processes in (I)–(IV) are calculated similarly. From the above arguments, we conclude that all the annihilation cross sections are almost determined by only the masses of the bino and the sleptons.
Finally, including all the cross sections described above, we solve the Boltzmann equation (3) numerically. Here we take $`m_{\stackrel{~}{B}}=1.1m_{\stackrel{~}{\tau }_1}`$ \[see Eq. (2)\]. In this case the (co)annihilation cross sections take their maximal values, which gives the smallest $`\stackrel{~}{\tau }_1`$ abundance. In Fig. 1 the result is shown in terms of $`Y`$, which is defined by $`Yn/s=_in_i/s`$ with the total entropy density of the universe $`s`$.
Here are some comments. We calculate both of the cases with and without coannihilation effects as mentioned before. If the mass difference between $`\stackrel{~}{\mu }_1(\stackrel{~}{e}_1)`$ and $`\stackrel{~}{\tau }_1`$ is large enough ($`\mathrm{\Delta }m_{}^>T_f`$), $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$ decay into $`\stackrel{~}{\tau }_1`$ and disappear from the thermal bath before $`\stackrel{~}{\tau }_1`$ freezes out. Therefore, the coannihilation processes become ineffective. In this case, the abundance becomes $`Y(n_{\stackrel{~}{\tau }_1}+n_{\stackrel{~}{\tau }_1^{}})/s`$ for $`T_{}^<T_f`$, since $`n_{\stackrel{~}{\mu }_1}n_{\stackrel{~}{e}_1}0`$. For $`T_{}^<T_f`$ the number density of $`\stackrel{~}{\tau }_1`$ and the entropy density decrease at the same rate as the universe expands, and $`Y`$ takes a constant value until $`\stackrel{~}{\tau }_1`$ decays. In fact, Fig. 1 shows the value of $`Y`$ just before the decay of $`\stackrel{~}{\tau }_1`$. On the other hand, when $`\mathrm{\Delta }m_{}^<T_f`$, $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$ are still in thermal equilibrium at $`TT_f`$ due to the effect of inverse decays, and thus we should consider the coannihilations. Note that $`Y`$ takes a constant value after sleptons decouple from the thermal bath even in this case. This is because $`Y`$ is an invariant parameter against the expansion of the universe and also because $`n=n_{\stackrel{~}{\tau }_1}+n_{\stackrel{~}{\tau }_1^{}}+n_{\stackrel{~}{\mu }_1}+n_{\stackrel{~}{\mu }_1^{}}+n_{\stackrel{~}{e}_1}+n_{\stackrel{~}{e}_1^{}}`$ does not change by the decays of $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$. Therefore, the final abundance for $`\stackrel{~}{\tau }_1`$ is given by $`Y`$ for both cases with and without coannihilation effects.
As shown in Fig. 1, the final abundance with coannihilation is larger than that without coannihilation. This can be understood as follows: Imagine when the coannihilation cross sections in (III) and (IV) are extremely large. In this case, the final abundance $`Y`$ would be smaller than that without coannihilation. On the other hand, if the coannihilation cross sections are zero, $`Y`$ increase by a factor of 3 because the relevant degrees of freedom is now 6, not 2. In our case, coannihilation cross sections in (III) and (IV) are smaller than those in (I) and (II), and the final abundance with coannihilation is about twice as large as that without coannihilation.
We present in Fig. 1 the results for the two extreme cases $`\mathrm{\Delta }m0`$ and $`\mathrm{\Delta }mT_f`$, for illustration and for comparison. In fact, we find that the final abundance of stau for the case $`0_{}^<\mathrm{\Delta }m_{}^<T_f`$ falls between the two lines in Fig. 1.
Before closing this section, we should comment on the case in which the mass difference between $`\stackrel{~}{\tau }_1`$ and $`\stackrel{~}{\mu }_1(\stackrel{~}{e}_1)`$ is extremely small. If the mass difference is smaller than the tau mass (of course we should include the coannihilation effects in this case), the decay channel $`\stackrel{~}{\mu }_1(\stackrel{~}{e}_1)\stackrel{~}{\tau }_1\overline{\tau }\mu (e)`$ is kinematically forbidden. Then, $`\stackrel{~}{\mu }_1`$ and $`\stackrel{~}{e}_1`$ may dominantly decay into gravitinos<sup>10</sup><sup>10</sup>10 The left-right mixing allows $`\stackrel{~}{\mu }_1`$ to decay through $`\stackrel{~}{\mu }_1\nu _\mu \stackrel{~}{\tau }_1\overline{\nu }_\tau `$. If this is the main decay channel, the following discussion does not change, since $`\stackrel{~}{\mu }_1`$ has already decayed into $`\stackrel{~}{\tau }_1`$ before the $`\stackrel{~}{\tau }_1`$ decays.and so we should consider the effect of the $`\stackrel{~}{\mu }_1(\stackrel{~}{e}_1)`$ decay on the BBN, as well as $`\stackrel{~}{\tau }_1`$. However, we simply neglect this possibility, since the final results do not change much.
## 3 BBN constraint on $`\stackrel{~}{\tau }_1`$ NLSP decay
We are now ready to investigate the cosmological consequence of the decay of the $`\stackrel{~}{\tau }_1`$ NLSP into the gravitino at the BBN epoch. The $`\stackrel{~}{\tau }_1`$ decay rate is estimated, for $`m_{\stackrel{~}{\tau }_1}m_{3/2}`$, as
$$\mathrm{\Gamma }_{\stackrel{~}{\tau }_1}\frac{1}{48\pi }\frac{m_{\stackrel{~}{\tau }_1}^5}{m_{3/2}^2M_{}^2},$$
(28)
with $`M_{}=2.4\times 10^{18}`$ GeV, and the lifetime is given by
$`\tau _{\stackrel{~}{\tau }_1}6\times 10^4\text{sec}\left({\displaystyle \frac{m_{3/2}}{1\text{GeV}}}\right)^2\left({\displaystyle \frac{m_{\stackrel{~}{\tau }_1}}{100\text{GeV}}}\right)^5.`$ (29)
Therefore, $`\stackrel{~}{\tau }_1`$ is found to be a long-lived particle and to decay during or even after the BBN epoch when $`m_{3/2}^{}{}_{}{}^{>}`$ 10 MeV for $`m_{\stackrel{~}{\tau }_1}`$ = 100 GeV, and hence we should seriously consider the effects of its decay on the BBN.
The energetic $`\tau `$’s produced by the $`\stackrel{~}{\tau }_1`$ decay trigger the electro-magnetic (EM) cascade processes and induce high energy photons. These photons may be abundant enough to destroy or overproduce various light elements (D, <sup>3</sup>He, <sup>4</sup>He, etc.) synthesized by the BBN. In order to keep the success of the BBN, the energy density (per the entropy density) of the extra EM particles emitted after the BBN era (i.e., $`t_{}^>10^4`$ sec) is severely constrained . Here it should be noted that not all the energy of tau contribute to these photo-dissociation processes. This is because the tau decays before it causes the EM cascade processes and the produced neutrinos give no effects on the photo-dissociation processes <sup>11</sup><sup>11</sup>11In fact, the high energy neutrinos produced by the tau decay also causes the EM cascade processes by scattering with the background neutrinos. However, this effect is less significant .. By using the result of the recent analysis in Ref. (see Fig.16 in Ref. ), we can obtain an upper bound on
$`{\displaystyle \frac{\rho _{\stackrel{~}{\tau }_1}}{s}}=m_{\stackrel{~}{\tau }_1}Y,`$ (30)
for a given lifetime of $`\stackrel{~}{\tau }_1`$. Here $`\rho _{\stackrel{~}{\tau }_1}`$ denotes the energy density of $`\stackrel{~}{\tau }_1`$. When one fixes the gravitino mass, this bound is translated into the lower bound on $`m_{\stackrel{~}{\tau }_1}`$, since both the abundance and the lifetime of $`\stackrel{~}{\tau }_1`$ are determined by its mass $`m_{\stackrel{~}{\tau }_1}`$. The obtained result is found in Fig. 2.
You can see that this BBN photo-dissociation constraint gives a more stringent lower bound on $`m_{\stackrel{~}{\tau }_1}`$ for $`m_{3/2}^{}{}_{}{}^{>}5`$ GeV, compared to the experimental bound $`m_{\stackrel{~}{\tau }_1}>90`$ GeV. This result will help us to estimate the upper bound on the reheating temperature in the next section.
## 4 Gravitino problem and constraint on $`T_R`$
Finally, we discuss the cosmological gravitino problem and obtain an upper bound on the reheating temperature of inflation. At the reheating epoch after the inflation ends, gravitinos are produced thermally by scatterings with particles in the hot plasma of the universe<sup>12</sup><sup>12</sup>12 As mentioned in Sec. 1, we discard the non-thermal production of gravitinos at the preheating epoch in order to obtain the most conservative result.. The relic abundance of the gravitino is given by
$$\mathrm{\Omega }_{3/2}^{\mathrm{th}}h^20.3\left(\frac{m_{3/2}}{1\text{GeV}}\right)^1\left(\frac{m_{\stackrel{~}{B}}}{100\text{GeV}}\right)^2\left(\frac{T_R}{10^8\text{GeV}}\right),$$
(31)
where $`h`$ is the present Hubble parameter in unit of 100km/sec/Mpc, and $`\mathrm{\Omega }_{3/2}^{\mathrm{th}}=\rho _{3/2}^{\mathrm{th}}/\rho _c`$ ($`\rho _{3/2}^{\mathrm{th}}`$ is the present energy density of the thermally produced gravitinos and $`\rho _c`$ is the critical density of the present universe). It is found from Eq. (31) that the overclosure limit of $`\mathrm{\Omega }_{3/2}^{\mathrm{th}}<1`$ puts an upper bound on $`T_R`$ as shown in Eq. (1). Notice that this upper bound on $`T_R`$, if one fixes the $`\stackrel{~}{B}`$ mass, becomes more severe for the lighter gravitino mass region, and the highest reheating temperature allowed in GMSB models is achieved for a relatively heavy gravitino mass.
Furthermore, it should be noted that gravitinos are also produced by the $`\stackrel{~}{\tau }_1`$ NLSP decays. Because one gravitino is produced per a $`\stackrel{~}{\tau }_1`$ decay, we find that
$$\mathrm{\Omega }_{3/2}^{\mathrm{NLSP}}\frac{\rho _{3/2}^{\mathrm{NLSP}}}{\rho _c}=\frac{m_{3/2}Y}{\rho _c/s_0},$$
(32)
where $`\rho _{3/2}^{\mathrm{NLSP}}`$ is the present energy density of the gravitinos produced by the $`\stackrel{~}{\tau }_1`$ decays, and $`s_0`$ denotes the entropy density of the present universe. In Fig. 2 we also show the upper bound on $`m_{\stackrel{~}{\tau }_1}`$ from the constraint $`\mathrm{\Omega }_{3/2}^{\mathrm{NLSP}}h^2<1`$. It is found that the gravitino mass is bounded from above as $`m_{3/2}^{}{}_{}{}^{<}1`$ TeV when combined with the lower bound on $`m_{\stackrel{~}{\tau }_1}`$ from the BBN photo-dissociation effects.
Therefore, the total relic abundance of gravitinos in GMSB models is now given by $`\mathrm{\Omega }_{3/2}^{\mathrm{tot}}=\mathrm{\Omega }_{3/2}^{\mathrm{th}}+\mathrm{\Omega }_{3/2}^{\mathrm{NLSP}}`$, as long as the decay of the $`\stackrel{~}{\tau }_1`$ NLSP takes place within the age of the universe. From the lower bound on the $`\stackrel{~}{\tau }_1`$ mass obtained in the previous section, we can estimate the most conservative upper bound on the reheating temperature in the following way. Since the abundance $`\mathrm{\Omega }_{3/2}^{\mathrm{th}}`$ depends on the bino mass $`m_{\stackrel{~}{B}}`$ as shown in Eq. (31), the weakest bound on $`T_R`$ is obtained by the lowest value of $`m_{\stackrel{~}{B}}`$. On the other hand, the lower bound on the mass of the $`\stackrel{~}{\tau }_1`$ NLSP represented in Fig. 2 is nothing but the lower bound on $`m_{\stackrel{~}{B}}`$. Now we are working in the parameter region of $`m_{\stackrel{~}{B}}`$ in Eq. (2) in order to forbid the coannihilation of $`\stackrel{~}{\tau }_1`$ with $`\stackrel{~}{B}`$. Therefore, the lowest value of $`m_{\stackrel{~}{B}}`$ coming from the $`m_{\stackrel{~}{\tau }_1}`$ constraint gives the most conservative upper bound on the reheating temperature<sup>13</sup><sup>13</sup>13 In fact, such a value of $`m_{\stackrel{~}{B}}`$ makes the $`\stackrel{~}{\tau }_1`$ abundance smallest. (See discussions in Sec. 2.) . The result is found in Fig. 3, where we take $`h=1`$ for simplicity<sup>14</sup><sup>14</sup>14 The upper bound on $`T_R`$ becomes more stringent for a smaller $`h`$. .
Note that the gravitino mass region of $`m_{3/2}^{}{}_{}{}^{>}1`$ TeV is excluded, since the gravitinos produced by the decay of $`\stackrel{~}{\tau }_1`$ overclose the present universe ($`\mathrm{\Omega }_{3/2}^{\mathrm{NLSP}}>1`$). On the other hand, $`\mathrm{\Omega }_{3/2}^{\mathrm{NLSP}}`$ plays no role in putting an upper bound on $`T_R`$ for a lighter gravitino mass region of $`m_{3/2}^{}{}_{}{}^{<}500`$ GeV.
You can see that the upper bound on $`T_R`$ is almost proportional to $`m_{3/2}`$ for $`m_{3/2}`$ $`_{}^<`$ 5 GeV. In this gravitino mass region, the experimental bound on $`m_{\stackrel{~}{\tau }_1}`$ gives the lowest value of $`m_{\stackrel{~}{B}}`$ (see Fig. 2). The upper bound on $`T_R`$, thus, is the conventional one Eq.(1) with $`m_{\stackrel{~}{B}}100`$ GeV. However, when $`m_{3/2}`$ $``$ 5 GeV–100 GeV, the upper bound on $`T_R`$ takes an almost constant value of $`10^9`$$`10^{10}`$ GeV. In this region the bino mass which gives the highest value of $`T_R`$ is determined by the lower bound on $`m_{\stackrel{~}{\tau }_1}`$ from the BBN photo-dissociation constraint. Since this lower bound on $`m_{\stackrel{~}{\tau }_1}`$ is more stringent than the experimental limit, the upper bound on $`T_R`$ also becomes more stringent than the conventional one. From Fig.2, the lower bound on the stau mass, i.e., the lower bound on the bino mass is almost proportional to $`(m_{3/2})^{1/2}`$ for $`m_{3/2}=`$ 5–100 GeV. Thus, it is found from Eq.(12) that the upper bound on $`T_R`$ becomes almost constant. Therefore, in GMSB models, the reheating temperature can be taken as high as $`T_R10^9`$$`10^{10}`$ GeV for $`m_{3/2}5`$–100 GeV.
## 5 Conclusions and discussions
In this letter, we have considered the cosmological gravitino problem in GMSB models with the $`\stackrel{~}{\tau }_1`$ NLSP. Especially, we have investigated in detail the cosmological consequence of the $`\stackrel{~}{\tau }_1`$ decay soon after the BBN epoch. Since the $`\stackrel{~}{\tau }_1`$ abundance when it decays is crucial for this discussion, we have solved numerically the Boltzmann equations for both cases with and without coannihilation effects, and obtained the final $`\stackrel{~}{\tau }_1`$ abundance. We have found that the BBN constraint on the $`\stackrel{~}{\tau }_1`$ abundance is translated into the lower bound on $`m_{\stackrel{~}{\tau }_1}`$, and that the obtained bound is more stringent for $`m_{3/2}^{}{}_{}{}^{>}5`$ GeV than the current experimental limit $`m_{\stackrel{~}{\tau }_1}>90`$ GeV. This gives some hints for GMSB models. If $`\stackrel{~}{\tau }_1`$ was detected at future collider experiments, the observed $`\stackrel{~}{\tau }_1`$ mass would enable us to set an upper bound on the gravitino mass, and hence on the SUSY breaking scale. It has also been found that the gravitino mass region of $`m_{3/2}^{}{}_{}{}^{>}`$ 1 TeV is excluded (although it might be marginal considering the SUSY flavor problem), since the energy density of gravitinos produced by the decays of $`\stackrel{~}{\tau }_1`$’s exceeds the present critical density.
By using the lower bound on $`m_{\stackrel{~}{\tau }_1}`$, we have obtained an upper bound on the reheating temperature $`T_R`$ of inflation in order to avoid the overclosure problem of the LSP gravitino produced in thermal scatterings and also in the $`\stackrel{~}{\tau }_1`$ decay. What we have found is that the most conservative upper bound on $`T_R`$ in GMSB models is $`T_{R}^{}{}_{}{}^{<}10^9`$$`10^{10}`$ GeV when $`m_{3/2}5`$–100 GeV. This upper bound on $`T_R`$ is weaker than those in the conventional hidden sector SUSY breaking models to solve the cosmological problem of unstable gravitinos ($`T_{R}^{}{}_{}{}^{<}10^6`$ GeV and $`10^8`$ GeV for $`m_{3/2}`$ 100–500 GeV and 500 GeV–1 TeV, respectively ). Therefore, it helps us a lot to build SUSY inflation models without the cosmological gravitino problem.
Such a high reheating temperature is also promising for baryogenesis. The relevant example is the leptogenesis via decays of heavy Majorana neutrinos, which is very attractive from the viewpoint of the observed tiny neutrino masses. The leptogenesis requires a high reheating temperature to generate the observed baryon asymmetry of the universe<sup>15</sup><sup>15</sup>15 In the leptogenesis scenarios where heavy Majorana neutrinos are thermally produced, the reheating temperature of $`T_{R}^{}{}_{}{}^{>}10^{10}`$ GeV is required to induce the sufficient baryon asymmetry. On the other hand, for the case when heavy Majorana neutrinos are non-thermally produced in the inflaton decays, the reheating temperature of $`T_{R}^{}{}_{}{}^{>}10^6`$ GeV is required .. Thus, the upper bound on the reheating temperature $`T_{R}^{}{}_{}{}^{<}10^9`$$`10^{10}`$ GeV obtained in this letter ensures some scenarios of the leptogenesis to work without the gravitino problem.
The results we have found here are almost independent on models of the GMSB, since the abundance of the $`\stackrel{~}{\tau }_1`$ NLSP is fixed by the $`\stackrel{~}{\tau }_1`$ mass. However, there are some loop-holes in which these constraints can be avoided. If one introduces the $`R`$-parity breaking, the NLSP can decay before the BBN (i.e., $`t_{}^<`$ 1 sec) and evade the BBN constraints, and also the LSP gravitino can decay within the age of the universe. Furthermore, if one assumes the late-time entropy production such as the thermal inflation in the thermal history of the universe, the abundances of the NLSP and also the LSP gravitino are diluted away so that we are free from these constraints.
In the present analysis, we have considered only the photo-dissociation effects of the decay of $`\stackrel{~}{\tau }_1`$ on the BBN to make a conservative analysis. As pointed out in Ref. the high energy hadrons produced by the $`\stackrel{~}{\tau }_1`$ decay might also be dangerous. Here we would like to briefly comment on their effects. The high energy hadrons, if they are produced at the time $`t1`$$`10^4`$ sec, delay the freeze-out of the $`p`$-$`n`$ conversion and raise the number ratio of $`n`$ to $`p`$, which leads to the overproduction of D and <sup>4</sup>He . Furthermore, if $`\stackrel{~}{\tau }_1`$ decays at the time $`t_{}^>10^4`$ sec, the produced hadrons destroy the light elements synthesized by the BBN and modify their abundances (e.g., increases the <sup>7</sup>Li abundance) . These effects give the upper bound on the $`\stackrel{~}{\tau }_1`$ abundance. However, there are some uncertainties in the BBN with high energy hadron injection. The relevant hadron scattering cross sections, especially those of Li, have not been observed experimentally in detail. Furthermore, the statistical treatment for the estimation of the errors has not completely settled yet. Therefore, in the present analysis we do not include the BBN constraints from the hadron injection.
## Acknowledgements
We would like to thank T. Yanagida for various suggestions and stimulating discussions and also K. Kohri and Y. Nomura for useful comments. This work was partially supported by the Japan Society for the Promotion of Science (T.A. and K.H.).
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# 1 Introduction
## 1 Introduction
Gauge symmetries have been established as a guiding principle that determines couplings among local fields. The fundamental interactions, electromagnetic, weak and strong interactions, are intermediated by gauge fields, and the couplings between the gauge fields and matter fields and the self-coupling of the gauge fields are determined by the gauge symmetries. In addition, the gauge symmetries lead to the gravitational interaction. Indeed, the gravitational interaction can be formulated within the framework of gauge theory based on a non-compact gauge group.
In supergravity theories, on the other hand, there is a non-trivial interaction which cannot be derived from gauge symmetries alone. In $`𝒩=1`$ supergravity theory (with $`𝒩=1`$ super Yang-Mills theory) in ten dimensions, the Lagrangian has a coupling between Yang-Mills fields and the abelian antisymmetric tensor field in the supergravity multiplet via the Chern-Simons 3-form. A non-trivial coupling was first introduced by de Wit et al. in the abelian case, and generalized to the non-abelian case by Chapline and Manton. Although this coupling is determined by the local supersymmetry in the supergravity theory, its derivation involves much tedious algebra. In this system, the antisymmetric tensor field must obey a deformed transformation rule including the Yang-Mills fields so as to maintain the invariance of the Lagrangian. This transformation rule can be determined uniquely, but only in a heuristic manner. The transformation rule plays an important role in proving the Green-Schwarz anomaly cancellation in superstring theory. The non-trivial coupling between Yang-Mills fields and the abelian antisymmetric tensor field is referred to as the Chapline-Manton coupling. The derivation of the coupling has also been discussed from the viewpoint of the BRST cohomology.
Recently, the Chapline-Manton coupling has also been derived on the basis of a gauge theory — a Yang-Mills theory in loop space. In this theory, the gauge field on loop space is given by the functional field on space-time, including an infinite number of local component fields. Yang-Mills fields (local Yang–Mills fields) are a part of the local component fields of the Yang-Mills fields on loop space, while an abelian antisymmetric tensor field of second rank is a part of the local component fields of the U(1) gauge field on loop space. The couplings among the local component fields of the Yang-Mills fields are determined by the symmetry of the loop gauge group. These couplings are essentially caused by the non-commutativety of the Lie algebra. In order to derive non-trivial couplings among the local component fields of the Yang-Mills fields and those of the U(1) gauge field, we need an extension of the loop gauge group. The suitable gauge group is the affine Lie group, which is a central extension of the loop group. Owing to the effect of the central extension, the extended gauge symmetry further leads to a new coupling between the local Yang-Mills fields and an abelian antisymmetric tensor field of second rank. This coincides with the Chapline-Manton coupling. In addition, the deformed transformation rule of the abelian antisymmetric tensor field of second rank is also derived from the transformation rule of the U(1) gauge field on loop space.
In addition to the abelian antisymmetric tensor field of second rank, for example, abelian (totally) antisymmetric tensor fields of higher rank appear in the Ramond-Ramond sector of superstring theories. These fields also contribute to a cancellation of the anomalies on D-branes under certain conditions. It is not difficult to extend the Chapline-Manton coupling for abelian antisymmetric tensor fields of higher rank. Indeed, such an antisymmetric tensor field can couple with the Chern-Simons (2n+1)-form. On the other hand, non-abelian antisymmetric tensor fields appear in the so-called BF term. The BF term is metric independent and takes the form of the product of a non-abelian antisymmetric tensor field $`B`$ and a field strength $`F`$ of Yang-Mills fields in the non-abelian case. The BF term is a generalization of the Chern-Simons term and is an important ingredient in BF Yang-Mills theories, topologically massive gauge theories, and so on. However, the extension of the Chapline-Manton coupling for non-abelian antisymmetric tensor fields is not known.
As is mentioned above, the Yang-Mills fields and U(1) gauge field on loop space possess an infinite number of local component fields. Non-abelian antisymmetric and symmetric tensor fields of second rank are the local component fields of the Yang-Mills fields on loop space, while an abelian tensor field of third rank with certain symmetric properties is the local component field of the U(1) gauge field on loop space. The interactions among these non-abelian tensor fields of second rank and local Yang-Mills fields can also be derived using the formalism of a non-linear realization developed for the loop gauge group applied to Yang-Mills theory in loop space. In this paper, we consider the Yang-Mills theory in loop space with the affine Lie gauge group and apply a non-linear realization method to the Yang-Mills theory. As we see below, the non-trivial interactions among the local Yang-Mills fields, non-abelian tensor fields of second rank, and an abelian tensor field of third rank can be systematically determined within the framework of the Yang-Mills theory. These local fields interact via a BF-like term.
This paper is organized as follows. In Sec. 2, we construct a Yang-Mills theory in loop space with the affine Lie gauge group. In Sec. 3, it is shown that the local field theory for Yang-Mills fields and an antisymmetric tensor field with the Chapline-Manton coupling is naturally derived on the basis of this theory. In Sec. 4, we apply a non-linear realization method developed for the affine Lie gauge group to the Yang-Mills theory in loop space. In Sec. 5, referring to the previous sections, we derive a local field theory for the non-abelian tensor fields, Yang-Mills fields and abelian tensor fields with non-trivial coupling. Section 6 is devoted to a summary and discussion of the possibilities for future development.
## 2 The Yang-Mills theory in loop space
We define a loop space $`\mathrm{\Omega }M^D`$ as the set of all loops in D-dimensional Minkowski space $`M^D`$. An arbitrary loop $`x^\mu =x^\mu (\sigma )[\mathrm{\hspace{0.17em}0}\sigma 2\pi ,x^\mu (0)=x^\mu (2\pi )]`$ in $`M^D`$ is represented as a point in $`\mathrm{\Omega }M^D`$ denoted by coordinates $`(x^{\mu \sigma })`$ with $`x^{\mu \sigma }x^\mu (\sigma )`$. <sup>1</sup><sup>1</sup>1<sup>)</sup> In the present paper the indices $`\mu `$, $`\nu `$, $`\kappa `$, $`\lambda `$, $`\xi `$ and $`\zeta `$ take the values 0, 1, 2, $`\mathrm{}`$, D-1, while the indices $`\rho `$, $`\sigma `$, $`\chi `$, and $`\omega `$ take continuous values from 0 to 2$`\pi `$.<sup>)</sup>
Let us consider a Yang-Mills theory in the loop space $`\mathrm{\Omega }M^D`$. We assume that a gauge group is an affine Lie group $`\widehat{G}_k`$ whose generators $`T_a(\sigma )`$ satisfy the commutation relation
$`[T_a(\rho ),T_b(\sigma )]=if_{ab}{}_{}{}^{c}T_{c}^{}(\rho )\delta (\rho \sigma )+ik\kappa _{ab}\delta ^{}(\rho \sigma )`$ (2.1)
and the hermiticity conditions $`T_a^{}(\sigma )=T_a(\sigma )`$. Here the $`f_{ab}^c`$ are the structure constants of the semisimple Lie group $`G`$ and $`\kappa _{ab}`$ is the Killing metric of $`G`$. <sup>2</sup><sup>2</sup>2 <sup>)</sup> The indices a, b, c, d, and e take the values 1, 2, 3, $`\mathrm{}`$, dim$`G`$.<sup>)</sup> The constant $`k`$ is called the ‘central charge’ of $`\widehat{G}_k`$ and takes an arbitrary value in this gauge theory. When we set $`k=0`$, the commutation relation (2.1) results in that for the loop group $`\widehat{G}_0`$. It is possible to make (non-trivial) central extensions to infinite-dimensional algebras. Hereafter, we refer to the Yang-Mills theory with an affine Lie gauge group as the ‘extended Yang-Mills theory’ (EYMT). It will become clear that the central extension of the gauge group leads to non-trivial couplings among non-abelian gauge fields and abelian gauge fields.
Let $`𝒜_{\mu \sigma }[x]`$ be a gauge field on $`\mathrm{\Omega }M^D`$. Owing to the central extension, the commutator given in (2.1) yields central terms without $`T_a(\sigma )`$, in addition to a linear combination of $`T_a(\sigma )`$. Consequently, the gauge field $`𝒜_{\mu \sigma }`$ needs extra terms without $`T_a(\sigma )`$ to allow for the consistency of the gauge transformation. We define the gauge field $`𝒜_{\mu \sigma }`$ as
$`𝒜_{\mu \sigma }[x]=𝒜_{\mu \sigma }^Y[x]+\stackrel{~}{𝒜}_{\mu \sigma }^U[x],`$ (2.2)
where $`𝒜_{\mu \sigma }^Y[x]`$ is a Yang-Mills field:
$`𝒜_{\mu \sigma }^Y[x]={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\rho }{2\pi }}𝒜_{\mu \sigma }{}_{}{}^{a\rho }[x]T_a(\rho ).`$ (2.3)
Here $`𝒜_{\mu \sigma }^{a\rho }`$ is a vector fields on $`\mathrm{\Omega }M^D`$, and $`\stackrel{~}{𝒜}_{\mu \sigma }^U`$ is a U(1) gauge field on $`\mathrm{\Omega }M^D`$ without $`T_a(\sigma )`$. As in the ordinary Yang-Mills theory, the infinitesimal gauge transformation for $`𝒜_{\mu \sigma }`$ is given by
$`\delta 𝒜_{\mu \sigma }[x]=_{\mu \sigma }\mathrm{\Lambda }[x]+i[𝒜_{\mu \sigma }[x],\mathrm{\Lambda }[x]],`$ (2.4)
where $`_{\mu \sigma }/x^{\mu \sigma }`$. An infinitesimal scalar function $`\mathrm{\Lambda }`$ on $`\mathrm{\Omega }M^D`$ is defined by
$`\mathrm{\Lambda }[x]=\mathrm{\Lambda }^Y[x]+\mathrm{\Lambda }^U[x],`$ (2.5)
where $`\mathrm{\Lambda }^Y[x]`$ is written
$`\mathrm{\Lambda }^Y[x]={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}\mathrm{\Lambda }^{a\sigma }[x]T_a(\sigma ),`$ (2.6)
with scalar functions $`\mathrm{\Lambda }^{a\sigma }`$ on $`\mathrm{\Omega }M^D`$, and $`\mathrm{\Lambda }^U`$ is a scalar function $`\mathrm{\Omega }M^D`$ without $`T_a(\sigma )`$. Since there is no relation between the gauge transformation (2.4) and a reparametrization $`\sigma \overline{\sigma }(\sigma )`$, $`\mathrm{\Lambda }^Y`$ and $`\mathrm{\Lambda }^U`$ obey the following reparametrization invariant conditions:
$`x^\mu (\sigma )_{\mu \sigma }\mathrm{\Lambda }^Y[x]={\displaystyle \frac{\mathrm{\Lambda }^{a\sigma }[x]}{\sigma }}T_a(\sigma ),`$ (2.7)
$`x^\mu (\sigma )_{\mu \sigma }\mathrm{\Lambda }^U[x]=0.`$ (2.8)
Here the prime denotes differentiation with respect to $`\sigma `$. Substituting (2.2) and (2.5) into (2.4) and considering the commutation relation (2.1), we obtain (the infinitesimal) gauge transformations for $`𝒜_{\mu \sigma }^Y`$ and $`\stackrel{~}{𝒜}_{\mu \sigma }^U`$:
$`\delta 𝒜_{\mu \sigma }^Y[x]=_{\mu \sigma }\mathrm{\Lambda }^Y[x]+i[𝒜_{\mu \sigma }^Y[x],\mathrm{\Lambda }^Y[x]]^Y,`$ (2.9)
$`\delta \stackrel{~}{𝒜}_{\mu \sigma }^U[x]=_{\mu \sigma }\mathrm{\Lambda }^U[x]+i[𝒜_{\mu \sigma }^Y[x],\mathrm{\Lambda }^Y[x]]^U.`$ (2.10)
Here, $`[,]^Y`$ denotes the part of a commutator $`[,]`$ written as a linear combination of $`T_a(\sigma )`$, while $`[,]^U`$ denotes the other part including the central charge $`k`$. In deriving these gauge transformations, we have used the fact that $`\stackrel{~}{𝒜}_{\mu \sigma }^U`$ and $`\mathrm{\Lambda }^U`$ are commutative. We note that the Yang-Mills fields $`𝒜_{\mu \sigma }^Y`$ appear in the gauge transformation of the U(1) gauge field $`\stackrel{~}{𝒜}_{\mu \sigma }^U`$. The transformation (2.10) is obviously different from an ordinary U(1) gauge transformation: $`\delta 𝒜_{\mu \sigma }^U=_{\mu \sigma }\mathrm{\Lambda }^U`$. The second term on the right-hand side of (2.10) is due to the central extension. Combining the reparametrization invariant condition (2.8) and the ordinary U(1) gauge transformation, we can obtain the condition $`x^\mu (\sigma )𝒜_{\mu \sigma }^U=0`$. In the present case, however, the U(1) gauge field $`\stackrel{~}{𝒜}_{\mu \sigma }^U`$ dose not satisfy the condition $`x^\mu (\sigma )\stackrel{~}{𝒜}_{\mu \sigma }^U=0`$ unless $`k=0`$.
Next, we consider the (naive) field strength of the gauge field $`𝒜_{\mu \sigma }`$ as
$`_{\mu \rho ,\nu \sigma }`$ $`=`$ $`_{\mu \rho }𝒜_{\nu \sigma }_{\nu \sigma }𝒜_{\mu \rho }+i[𝒜_{\mu \rho },𝒜_{\nu \sigma }]`$ (2.11)
$`=`$ $`_{\mu \rho ,\nu \sigma }^Y+_{\mu \rho ,\nu \sigma }^U,`$
with
$`_{\mu \rho ,\nu \sigma }^Y_{\mu \rho }𝒜_{\nu \sigma }^Y_{\nu \sigma }𝒜_{\mu \rho }^Y+i[𝒜_{\mu \rho }^Y,𝒜_{\nu \sigma }^Y]^Y,`$ (2.12)
$`_{\mu \rho ,\nu \sigma }^U_{\mu \rho }\stackrel{~}{𝒜}_{\nu \sigma }^U_{\nu \sigma }\stackrel{~}{𝒜}_{\mu \rho }^U+i[𝒜_{\mu \rho }^Y,𝒜_{\nu \sigma }^Y]^U.`$ (2.13)
Note that (2.13) is different from the ordinary field strength of the U(1) gauge field $`\stackrel{~}{𝒜}_{\mu \sigma }^U`$. The (naive) field strength (2.11) obeys the ordinary gauge transformation rule: $`\delta _{\mu \rho ,\nu \sigma }=i[_{\mu \rho ,\nu \sigma },\mathrm{\Lambda }]`$. Because of the central extension, however, we can immediately find that $`_{\mu \rho ,\nu \sigma }^Y`$ obeys the homogeneous gauge-transformation rule $`\delta _{\mu \rho ,\nu \sigma }^Y=i[_{\mu \rho ,\nu \sigma }^Y,\mathrm{\Lambda }^Y]^Y`$, while $`_{\mu \rho ,\nu \sigma }^U`$ obeys the inhomogeneous gauge-transformation rule $`\delta _{\mu \rho ,\nu \sigma }^U=i[_{\mu \rho ,\nu \sigma }^Y,\mathrm{\Lambda }^Y]^U`$. Therefore, (2.11) is not suitable for the field strength under the affine Lie gauge group. For this reason, we modify (2.11) as
$`_{\mu \rho ,\nu \sigma }_{\mu \rho ,\nu \sigma }^Y+_{\mu \rho ,\nu \sigma }^U,`$ (2.14)
with
$`_{\mu \rho ,\nu \sigma }^U_{\mu \rho ,\nu \sigma }^U+k{\displaystyle \frac{d\omega }{2\pi }x^\lambda (\omega )\mathrm{Tr}\left[𝒜_{\lambda \omega }^Y_{\mu \rho ,\nu \sigma }^Y\right]}.`$ (2.15)
Here $`\mathrm{`}\mathrm{`}\mathrm{Tr}\mathrm{"}`$ denotes the inner product of two elements of the affine Lie algebra $`\widehat{G}`$: $`\mathrm{Tr}[VW]=_{a,b}_0^{2\pi }\frac{d\sigma }{2\pi }\kappa _{ab}V^{a\sigma }W^{b\sigma }`$. We can confirm that $`_{\mu \rho ,\nu \sigma }^U`$ is gauge invariant under the reparametrization invariant condition (2.7). Note that the gauge invariance is still maintained without the central extension. In order for the right-hand side of (2.15) to transform in the same manner as $`_{\mu \rho ,\nu \sigma }^U`$ under the reparametrization, however, it is necessary to restrict $`𝒜_{\mu \sigma }^{a\rho }`$ in (2.4) to the form
$`𝒜_{\mu \sigma }{}_{}{}^{a\rho }[x]=\delta (\rho \sigma )𝒜_\mu {}_{}{}^{a\rho }[x].`$ (2.16)
Here $`𝒜_\mu ^{a\rho }`$ are the fields on $`\mathrm{\Omega }M^D`$ that behave as vector functionals on $`M^D`$.
The action for $`𝒜_{\mu \sigma }`$ is defined as
$`S_\mathrm{R}={\displaystyle \frac{1}{V_R}}{\displaystyle [dx]\left(^Y+^U\right)\mathrm{exp}\left(\frac{L}{l^2}\right)},`$ (2.17)
with $`V_R\{dx\}\mathrm{exp}(L/l^2)`$, and the Lagrangians
$`^Y={\displaystyle \frac{1}{4}}N^Y𝒢^{\kappa \rho ,\lambda \sigma }𝒢^{\mu \chi ,\nu \omega }\mathrm{Tr}[_{\kappa \rho ,\mu \chi }^Y_{\lambda \sigma ,\nu \omega }^Y],`$ (2.18)
$`^U={\displaystyle \frac{1}{4}}N^U𝒢^{\kappa \rho ,\lambda \sigma }𝒢^{\mu \chi ,\nu \omega }_{\kappa \rho ,\mu \chi }^U_{\lambda \sigma ,\nu \omega }^U,`$ (2.19)
where $`N^Y`$ and $`N^U`$ are arbitrary constants.
Here, the measures $`[dx]`$ and $`\{dx\}`$ are given by $`[dx]_{\mu =0}^{D1}_{n=\mathrm{}}^{\mathrm{}}dx^{\mu n}`$ and $`\{dx\}_{\mu =0}^{D1}_{n=\mathrm{}}^{\mathrm{},n0}dx^{\mu n}`$, where the $`x^{\mu n}`$ are the coefficients of the Fourier expansion $`x^\mu (\sigma )=_{n=\mathrm{}}^{\mathrm{}}x^{\mu n}e^{in\sigma }`$. These measures are invariant under a reparametrization. The (inverse) metric tensor $`𝒢^{\mu \rho ,\nu \sigma }`$ on $`\mathrm{\Omega }M^D`$ is defined by $`𝒢^{\mu \rho ,\nu \sigma }\eta ^{\mu \nu }\delta (\rho \sigma )`$, where $`\eta _{\mu \nu }`$ is the metric tensor on $`M^D`$. <sup>3</sup><sup>3</sup>3 <sup>)</sup> diag $`\eta _{\mu \nu }=(1,1,1,\mathrm{},1)`$ .<sup>)</sup> The damping factor $`\mathrm{exp}(L/l^2)`$ with $`L_0^{2\pi }\frac{d\sigma }{2\pi }\eta _{\mu \nu }x^\mu (\sigma )x^\nu (\sigma )`$ is inserted into the action so that it becomes well defined, where $`l(>0)`$ is a constant with the dimension of length giving the size of loops.
We would like to focus attention on the fact that there is coupling between the Yang-Mills fields $`𝒜_{\mu \sigma }^Y`$ and the U(1) gauge field $`𝒜_{\mu \sigma }^U`$ in the Lagrangian (2.19). It is obvious that the coupling is due to the central extension of the gauge group.
The Lagrangians $`^Y`$ and $`^U`$ are gauge invariant, while they are not reparametrization invariant, due to the definition of the inner product $`\mathrm{`}\mathrm{`}\mathrm{Tr}\mathrm{"}`$ and the metric tensor $`𝒢^{\mu \rho ,\nu \sigma }`$. If necessary, we can indeed define an inner product and metric tensor to maintain reparametrization invariance as well as gauge-invariance. The metric $`\eta ^{\mu \nu }\delta (\rho \sigma )`$ and the inner product $`\mathrm{`}\mathrm{`}\mathrm{Tr}\mathrm{"}`$ that we employ can be shown in a concrete calculation to be forms of the reparametrization invariant inner product and metric tensor in a certain gauge of reparametrization.
## 3 The Chapline-Manton coupling
In this section, we derive a local field theory with a coupling between local Yang-Mills fields and an abelian antisymmetric tensor field of second rank on the basis of the Yang-Mills theory in loop space. Let us consider the simplest solutions of (2.7) consisting of local functions on $`M^D`$,
$`\mathrm{\Lambda }^{Y(0)}[x]={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}g_0\mathrm{\Lambda }^a(x(\sigma ))T_a(\sigma ),`$ (3.1)
where $`\mathrm{\Lambda }^a`$ is an infinitesimal scalar function on $`M^D`$ and $`g_0`$ is a constant of dimension $`[\mathrm{length}]^{\frac{D4}{2}}`$. On the other hand, the simplest solution of (2.8) consisting of a local function on $`M^D`$ is given by
$`\mathrm{\Lambda }^{U(0)}[x]={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}q_0x^\mu (\sigma )\lambda _\mu (x(\sigma )),`$ (3.2)
where $`\lambda _\mu `$ is a vector function on $`M^D`$ and $`q_0`$ is a constant of dimension $`[\mathrm{length}]^{\frac{D6}{2}}`$.
Corresponding to (3.1) and (3.2), we consider the (restricted) Yang-Mills field $`𝒜_{\mu \sigma }^Y`$ and the U(1) gauge field $`𝒜_{\mu \sigma }^U`$ written in terms of local fields as
$`𝒜_{\mu \sigma }^{Y(0)}[x]=g_0A_\mu {}_{}{}^{a}(x(\sigma ))T_a(\sigma ),`$ (3.3)
$`\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}[x]=q_0x^\nu (\sigma )\{B_{\mu \nu }(x(\sigma ))+C_{\mu \nu }(x(\sigma ))\},`$ (3.4)
respectively, where the $`A_\mu {}_{}{}^{a}(x)`$ are vector fields on $`M^D`$, and the $`B_{\mu \nu }(x)`$ and $`C_{\mu \nu }(x)`$ are antisymmetric and symmetric tensor fields of second rank on $`M^D`$. Obviously, the right-hand sides of (3.3) and (3.4) transform in the same manner as the left-hand sides of (3.3) and (3.4) under the reparametrization. Substituting (3.1) and (3.3) into (2.9), we obtain the transformation rule of Yang-Mills fields for $`A_\mu ^a`$ as
$`\delta A_\mu {}_{}{}^{a}(x)=D_\mu \mathrm{\Lambda }^a(x)_\mu \mathrm{\Lambda }^a(x)q_0A_\mu {}_{}{}^{b}(x)\mathrm{\Lambda }^c(x)f_{bc}{}_{}{}^{a},`$ (3.5)
by virtue of (2.1). On the other hand, substitution of (3.1), (3.2), (3.3) and (3.4) into (2.10) yields the transformation rules of $`B_{\mu \nu }`$ and $`C_{\mu \nu }`$ as
$`\delta B_{\mu \nu }(x)=_{[\mu }\lambda _{\nu ]}(x)\stackrel{~}{k}_0A_{[\mu }{}_{}{}^{a}(x)_{\nu ]}\lambda _a(x),`$ (3.6)
$`\delta C_{\mu \nu }(x)=\stackrel{~}{k}_0A_{(\mu }{}_{}{}^{a}(x)_{\nu )}\lambda _a(x),^)`$ (3.7)
where the lowering of the index $`a`$ has been carried out with $`\kappa _{ab}`$. Here, $`\stackrel{~}{k}_0kg_0/2q_0`$ is a constant of dimension \[ length \]. Owing to the central extension, the local Yang-Mills fields $`A_\mu ^a`$ appear in the transformation rules of $`B_{\mu \nu }`$ and $`C_{\mu \nu }`$. However, the transformation rules of $`C_{\mu \nu }`$ and $`A_\mu ^a`$ are not independent. Indeed, the following symmetric tensor field of second rank is invariant under the transformation rules (3.5) and (3.7):
$`\stackrel{~}{C}_{\mu \nu }(x)C_{\mu \nu }(x)+{\displaystyle \frac{1}{2}}\stackrel{~}{k}_0A_{(\mu }{}_{}{}^{a}(x)A_{\nu ),a}(x).`$ (3.8)
Substituting $`C_{\mu \nu }=\stackrel{~}{C}_{\mu \nu }\frac{1}{2}\stackrel{~}{k}_0A_{(\mu }{}_{}{}^{a}A_{\nu ),a}^{}`$ into (3.4), we obtain $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}`$ written in terms of the local fields $`A_\mu ^a`$, $`B_{\mu \nu }`$ and $`\stackrel{~}{C}_{\mu \nu }`$. As we shall see, the couplings of the local fields are uniquely determined in the Yang-Mills theory in loop space. On examination of the couplings of these local fields, however, we find that the gauge invariant tensor field $`\stackrel{~}{C}_{\mu \nu }`$ is free of $`A_\mu ^a`$ and $`B_{\mu \nu }`$. Since we are interested in the couplings of local fields, we omit $`\stackrel{~}{C}_{\mu \nu }`$ from $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}`$, written in terms of the local fields $`A_\mu ^a`$, $`B_{\mu \nu }`$ and $`\stackrel{~}{C}_{\mu \nu }`$ for simplicity. In other words, we replace $`C_{\mu \nu }`$ with $`\frac{1}{2}\stackrel{~}{k}_0A_{(\mu }{}_{}{}^{a}A_{\nu ),a}^{}`$ in (3.4). Then (3.4) is rewritten as
$`\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}[x]=𝒜_{\mu \sigma }^{U(0)}[x]{\displaystyle \frac{1}{2}}\stackrel{~}{k}_0q_0x^\nu (\sigma )A_{(\mu }{}_{}{}^{a}(x(\sigma ))A_{\nu ),a}(x(\sigma )),`$ (3.9)
with
$`𝒜_{\mu \sigma }^{U(0)}[x]=q_0x^\nu (\sigma )B_{\mu \nu }(x(\sigma )),`$ (3.10)
where $`𝒜_{\mu \sigma }^{U(0)}`$ is written in terms of the abelian local field only. Since the constant $`\stackrel{~}{k}_0`$ is proportional to the central charge $`k`$, the U(1) gauge field $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}`$ is reduced to $`𝒜_{\mu \sigma }^{U(0)}`$ by setting $`k=0`$. Note that $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}`$ does not satisfy the condition $`x^\mu (\sigma )\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}=0`$ for $`k0`$, while $`𝒜_{\mu \sigma }^{U(0)}`$ satisfies the condition $`x^\mu (\sigma )𝒜_{\mu \sigma }^{U(0)}=0`$. As we have discussed in Sec. 2, this fact is consistent with the fact that the U(1) gauge field $`\stackrel{~}{𝒜}_{\mu \sigma }^U`$ does not satisfy the condition $`x^\mu (\sigma )\stackrel{~}{𝒜}_{\mu \sigma }^U=0`$ unless $`k=0`$.
Next, let us consider the field strength. Substituting (3.3) into (2.12), by virtue of (2.1) we obtain $`_{\mu \rho ,\nu \sigma }^{Y(0)}`$ written in terms of the field strength of $`A_\mu ^a`$:
$`_{\mu \rho ,\nu \sigma }^{Y(0)}=g_0F_{\mu \nu }{}_{}{}^{a}(x(\sigma ))T_a(\sigma )\delta (\rho \sigma ),`$ (3.11)
with
$`F_{\mu \nu }{}_{}{}^{a}(x)=_\mu A_\nu {}_{}{}^{a}(x)_\nu A_\mu {}_{}{}^{a}(x)A_\mu {}_{}{}^{b}(x)A_\nu {}_{}{}^{c}(x)f_{bc}{}_{}{}^{a}.`$ (3.12)
Also the substitution of (3.3), (3.4) and (3.11) into (2.15) yields
$`_{\mu \rho ,\nu \sigma }^{U(0)}=q_0x^\lambda (\sigma )H_{\mu \nu \lambda }(x(\sigma ))`$ (3.13)
with
$`H_{\mu \nu \kappa }(x)=F_{\mu \nu \kappa }(x)+\stackrel{~}{k}_0\mathrm{\Omega }_{\mu \nu \kappa }(x),`$ (3.14)
and
$`F_{\mu \nu \lambda }(x)_\mu B_{\nu \lambda }(x)+_\nu B_{\lambda \mu }(x)+_\lambda B_{\mu \nu }(x),`$ (3.15)
$`\mathrm{\Omega }_{\mu \nu \lambda }(x)=A_{[\mu }{}_{}{}^{a}(x)_\nu A_{\lambda ],a}(x){\displaystyle \frac{g_0}{3}}A_{[\mu }{}_{}{}^{a}(x)A_\nu {}_{}{}^{b}(x)A_{\lambda ]}{}_{}{}^{c}(x)f_{abc}.`$ (3.16)
Here, $`\mathrm{\Omega }_{\mu \nu \lambda }`$ occurring in $`H_{\mu \nu \lambda }`$ is a Chern-Simons 3-form. Reflecting the fact that $`\stackrel{~}{}_{\mu \rho ,\nu \sigma }^{U(0)}`$ is gauge invariant, $`H_{\mu \nu \lambda }`$ becomes also invariant under the transformation rules of (3.5) and (3.6).
Finally, let us derive the action $`S_R^{(0)}`$ of the local fields $`A_\mu ^a`$ and $`B_{\mu \nu }`$. Substituting (3.11) and (3.13) into (2.18) and (2.19), respectively, and integrating them over $`\rho `$, $`\chi `$ and $`\omega `$, we obtain the Lagrangians $`^{Y(0)}`$ and $`^{U(0)}`$ expressed as integrals over $`\sigma `$:
$`^{Y(0)}={\displaystyle \frac{1}{4}}N^Yg_0^2\delta (0)^2{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}F_{\mu \nu ,a}(x(\sigma ))F^{\mu \nu ,a}(x(\sigma )),`$ (3.17)
$`^{U(0)}={\displaystyle \frac{1}{4}}N^Uq_0^2\delta (0){\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}x^\kappa (\sigma )x_\lambda ^{}(\sigma )H_{\kappa \mu \nu }(x(\sigma ))H^{\lambda \mu \nu }(x(\sigma )).`$ (3.18)
We next insert (3.17) and (3.18) into (2.17) and expand the functions of $`x^\mu (\sigma )`$, $`F_{\mu \nu ,a}(x(\sigma ))F^{\mu \nu ,a}(x(\sigma ))`$ and $`H_{\mu \nu \lambda }(x(\sigma ))H^{\mu \nu \lambda }(x(\sigma ))`$, about $`x^{\mu 0}`$. Then, all the differential coefficients at $`x^{\mu 0}`$ in each Taylor series become total derivatives with respect to $`x^{\mu 0}`$ and vanish under the boundary conditions with $`|x^{\mu 0}|\mathrm{}`$. As a result, we obtain an action in which the argument $`x^\mu (\sigma )`$ of the functions is replaced with $`x^{\mu 0}`$.
Carrying out the integrations with respect to $`x^{\mu n}`$ after the Wick rotations $`x^{\mu n}ix^{\mu n}(n0)`$, we obtain the action $`S_\mathrm{R}^{(0)}`$ as
$`S_\mathrm{R}^{(0)}={\displaystyle }d^Dx\{{\displaystyle \frac{1}{4}}F_{\mu \nu }{}_{}{}^{a}(x)F^{\mu \nu }{}_{a}{}^{}(x)+{\displaystyle \frac{1}{12}}H_{\mu \nu \lambda }(x)H^{\mu \nu \lambda }(x)\}.`$ (3.19)
Here, we set the normalization conditions as $`N^Yg_0^2\delta (0)=1`$ and $`3N^Uq_0^2l^2\delta ^2(0)/2=1`$. Thus, we obtain the action describing the system with the coupling between the local Yang-Mills fields $`A_\mu ^a`$ and the abelian antisymmetric tensor field $`B_{\mu \nu }`$ via the Chern-Simons 3-form. Setting the structure constants $`f_{ab}^c`$ to $`0`$ in (2.1), we can obtain the action (3.19) in the abelian version. It describes the system with the coupling between the U(1) gauge field $`A_\mu `$ and the abelian antisymmetric tensor field $`B_{\mu \nu }`$ via the abelian Chern-Simons 3-form. Such couplings have been introduced by de Wit et al. in the abelian case and Chapline and Manton in the non-abelian case. They assign the transformation rule (3.6) to the antisymmetric tensor field $`B_{\mu \nu }`$ in order for (3.14) to be gauge invariant. In contrast, we can naturally derive (3.6) in the framework of the gauge theory.
## 4 Application of the non-linear realization to the EYMT
In this section, we discuss the local field theories for higher rank tensor fields based on the EYMT.
Let us consider the solutions of (2.7) consisting of local functions on $`M^D`$:
$`\mathrm{\Lambda }^{Y(p)}[x]={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}g_pQ^{\mu _1}(\sigma )Q^{\mu _2}(\sigma )\mathrm{}Q^{\mu _p}(\sigma )\lambda _{\mu _1\mu _2,\mathrm{},\mu _p}{}_{}{}^{a}(x(\sigma ))T_a(\sigma ).`$ (4.1)
Here, $`Q_\mu (\sigma )x^\mu (\sigma )/\sqrt{x^{\mathrm{\hspace{0.17em}2}}(\sigma )}`$, <sup>5</sup><sup>5</sup>5<sup>)</sup> $`x^{\mathrm{\hspace{0.17em}2}}(\sigma )x_\mu ^{}(\sigma )x^\mu (\sigma )`$ . <sup>)</sup> where the $`\lambda _{\mu _1\mu _2,\mathrm{},\mu _p}{}_{}{}^{a}(x)`$ are infinitesimal tensor function of rank $`p`$ ($`p=0,1,2,\mathrm{},`$) on $`M^D`$ and $`g_p`$ is a constant of dimension $`[\mathrm{length}]^{\frac{D4}{2}}`$. Similarly, the solutions of (2.8) are given by
$`\mathrm{\Lambda }^{U(p)}[x]`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}q_p\sqrt{x^{\mathrm{\hspace{0.17em}2}}(\sigma )}Q^{\mu _1}(\sigma )Q^{\mu _2}(\sigma )\mathrm{}Q^{\mu _p}(\sigma )Q^{\mu _{p+1}}(\sigma )\lambda _{\mu _1\mu _2,\mathrm{},\mu _p\mu _{p+1}}(x(\sigma ))`$ (4.2)
$`+`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}e_pQ^{\mu _1}(\sigma )Q^{\mu _2}(\sigma )\mathrm{}Q^{\mu _p}(\sigma )\kappa _{\mu _1\mu _2,\mathrm{},\mu _p}(x(\sigma )),`$
where $`\lambda _{\mu _1\mu _2,\mathrm{},\mu _p\mu _{p+1}}(x)`$ and $`\kappa _{\mu _1\mu _2,\mathrm{},\mu _p}(x)`$ are infinitesimal tensor functions of rank ($`p`$ +1) and $`p`$ on $`M^D`$, and $`q_p`$ and $`e_p`$ are constants of dimension $`[\mathrm{length}]^{\frac{D6}{2}}`$ and $`[\mathrm{length}]^{\frac{D4}{2}}`$, respectively. Setting $`p=0`$, the infinitesimal functions (4.1) and (4.2) correspond to (3.1) and (3.2), respectively.
Any general solution of (2.7) is given as a linear combination of $`\mathrm{\Lambda }^{Y(p)}`$, while any general solution of (2.8) is given as a linear combination of $`\mathrm{\Lambda }^{U(p)}`$. Explicitly, $`\mathrm{\Lambda }^Y[x]`$ and $`\mathrm{\Lambda }^U[x]`$ can be expressed as $`\mathrm{\Lambda }^Y[x]_{p=0}^{\mathrm{}}\mathrm{\Lambda }^{Y(p)}[x]`$ and $`\mathrm{\Lambda }^U[x]_{p=0}^{\mathrm{}}\mathrm{\Lambda }^{U(p)}[x]`$. (The coefficients of the terms in these sums are absorbed into $`g_p`$ and $`q_p`$.) Consequently, any infinitesimal function $`\mathrm{\Lambda }`$ consisting of local functions on $`M^D`$ is given in the form of a linear combination:
$`\mathrm{\Lambda }[x]=\mathrm{\Lambda }^Y[x]+\mathrm{\Lambda }^U[x]{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}\mathrm{\Lambda }^{(p)}[x],`$ (4.3)
with
$`\mathrm{\Lambda }^{(p)}[x]\mathrm{\Lambda }^{Y(p)}[x]+\mathrm{\Lambda }^{U(p)}[x].`$ (4.4)
Corresponding to (4.3), we express the gauge field $`𝒜_{\mu \sigma }`$ as a linear combination of $`𝒜_{\mu \sigma }^{(p)}`$ associated with $`\mathrm{\Lambda }^{(p)}`$ :
$`𝒜_{\mu \sigma }[x]=𝒜_{\mu \sigma }^Y[x]+\stackrel{~}{𝒜}_{\mu \sigma }^U[x]{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}𝒜_{\mu \sigma }^{(p)}[x],`$ (4.5)
with
$`𝒜_{\mu \sigma }^{(p)}[x]𝒜_{\mu \sigma }^{Y(p)}[x]+\stackrel{~}{𝒜}_{\mu \sigma }^{U(p)}[x],`$ (4.6)
where $`𝒜_{\mu \sigma }^{Y(p)}`$ are Yang-Mills fields consisting of the local tensor fields of rank $`p`$ and $`(p+1)`$ on $`M^D`$, and $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(p)}`$ are U(1) gauge fields consisting of the local tensor fields of rank $`p`$, $`(p+1)`$ and $`(p+2)`$ on $`M^D`$. (As an exceptional case, $`𝒜_{\mu \sigma }^{Y(0)}`$ consists of local vector fields only, and $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(0)}[x]`$ consists of local fields of second rank only.)
Substituting (4.3) and (4.5) into (2.2), and comparing the two sides of the resulting equation, we conclude that
$`\delta 𝒜_{\mu \sigma }^{(p)}=_{\mu \sigma }\mathrm{\Lambda }^{(p)}+i{\displaystyle \underset{k=0}{\overset{p}{}}}[𝒜_{\mu \sigma }^{(k)},\mathrm{\Lambda }^{(pk)}].`$ (4.7)
We note that $`𝒜_{\mu \sigma }^{(0)}`$ obeys the same gauge transformation as (2.2), while the $`𝒜_{\mu \sigma }^{(p)}(p=1,2,3,\mathrm{})`$ do not obey the gauge transformation as (2.2). Indeed, the gauge transformation of $`𝒜_{\mu \sigma }^{(p)}(p=1,2,3,\mathrm{})`$ depends on other gauge fields $`𝒜_{\mu \sigma }^{(k)}`$ $`[k(<p)=0,1,2,3,\mathrm{}]`$.
Next, we consider the (naive) field strength of $`𝒜_{\mu \sigma }^{(p)}`$. We substitute (4.5) into (2.11). Then we can decompose $`_{\mu \rho ,\nu \sigma }`$ as
$`_{\mu \rho ,\nu \sigma }=_{\mu \rho ,\nu \sigma }^Y+_{\mu \rho ,\nu \sigma }^U={\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}_{\mu \rho ,\nu \sigma }^{(p)},`$ (4.8)
with
$`_{\mu \rho ,\nu \sigma }^{(p)}_{\mu \rho }𝒜_{\nu \sigma }^{(p)}_{\nu \sigma }𝒜_{\mu \rho }^{(p)}+{\displaystyle \underset{k=0}{\overset{p}{}}}[𝒜_{\mu \rho }^{(k)},𝒜_{\nu \sigma }^{(pk)}].`$ (4.9)
Under the gauge transformation (4.7), the $`_{\mu \rho ,\nu \sigma }^{(p)}`$ transform as
$`\delta _{\mu \rho ,\nu \sigma }^{(p)}=i{\displaystyle \underset{k=0}{\overset{p}{}}}[_{\mu \rho ,\nu \sigma }^{(k)},\mathrm{\Lambda }^{(pk)}].`$ (4.10)
From (4.10), we see that the $`_{\mu \rho ,\nu \sigma }^{(p)}(p=1,2,3,\mathrm{})`$ do not obey the same transformation rule $`\delta _{\mu \rho ,\nu \sigma }=i[_{\mu \rho ,\nu \sigma },\mathrm{\Lambda }]`$, except $`_{\mu \rho ,\nu \sigma }^{(0)}`$. Accordingly, we cannot construct the “modified” field strengths corresponding to (2.14) for $`𝒜_{\mu \sigma }^{(p)}(p=1,2,3,\mathrm{},)`$ in the manner discussed in Sec. 2. To begin with, we must find the suitable field strengths of $`𝒜_{\mu \sigma }^{(p)}(p=1,2,3,\mathrm{},)`$ that obey the same transformation rule as $`_{\mu \rho ,\nu \sigma }`$. Such field strengths can be systematically derived by using a non-linear realization method.
Let us consider the linear subspace $`\widehat{𝐠}_k^{(p)}`$ all of whose elements have the form
$`\mathrm{\Xi }^{(p)}[x]\mathrm{\Xi }^{Y(p)}[x]+\mathrm{\Xi }^{U(p)}[x],`$ (4.11)
with
$`\mathrm{\Xi }^{Y(p)}[x]`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}g_pQ^{\mu _1}(\sigma )Q^{\mu _2}(\sigma )\mathrm{}Q^{\mu _p}(\sigma )\xi _{\mu _1\mu _2,\mathrm{},\mu _p}{}_{}{}^{a}(x(\sigma ))T_a(\sigma ),`$ (4.12)
$`\mathrm{\Xi }^{U(p)}[x]`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}q_p\sqrt{x^{\mathrm{\hspace{0.17em}2}}(\sigma )}Q^{\mu _1}(\sigma )Q^{\mu _2}(\sigma )\mathrm{}Q^{\mu _p}(\sigma )Q^{\mu _{p+1}}(\sigma )\xi _{\mu _1\mu _2,\mathrm{},\mu _p\mu _{p+1}}(x(\sigma ))`$ (4.13)
$`+`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}e_pQ^{\mu _1}(\sigma )Q^{\mu _2}(\sigma )\mathrm{}Q^{\mu _p}(\sigma )\zeta _{\mu _1\mu _2,\mathrm{},\mu _p}(x(\sigma )).`$
Here, $`\xi _{\mu _1\mu _2,\mathrm{},\mu _p}{}_{}{}^{a}(x)`$ and $`\zeta _{\mu _1\mu _2,\mathrm{},\mu _p}(x)`$ are arbitrary tensor functions of rank $`p`$ and $`\xi _{\mu _1\mu _2,\mathrm{},\mu _p\mu _{p+1}}(x)`$ is an arbitrary tensor function of rank ($`p`$ +1) on $`M^D`$. Using (2.1), we obtain a commutation relation for every $`\mathrm{\Xi }^{(p)}[x]\widehat{𝐠}_k^{(p)}`$ and $`\mathrm{\Xi }^{(q)}[x]\widehat{𝐠}_k^{(q)}`$ as
$`[\mathrm{\Xi }^{(p)},\mathrm{\Xi }^{(q)}]\widehat{𝐠}_k^{(p+q)}.`$ (4.14)
This commutation relation shows that the linear subspace $`\widehat{𝐠}_k^{(0)}`$ is a Lie algebra, whereas $`\widehat{𝐠}_k^{(p)}(p=1,2,3,\mathrm{})`$ is not a Lie algebra. The direct sum $`\widehat{𝐠}_k`$ $`_{p=0}^{\mathrm{}}`$ $`\widehat{𝐠}_k^{(p)}`$ is obviously a Lie algebra. Thus, the linear subspace $`\widehat{𝐠}_k^{(0)}`$ forms a subalgebra of $`\widehat{𝐠}_k`$. We consider the Lie groups $`\widehat{𝐆}_k`$ and $`\widehat{𝐆}_k^{(0)}`$ associated with $`\widehat{𝐠}_k`$ and $`\widehat{𝐠}_k^{(0)}`$, respectively. Here, both $`\widehat{𝐆}_k`$ and $`\widehat{𝐆}_k^{(0)}`$ are Lie subgroups of the affine Lie group $`\widehat{G}_k`$. Since $`\widehat{𝐆}_k^{(0)}`$ is a subgroup of $`\widehat{𝐆}_k`$, we can consider the coset manifold $`\widehat{𝐆}_k`$/ $`\widehat{𝐆}_k^{(0)}`$. We introduce the scalar field $`\mathrm{\Phi }^{(p)}[x]`$ on loop space so as to parameterize the coset representative of $`\widehat{𝐆}_k`$/$`\widehat{𝐆}_k^{(0)}`$:
$`𝒱[\mathrm{\Phi }]=\mathrm{exp}\left(i{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}\mathrm{\Phi }^{(p)}[x]\right).`$ (4.15)
Here $`\mathrm{\Phi }^{(p)}[x]`$ is an element of $`\widehat{𝐠}_k^{(p)}`$:
$`\mathrm{\Phi }^{(p)}[x]=\mathrm{\Phi }^{Y(p)}[x]+\mathrm{\Phi }^{U(p)}[x].`$ (4.16)
The (finite) transformation rule of $`𝒱[\mathrm{\Phi }]𝒱[\overline{\mathrm{\Phi }}]`$ for $`𝒳\widehat{𝐆}_k`$ is given by
$`𝒱[\mathrm{\Phi }]𝒱[\overline{\mathrm{\Phi }}]=𝒳𝒱[\mathrm{\Phi }]𝒴^1[\mathrm{\Phi },𝒳],`$ (4.17)
where $`𝒳=\mathrm{exp}(i_{p=0}^{\mathrm{}}\mathrm{\Xi }^{(p)})`$ and $`𝒴=\mathrm{exp}(i\mathrm{\Xi }^{(0)})`$. (Here, $`𝒴`$ depends on $`𝒳`$ and $`\mathrm{\Phi }`$.)
Next, we define the vector field $`\widehat{𝒜}_{\mu \sigma }`$ by using $`𝒜_{\mu \sigma }`$ and $`𝒱[\mathrm{\Phi }]`$:
$`\widehat{𝒜}_{\mu \sigma }𝒱^1𝒜_{\mu \sigma }𝒱i𝒱^1_{\mu \sigma }𝒱.`$ (4.18)
Substituting (4.5) and (4.15) into (4.18), we can express $`\widehat{𝒜}_{\mu \sigma }`$ as a linear combination: $`\widehat{𝒜}_{\mu \sigma }=_{p=0}^{\mathrm{}}\widehat{𝒜}_{\mu \sigma }^{(p)}`$. The concrete forms of $`\widehat{𝒜}_{\mu \sigma }^{(0)}`$ and $`\widehat{𝒜}_{\mu \sigma }^{(1)}`$ are given by
$`\widehat{𝒜}_{\mu \sigma }^{(0)}=𝒜_{\mu \sigma }^{(0)},`$ (4.19)
$`\widehat{𝒜}_{\mu \sigma }^{(1)}=𝒜_{\mu \sigma }^{(1)}(_{\mu \sigma }\mathrm{\Phi }^{(1)}+i[𝒜_{\mu \sigma }^{(0)},\mathrm{\Phi }^{(1)}]),`$ (4.20)
respectively. Under the transformation (4.17) and the finite gauge transformation $`𝒜_{\mu \sigma }\overline{𝒜}_{\mu \sigma }=𝒳𝒜_{\mu \sigma }𝒳^1i𝒳_{\mu \sigma }𝒳^1`$, we obtain the finite gauge transformation rule of $`\widehat{𝒜}_{\mu \sigma }`$:
$`\widehat{𝒜}_{\mu \sigma }\overline{\widehat{𝒜}}_{\mu \sigma }=𝒴\widehat{𝒜}_{\mu \sigma }𝒴^1i𝒴_{\mu \sigma }𝒴^1.`$ (4.21)
Consequently, the (naive) field strength $`\widehat{}_{\mu \rho ,\nu \sigma }_{\mu \rho }\widehat{𝒜}_{\nu \sigma }_{\nu \sigma }\widehat{𝒜}_{\mu \rho }+i[\widehat{𝒜}_{\mu \rho },\widehat{𝒜}_{\nu \sigma }]`$ obeys the transformation rule
$`\widehat{}_{\mu \rho ,\nu \sigma }\overline{\widehat{}}_{\mu \rho ,\nu \sigma }=𝒴\widehat{}_{\mu \rho ,\nu \sigma }𝒴^1.`$ (4.22)
Replacing $`\mathrm{\Xi }^{(0)}`$ with the infinitesimal function $`\mathrm{\Lambda }^{(0)}`$ in $`𝒴`$, we obtain the infinitesimal transformation rule of $`\widehat{}_{\mu \rho ,\nu \sigma }`$ from (4.22):
$`\delta \widehat{}_{\mu \rho ,\nu \sigma }=i[\widehat{}_{\mu \rho ,\nu \sigma },\mathrm{\Lambda }^{(0)}].`$ (4.23)
As in (4.8), we express $`\widehat{}_{\mu \rho ,\nu \sigma }`$ as a linear combination, $`\widehat{}_{\mu \rho ,\nu \sigma }=_{p=0}^{\mathrm{}}\widehat{}_{\mu \rho ,\nu \sigma }^{(p)}`$, with
$`\widehat{}_{\mu \rho ,\nu \sigma }^{(p)}=\widehat{}_{\mu \rho ,\nu \sigma }^{Y(p)}+\widehat{}_{\mu \rho ,\nu \sigma }^{U(p)},`$ (4.24)
and
$`\widehat{}_{\mu \rho ,\nu \sigma }^{Y(p)}=_{\mu \rho }\widehat{𝒜}_{\nu \sigma }^{Y(p)}_{\nu \sigma }\widehat{𝒜}_{\mu \rho }^{Y(p)}+i{\displaystyle \underset{k=0}{\overset{p}{}}}[\widehat{𝒜}_{\mu \rho }^{Y(k)},\widehat{𝒜}_{\nu \sigma }^{Y(pk)}]^Y,`$ (4.25)
$`\widehat{}_{\mu \rho ,\nu \sigma }^{U(p)}=_{\mu \rho }\widehat{𝒜}_{\nu \sigma }^{U(p)}_{\nu \sigma }\widehat{𝒜}_{\mu \rho }^{U(p)}+i{\displaystyle \underset{k=0}{\overset{p}{}}}[\widehat{𝒜}_{\mu \rho }^{Y(k)},\widehat{𝒜}_{\nu \sigma }^{Y(pk)}]^U.`$ (4.26)
Here, $`\widehat{}_{\mu \rho ,\nu \sigma }^{(0)}`$ is identical with $`_{\mu \rho ,\nu \sigma }^{(0)}`$. From (4.23), we see that the infinitesimal transformation rules of $`\widehat{}_{\mu \rho ,\nu \sigma }^{(p)}(p=0,1,2,\mathrm{})`$ are given by $`\delta \widehat{}_{\mu \rho ,\nu \sigma }^{(p)}=i[\widehat{}_{\mu \rho ,\nu \sigma }^{(p)},\mathrm{\Lambda }^{(0)}]`$. Note that every $`\widehat{}_{\mu \rho ,\nu \sigma }^{(p)}(p=0,1,2,\mathrm{})`$ obeys the same transformation rule as $`\widehat{}_{\mu \rho ,\nu \sigma }`$. From $`\widehat{}_{\mu \rho ,\nu \sigma }^{(p)}`$, we can construct the modified field strengths of $`𝒜_{\mu \sigma }^{(p)}(p=1,2,3,\mathrm{})`$ in the manner discussed in Sec. 2.
We define the field strength of $`𝒜_{\mu \sigma }^{(p)}`$ in the same way as (2.14):
$`\widehat{}_{\mu \rho ,\nu \sigma }^{(p)}=\widehat{}_{\mu \rho ,\nu \sigma }^{Y(p)}+\widehat{}_{\mu \rho ,\nu \sigma }^{U(p)},`$ (4.27)
with
$`\widehat{}_{\mu \rho ,\nu \sigma }^{U(p)}`$ $``$ $`\widehat{}_{\mu \rho ,\nu \sigma }^{U(p)}+k{\displaystyle \frac{d\omega }{2\pi }x^\lambda (\omega )\mathrm{Tr}\left[𝒜_{\lambda \omega }^{Y(0)}\widehat{}_{\mu \rho ,\nu \sigma }^{Y(p)}\right]},`$ (4.28)
where $`\widehat{}_{\mu \rho ,\nu \sigma }^{U(p)}`$ becomes gauge-invariant owing to the condition (2.7).
Finally, we obtain the action for $`𝒜_{\mu \sigma }^{(p)}`$ as
$`S_\mathrm{R}^{(p)}={\displaystyle \frac{1}{V_R}}{\displaystyle [dx]\left(^{Y(p)}+^{U(p)}\right)\mathrm{exp}\left(\frac{L}{l^2}\right)},`$ (4.29)
with the Lagrangians
$`^{Y(p)}={\displaystyle \frac{1}{4}}N^Y𝒢^{\kappa \rho ,\lambda \sigma }𝒢^{\mu \chi ,\nu \omega }\mathrm{Tr}[\widehat{}_{\kappa \rho ,\mu \chi }^{Y(p)}\widehat{}_{\lambda \sigma ,\nu \omega }^{Y(p)}],`$ (4.30)
$`^{U(p)}={\displaystyle \frac{1}{4}}N^U𝒢^{\kappa \rho ,\lambda \sigma }𝒢^{\mu \chi ,\nu \omega }\widehat{}_{\kappa \rho ,\mu \chi }^{U(p)}\widehat{}_{\lambda \sigma ,\nu \omega }^{U(p)},`$ (4.31)
where the same damping factor is introduced as in (2.17). Obviously, the Lagrangians $`^{Y(p)}`$ and $`^{U(p)}`$ are gauge invariant. Also, the action $`S_\mathrm{R}^{(p=0)}`$ is identical with the action given in Sec. 3. In this way, the suitable action for $`𝒜_{\mu \sigma }^{(p)}`$ can be systematically derived using the formalism of the non-linear realization.
## 5 The Chapline-Manton coupling for higher rank tensor fields
Referring to the previous section, we now derive the local field theory with couplings among non-abelian tensor fields of second rank, local Yang-Mills fields and abelian tensor fields of third rank. Let us consider the next simplest solutions (2.7) and (2.8) consisting of local functions on $`M^D`$,
$`\mathrm{\Lambda }^{Y(1)}[x]`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}g_1Q^\mu (\sigma )\lambda _\mu {}_{}{}^{a}(x(\sigma ))T_a(\sigma ),`$ (5.1)
$`\mathrm{\Lambda }^{U(1)}[x]`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}q_1\sqrt{x^{\mathrm{\hspace{0.17em}2}}(\sigma )}Q^\mu (\sigma )Q^\nu (\sigma )\xi _{\mu \nu }(x(\sigma ))`$ (5.2)
$`+`$ $`{\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}e_1Q^\mu (\sigma )\kappa _\mu (x(\sigma )),`$
where $`\lambda _\mu {}_{}{}^{a}(x)`$ and $`\kappa _\mu (x)`$ are infinitesimal vector functions, and $`\xi _{\mu \nu }(x)`$ is an infinitesimal tensor function of second rank. The constants $`g_1`$ and $`e_1`$ are of dimension $`[\mathrm{length}]^{\frac{D4}{2}}`$, and $`q_1`$ is of dimension $`[\mathrm{length}]^{\frac{D6}{2}}`$. The infinitesimal functions (5.1) and (5.2) correspond to (4.1) and (4.2) with $`p=1`$, respectively. In this case, however, we can rewrite (5.2) in a more simple form as
$`\mathrm{\Lambda }^{U(1)}[x]={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}q_1\sqrt{x^{\mathrm{\hspace{0.17em}2}}(\sigma )}Q^\mu (\sigma )Q^\nu (\sigma )\lambda _{\mu \nu }(x(\sigma ))`$ (5.3)
with a redefinition of the infinitesimal tensor function: $`\lambda _{\mu \nu }(x)\xi _{\mu \nu }(x)(e_1/2q_1)_{(\mu }\kappa _{\nu )}(x)`$. In deriving (5.3), we carried out a partial integration over $`\sigma `$ for the second term on the right-hand side of (5.2). Corresponding to (5.1) and (5.3), we take the Yang-Mills field $`𝒜_{\mu \sigma }^{Y(1)}`$ and the U(1) gauge field $`𝒜_{\mu \sigma }^{U(1)}`$ consisting of local fields on $`M^D`$ as
$`𝒜_{\mu \sigma }^{Y(1)}[x]`$ $`=`$ $`g_1Q^\nu (\sigma )B_{\mu \nu }{}_{}{}^{a}(x(\sigma ))T_a(\sigma ){\displaystyle \frac{1}{2}}g_1Q_\mu (\sigma )Q^\nu (\sigma )Q^\kappa (\sigma )C_{\nu \kappa }{}_{}{}^{a}(x(\sigma ))T_a(\sigma )`$ (5.4)
$``$ $`\stackrel{~}{g}_1\varphi _\nu {}_{}{}^{a}(x(\sigma ))D_\sigma \left(\mathrm{\Pi }_\mu {}_{}{}^{\nu }(\sigma ){\displaystyle \frac{T_a(\sigma )}{\sqrt{x^2(\sigma )}}}\right),`$
$`\stackrel{~}{𝒜}_{\mu \sigma }^{U(1)}[x]`$ $`=`$ $`q_1\sqrt{x^2(\sigma )}Q^\nu (\sigma )Q^\kappa (\sigma )U_{\mu \nu \kappa }(x(\sigma ))`$ (5.5)
$`+`$ $`q_1\sqrt{x^2(\sigma )}Q_\mu (\sigma )Q^\nu (\sigma )Q^\kappa (\sigma )Q^\lambda (\sigma )V_{\nu \kappa \lambda }(x(\sigma ))`$
$``$ $`\stackrel{~}{q}_1Q^\nu (\sigma )\{B_{\mu \nu }(x(\sigma ))+C_{\mu \nu }(x(\sigma ))\}`$
$``$ $`\stackrel{~}{q}_1{\displaystyle \frac{d}{d\sigma }}\left(Q_\mu (\sigma )Q^\nu (\sigma )Q^\kappa (\sigma )\right)\varphi _{\nu \kappa }(x(\sigma )),`$
where $`\mathrm{\Pi }^\mu {}_{\nu }{}^{}(\sigma )\delta ^\mu {}_{\nu }{}^{}+Q^\mu (\sigma )Q_\nu (\sigma )`$. The differential $`D_\sigma `$ is defined as $`D_\sigma P_a(\sigma )dP_a(\sigma )/d\sigma +g_0x^\mu (\sigma )A_\mu {}_{}{}^{b}(x(\sigma ))f_{ba}{}_{}{}^{c}P_{c}^{}(\sigma )`$ with $`P_a(\sigma )`$ an arbitrary function of $`\sigma `$, where the $`A_\mu {}_{}{}^{a}(x)`$ are the local Yang-Mills fields. Here, $`U_{\mu \nu \lambda }(x)`$ is a tensor field of third rank with the symmetric property $`U_{\lambda \mu \nu }=U_{\lambda \nu \mu }`$, $`V_{\mu \nu \lambda }(x)`$ is a totally symmetric tensor field of third rank, $`B_{\mu \nu }{}_{}{}^{a}(x)`$ and $`B_{\mu \nu }(x)`$ are antisymmetric tensor fields of second rank, $`C_{\mu \nu }{}_{}{}^{a}(x)`$, $`C_{\mu \nu }(x)`$ and $`\varphi _{\mu \nu }(x)`$ are symmetric tensor fields of second rank, and $`\varphi _\mu {}_{}{}^{a}(x)`$ are vector fields. The constants $`\stackrel{~}{g}_1`$ and $`\stackrel{~}{q}_1`$ are of dimension $`[\mathrm{length}]^{\frac{D6}{2}}`$ and $`[\mathrm{length}]^{\frac{D2}{2}}`$, respectively.
From (4.7), the gauge-transformation rules of $`𝒜_{\mu \sigma }^{Y(1)}`$ and $`𝒜_{\mu \sigma }^{U(1)}`$ are given by
$`\delta 𝒜_{\mu \sigma }^{Y(1)}=_{\mu \sigma }\mathrm{\Lambda }^{Y(1)}+i[𝒜_{\mu \sigma }^{Y(1)},\mathrm{\Lambda }^{Y(0)}]^Y+i[𝒜_{\mu \sigma }^{Y(0)},\mathrm{\Lambda }^{Y(1)}]^Y,`$ (5.6)
$`\delta \stackrel{~}{𝒜}_{\mu \sigma }^{U(1)}=_{\mu \sigma }\mathrm{\Lambda }^{U(1)}+i[𝒜_{\mu \sigma }^{Y(1)},\mathrm{\Lambda }^{Y(0)}]^U+i[𝒜_{\mu \sigma }^{Y(0)},\mathrm{\Lambda }^{Y(1)}]^U.`$ (5.7)
Substituting (3.1), (3.3), (5.1) and (5.4) into (5.6), we obtain the infinitesimal transformation rules of the local fields $`B_{\mu \nu }{}_{}{}^{a}(x)`$, $`C_{\mu \nu }{}_{}{}^{a}(x)`$ and $`\varphi _\mu {}_{}{}^{a}(x)`$ as
$`\delta B_{\mu \nu }{}_{}{}^{a}(x)=D_{[\mu }\lambda _{\nu ]}{}_{}{}^{a}(x)g_0B_{\mu \nu }{}_{}{}^{b}(x)\lambda ^c(x)f_{bc}{}_{}{}^{a},`$ (5.8)
$`\delta C_{\mu \nu }{}_{}{}^{a}(x)=D_{(\mu }\lambda _{\nu )}{}_{}{}^{a}(x)g_0C_{\mu \nu }{}_{}{}^{b}(x)\lambda ^c(x)f_{bc}{}_{}{}^{a},`$ (5.9)
$`\delta \varphi _\mu {}_{}{}^{a}(x)=m_{g1}\lambda _\mu {}_{}{}^{a}(x)g_0\varphi _\mu {}_{}{}^{b}(x)\lambda ^c(x)f_{bc}{}_{}{}^{a},`$ (5.10)
where $`m_{g1}g_1/\stackrel{~}{g}_1`$ is a constant of dimension $`[\mathrm{length}]^1`$ and $`D_\mu `$ denotes the covariant derivative given by (3.5). We see that the local fields $`B_{\mu \nu }{}_{}{}^{a}(x)`$, $`C_{\mu \nu }{}_{}{}^{a}(x)`$ and $`\varphi _\mu {}_{}{}^{a}(x)`$ obey non-abelian gauge transformation rules . Similarly, substituting (3.1), (3.2), (3.3), (5.1), (5.3), (5.4) and (5.5) into (5.7) yields the infinitesimal transformation rules of the local fields $`U_{\mu \nu \lambda }(x)`$, $`V_{\mu \nu \lambda }(x)`$, $`B_{\mu \nu }(x)`$, $`C_{\mu \nu }(x)`$ and $`\varphi _{\mu \nu }(x)`$:
$`\delta B_{\mu \nu }(x)`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{k}_1}{2}}(m_{q1}A_{[\mu }{}_{}{}^{a}(x)\lambda _{\nu ],a}(x)\varphi _{[\mu }{}_{}{}^{a}(x)_{\nu ]}\lambda _a(x)),`$ (5.11)
$`\delta C_{\mu \nu }(x)`$ $`=`$ $`2m_{q1}\lambda _{\mu \nu }(x)+{\displaystyle \frac{\stackrel{~}{k}_1}{2}}(m_{q1}A_{(\mu }{}_{}{}^{a}(x)\lambda _{\nu ),a}(x)\varphi _{(\mu }{}_{}{}^{a}(x)_{\nu )}\lambda _a(x)),`$ (5.12)
$`\delta U_{\mu \nu \kappa }(x)`$ $`=`$ $`_\mu \lambda _{\nu \kappa }(x)_{(\nu }\lambda _{\kappa )\mu }(x)`$
$``$ $`{\displaystyle \frac{\stackrel{~}{k}_1}{2}}(A_\mu {}_{}{}^{a}(x)_{(\nu }\lambda _{\kappa ),a}(x)+B_{\mu (\nu }{}_{}{}^{a}_{\kappa )}^{}\lambda _a(x){\displaystyle \frac{1}{m_{g1}}}\varphi _\mu {}_{}{}^{a}(x)D_{(\nu }_{\kappa )}\lambda _a(x)),`$
$`\delta V_{\mu \nu \kappa }(x)`$ $`=`$ $`{\displaystyle \frac{1}{6}}_{(\mu }\lambda _{\nu \kappa )}(x)`$ (5.14)
$`+`$ $`{\displaystyle \frac{\stackrel{~}{k}_1}{12}}(C_{(\mu \nu }{}_{}{}^{a}_{\kappa )}^{}\lambda _a(x)+{\displaystyle \frac{2}{m_{g1}}}\varphi _{(\mu }{}_{}{}^{a}(x)D_\nu _{\kappa )}\lambda _a(x)),^)`$
$`\delta \varphi _{\mu \nu }(x)`$ $`=`$ $`m_{q1}\lambda _{\mu \nu }(x){\displaystyle \frac{\stackrel{~}{k}_1}{2}}{\displaystyle \frac{m_{q1}}{m_{g1}}}\varphi _{(\mu }{}_{}{}^{a}(x)_{\nu )}\lambda _a(x).`$ (5.15)
Here $`m_{q1}q_1/\stackrel{~}{q}_1`$ and $`\stackrel{~}{k}_1kg_0g_1/q_1`$ are constants of dimension $`[\mathrm{length}]^1`$ and $`[\mathrm{length}]`$, respectively. The local fields $`U_{\mu \nu \lambda }(x)`$, $`V_{\mu \nu \lambda }(x)`$, $`B_{\mu \nu }(x)`$, $`C_{\mu \nu }(x)`$ and $`\varphi _{\mu \nu }(x)`$ obey abelian gauge transformation rules. As we expected, the non-abelian local fields $`A_\mu ^a`$, $`B_{\mu \nu }^a`$, $`C_{\mu \nu }^a`$ and $`\varphi _\mu ^a`$ appear in the infinitesimal transformation rules of these abelian local fields.
As in the case of (3.8), we can find the gauge invariant tensor fields under the transformation rules (5.8)–(5.15) and (3.5). We obtain <sup>6</sup><sup>6</sup>footnotetext: <sup>)</sup> $`X_{(\mu }Y_\nu Z_{\lambda )}X_\mu Y_\nu Z_\lambda +X_\nu Y_\lambda Z_\nu +X_\lambda Y_\mu Z_\nu +X_\lambda Y_\nu Z_\mu +X_\nu Y_\mu Z_\lambda +X_\mu Y_\lambda Z_\nu `$ .
$`\stackrel{~}{B}_{\mu \nu }(x)B_{\mu \nu }(x){\displaystyle \frac{\stackrel{~}{k}_1}{2}}{\displaystyle \frac{m_{q1}}{m_{g1}}}A_{[\mu }{}_{}{}^{a}(x)\varphi _{\nu ],a}(x),`$ (5.16)
$`\stackrel{~}{C}_{\mu \nu }(x)C_{\mu \nu }(x)2\varphi _{\mu \nu }(x){\displaystyle \frac{\stackrel{~}{k}_1}{2}}{\displaystyle \frac{m_{q1}}{m_{g1}}}A_{(\mu }{}_{}{}^{a}(x)\varphi _{\nu ),a}(x),`$ (5.17)
$`\stackrel{~}{V}_{\mu \nu \lambda }(x)V_{\mu \nu \lambda }(x)S_{\mu \nu \lambda }(x){\displaystyle \frac{\stackrel{~}{k}_1}{12}}C_{(\mu \nu }{}_{}{}^{a}(x)A_{\lambda ),a}(x),`$ (5.18)
where $`S_{\mu \nu \lambda }(x)`$ is the irreducible component of $`U_{\mu \nu \lambda }(x)`$ given by $`S_{\mu \nu \lambda }(U_{\mu \nu \lambda }+U_{\nu \lambda \mu }+U_{\lambda \mu \nu })/3`$.
Eliminating $`B_{\mu \nu }`$, $`C_{\mu \nu }`$ and $`V_{\mu \nu \lambda }`$ from (5.5) by using (5.16), (5.17) and (5.18), we can rewrite (5.5) in terms of the abelian local fields $`\stackrel{~}{B}_{\mu \nu }`$, $`\stackrel{~}{C}_{\mu \nu }`$, $`\stackrel{~}{V}_{\mu \nu \lambda }`$, $`U_{\mu \nu \lambda }`$, $`S_{\mu \nu \lambda }`$ and $`\varphi _{\mu \nu }`$ and the non-abelian local fields $`C_{\mu \nu }^a`$, $`\varphi _\mu ^a`$ and $`A_\mu ^a`$. As in Sec. 3, we next remove the gauge invariant tensor fields $`\stackrel{~}{B}_{\mu \nu }`$, $`\stackrel{~}{C}_{\mu \nu }`$ and $`\stackrel{~}{V}_{\mu \nu \lambda }`$ from (5.5). Then (5.5) can be rewritten as
$`\stackrel{~}{𝒜}_{\mu \sigma }^{U(1)}[x]`$ $`=`$ $`𝒜_{\mu \sigma }^{U(1)}[x]\stackrel{~}{k}_1q_1{\displaystyle \frac{1}{m_{g1}}}Q^\nu (\sigma )A_\mu {}_{}{}^{a}(x(\sigma ))\varphi _{\nu ,a}(x(\sigma ))`$
$`+`$ $`{\displaystyle \frac{\stackrel{~}{k}_1}{2}}q_1\sqrt{x^2(\sigma )}Q_\mu (\sigma )Q^\nu (\sigma )Q^\kappa (\sigma )Q^\lambda (\sigma )C_{\nu \kappa }{}_{}{}^{a}(x(\sigma ))A_{\lambda ,a}(x(\sigma )),`$
with
$`𝒜_{\mu \sigma }^{U(1)}[x]`$ $`=`$ $`q_1\sqrt{x^2(\sigma )}\{\delta _\mu {}_{}{}^{[\nu }Q_{}^{\kappa ]}(\sigma )Q^\zeta (\sigma )\mathrm{\Pi }_\mu {}_{}{}^{\zeta }(\sigma )Q^\nu (\sigma )Q^\kappa (\sigma )\}A_{\nu \kappa \zeta }(x(\sigma ))`$
$``$ $`\stackrel{~}{q}_1\sqrt{x^2(\sigma )}\{Q_\mu ^{}(\sigma )Q^\nu (\sigma )Q^\kappa (\sigma )+2\mathrm{\Pi }_\mu {}_{}{}^{\nu }(\sigma )Q^\kappa (\sigma )\}\varphi _{\nu \kappa }(x(\sigma )),`$
where $`A_{\mu \nu \lambda }(x)`$ is a tensor field of third rank defined by $`A_{\mu \nu \lambda }(U_{\mu \nu \lambda }3S_{\mu \nu \lambda })/2`$ and has the symmetry property $`A_{\lambda \mu \nu }=A_{\lambda \nu \mu }`$. The infinitesimal transformation rule of $`A_{\mu \nu \lambda }(x)`$ is given by
$`\delta A_{\mu \nu \kappa }(x)`$ $`=`$ $`_\mu \lambda _{\nu \kappa }(x)`$ (5.21)
$`+`$ $`{\displaystyle \frac{\stackrel{~}{k}_1}{4}}(A_{(\nu }{}_{}{}^{a}(x)_{\kappa )}\lambda _{\mu ,a}(x)+_\mu \lambda _{(\nu }{}_{}{}^{a}(x)A_{\kappa ),a}(x)B_{\mu (\nu }{}_{}{}^{a}(x)_{\kappa )}\lambda _a(x).`$
$`.{\displaystyle \frac{1}{m_{g1}}}\varphi _{(\nu }{}_{}{}^{a}(x)D_{\kappa )}_\mu \lambda _a(x){\displaystyle \frac{1}{m_{g1}}}D_\mu _{(\nu }\lambda ^a(x)\varphi _{\kappa ),a}(x)).`$
The U(1) gauge field $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(1)}`$ is reduced to $`𝒜_{\mu \sigma }^{U(1)}`$ by setting $`k=0`$. As we have seen in Sec. 3, we can also check that $`\stackrel{~}{𝒜}_{\mu \sigma }^{U(1)}`$ does not satisfy the condition $`x^\mu (\sigma )\stackrel{~}{𝒜}_{\mu \sigma }^{U(1)}=0`$ if $`k0`$, while $`𝒜_{\mu \sigma }^{U(1)}`$, which is written in terms of the abelian local tensor fields, satisfies the condition $`x^\mu (\sigma )𝒜_{\mu \sigma }^{U(1)}=0`$.
Next, let us consider the scalar field $`\mathrm{\Phi }^{(1)}=\mathrm{\Phi }^{Y(1)}+\mathrm{\Phi }^{U(1)}`$, where $`\mathrm{\Phi }^{(1)}`$ corresponds to (4.16) with $`p=1`$. From (4.17), we see that $`\mathrm{\Phi }^{Y(1)}`$ and $`\mathrm{\Phi }^{U(1)}`$ obey the infinitesimal transformation rules
$`\delta \mathrm{\Phi }^{Y(1)}=\mathrm{\Lambda }^{Y(1)}+i[\mathrm{\Phi }^{Y(1)},\mathrm{\Lambda }^{Y(0)}]^Y,`$ (5.22)
$`\delta \mathrm{\Phi }^{U(1)}=\mathrm{\Lambda }^{U(1)}+i[\mathrm{\Phi }^{Y(1)},\mathrm{\Lambda }^{Y(0)}]^U.`$ (5.23)
We express the scalar fields $`\mathrm{\Phi }^{Y(1)}`$ and $`\mathrm{\Phi }^{U(1)}`$ in terms of the local fields on $`M^D`$. Since the infinitesimal scalar functions $`\mathrm{\Lambda }^{Y(1)}`$, $`\mathrm{\Lambda }^{U(1)}`$ and $`\mathrm{\Lambda }^{Y(0)}`$ are reparametrization invariant, $`\mathrm{\Phi }^{Y(1)}`$ and $`\mathrm{\Phi }^{U(1)}`$ also have to be reparametrization invariant. Taking account of this, we express $`\mathrm{\Phi }^{Y(1)}`$ and $`\mathrm{\Phi }^{U(1)}`$ in terms of the local fields $`\varphi _\mu ^a`$ and $`\varphi _{\mu \nu }`$, respectively, as
$`\mathrm{\Phi }^{Y(1)}[x]={\displaystyle _0^{2\pi }}{\displaystyle \frac{d\sigma }{2\pi }}\stackrel{~}{g}_1Q^\mu (\sigma )\varphi _\mu {}_{}{}^{a}(x(\sigma ))T_a(\sigma ),`$ (5.24)
$`\mathrm{\Phi }^{U(1)}[x]={\displaystyle \frac{d\sigma }{2\pi }\stackrel{~}{q}_1\sqrt{x^2(\sigma )}Q^\mu (\sigma )Q^\nu (\sigma )\varphi _{\mu \nu }(x(\sigma ))},`$ (5.25)
where $`\stackrel{~}{g}_1`$ and $`\stackrel{~}{q}_1`$ are the constants appearing in (5.4) and (5.5). The transformation rules (5.22) and (5.23) lead to infinitesimal transformation rules of $`\varphi _{\mu \nu }`$ and $`\varphi _\mu ^a`$. However, we find that the transformation rules of $`\varphi _\mu ^a`$ and $`\varphi _{\mu \nu }`$ are compatible with (5.10) and (5.15).
From (4.20), the vector fields $`\widehat{𝒜}_{\mu \sigma }^{Y(1)}`$ and $`\widehat{𝒜}_{\mu \sigma }^{U(1)}`$ are given by
$`\widehat{𝒜}_{\mu \sigma }^{Y(1)}𝒜_{\mu \sigma }^{Y(1)}_{\mu \sigma }\mathrm{\Phi }^{Y(1)}+i[𝒜_{\mu \sigma }^{Y(0)},\mathrm{\Phi }^{Y(1)}]^Y,`$ (5.26)
$`\widehat{𝒜}_{\mu \sigma }^{U(1)}\stackrel{~}{𝒜}_{\mu \sigma }^{U(1)}_{\mu \sigma }\mathrm{\Phi }^{U(1)}+i[𝒜_{\mu \sigma }^{Y(0)},\mathrm{\Phi }^{Y(1)}]^U.`$ (5.27)
Substituting (5.4), (5.24), and (3.3) into (5.26), we obtain
$`\widehat{𝒜}_{\mu \sigma }^{Y(1)}`$ $`=`$ $`g_1Q^\nu (\sigma )\widehat{B}_{\mu \nu }{}_{}{}^{a}(x(\sigma ))T_a(\sigma )`$ (5.28)
$``$ $`{\displaystyle \frac{1}{2}}g_1Q_\mu (\sigma )Q^\nu (\sigma )Q^\lambda (\sigma )\widehat{C}_{\nu \lambda }{}_{}{}^{a}(x(\sigma ))T_a(\sigma ),`$
with
$`\widehat{B}_{\mu \nu }{}_{}{}^{a}(x)`$ $``$ $`B_{\mu \nu }{}_{}{}^{a}(x){\displaystyle \frac{1}{m_{g1}}}D_{[\mu }\varphi _{\nu ]}{}_{}{}^{a}(x),`$ (5.29)
$`\widehat{C}_{\mu \nu }{}_{}{}^{a}(x)`$ $``$ $`C_{\mu \nu }{}_{}{}^{a}(x){\displaystyle \frac{1}{m_{g1}}}D_{(\mu }\varphi _{\nu )}{}_{}{}^{a}(x).`$ (5.30)
Under the transformation rules (5.8), (5.9), (5,10) and (3.5), $`\widehat{B}_{\mu \nu }^a`$ and $`\widehat{C}_{\mu \nu }^a`$ transform homogeneously:
$`\delta \widehat{B}_{\mu \nu }{}_{}{}^{a}(x)=g_0\widehat{B}_{\mu \nu }{}_{}{}^{b}(x)\lambda ^c(x)f_{bc}{}_{}{}^{a},`$ (5.31)
$`\delta \widehat{C}_{\mu \nu }{}_{}{}^{a}(x)=g_0\widehat{C}_{\mu \nu }{}_{}{}^{b}(x)\lambda ^c(x)f_{bc}{}_{}{}^{a}.`$ (5.32)
The transformation rules (5.31) and (5.32) are compatible with the gauge transformation rules $`\delta \widehat{𝒜}_{\mu \sigma }^{Y(1)}=i[\widehat{𝒜}_{\mu \sigma }^{Y(1)},\mathrm{\Lambda }^{Y(0)}]^Y`$. Also, substitution of (5,19), (5,24), (5.25) and (3.3) into (5.27) yields
$`\widehat{𝒜}_{\mu \sigma }^{U(1)}[x]`$ $`=`$ $`q_1\sqrt{x^2(\sigma )}Q^\nu (\sigma )Q^\lambda (\sigma )B_{\mu \nu \lambda }(x(\sigma ))`$ (5.33)
$``$ $`{\displaystyle \frac{1}{2}}q_1\sqrt{x^2(\sigma )}Q_\mu (\sigma )Q^\nu (\sigma )Q^\kappa (\sigma )Q^\lambda (\sigma )C_{\nu \kappa \lambda }(x(\sigma )).`$
with
$`B_{\mu \nu \lambda }(x)`$ $``$ $`\stackrel{~}{A}_{\mu \nu \lambda }(x)\stackrel{~}{A}_{(\nu \lambda )\mu }(x)+{\displaystyle \frac{\stackrel{~}{k}_1}{2m_{g1}}}A_\mu {}_{}{}^{a}(x)_{(\nu }\varphi _{\lambda ),a}(x),`$ (5.34)
$`C_{\mu \nu \lambda }(x)`$ $``$ $`{\displaystyle \frac{1}{3}}\stackrel{~}{A}_{(\mu \nu \lambda )}(x){\displaystyle \frac{\stackrel{~}{k}_1}{6}}C_{(\mu \nu }{}_{}{}^{a}(x)A_{\lambda ),a}(x),`$ (5.35)
and $`\stackrel{~}{A}_{\mu \nu \lambda }(x)A_{\mu \nu \lambda }(x)(1/m_{q1})_\mu \varphi _{\nu \lambda }(x)`$. Here, $`\stackrel{~}{A}_{\mu \nu \lambda }`$ is a third-rank tensor field that is gauge invariant under $`k=0`$. By setting $`k=0`$, the abelian tensor fields $`B_{\mu \nu \lambda }`$ and $`C_{\mu \nu \lambda }`$ are reduced to the (irreducible) components of $`\stackrel{~}{A}_{\mu \nu \lambda }`$. Under the transformation rules of (5.9), (5.10), (5.15), (5,21) and (3.5), we see that $`B_{\mu \nu \lambda }`$ and $`C_{\mu \nu \lambda }`$ obey simple transformation rules:
$`\delta B_{\mu \nu \kappa }(x)={\displaystyle \frac{\stackrel{~}{k}_1}{2}}\widehat{B}_{\mu (\nu }{}_{}{}^{a}(x)_{\kappa )}\lambda _a(x),`$ (5.36)
$`\delta C_{\mu \nu \kappa }(x)={\displaystyle \frac{\stackrel{~}{k}_1}{6}}\widehat{C}_{(\mu \nu }{}_{}{}^{a}(x)_{\kappa )}\lambda _a(x).`$ (5.37)
These transformation rules are also compatible with the gauge transformation rule $`\delta \widehat{𝒜}_{\mu \sigma }^{U(1)}=i[\widehat{𝒜}_{\mu \sigma }^{Y(1)},\mathrm{\Lambda }^{Y(0)}]^U`$.
We next express the field strength $`\widehat{}_{\mu \rho ,\nu \sigma }^{(1)}`$ in terms of the local fields, where $`\widehat{}_{\mu \rho ,\nu \sigma }^{(1)}`$ corresponds to (4.27) with $`p=1`$. The concrete forms of $`\widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)}`$ and $`\widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}`$ are given by
$`\widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)}`$ $``$ $`_{\mu \rho }\widehat{𝒜}_{\nu \sigma }^{Y(1)}_{\nu \sigma }\widehat{𝒜}_{\mu \rho }^{Y(1)}+i[𝒜_{\mu \rho }^{Y(0)},\widehat{𝒜}_{\nu \sigma }^{Y(1)}]^Y+i[\widehat{𝒜}_{\mu \rho }^{Y(1)},𝒜_{\nu \sigma }^{Y(0)}]^Y,`$
$`\widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}`$ $``$ $`_{\mu \rho }\widehat{𝒜}_{\nu \sigma }^{U(1)}_{\nu \sigma }\widehat{𝒜}_{\mu \rho }^{U(1)}+i[𝒜_{\mu \rho }^{Y(0)},\widehat{𝒜}_{\nu \sigma }^{Y(1)}]^U+i[\widehat{𝒜}_{\mu \rho }^{Y(1)},𝒜_{\nu \sigma }^{Y(0)}]^U`$ (5.39)
$`+`$ $`k{\displaystyle \frac{d\omega }{2\pi }x^\lambda (\omega )\mathrm{Tr}\left[𝒜_{\lambda \omega }^{Y(0)}\widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)}\right]}.`$
As was mentioned in Sec. 4, $`\widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)}`$ obeys the homogeneous gauge transformation rule $`\delta \widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)}=i[\widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)},\mathrm{\Lambda }^{Y(0)}]^Y`$, while $`\widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}`$ is gauge invariant: $`\delta \widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}=0`$. Substituting (5.28) and (3.3) into (5.38), we obtain $`\widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)}`$ in terms of the local fields:
$`\widehat{}_{\mu \rho ,\nu \sigma }^{Y(1)}[x]`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_1\delta (\rho \sigma )[.Q^\lambda (\rho )\widehat{H}_{\lambda \mu \nu }{}_{}{}^{a}(x(\rho ))T_a(\rho )+\widehat{B}_{\mu \nu }{}_{}{}^{a}(x(\rho ))D_\rho \left({\displaystyle \frac{T_a(\rho )}{\sqrt{x^2(\rho )}}}\right)`$ (5.40)
$`{\displaystyle \frac{1}{2}}\widehat{I}_{\lambda [\mu }{}_{}{}^{a}(x(\rho ))D_\rho \left(Q_{\nu ]}(\rho )Q^\lambda (\rho ){\displaystyle \frac{T_a(\rho )}{\sqrt{x^2(\rho )}}}\right)`$
$`{\displaystyle \frac{1}{2}}Q^\kappa (\rho )Q^\lambda (\rho )Q_{[\mu }(\rho )\widehat{J}_{\nu ]\kappa \lambda }{}_{}{}^{a}(x(\rho ))T_a(\rho ).]`$
$`{\displaystyle \frac{1}{2}}g_1\delta ^{}(\rho \sigma )[.\{Q^\lambda (\rho )Q_{(\mu }(\rho )\widehat{I}_{\nu )\lambda }{}_{}{}^{a}(x(\rho ))`$
$`(3Q_\mu (\rho )Q_\nu (\rho )+\eta _{\mu \nu })Q^\kappa (\rho )Q^\lambda (\rho )\widehat{C}_{\kappa \lambda }{}_{}{}^{a}(x(\rho ))\}{\displaystyle \frac{T_a(\rho )}{\sqrt{x^2(\rho )}}}.]`$
$`(\text{all of the above terms with }\mu \nu \text{ and }\rho \sigma \text{}),`$
with
$`\widehat{H}_{\lambda \mu \nu }{}_{}{}^{a}(x)`$ $``$ $`D_\lambda \widehat{B}_{\mu \nu }{}_{}{}^{a}(x)+D_\mu \widehat{B}_{\nu \lambda }{}_{}{}^{a}(x)+D_\nu \widehat{B}_{\lambda \mu }{}_{}{}^{a}(x),`$ (5.41)
$`\widehat{I}_{\mu \nu }{}_{}{}^{a}(x)`$ $``$ $`\widehat{B}_{\mu \nu }{}_{}{}^{a}(x)\widehat{C}_{\mu \nu }{}_{}{}^{a}(x),`$ (5.42)
$`\widehat{J}_{\lambda \mu \nu }{}_{}{}^{a}(x)`$ $``$ $`D_\mu \widehat{B}_{\lambda \nu }{}_{}{}^{a}(x)+D_\mu \widehat{C}_{\lambda \nu }{}_{}{}^{a}(x)+D_\lambda \widehat{C}_{\mu \nu }{}_{}{}^{a}(x).`$ (5.43)
It is obvious that $`\widehat{H}_{\lambda \mu \nu }{}_{}{}^{a}(x)`$, $`\widehat{I}_{\mu \nu }{}_{}{}^{a}(x)`$ and $`\widehat{J}_{\lambda \mu \nu }{}_{}{}^{a}(x)`$ transform homogeneously. We next substitute (5.28), (5.33), (5.40) and (3.3) into (5.39). After a little tedious calculation, we obtain $`\widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}`$ in terms of the local fields:
$`\widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}={\displaystyle \frac{1}{4}}q_1\delta (\rho \sigma )`$
$`\times [.\sqrt{x^2(\rho )}\{.2Q^\kappa (\rho )Q^\lambda (\rho )\delta _{[\dot{\mu }}{}_{}{}^{\zeta }+2Q^\zeta (\rho )Q^\lambda (\rho )\delta _{[\dot{\mu }}^\kappa `$
$`+Q^\kappa (\rho )Q^\lambda (\rho )Q^\zeta (\rho )Q_{[\dot{\mu }}(\rho ).\}_\zeta \stackrel{~}{B}_{\dot{\nu }]\kappa \lambda }(x(\rho ))`$
$`+\sqrt{x^2(\sigma )}Q^\lambda (\rho )Q^\zeta (\rho )\{.Q^\kappa (\rho )Q_{[\mu }(\rho )\delta _{\nu ]}^\omega `$
$`+{\displaystyle \frac{3}{2}}Q^\omega (\rho )Q_{[\mu }(\rho )\delta _{\nu ]}{}_{}{}^{\kappa }.\}_\omega \stackrel{~}{C}_{\kappa \lambda \zeta }(x(\rho ))`$
$`+2\stackrel{~}{k}_1\sqrt{x^2(\rho )}Q^\kappa (\rho )Q^\lambda (\rho )\widehat{B}_{\kappa [\mu }{}_{}{}^{a}(x(\rho ))F_{\nu ]\lambda ,a}(x(\rho ))`$
$`+\stackrel{~}{k}_1\sqrt{x^2(\rho )}Q^\kappa (\rho )Q^\lambda (\rho )Q^\zeta (\rho )Q_{[\dot{\mu }}(\rho )\widehat{C}_{\kappa \lambda }{}_{}{}^{a}(x(\rho ))F_{\dot{\nu }]\zeta ,a}(x(\rho ))`$
$`+{\displaystyle \frac{d}{d\rho }}\{.2Q^\lambda (\rho )\delta _{[\mu }{}_{}{}^{\kappa }+Q^\kappa (\rho )Q^\lambda (\rho )Q_{[\mu }(\rho ).\}\stackrel{~}{B}_{\nu ]\kappa \lambda }(x(\rho ))`$
$`+{\displaystyle \frac{3}{2}}{\displaystyle \frac{d}{d\rho }}\{.Q^\kappa (\rho )Q^\lambda (\rho )Q_{[\mu }(\rho )\delta _{\nu ]}{}_{}{}^{\zeta }.\}\stackrel{~}{C}_{\kappa \lambda \zeta }(x(\rho )).]`$
$`+{\displaystyle \frac{1}{2}}q_1\delta ^{}(\rho \sigma )`$
$`\times [.\{.2Q^\lambda (\rho )\delta _{(\mu }{}_{}{}^{\kappa }+Q^\kappa (\rho )Q^\lambda (\rho )Q_{(\mu }(\rho ).\}\stackrel{~}{B}_{\nu )\kappa \lambda }(x(\rho ))`$
$`+\{.(3Q_\mu (\rho )Q_\nu (\rho )+\eta _{\mu \nu })Q^\kappa (\rho )Q^\lambda (\rho )Q^\zeta (\rho )`$
$`+{\displaystyle \frac{3}{2}}Q^\kappa (\rho )Q^\lambda (\rho )Q_{(\mu }(\rho )\delta _{\nu )}{}_{}{}^{\zeta }.\}\stackrel{~}{C}_{\kappa \lambda \zeta }(x(\rho )).]^)`$
$`(\text{all of the above terms with }\mu \nu \text{ and }\rho \sigma \text{}),`$ (5.44)
with
$`\stackrel{~}{B}_{\mu \nu \lambda }(x)`$ $``$ $`B_{\mu \nu \lambda }(x)+{\displaystyle \frac{1}{2}}\stackrel{~}{k}_1\widehat{B}_{\mu (\nu }{}_{}{}^{a}(x)A_{\lambda ),a}(x),`$ (5.45)
$`\stackrel{~}{C}_{\mu \nu \lambda }(x)`$ $``$ $`C_{\mu \nu \lambda }(x)+{\displaystyle \frac{1}{6}}\stackrel{~}{k}_1\widehat{C}_{(\mu \nu }{}_{}{}^{a}(x)A_{\lambda ),a}(x).`$ (5.46)
Here $`F_{\mu \nu }^a`$ is the field strength of $`A_\mu ^a`$. Obviously, $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ are invariant under the transformation rules (5.31), (5.32), (5.36), (5.37) and (3.5). Besides $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ and their derivatives, the two distinctive terms $`\widehat{B}_{\mu \nu }{}_{}{}^{a}F_{\lambda \kappa ,a}^{}`$ and $`\widehat{C}_{\mu \nu }{}_{}{}^{a}F_{\lambda \kappa ,a}^{}`$ occur in (5.44). These terms are also invariant under the transformation rules of (5.31), (5.32) and (3.5). Consequently, $`\widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}`$ remains invariant under the transformation of the local fields. This invariance is compatible with the fact that $`\delta \widehat{}_{\mu \rho ,\nu \sigma }^{U(1)}=0`$. Note that the terms $`\widehat{B}_{\mu \nu }{}_{}{}^{a}F_{\lambda \kappa ,a}^{}`$ and $`\widehat{C}_{\mu \nu }{}_{}{}^{a}F_{\lambda \kappa ,a}^{}`$ take the form of products of the field strengths $`F_{\mu \nu }^a`$ and the non-abelian tensor fields $`\widehat{B}_{\mu \nu }^a`$ or $`\widehat{C}_{\mu \nu }^a`$. These terms take the same form as the (non-abelian) BF-term, except for a totally antisymmetric property. Accordingly, we refer to these terms as “BF-like terms” hereafter. We can regard the BF-like terms in (5.44) as a kind of generalization of the Chern-Simons terms $`\mathrm{\Omega }_{\mu \nu \lambda }`$ in (3.13) for the non-abelian tensor fields $`\widehat{B}_{\mu \nu }^a`$ and $`\widehat{C}_{\mu \nu }^a`$.
Finally, we consider the action $`S_\mathrm{R}^{(1)}`$ corresponding to (4.29) with $`p=1`$. We divide $`S_\mathrm{R}^{(1)}`$ into $`S_\mathrm{R}^{Y(1)}`$ and $`S_\mathrm{R}^{U(1)}`$ as
$`S_\mathrm{R}^{Y(1)}`$ $`=`$ $`{\displaystyle \frac{1}{V_R}}{\displaystyle [dx]^{Y(1)}\mathrm{exp}\left(\frac{L}{l^2}\right)},`$ (5.47)
$`S_\mathrm{R}^{U(1)}`$ $`=`$ $`{\displaystyle \frac{1}{V_R}}{\displaystyle [dx]^{U(1)}\mathrm{exp}\left(\frac{L}{l^2}\right)},`$ (5.48)
with the Lagrangians
$`^{Y(1)}={\displaystyle \frac{1}{4}}N^Y𝒢^{\kappa \rho ,\lambda \sigma }𝒢^{\mu \chi ,\nu \omega }\mathrm{Tr}[\widehat{}_{\kappa \rho ,\mu \chi }^{Y(1)}\widehat{}_{\lambda \sigma ,\nu \omega }^{Y(1)}],`$ (5.49)
$`^{U(1)}={\displaystyle \frac{1}{4}}N^U𝒢^{\kappa \rho ,\lambda \sigma }𝒢^{\mu \chi ,\nu \omega }\widehat{}_{\kappa \rho ,\mu \chi }^{U(1)}\widehat{}_{\lambda \sigma ,\nu \omega }^{U(1)}.`$ (5.50)
From the definition of (5.47), it is obvious that $`S_\mathrm{R}^{Y(1)}`$ does not include the central charge $`k`$. This means that the action $`S_\mathrm{R}^{Y(1)}`$ is not affected by the central extension of the gauge group. In other words, the action $`S_\mathrm{R}^{Y(1)}`$ is coincident with the action whose gauge group is the loop group. The action $`S_\mathrm{R}^{Y(1)}`$ written in terms of the local fields is derived in Ref. 14). It describes a massive tensor field theory as non-abelian Stückelberg formalism for the tensor fields $`\widehat{B}_{\mu \nu }^a`$ and $`\widehat{C}_{\mu \nu }^a`$. The local interactions among $`\widehat{B}_{\mu \nu }^a`$, $`\widehat{C}_{\mu \nu }^a`$ and the local Yang-Mills fields $`A_\mu ^a`$ are determined by the non-linear realization method developed for the loop gauge group. The central charge $`k`$ does not appear in these interactions.
Let us now derive the action $`S_\mathrm{R}^{U(1)}`$ written in terms of the local fields. As was done in Sec. 3, substituting (5.44) into (5.50), we obtain $`^{U(1)}`$ in terms of the local fields. Next, inserting the Lagrangian into (5.48), and carrying out the integrations over $`x^{\mu n}(n0)`$ for (5.48), we obtain the action $`S_\mathrm{R}^{U(1)}`$ written in terms of the local fields. The concrete form of $`S_\mathrm{R}^{U(1)}`$ is given by
$`S_\mathrm{R}^{U(1)}={\displaystyle }d^Dx[.{\displaystyle \frac{1}{4}}\{.R_{\mu \gamma \xi \eta }{}_{}{}^{\nu \zeta \kappa \lambda }_{}^{\mu }\stackrel{~}{B}^{\gamma \xi \eta }(x)_\nu \stackrel{~}{B}_{\zeta \kappa \lambda }(x)`$
$`+a_1R_{\xi \eta }{}_{}{}^{\kappa \lambda }(_\kappa \stackrel{~}{B}_{[\mu \nu ]\lambda }(x)+_{[\mu }\stackrel{~}{B}_{\nu ]\kappa \lambda }(x))`$
$`\times \left(^\xi \stackrel{~}{B}^{\mu \nu \eta }(x)+^\mu \stackrel{~}{B}^{\nu \xi \eta }(x)\right)`$
$`+4a_2R_{\xi \eta \gamma }{}_{}{}^{\mu \kappa \lambda }_{}^{\gamma }\stackrel{~}{B}^{\nu \xi \eta }(x)`$
$`\times \left(4_\kappa \stackrel{~}{B}_{[\mu \nu ]\lambda }(x)+4_{[\mu }\stackrel{~}{B}_{\nu ]\kappa \lambda }(x)+_\mu \stackrel{~}{B}_{\nu \kappa \lambda }(x)\right)`$
$`+a_3R_{\mu \gamma \xi \eta }{}_{}{}^{\nu \zeta \kappa \lambda }_{}^{\mu }\stackrel{~}{C}^{\gamma \xi \eta }(x)_\nu \stackrel{~}{C}_{\zeta \kappa \lambda }(x)`$
$`+a_2R_{\xi \eta \gamma }{}_{}{}^{\kappa \lambda \zeta }(2_\mu \stackrel{~}{C}_{\kappa \lambda \zeta }(x)3_\kappa \stackrel{~}{C}_{\mu \lambda \zeta }(x))`$
$`\times \left(2^\mu \stackrel{~}{C}^{\xi \eta \gamma }(x)3^\xi \stackrel{~}{C}^{\mu \eta \gamma }(x)\right)`$
$`2a_2R_{\xi \eta \gamma }{}_{}{}^{\mu \kappa \lambda }(2^\nu \stackrel{~}{C}^{\xi \eta \gamma }(x)3^\xi \stackrel{~}{C}^{\nu \eta \gamma }(x))`$
$`\times \left(2_\kappa \stackrel{~}{B}_{[\mu \nu ]\lambda }(x)+2_{[\mu }\stackrel{~}{B}_{\nu ]\kappa \lambda }(x)+_\mu \stackrel{~}{B}_{\nu \kappa \lambda }(x)\right)`$
$`2a_3R_{\mu \gamma \xi \eta }{}_{}{}^{\nu \zeta \kappa \lambda }_{\nu }^{}\stackrel{~}{B}_{\zeta \kappa \lambda }(x)^\mu \stackrel{~}{C}^{\gamma \xi \eta }(x)`$
$`+2\stackrel{~}{k}_1a_1R_{\xi \eta }{}_{}{}^{\kappa \lambda }\widehat{B}_{}^{\xi \mu }{}_{a}{}^{}(x)F^{\nu \eta ,a}(x)(_\kappa \stackrel{~}{B}_{[\mu \nu ]\lambda }(x)+_{[\mu }\stackrel{~}{B}_{\nu ]\kappa \lambda }(x))`$
$`4\stackrel{~}{k}_1a_2R_{\xi \eta \gamma }{}_{}{}^{\mu \kappa \lambda }\widehat{C}_{}^{\xi \eta }{}_{a}{}^{}(x)F^{\nu \gamma ,a}(x)`$
$`\times \left(2_\kappa \stackrel{~}{B}_{[\mu \nu ]\lambda }(x)+2_{[\mu }\stackrel{~}{B}_{\nu ]\kappa \lambda }(x)+_\mu \stackrel{~}{B}_{\nu \kappa \lambda }(x)\right)`$
$`+2\stackrel{~}{k}_1a_2R_{\xi \eta \gamma }{}_{}{}^{\kappa \lambda \zeta }\widehat{C}_{}^{\xi \eta }{}_{a}{}^{}(x)F^{\mu \gamma ,a}(x)`$
$`\times \left(2_\mu \stackrel{~}{C}_{\kappa \lambda \zeta }(x)3_\kappa \stackrel{~}{C}_{\mu \lambda \zeta }(x)\right)`$
$`+\stackrel{~}{k}_1{}_{}{}^{2}a_{1}^{}R_{\xi \gamma }{}_{}{}^{\kappa \lambda }\widehat{B}_{\kappa [\mu }^{}{}_{}{}^{a}(x)F_{\nu ]\lambda ,a}(x)\widehat{B}^{\xi \mu }{}_{b}{}^{}(x)F^{\nu \gamma ,b}(x)`$
$`+4\stackrel{~}{k}_1{}_{}{}^{2}a_{2}^{}R_{\xi \eta \gamma }{}_{}{}^{\kappa \lambda \zeta }\widehat{C}_{\kappa \lambda }^{}{}_{}{}^{a}(x)F_{\mu \zeta ,a}(x)\widehat{C}^{\xi \eta ,b}(x)F^{\mu \gamma }{}_{b}{}^{}(x).\}`$
$`+{\displaystyle \frac{1}{2}}m_{q1}{}_{}{}^{2}\{.\stackrel{~}{B}_{\mu \nu \kappa }(x)\stackrel{~}{B}^{\mu \nu \kappa }(x)+b_1\stackrel{~}{B}_{\mu \nu \kappa }(x)\stackrel{~}{B}^{\nu \mu \kappa }(x)`$
$`+b_2\stackrel{~}{B}_{\mu \kappa }{}_{}{}^{\kappa }(x)\stackrel{~}{B}^{\mu \lambda }{}_{\lambda }{}^{}(x)+b_3\stackrel{~}{B}^\kappa {}_{\kappa \mu }{}^{}(x)\stackrel{~}{B}^{\mu \lambda }{}_{\lambda }{}^{}(x)`$
$`+b_4\stackrel{~}{B}^\kappa {}_{\kappa \mu }{}^{}(x)\stackrel{~}{B}_\lambda {}_{}{}^{\lambda \mu }(x)+b_5\stackrel{~}{C}_{\mu \nu \kappa }(x)\stackrel{~}{C}^{\mu \nu \kappa }(x)`$
$`+b_6\stackrel{~}{C}^\kappa {}_{\kappa \mu }{}^{}(x)\stackrel{~}{C}^{\mu \lambda }{}_{\lambda }{}^{}(x)+b_7\stackrel{~}{B}_{\mu \nu \kappa }(x)\stackrel{~}{C}^{\mu \nu \kappa }(x)`$
$`+b_8\stackrel{~}{B}_{\mu \kappa }{}_{}{}^{\kappa }(x)\stackrel{~}{C}^{\mu \kappa }{}_{\kappa }{}^{}(x)+b_9\stackrel{~}{B}^\kappa {}_{\kappa \mu }{}^{}(x)\stackrel{~}{C}^{\mu \lambda }{}_{\lambda }{}^{}(x).\}.],`$ (5.51)
with
$`R_{\mu \nu }{}_{}{}^{\kappa \lambda }=\delta _{(\mu }{}_{}{}^{\kappa }\delta _{\nu )}^{}{}_{}{}^{\lambda }+\eta _{\mu \nu }\eta ^{\kappa \lambda },`$ (5.52a)
$`R_{\mu \nu \xi }{}_{}{}^{\kappa \lambda \gamma }=\delta _{(\mu }{}_{}{}^{\kappa }\delta _{\nu }^{}{}_{}{}^{\lambda }\delta _{\xi )}^{}{}_{}{}^{\gamma }+{\displaystyle \frac{1}{4}}\eta _{(\mu \nu }\eta ^{(\kappa \lambda }\delta _{\xi )}{}_{}{}^{\gamma )},`$ (5.52b)
$`R_{\mu \nu \xi \eta }{}_{}{}^{\kappa \lambda \gamma \zeta }=\delta _{(\mu }{}_{}{}^{\kappa }\delta _{\nu }^{}{}_{}{}^{\lambda }\delta _{\xi }^{}{}_{}{}^{\gamma }\delta _{\eta )}^{}{}_{}{}^{\zeta }+{\displaystyle \frac{1}{8}}\eta _{(\mu \nu }\eta ^{(\kappa \lambda }\delta _\xi {}_{}{}^{\gamma }\delta _{\eta )}^{}^{\zeta )}`$
$`+{\displaystyle \frac{1}{64}}\eta _{(\mu \nu }\eta _{\xi \eta )}\eta ^{(\kappa \lambda }\eta ^{\gamma \zeta )}.`$ (5.52c)
Here $`a_i(i=1,2,3)`$ and $`b_i(i=1,2,\mathrm{},9)`$ are constants given by
$`a_1=(D+4)(D+6),a_2={\displaystyle \frac{D+6}{16}},a_3={\displaystyle \frac{1}{16}},`$
$`b_1={\displaystyle \frac{1}{2}}J\{4K(2D1)(D+3)\delta (0)16(D1)(D^2+4D+1)\delta ^{\prime \prime }(0)\},`$
$`b_2={\displaystyle \frac{1}{2}}J\{K(3D^2+22D41)\delta (0)+4(D1)(3D+11)\delta ^{\prime \prime }(0)\},`$
$`b_3=J\{4K(D+1)^2\delta (0)+16(D1)(D+3)\delta ^{\prime \prime }(0)\},`$
$`b_4=J\{12K(D+1)\delta (0)+16(D1)\delta ^{\prime \prime }(0)\},`$
$`b_5={\displaystyle \frac{1}{2}}J\{27K(D+1)^2\delta (0)+12(D1)(D+21)\delta ^{\prime \prime }(0)\},`$
$`b_6={\displaystyle \frac{1}{4}}J\{9K(D^218D7)\delta (0)36(D1)(D13)\delta ^{\prime \prime }(0)\},`$
$`b_7=6J\{3K(D+1)^2\delta (0)2(D1)(3D+7)\delta ^{\prime \prime }(0)\},`$
$`b_8=3J\{K(D^2+10D3)\delta (0)+2(D1)^2\delta ^{\prime \prime }(0)\},`$
$`b_9=3J\{2K(D+1)(D5)\delta (0)4(D1)(3D+7)\delta ^{\prime \prime }(0)\},`$
with
$`K{\displaystyle \frac{1}{V_\mathrm{R}}}{\displaystyle \{dx\}_0^{2\pi }\frac{d\sigma }{2\pi }Q_\mu ^{}(\sigma )Q^\mu (\sigma )\mathrm{exp}\left(\frac{L}{l^2}\right)}(>0),`$
$`J^12K(4D^3+19D^220D15)\delta (0)8(D1)(2D^2+9D+5)\delta ^{\prime \prime }(0).`$
In deriving the action (5.51), we have set the free parameters $`k_u`$, $`q_1`$, and $`\stackrel{~}{q}_1`$ so as to satisfy the normalization conditions
$`{\displaystyle \frac{k_uq_1{}_{}{}^{2}l_{}^{2}}{4(D+2)(D+4)(D+6)}}\delta (0)^2=1,`$ (5.53a)
$`{\displaystyle \frac{k_u\stackrel{~}{q}_1^2}{D(D1)(D+2)(D+4)}}\times \{2K(4D^3+19D^220D15)\delta (0)`$
$`8(D1)(2D^2+9D+5)\delta ^{\prime \prime }(0)\}=1.`$ (5.53b)
The action (5.51) describes the massive tensor fields theory for $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ without spoiling the gauge invariance. This property is also possessed by Stückelberg formalism. Reflecting the non-abelian gauge theory, however, the action (5.51) includes non-trivial couplings via the BF-like terms $`\widehat{B}_{\mu \nu }{}_{}{}^{a}F_{\lambda \kappa ,a}^{}`$ and $`\widehat{C}_{\mu \nu }{}_{}{}^{a}F_{\lambda \kappa ,a}^{}`$. It is obvious that these couplings are due to the central extension of the gauge group. Indeed, all the interaction terms occurring in (5.51) include the central charge $`k`$. By setting $`k=0`$, we find that the gauge invariant tensor fields $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ reduce to the components of $`\stackrel{~}{A}_{\mu \nu \lambda }`$, and all the interaction terms occurring in (5.51) vanish. Hence, (5.51) becomes the action for the massive tensor field $`\stackrel{~}{A}_{\mu \nu \lambda }`$ without interactions. The gauge invariance still holds, because $`\stackrel{~}{A}_{\mu \nu \lambda }`$ is invariant under the transformation rules of (5.15) and (5.21) with $`k=0`$. Therefore, (5.51) results in the action of the Stückelberg formalism for the abelian tensor field of third rank $`\stackrel{~}{A}_{\mu \nu \lambda }`$. Consequently, we can regard (5.51) as the action of the “generalized” Stückelberg formalism for the tensor fields of third rank $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ in a broad sense.
We next comment on the interactions in (5.51). Although the types of interactions in (5.51) are somewhat complicated, we can find some features of the interactions. First, the abelian tensor fields $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ couple with the non-abelian tensor fields $`\widehat{B}_{\mu \nu }^a`$ and $`\widehat{C}_{\mu \nu }^a`$ and the local Yang-Mills fields $`A_\mu ^a`$ via the BF-like terms. (Here, $`\stackrel{~}{C}_{\mu \nu \lambda }`$ does not couple with $`\widehat{B}_{\mu \nu }{}_{}{}^{a}F_{\lambda \kappa ,a}^{}`$.) Second, the second power of the BF-like terms give couplings among the non-abelian fields $`\widehat{B}_{\mu \nu }^a`$, $`\widehat{C}_{\mu \nu }^a`$ and $`A_\mu ^a`$ that are obviously different from the minimal interactions resulting from the covariant derivative. We would like to emphasize that these features are analogous to those of the couplings in the action (3.19). Instead of the Chern-Simons term, the BF-like terms contribute to the non-trivial couplings in the action (5.51). We may regard the couplings between the abelian tensor fields and the BF-like terms as a kind of generalization of the Chapline-Manton coupling.
## 6 Conclusion and discussion
In this paper, we have considered the EYMT in loop space whose gauge group is the affine Lie group. In the EYMT, central extension of the gauge group leads to a coupling between the Yang-Mills fields $`𝒜_{\mu \sigma }^Y`$ and the U(1) gauge field $`𝒜_{\mu \sigma }^U`$. The coupling is different from the minimal coupling and the coupling via the Pauli terms existing in the Standard Model. The coupling yields non-trivial couplings between non-abelian local fields included in the Yang-Mills fields $`𝒜_{\mu \sigma }^Y`$ and an abelian local field included in the U(1) gauge field, $`𝒜_{\mu \sigma }^U`$. The Chapline-Manton coupling, which was originally introduced in order to combine a supergravity and a super Yang-Mills system, can be systematically derived within the framework of the Yang-Mills theory. In the supergravity theory, the Chapline-Manton coupling is derived using the local supersymmetry. It is interesting to study the relation of the central extension of the gauge group and local supersymmetry.
By using the formalism of the non-linear realization developed for the affine Lie gauge group, furthermore, we can derive the “generalized” Chapline-Manton coupling for higher-rank tensor fields. This coupling is given by the couplings among the local Yang-Mills fields $`A_\mu ^a`$, the non-abelian tensor fields of second rank $`\widehat{B}_{\mu \nu }^a`$ and $`\widehat{C}_{\mu \nu }^a`$, and the abelian tensor fields of third rank $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ via the BF-like terms. In the (bosonic) string theory, an abelian antisymmetric tensor field of second rank appears as massless excited states, while an abelian tensor field of third rank having the same symmetric property as $`\stackrel{~}{A}_{\mu \nu \lambda }`$ ($`\stackrel{~}{A}_{\mu \nu \lambda }=\stackrel{~}{A}_{\mu \lambda \nu }`$) appears as massive excited states. The Chapline-Manton coupling is realized in the type I supergravity theory. This theory is obtained as the low energy effective theory of the type I (or heterotic) superstring theory. The generalized Chapline-Manton coupling including the abelian tensor field of third rank might be realized in a massive mode in string theory.
If an abelian tensor field couples with the BF term, then it must have the totally antisymmetric property, because BF terms have this property. However, both the abelian tensor fields of third rank $`\stackrel{~}{B}_{\mu \nu \lambda }`$ and $`\stackrel{~}{C}_{\mu \nu \lambda }`$ have some specific symmetric properties. For this reason, these abelian tensor fields cannot couple with the BF term. Such a difficulty as this might be settled by considering the EYMT in closed p-manifold space $`\mathrm{\Omega }^pM^D`$, which is the configuration space for closed p-branes. A U(1) gauge field $`𝒜_{\mu \sigma }^{U(0)}[x]`$ on $`\mathrm{\Omega }^pM^D`$ consisting of an abelian local tensor fields is given by
$`𝒜_{\mu \stackrel{}{\sigma }}^{U(0)}[x]={\displaystyle \frac{q_0}{p!}}\mathrm{\Sigma }^{\nu _1\nu _2\mathrm{}\nu _p}(\stackrel{}{\sigma })B_{\mu \nu _1\nu _2\mathrm{}\nu _p}(x(\stackrel{}{\sigma })),`$ (6.1)
where $`B_{\mu \nu _1\nu _2\mathrm{}\nu _p}(x)`$ is an (abelian) totally antisymmetric tensor field of rank $`(p+1)`$ on $`M^D`$ and $`\mathrm{\Sigma }^{\nu _1\nu _2\mathrm{}\nu _p}(\stackrel{}{\sigma })x_1^{}{}_{}{}^{[\nu _1}(\stackrel{}{\sigma })x_2^{}{}_{}{}^{\nu _2}(\stackrel{}{\sigma })\mathrm{}x_p^{}{}_{}{}^{\nu _p]}(\stackrel{}{\sigma })`$. \[Here, $`\stackrel{}{\sigma }=(\sigma _1,\sigma _2,\mathrm{},\sigma _p)`$ represents the parameters describing a closed p-brane and $`x_n^\mu (\stackrel{}{\sigma })x^\mu (\stackrel{}{\sigma })/\sigma _n(1np)`$.\] We can indeed obtain the local field theory of $`B_{\mu \nu _1\nu _2\mathrm{}\nu _p}(x)`$ from the U(1) gauge theory in closed p-manifold space. In order to carry out a similar extension to the Yang-Mills theory, we have to find a suitable gauge group other than the affine Lie gauge group. It is conceivable that the suitable gauge group for the Yang-Mills theory in closed p-manifold space is the Mickelson-Faddeev group (and its generalization to higher dimensions). The commutation relations of the generators of the (generalized) Mickelson-Faddeev group are given by
$`[T_a(\stackrel{}{\rho }),T_b(\stackrel{}{\sigma })]=if_{ab}{}_{}{}^{c}T_{c}^{}(\stackrel{}{\sigma })\delta ^p(\stackrel{}{\rho }\stackrel{}{\sigma })`$
$`+kϵ^{j_1j_2\mathrm{}j_p}_{j_1}T_{(ab)}{}_{j_2\mathrm{}j_{p1}}{}^{}(\stackrel{}{\sigma })_{j_p}\delta ^{(p)}(\stackrel{}{\rho }\stackrel{}{\sigma }).`$ (6.2)
Setting $`p=1`$, we find that (6.2) coincides with (2.1). The commutation relations (6.2) are a natural extension of (2.1). We hope to discuss this subject in the future.
## Acknowledgements
I am grateful to S. Deguchi for valuable discussions and helpful suggestions. I would also like to thank A. Sugamoto for his careful reading of the manuscript. Thanks are also due to S. L. Shatashvili for information regarding Ref. 23)
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# 1 Introduction
## 1 Introduction
The production cross sections for all the processes at hadron-collider experiments are controlled by strong interaction physics and, hence, by its underlying field theory, QCD (see recent overviews in Refs. ). Studies of QCD at the Tevatron and the LHC have two main purposes . First, they are important to test the predictions of QCD, to measure its fundamental parameters (e.g. the strong coupling $`\alpha _\mathrm{S}`$) and to extract quantitative information on its non-perturbative dynamics (e.g. the distribution of partons in the proton). Second, they are relevant to a precise estimate of the background to other Standard Model processes and to signals of new physics.
This contribution is not a comprehensive review of QCD at high-energy hadron colliders. It is based on a selection of the topics presented in my introductory lectures at this Workshop. The selection highlights the QCD subjects that were most discussed during the Workshop and includes a pedagogical overview of some of the corresponding theoretical tools.
After the introduction of the general theoretical framework, I summarize in Sect. 2 the present knowledge on the parton densities and its impact on QCD predictions for hard-scattering processes at the Tevatron and the LHC. In Sect. 3, I then discuss some issues related to processes that are sensitive to the gluon density and, hence, to its determination. Section 4 presents a dictionary of different approaches (fixed-order expansions, resummed calculations, parton showers) to perturbative QCD calculations. The dictionary continues in Sect. 5, where I review soft-gluon resummation and discuss some recent phenomenological applications of threshold resummation to hadron collisions.
The QCD framework to describe any inclusive hard-scattering process,
$$h_1(p_1)+h_2(p_2)H(Q,\{\mathrm{}\})+X,$$
(1)
in hadron–hadron collisions is based on perturbation theory and on the factorization theorem of mass singularities. The corresponding cross section is computed by using the factorization formula
$`\sigma (p_1,p_2;Q,\{\mathrm{}\})`$ $`=`$ $`{\displaystyle \underset{a,b}{}}{\displaystyle _{x_{\mathrm{min}}}^1}𝑑x_1𝑑x_2f_{a/h_1}(x_1,\mu _F^2)f_{b/h_2}(x_2,\mu _F^2)\widehat{\sigma }_{ab}(x_1p_1,x_2p_2;Q,\{\mathrm{}\};\mu _F^2)`$ (2)
$`+`$ $`𝒪\left((\mathrm{\Lambda }_{QCD}/Q)^p\right).`$
The colliding hadrons $`h_1`$ and $`h_2`$ have momenta $`p_1`$ and $`p_2`$, $`H`$ denotes the triggered hard probe (vector bosons, jets, heavy quarks, Higgs bosons, SUSY particles and so on) and $`X`$ stands for any unobserved particle produced by the collision. The typical scale $`Q`$ of the scattering process is set by the invariant mass or the transverse momentum of the hard probe, and the notation $`\{\mathrm{}\}`$ stands for any other relevant scale and kinematic variable of the process. For instance, in the case of $`W`$ production we have $`Q=M_W`$ and $`\{\mathrm{}\}=\{Q_{},y,\mathrm{}\}`$, where $`M_W,Q_{}`$ and $`y`$ are the mass of the vector boson, its transverse momentum and its rapidity, respectively.
The factorization formula (2) involves the convolution of the partonic cross sections $`\widehat{\sigma }_{ab}`$ (where $`a,b=q,\overline{q},g)`$ and the parton distributions $`f_{a/h}(x,\mu _F^2)`$ of the colliding hadrons. If the hard probe $`H`$ is a hadron or a photon, the factorization formula has to include an additional convolution with the corresponding parton fragmentation function $`d_{a/H}(z,\mu _F^2)`$.
The term $`𝒪\left((\mathrm{\Lambda }_{QCD}/Q)^p\right)`$ on the right-hand side of Eq. (2) generically denotes non-perturbative contributions (hadronization effects, multiparton interactions, contributions of the soft underlying event, and so on). Provided the hard-scattering process (1) is sufficiently inclusive<sup>1</sup><sup>1</sup>1 More precisely, it has to be defined in an infrared- and collinear-safe manner., $`\widehat{\sigma }_{ab}`$ is computable as a power series expansion in $`\alpha _\mathrm{S}(Q^2)`$ and the non-perturbative contributions are (small) power-suppressed corrections (i.e. the power $`p`$ is positive) as long as the hard-scattering scale $`Q`$ is larger than few hundred MeV, the typical size of the QCD scale $`\mathrm{\Lambda }_{QCD}`$.
The parton densities $`f_{a/h}(x,\mu _F^2)`$ are phenomenological distributions that describe how partons are bounded in the colliding hadrons. Although they are not calculable in QCD perturbation theory, the parton densities are universal (process-independent) quantities. The scale $`\mu _F`$ is a factorization scale introduced in Eq. (2) to separate the bound-state effects from the perturbative interactions of the partons. The physical cross section $`\sigma (p_1,p_2;Q,\{\mathrm{}\})`$ does not depend on this arbitrary scale, but parton densities and partonic cross sections separately depend on $`\mu _F`$. In particular, higher-order contributions to $`\widehat{\sigma }_{ab}(x_1p_1,x_2p_2;Q,\{\mathrm{}\};\mu _F^2)`$ contain corrections of relative order $`(\alpha _\mathrm{S}(Q^2)\mathrm{ln}Q^2/\mu _F^2)^n`$. If $`\mu _F`$ is very different from $`Q`$, these corrections become large and spoil the reliability of the perturbative expansion. Thus, in practical applications of the factorization formula (2), the scale $`\mu _F`$ is set approximately equal to the hard scale $`Q`$ and variations of $`\mu _F`$ around this central value are used to estimate the uncertainty of the perturbative expansion.
The lower limit $`x_{\mathrm{min}}`$ of the integrations over the parton momentum fractions $`x_1`$ and $`x_2`$, as well as the values of $`x_1`$ and $`x_2`$ that dominate the convolution integral in Eq. (2), are controlled by the kinematics of the hard-scattering process. Typically we have $`x_{\mathrm{min}}>Q^2/S`$, where $`S=(p_1+p_2)^2`$ is the square of the centre-of-mass energy of the collision. If the hard probe is a state of invariant mass $`M`$ and rapidity $`y`$, the dominant values of the momentum fractions are $`x_{1,2}(Me^{\pm y})/\sqrt{S}`$ (see Fig. 1). Thus varying $`M`$ and $`y`$ at fixed $`\sqrt{S}`$, we are sensitive to partons with different momentum fractions. Increasing $`\sqrt{S}`$ the parton densities are probed in a kinematic range that extends towards larger values of $`Q`$ and smaller values of $`x_{1,2}`$.
## 2 Parton densities
The parton densities are an essential ingredient to study hard-scattering collisions. Once the partonic cross sections have been perturbatively computed, cross section measurements can be used to determine the parton densities. Then, they can in turn be used to predict cross sections for other hard-scattering processes.
The dependence of the parton densities<sup>2</sup><sup>2</sup>2In the following the parton densities of the proton $`f_{a/p}`$ are simply denoted by $`f_a`$ and those of the antiproton are obtained by using charge-conjugation invariance, i.e. $`f_{a/\overline{p}}=f_{\overline{a}/p}=f_{\overline{a}}`$. $`f_a(x,\mu ^2)`$ on the momentum fraction $`x`$ and their absolute value at any fixed scale $`\mu `$ are not computable in perturbation theory. However, the scale dependence is perturbatively controlled by the DGLAP evolution equation
$$\frac{df_a(x,\mu ^2)}{d\mathrm{ln}\mu ^2}=\underset{b}{}_x^1\frac{dz}{z}P_{ab}(\alpha _\mathrm{S}(\mu ^2),z)f_a(x/z,\mu ^2).$$
(3)
The kernels $`P_{ab}(\alpha _\mathrm{S},z)`$ are the Altarelli–Parisi (AP) splitting functions. As the partonic cross sections in Eq. (2), the AP splitting functions can be computed as a power series expansion in $`\alpha _\mathrm{S}`$:
$$P_{ab}(\alpha _\mathrm{S},z)=\alpha _\mathrm{S}P_{ab}^{(LO)}(z)+\alpha _\mathrm{S}^2P_{ab}^{(NLO)}(z)+\alpha _\mathrm{S}^3P_{ab}^{(NNLO)}(z)+𝒪(\alpha _\mathrm{S}^4).$$
(4)
The leading order (LO) and next-to-leading order (NLO) terms $`P_{ab}^{(LO)}(z)`$ and $`P_{ab}^{(NLO)}(z)`$ in the expansion are known . These first two terms are used in most of the QCD studies. Having determined $`f_a(x,Q_0^2)`$ at a given input scale $`\mu =Q_0`$, the evolution equation (3) can be used to compute the parton densities at different perturbative scales $`\mu `$ and larger values of $`x`$.
The parton densities are determined by performing global fits to data from deep-inelastic scattering (DIS), Drell–Yan (DY), prompt-photon and jet production. The method consists in parametrizing the parton densities at some input scale $`Q_0`$ and then adjusting the parameters to fit the data. The parameters are usually constrained by imposing the positivity of the parton densities $`(f_a(x,\mu ^2)0)`$ and the momentum sum rule $`(_a_0^1𝑑xxf_a(x,\mu ^2)=1)`$.
The present knowledge on the parton densities of the proton is reviewed in Refs. . Their typical behaviour is shown in Fig. 2. All densities decrease at large $`x`$. At small $`x`$ the valence quark densities vanish and the gluon density dominates. The sea-quark densities also increase at small $`x`$ because they are driven by the strong rise of the gluon density and the splitting of gluons in $`q\overline{q}`$ pairs. Note that the quark densities are not flavour-symmetric either in the valence sector $`(u_vd_v)`$ or in the sea sector $`(\overline{u}\overline{d})`$.
In addition to having the best estimate of the parton densities, it is important to quantify the corresponding uncertainty. This is a difficult issue. The uncertainty depends on the kinematic range in $`x`$ and $`Q^2`$. Moreover, it cannot be reliably estimated by simply comparing the parton densities obtained by different global fits. In fact, a lot of common systematic and common assumptions affect the global-fit procedures. Recent attempts to obtain parton densities with error bands that take into account correlations between experimental errors are described in Refs. . Some important theoretical uncertainties that are still to be understood are also discussed in Ref. .
The overall conclusion is that the quark densities<sup>3</sup><sup>3</sup>3Uncertainties on the determination of the quark densities at very high $`x`$ are discussed in Refs. . are reasonably well constrained and determined by DIS and DY processes, while the gluon density is certainly more uncertain . At small $`x`$ $`(x<10^3)`$, the gluon density $`f_g`$ is at present constrained by a single process, namely DIS at HERA. Thus, large higher-order corrections of the type $`(\alpha _\mathrm{S}\mathrm{ln}1/x)^n`$ could possibly affect the extraction of $`f_g`$. Assuming that $`f_g`$ is well determined at small $`x`$, the momentum sum rule reasonably constrains $`f_g`$ at intermediate values of $`x`$ $`(x10^2)`$. Jet production at the Tevatron at low to moderate values of the jet transverse energy $`E_T`$ can also be useful in constraining the gluon distribution in the range $`0.05<x<0.2`$. At large $`x`$ $`(x>10^1)`$, the most sensitive process to $`f_g`$ is prompt-photon production. Since, at present, prompt-photon data are not well described/predicted by perturbative QCD calculations, they cannot be used for a precise determination of $`f_g`$. Further discussion on these points is given in Sect. 3.
The conclusion that the gluon density is not well known can also be drawn by inspection (see Fig. 3) of the differences between the most updated analyses performed by the CTEQ Collaboration and the MRST group.
The differences between the MRST gluons and the CTEQ ones are due to the fact that the two groups used different data sets. The various gluon densities are very similar at small $`x`$, because in this region both groups used the HERA data. The MRST group includes prompt-photon data in the global fit: these data constrain the gluon directly at $`x>10^1`$ and indirectly (by the momentum sum rule) at $`x10^2`$. The CTEQ group does not use prompt-photon data, but it includes Tevatron data on the one-jet inclusive cross section. These data give a good constraint on $`f_g`$ in the region $`0.05<x<0.2`$.
There are also differences within the MRST and CTEQ sets. The various gluon densities of the MRST set correspond to different values of the non-perturbative transverse-momentum smearing that can be introduced to describe the differences among the prompt-photon data that are available at several centre-of-mass energies. The CTEQ5M and CTEQ5HJ gluons correspond to different assumptions on the parametrization of the functional form of $`f_g(x,Q_0^2)`$ at large $`x`$; the CTEQ5M set corresponds to the minimum-$`\chi ^2`$ solution of the fit while the CTEQ5HJ set (with a slightly higher $`\chi ^2`$) provides the best fit to the high-$`E_T`$ tail of the CDF and D0 jet cross sections.
This brief illustration shows that the differences in the most recent parton densities are mainly due to either inconsistencies between data sets and/or poor theoretical understanding of them. A more quantitative picture of the dependence on $`x`$ and $`Q^2`$ of the gluon density uncertainty is presented in Fig. 4.
We can see that the DIS and DY data sets weakly constrain $`f_g`$ for $`x>10^1`$. Since the AP splitting functions lead to negative scaling violation at large $`x`$, when $`f_g(x,Q^2)`$ is evolved at larger scales $`Q`$ according to Eq. (3) the gluon uncertainty is diluted: it propagates at smaller values of $`x`$ and its size is reduced at fixed $`x`$.
Figure 5 shows the typical predictions for hard-scattering cross sections at the Tevatron and the LHC, as obtained by using the parton densities of the MRST set. These predictions have to be supplemented with the corresponding uncertainties coming from the determination of the parton densities and from perturbative corrections beyond the NLO. Owing to the increased centre-of-mass energy and to QCD scaling violation (see Fig. 4), the kinematic region with small uncertainties is larger at the LHC than at the Tevatron.
For most of the QCD processes at the LHC, the uncertainty from the parton densities is smaller than $`\pm 10\%`$ and, in particular, it is smaller than the uncertainty from higher-order corrections. Some relevant exceptions are the single-jet, $`W/Z`$ and top quark cross sections. In the case of the single-jet inclusive cross section at high $`E_T`$ $`(E_T>2`$ TeV), the uncertainty from the poorly known gluon density at high $`x`$ is larger than that $`(\pm 10\%)`$ from higher-order corrections. The $`W`$ and $`Z`$ production cross sections are dominated by $`q\overline{q}`$ annihilation. Since the quark densities are well known, the ensuing uncertainty on the $`W/Z`$ cross section is small $`(\pm 5\%)`$. Nonetheless, in this case the uncertainty from higher-order corrections is even smaller, since the partonic cross sections for the DY process are known at the next-to-next-to-leading order (NNLO) in perturbation theory. In the case of top-quark production at the LHC, the gluon channel dominates and leads to an uncertainty of $`\pm 10\%`$ on the total cross section. Also for this process, however, the perturbative component is known beyond the NLO. Including all-order resummation of soft-gluon contributions , the estimated uncertainty from unknown higher-order corrections is approximately $`\pm 5\%`$ .
## 3 The gluon density issue
At present, the processes<sup>4</sup><sup>4</sup>4The rôle of jet production at the Tevatron has briefly been recalled in Sect. 2, and it is discussed in detail in Ref. . that are, in principle, most sensitive to the gluon density are DIS at HERA, $`b`$-quark production at the Tevatron, and prompt-photon production at fixed-target experiments. These processes constrain $`f_g`$ for $`x<10^3`$, $`x10^3`$$`10^2`$ and $`x>10^1`$, respectively. Nonetheless, the gluon density is, in practice, not well determined. The issue (or, perhaps, the puzzle) is that from a phenomenological viewpoint the standard theory, namely perturbative QCD at NLO, works pretty well for $`x<10^3`$ but not so well at larger values of $`x`$, while from theoretical arguments we should expect just the opposite to happen. This issue is discussed below mainly in its perturbative aspects. We should however keep it in mind that all these processes are dominated by hard-scattering scales $`Q`$ of the order of few GeV. Different types of non-perturbative contributions can thus be important.
From the study of DIS at HERA we can extract information on the gluon and sea-quark densities of the proton. The main steps in the QCD analysis of the structure functions at small values of the Bjorken variable $`x`$ are the following. The measurement of the proton structure function $`F_2(x,Q^2)q_S(x,Q^2)`$ directly determines the sea-quark density $`q_S=x(f_q+f_{\overline{q}})`$. Then, the DGLAP evolution equation (3) or, more precisely, the following equations (the symbol $``$ denotes the convolution integral with respect to $`x`$):
$`dF_2(x,Q^2)/d\mathrm{ln}Q^2`$ $``$ $`P_{qq}q_S+P_{qg}g,`$ (5)
$`dg(x,Q^2)/d\mathrm{ln}Q^2`$ $``$ $`P_{gq}q_S+P_{gg}g,`$ (6)
are used to extract a gluon density $`g(x,Q^2)=xf_g(x,Q^2)`$ that agrees with the measured scaling violation in $`dF_2(x,Q^2)/d\mathrm{ln}Q^2`$ (according to Eq. (5)) and fulfils the self-consistency equation (6).
The perturbative-QCD ingredients in this analysis are the AP splitting functions $`P_{ab}(\alpha _\mathrm{S},x)`$. Once they are known (and only then), the non-perturbative gluon density can be determined.
The standard perturbative-QCD framework to extract $`g(x,Q^2)`$ consists in using the truncation of the AP splitting functions at the NLO. This approach has been extensively compared with structure function data over the last few years and it gives a good description of the HERA data, down to low values of $`Q^22\mathrm{GeV}^2`$. The NLO QCD fits simply require a slightly steep input gluon density at these low momentum scales. Typically , we have $`g(x,Q_0^2)x^\lambda `$, with $`\lambda 0.2`$ at $`Q_0^22\mathrm{GeV}^2`$, and the data constrain $`g(x,Q_0^2)`$ with an uncertainty of approximately $`\pm 20\%`$.
Although it is phenomenologically successful, the NLO approach is not fully satisfactory from a theoretical viewpoint. The truncation of the splitting functions at a fixed perturbative order is equivalent to assuming that the dominant dynamical mechanism leading to scaling violations is the evolution of parton cascades with strongly-ordered transverse momenta. However, at high energy this evolution takes place over large rapidity intervals $`(\mathrm{\Delta }y\mathrm{ln}1/x)`$ and diffusion in transverse momentum becomes relevant. Formally, this implies that higher-order corrections to $`P_{ab}(\alpha _\mathrm{S},x)`$ are logarithmically enhanced:
$$P_{ab}(\alpha _\mathrm{S},x)\frac{\alpha _\mathrm{S}}{x}+\frac{\alpha _\mathrm{S}}{x}(\alpha _\mathrm{S}\mathrm{ln}x)+\mathrm{}+\frac{\alpha _\mathrm{S}}{x}(\alpha _\mathrm{S}\mathrm{ln}x)^n+\mathrm{}.$$
(7)
At asymptotically small values of $`x`$, resummation of these corrections is mandatory to obtain reliable predictions.
Small-$`x`$ resummation is, in general, accomplished by the BFKL equation . In the context of structure-function calculations, the BFKL equation provides us with improved expressions of the AP splitting functions $`P_{ab}(\alpha _\mathrm{S},x)`$, in which the leading logarithmic (LL) terms $`(\alpha _\mathrm{S}\mathrm{ln}x)^n`$, the next-to-leading logarithmic (NLL) terms $`\alpha _\mathrm{S}(\alpha _\mathrm{S}\mathrm{ln}x)^n`$, and so forth, are systematically summed to all orders $`n`$ in $`\alpha _\mathrm{S}`$. The present theoretical status of small-$`x`$ resummation is discussed in Ref. . Since in the small-$`x`$ region the gluon channel dominates, only the gluon splitting functions $`P_{gg}`$ and $`P_{gq}`$ contain LL contributions. These are known to be positive but numerically smaller than naively expected (the approach to the asymptotic regime is much delayed by cancellations of logarithmic corrections that occur at the first perturbative orders in $`P_{gg}`$ and $`P_{gq}`$). The NLL terms in the quark splitting functions $`P_{qg}`$ and $`P_{qq}`$ are known and turn out to be positive and large. A very important progress is the recent calculation of the NLL terms in $`P_{gg}`$, which are found to be negative and large. The complete NLL terms in $`P_{gq}`$ are still unknown.
The results of Refs. , the large size of the NLL terms and the alternating sign (from the LL to the NLL order and from the gluon to the quark channel) of the resummed small-$`x`$ contributions have prompted a lot of activity (see the list of references in Ref. ) on the conceptual basis and the phenomenological implications of small-$`x`$ resummation. This activity is still in progress and definite quantitative conclusions on the impact of small-$`x`$ resummation at HERA cannot be drawn yet.
At the same time, the capability of the fixed-order approach to produce a good description of the proton structure function $`F_2(x,Q^2)`$ at HERA cannot be used to conclude that the small-$`x`$ behaviour of the gluon density is certainly well determined. In fact, by comparing LO and NLO results, we could argue that the ensuing theoretical uncertainty on $`f_g`$ is sizeable . Going from LO to NLO, we can obtain stable predictions for $`F_2`$, but we have to vary the gluon density a lot. As shown in Fig. 6, the NLO gluon density sizeably differs from its LO parametrization, not only in absolute normalization but also in $`x`$-shape. For instance, at $`x=10^4`$ and $`Q^2=5\mathrm{GeV}^2`$ the NLO gluon is a factor of 2 smaller than the LO gluon. This can be understood from the fact that the scaling violation of $`F_2`$ is produced by the convolution $`P_{qg}g`$ (see the right-hand side of Eq. (5)). The quark splitting function $`P_{qg}`$ behaves as
$$P_{qg}(\alpha _\mathrm{S},x)\alpha _\mathrm{S}P_{qg}^{(LO)}(x)\left[1+2.2\frac{C_A\alpha _\mathrm{S}}{\pi }\frac{1}{x}+\mathrm{}\right],$$
(8)
where the LO term $`P_{qg}^{(LO)}(x)`$ is flat at small $`x`$, whereas the NLO correction is steep. To obtain a stable evolution of $`F_2`$, the NLO steepness of $`P_{qg}`$ has to be compensated by a gluon density that is less steep at NLO than at LO. This has to be kept in mind when concluding on the importance of small-$`x`$ resummation because the NLO steepness of $`P_{qg}`$ is the lowest-order manifestation of BFKL dynamics in the quark channel.
In the large-$`x`$ region, there is a well-known correlation between $`\alpha _\mathrm{S}`$ and $`f_g`$. At small $`x`$, there is an analogous strong correlation between the $`x`$-shapes of $`P_{qg}`$ and $`f_g`$. In the fixed-order QCD analysis of $`F_2`$, large NLO perturbative corrections at small $`x`$ can be balanced by the extreme flexibility of parton density parametrizations. It is difficult to disentangle this correlation between process-dependent perturbative contributions and non-perturbative parton densities from the study of a single quantity, as in the case of $`F_2`$ at HERA. The uncertainty on the gluon density at small $`x`$, as estimated from the NLO QCD fits of the HERA data, is evidently only a lower limit on the actual uncertainty on $`f_g`$.
The production of $`b`$ quarks at the Tevatron is also sensitive to the gluon density at relatively small values of $`x`$. The comparison between Tevatron data and perturbative-QCD predictions at NLO is shown in Fig. 7. Using standard sets of parton densities, the theoretical predictions typically underestimate the measured cross section by a factor of 2. This certainly is disappointing, although justifiable by the large theoretical uncertainty of the perturbative calculation . A lower limit on this uncertainty can be estimated by studying the scale dependence and the convergence of the perturbative expansion. Varying the factorization and renormalization scales by a factor of four around the $`b`$-quark mass $`m_b`$, the NLO cross section varies by a factor of almost 2 at the Tevatron and by a factor of 4–5 at the LHC . Similar factors are obtained by considering the ratio of the NLO and LO cross sections.
The present theoretical predictions for $`b`$-quark production at hadron colliders certainly need to be improved . Since the hard scale $`Qm_b`$ is not very large, a possible improvement regards estimates of non-perturbative contributions (for instance, effects of the fragmentation of the $`b`$-quark and of the intrinsic transverse momentum of the colliding partons). As for the evaluation of perturbative contributions at higher orders, the resummation of logarithmic terms of the type $`\alpha _\mathrm{S}^n\mathrm{ln}^n(p_t/m_b)`$ is important when the transverse momentum $`p_t`$ of the $`b`$ quark is much larger than $`m_b`$. The resummation of small-$`x`$ logarithmic contributions $`\alpha _\mathrm{S}^n\mathrm{ln}^nx`$ can also be relevant, because $`x2m_b/\sqrt{S}`$ is as small as $`10^3`$ at the Tevatron and as $`10^4`$ at the LHC. The theoretical tool to perform this resummation, namely the $`k_{}`$-factorization approach , is available. Updated phenomenological studies based on this tool and on the information from small-$`x`$ DIS at HERA would be interesting.
Prompt-photon production at fixed-target experiments is sensitive to the behaviour of the gluon density at large $`x`$ $`(x>0.1)`$. The theoretical predictions for this process, however, are not very accurate. Figure 8 shows the factorization- and renormalization-scale dependence of the perturbative cross section for the case of the E706 kinematics. If the scale is varied by a factor of 4 around the transverse energy $`E_T`$ of the prompt photon, the LO cross section varies by a factor of almost 4. Going to NLO the situation improves, but not very much, because the NLO cross section still varies by a factor of about 2.
A detailed comparison between NLO QCD calculations and data from the ISR and fixed-target experiments has recently been performed in Ref. . As shown in Fig. 9, the overall agreement with the theory is not satisfactory, even taking into account the uncertainty coming from scale variations in the theoretical predictions. Modifications of the gluon density can improve the agreement with some data sets only at the expense of having larger disagreement with other data sets. The differences between experiments at similar centre-of-mass energies (see, for instance, E706 pBe/530 at $`\sqrt{S}=31.6\mathrm{GeV}`$ and WA70 pp at $`\sqrt{S}=23\mathrm{GeV})`$ are much larger than expected from perturbative scaling violations. This can possibly suggest inconsistencies of experimental origin.
Another (not necessarily alternative) origin of the differences between data and theory could be the presence of non-perturbative effects that are not included in the NLO perturbative calculation. This explanation has been put forward in Refs. by introducing some amount of intrinsic<sup>5</sup><sup>5</sup>5To be precise, in Ref. the $`k_{}`$ of the colliding partons is not called ‘intrinsic’, but it is more generically called the $`k_{}`$ ‘from initial-state soft-gluon radiation’. transverse momentum $`k_{}`$ of the colliding partons. Owing to the steeply falling $`E_T`$ distribution $`(d\sigma /dE_T1/E_T^7)`$ of the prompt photon, even a small transverse-momentum kick<sup>6</sup><sup>6</sup>6The $`E_T`$ distribution of the single-photon is not calculable down to $`E_T=0`$ or, in other words, $`d\sigma /dE_T`$ is not integrable in the entire kinematic range of $`E_T`$. Thus, the intrinsic $`k_{}`$ of the incoming partons does not simply produce a shift of events from the low-$`E_T`$ to the high-$`E_T`$ region. For this reason, the terminology ‘$`k_{}`$ kick’ seems to be more appropriate than ‘$`k_{}`$ smearing’. can indeed produce a large effect on the cross section, in particular, at small values of $`E_T`$. Phenomenological investigations show that this additional $`k_{}`$ kick can lead to a better agreement between calculations and data. The E706 data suggest the value $`k_{}1.2\mathrm{GeV}`$, the WA70 data prefer no $`k_{}`$, and the UA6 data in the intermediate range of centre-of-mass energy $`(\sqrt{S}=24.3\mathrm{GeV}`$) may prefer an intermediate value of $`k_{}`$. Similar conclusions are obtained in the analysis by the MRST group .
A precise physical understanding of $`k_{}`$ effects is still missing. On one side, since the amount of $`k_{}`$ suggested by prompt-photon data varies with $`\sqrt{S}`$, it is difficult to argue that the transverse momentum is really ‘intrinsic’ and has an entirely non-perturbative origin. On the other side, in the case of the inclusive production of a single photon, a similar effect cannot be justified by higher-order logarithmic corrections produced by perturbative soft-gluon radiation (see Sect. 5). A lot of model-dependent assumptions (and ensuing uncertainties) certainly enter in the present implementations of the $`k_{}`$ kick. A general framework to consistently include non-perturbative transverse-momentum effects in perturbative calculations is not yet available. Recent proposals with this aim are presented in Refs. and .
Further studies on the consistency between different prompt-photon experiments and on the issue of intrinsic-$`k_{}`$ effects in hadron–hadron collisions are necessary. Owing to the present theoretical (and, possibly, experimental) uncertainties, it is difficult to use prompt-photon data to accurately determine the gluon density at large $`x`$. Other recent theoretical improvements, such as soft-gluon resummation, of the perturbative calculations for prompt-photon production at large $`x_T=2E_T/\sqrt{S}`$ are discussed in Sect. 5.
Studies of other single-particle inclusive cross sections, such as $`\pi ^0`$ cross sections , can be valuable to constrain the parton densities and could possibly help to clarify some of the experimental and theoretical issues arisen by prompt-photon production.
## 4 Partonic cross sections: fixed-order expansions, <br>resummed calculations, parton showers
The calculation of hard-scattering cross sections according to the factorization formula (2) requires the knowledge of the partonic cross sections $`\widehat{\sigma }`$, besides that of the parton densities. The partonic cross sections are usually computed by truncating their perturbative expansion at a fixed order in $`\alpha _\mathrm{S}`$:
$`\widehat{\sigma }(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _F^2)`$ $`=`$ $`\alpha _\mathrm{S}^k(\mu _R^2)\{\widehat{\sigma }^{(LO)}(p_1,p_2;Q,\{Q_1,\mathrm{}\})`$
$`+\alpha _\mathrm{S}(\mu _R^2)\widehat{\sigma }^{(NLO)}(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _R^2;\mu _F^2)`$
$`+\alpha _\mathrm{S}^2(\mu _R^2)\widehat{\sigma }^{(NNLO)}(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _R^2;\mu _F^2)+\mathrm{}\}.`$
The scale $`\mu _R`$ is the arbitrary renormalization scale introduced to define the perturbative expansion. Although the ‘exact’ partonic cross section on the left-hand side of Eq. (4) does not depend on $`\mu _R`$, each term on the right-hand side (and, hence, any fixed-order truncation) separately depends on it.
The LO (or tree-level) term $`\widehat{\sigma }^{(LO)}`$ gives only an estimate of the order of magnitude of the partonic cross section, because at this order $`\alpha _\mathrm{S}`$ is not unambiguously defined. Equivalently, we can say that since $`\widehat{\sigma }^{(LO)}`$ does not depend on $`\mu _R`$, the size of its contribution can be varied quite arbitrarily by changing $`\mu _R`$ in its coefficient $`\alpha _\mathrm{S}^k(\mu _R^2)`$. The strong coupling $`\alpha _\mathrm{S}`$ can be precisely defined only starting from NLO. A ‘reliable’ estimate of the central value of $`\widehat{\sigma }`$ thus requires the knowledge of (at least) the NLO term $`\widehat{\sigma }^{(NLO)}`$. This term explicitly depends on $`\mu _R`$ and this dependence begins to compensate that of $`\alpha _\mathrm{S}(\mu _R^2)`$.
In general, the $`n`$-th term in the curly bracket of Eq. (4) contains contributions of the type $`(\alpha _\mathrm{S}(\mu _R^2)\mathrm{ln}Q/\mu _R)^n`$. If $`\mu _R`$ is very different from the hard scale $`Q`$, these contributions become large and spoil the reliability of the truncated expansion (4). Thus, in practical applications the scale $`\mu _R`$ should be set approximately equal to the hard scale $`Q`$. As mentioned in Sect. 3, variations of $`\mu _R`$ around this central value are typically used to set a lower limit on the theoretical uncertainty of the perturbative calculation.
A better estimate of the accuracy of any perturbative expansion is obtained by considering the effect of removing the last perturbative term that has been computed. Since $`\alpha _\mathrm{S}`$ can be precisely defined only at NLO, this procedure can consistently be applied to Eq. (4) only as from its NNLO term. A ‘reliable’ estimate of the theoretical error on $`\widehat{\sigma }`$ thus requires the knowledge of the NNLO term $`\widehat{\sigma }^{(NNLO)}`$ in Eq. (4).
The LO and NLO approximations of $`\widehat{\sigma }`$ are used at present in (most of) the fixed-order QCD calculations. Prospects towards NNLO calculations of partonic cross sections and AP splitting functions are reviewed in Refs. .
The fixed-order expansion (4) provides us with a well-defined and systematic framework to compute the partonic cross section $`\widehat{\sigma }(p_1,p_2;Q,\{Q_1,\mathrm{}\};\mu _F^2)`$ of any hard-scattering process that is sufficiently inclusive or, more precisely, that is defined in an infrared- and collinear-safe manner. However, the fixed-order expansion is reliable only when all the kinematical scales $`Q,\{Q_1,\mathrm{}\}`$ are of the same order of magnitude. When the hard-scattering process involves two (or several) very different scales, say $`Q_1Q`$, the $`\mathrm{N}^n\mathrm{LO}`$ term in Eq. (4) can contain double- and single-logarithmic contributions of the type $`(\alpha _\mathrm{S}L^2)^n`$ and $`(\alpha _\mathrm{S}L)^n`$ with $`L=\mathrm{ln}(Q_1/Q)1`$. These terms spoil the reliability of the fixed-order expansion and have to be summed to all orders by systematically defining logarithmic expansions (resummed calculations).
Typical large logarithms, $`L=\mathrm{ln}Q/Q_0`$, are those related to the evolution of the parton densities from a low input scale $`Q_0`$ to the hard-scattering scale $`Q`$. These logarithms are produced by collinear radiation from the colliding partons and give single-logarithmic contributions. They never explicitly appear in the calculation of the partonic cross section, because they are systematically (LO, NLO and so forth) resummed in the evolved parton densities $`f(x,Q^2)`$ by using the DGLAP equation (3).
Different large logarithms, $`L=\mathrm{ln}Q/\sqrt{S}`$, appear when the centre-of-mass energy $`\sqrt{S}`$ of the collision is much larger than the hard scale $`Q`$. These small-$`x`$ $`(x=Q/\sqrt{S})`$ logarithms are produced by multiple radiation over the wide rapidity range that is available at large energy. They usually give single-logarithmic contributions that can be resummed by using the BFKL equation. BFKL resummation is relevant to DIS structure functions at small values of the Bjorken variable $`x`$ (see Sect. 3) and it can also be important at the LHC for the production of $`b`$ quarks and of prompt photons at relatively low $`E_T`$.
Another class of large logarithms is associated to the bremsstrahlung spectrum of soft gluons. Since soft gluons can be radiated collinearly, they give rise to double-logarithmic contributions to the partonic cross section:
$$\widehat{\sigma }\alpha _\mathrm{S}^k\widehat{\sigma }^{(LO)}\left\{1+\underset{n=1}{\overset{\mathrm{}}{}}\alpha _\mathrm{S}^n\left(C_{2n}^{(n)}L^{2n}+C_{2n1}^{(n)}L^{2n1}+C_{2n2}^{(n)}L^{2n2}+\mathrm{}\right)\right\}.$$
(10)
Soft-gluon resummation is discussed in Sect. 5.
A related approach to evaluate higher-order contributions to the partonic cross sections is based on Monte Carlo parton showers (see and the updated list of references in ). Rather than computing exactly $`\widehat{\sigma }^{(NLO)}`$, $`\widehat{\sigma }^{(NNLO)}`$ and so forth, the parton shower gives an all-order approximation of the partonic cross section in the soft and collinear regions. In this respect, the computation of the partonic cross sections performed by parton showers is somehow similar to that obtained by soft-gluon resummed calculations. There is, however, an important conceptual difference between the two approaches. This difference and the limits of applicability of the parton-shower method are briefly recalled below. Apart from these limits, parton-shower calculations can give some advantages. Multiparton kinematics can be treated exactly. The parton shower can be supplemented with models of non-perturbative effects (hadronization, intrinsic $`k_{}`$, soft underlying event) to provide a complete description of the hard-scattering process at the hadron level.
For a given cross section, resummed calculations can in principle be performed to any logarithmic accuracy. The logarithmic accuracy achievable by parton showers is instead intrinsically limited by quantum mechanics. The parton-shower algorithms are probabilistic. Starting from the LO cross section, the parton shower generates multiparton final states according to a probability distribution that approximates the square of the QCD matrix elements. The approximation is based on the universal (process-independent) factorization properties of multiparton matrix elements in the soft and collinear limits. Although the matrix element does factorize, its square contains quantum interferences, which are not positive-definite and, in general, cannot be used to define probability distributions. To leading infrared accuracy, this problem is overcome by exploiting QCD coherence (see Refs. and referencees therein): soft gluons radiated at large angle from the partons involved in the LO subprocess destructively interfere. This quantum mechanical effect can be simply implemented by enforcing an angular-ordering constraint on the phase space available for the parton shower evolution. Thus, angular-ordered parton showers can consistently compute the first two dominant towers ($`\alpha _\mathrm{S}^nL^{2n}`$ and $`\alpha _\mathrm{S}^nL^{2n1}`$) of logarithmic contributions in Eq. (10). However, parton showers contain also some subleading logarithmic contributions. For instance, they correctly compute the single-logarithmic terms $`\alpha _\mathrm{S}^nL^n`$ of purely collinear origin that lead to the LO evolution of the parton densities. Moreover, as discussed in Ref. by a comparison with resummed calculations, in the case of hard-scattering processes whose LO subprocess involves two coloured partons (e.g. DIS or DY production), angular-ordered parton showers have a higher logarithmic accuracy: they can consistently evaluate the LL and NLL terms in Eq. (15). The extension of parton-shower algorithms to higher logarithmic accuracy is not necessarily feasible and is, in any case, challenging.
Of course, because of quantum interferences and quantum fluctuations, the probabilistic parton-shower approach cannot be used to systematically perform exact calculations at NLO, NNLO and so forth. Nonetheless, important progress has been made to include matrix element corrections in parton shower algorithms . The purpose is to consider the multiparton configurations generated by parton showering from the LO matrix element and to correct them in the hard (non-soft and non-collinear) region by using the exact expressions of the higher-order matrix elements. Hard matrix element corrections to parton showers have been implemented for top quark decay and for production of $`W,Z`$ and DY lepton pairs . The same techniques could be applied to other processes, as, for instance, production of Higgs boson and vector-boson pairs .
Note also that, at present, angular-ordered parton showers cannot be considered as true ‘next-to-leading’ tools, even where their logarithmic accuracy is concerned. The consistent computation of the first two towers of logarithmic contributions in Eq. (10) is not sufficient for this purpose. For instance, to precisely introduce an NLO definition of $`\alpha _\mathrm{S}`$, we should control all the terms obtained by the replacement $`\alpha _\mathrm{S}\alpha _\mathrm{S}+c\alpha _\mathrm{S}^2+𝒪(\alpha _\mathrm{S}^3)`$. When it is introduced in the towers of double-logarithmic terms $`\alpha _\mathrm{S}^nL^{2n}`$ of Eq. (10), this replacement leads to contributions of the type $`\alpha _\mathrm{S}^{n+1}L^{2n}\alpha _\mathrm{S}^nL^{2n2}`$. Since these contributions are not fully computable at present, the parameter $`\alpha _\mathrm{S}`$ used in the parton showers corresponds to a simple LO parametrization of QCD running coupling.
## 5 Soft-gluon resummation
Double-logarithmic contributions due to soft gluons arise in all the kinematic configurations where radiation of real and virtual partons is highly unbalanced (see Ref. and references therein). For instance, this happens in the case of transverse-momentum distributions at low transverse momentum, in the case of hard-scattering production near threshold or when the structure of the final state is investigated with high resolution (internal jet structure, shape variables).
Soft-gluon resummation for jet shapes has been extensively studied and applied to hadronic final states produced by $`e^+e^{}`$ annihilation . Applications to hadron–hadron collisions have just begun to appear and have a large, yet uncovered, potential (from $`\alpha _\mathrm{S}`$ determinations to studies of non-perturbative dynamics).
Transverse-momentum logarithms, $`L=\mathrm{ln}Q^2/𝑸_{}^2`$, occur in the distribution of transverse momentum $`𝑸_{}`$ of systems with high mass $`Q`$ $`(QQ_{})`$ that are produced with a vanishing $`𝑸_{}`$ in the LO subprocess. Examples of such systems are DY lepton pairs, lepton pairs produced by $`W`$ and $`Z`$ decay, heavy quark–antiquark pairs, photon pairs and Higgs bosons. In these processes the LO transverse-momentum distribution is sharply peaked around $`𝑸_{}=0`$ ($`d\widehat{\sigma }/d^2𝑸_{}\delta ^{(2)}(𝑸_{}`$)). If the heavy system is produced with $`𝑸_{}^2Q^2`$, the emission of real radiation at higher orders is strongly suppressed and cannot balance the virtual contributions. The ensuing logarithms, $`L=\mathrm{ln}Q^2/𝑸_{}^2`$, diverge order by order when $`𝑸_{}0`$, but after all-order resummation they leads to a finite smearing of the LO distribution.
Threshold logarithms, $`L=\mathrm{ln}(1x)`$, occur when the tagged final state produced by the hard scattering is forced to carry a very large fraction $`x`$ ($`x1`$) of the available centre-of-mass energy $`\sqrt{S}`$. Also in this case, the radiative tail of real emission is stronly suppressed at higher perturbative orders. Oustanding examples of hard processes near threshold are DIS at large $`x`$ (here $`x`$ is the Bjorken variable), production of DY lepton pairs with large invariant mass $`Q`$ ($`x=Q/\sqrt{S}`$), production of heavy quark–antiquark pairs ($`x=2m_Q/\sqrt{S}`$), production of single jets and single photons at large transverse energy $`E_T`$ ($`x=2E_T/\sqrt{S}`$).
To emphasize the difference between transverse-momentum logarithms and threshold logarithms generated by soft gluons, it can be instructive to consider prompt-photon production. In the case of production of a photon pair<sup>7</sup><sup>7</sup>7The same discussion applies to the production of a DY lepton pair. with invariant mass squared $`Q^2=(p_1^{(\gamma )}+p_2^{(\gamma )})^2`$ and total transverse momentum $`𝑸_{}=𝒑_1^{(\gamma )}+𝒑_2^{(\gamma )}`$, transverse-momentum logarithms and threshold logarithms appear when $`𝑸_{}^2Q^2`$ and $`𝑸_{}^2(S/4Q^2)`$, respectively. However, in the case of production of a single photon with transverse energy (or, equivalently, transverse momentum) $`E_T`$, soft gluons can produce logarithms only in the threshold region $`x_T=2E_T/\sqrt{S}1`$. If the prompt photon has a transverse energy that is not close<sup>8</sup><sup>8</sup>8Eventually, when $`x_T1`$, higher-order corrections are single-logarithmically enhanced. This small-$`x`$ logarithms, $`(\alpha _\mathrm{S}\mathrm{ln}x_T)^n`$, have to be taken into account by BFKL resummation. to its threshold value, the emission of accompanying radiation is not kinematically suppressed and there are no soft logarithms analogous to those in the transverse-momentum distribution of a photon pair. In particular, there are no double-logarithmic contributions of the type $`(\alpha _\mathrm{S}\mathrm{ln}^2E_T^2/S)^n`$, and perturbative soft gluons are not distinguishable from perturbative hard gluons.
Studies of soft-gluon resummation for transverse-momentum distributions at low transverse momentum and for hard-scattering production near threshold started two decades ago . The physical bases for a systematic all-order summation of the soft-gluon contributions are dynamics and kinematics factorizations . The first factorization follows from gauge invariance and unitarity: in the soft limit, multigluon amplitudes fulfil factorization formulae given in terms of universal (process-independent) soft contributions. The second factorization regards kinematics and strongly depends on the actual cross section to be evaluated. If, in the appropriate soft limit, the multiparton phase space for this cross section can be written in a factorized way, resummation is analytically feasible in form of generalized exponentiation of the universal soft contributions that appear in the factorization formulae of the QCD amplitudes.
Note that the phase space depends in a non-trivial way on multigluon configurations and, in general, is not factorizable in single-particle contributions<sup>9</sup><sup>9</sup>9In the case of jet cross sections, for instance, phase-space factorization depends on the detailed definition of jets and it can easily be violated . Some jet algorithms, such as the $`k_{}`$-algorithm , have better factorization properties.. Moreover, even when phase-space factorization is achievable, it does not always occur in the space of the kinematic variables where the cross section is defined. Usually, it is necessary to introduce a conjugate space to overcome phase-space constraints. This is the case for transverse-momentum distributions and hard-scattering production near threshold. The relevant kinematical constraint for $`𝑸_{}`$-distributions is (two-dimensional) transverse-momentum conservation and it can be factorized by performing a Fourier transformation. Soft-gluon resummation for $`𝑸_{}`$-distributions is thus carried out in $`𝒃`$-space , where the impact parameter $`𝒃`$ is the variable conjugate to $`𝑸_{}`$ via the Fourier transformation. Analogously, the relevant kinematical constraint for hard-scattering production near threshold is (one-dimensional) energy conservation and it can be factorized by working in $`N`$-moment space , $`N`$ being the variable conjugate to the threshold variable $`x`$ (energy fraction) via a Mellin or Laplace transformation.
Using a short-hand notation, the general structure of the partonic cross section $`\widehat{\sigma }`$ after summation of soft-gluon contributions is
$$\widehat{\sigma }=\widehat{\sigma }_{\mathrm{res}.}+\widehat{\sigma }_{\mathrm{rem}.}.$$
(11)
The term $`\widehat{\sigma }_{\mathrm{res}.}`$ embodies the all-order resummation, while the remainder $`\widehat{\sigma }_{\mathrm{rem}.}`$ contains no large logarithmic contributions. The latter has the form
$$\widehat{\sigma }_{\mathrm{rem}.}=\widehat{\sigma }^{(\mathrm{f}.\mathrm{o}.)}\left[\widehat{\sigma }_{\mathrm{res}.}\right]^{(\mathrm{f}.\mathrm{o}.)},$$
(12)
and it is obtained from $`\widehat{\sigma }^{(\mathrm{f}.\mathrm{o}.)}`$, the truncation of the perturbative expansion for $`\widehat{\sigma }`$ at a given fixed order (LO, NLO, …), by subtracting the corresponding truncation $`\left[\widehat{\sigma }_{\mathrm{res}.}\right]^{(\mathrm{f}.\mathrm{o}.)}`$ of the resummed part. Thus, the expression on the right-hand side of Eq. (11) includes soft-gluon logarithms to all orders and it is matched to the exact (with no logarithmic approximation) fixed-order calculation. It represents an improved perturbative calculation that is everywhere as good as the fixed-order result, and much better in the kinematics regions where the soft-gluon logarithms become large ($`\alpha _\mathrm{S}L1`$). Eventually, when $`\alpha _\mathrm{S}L>1`$, the resummed perturbative contributions are of the same size as the non-perturbative contributions and the effect of the latter has to be implemented in the resummed calculation.
The resummed cross section has the following typical form:
$$\widehat{\sigma }_{\mathrm{res}.}=\alpha _\mathrm{S}^k_{\mathrm{inv}.}\widehat{\sigma }^{(LO)}CS,$$
(13)
where the integral $`_{\mathrm{inv}.}`$ denotes the inverse tranformation from the conjugate space where resummation is actually carried out. Methods to perform the inverse transformation are discussed in Refs. and for $`Q_{}`$-resummation and threshold resummation, respectively. The $`C`$ term has the perturbative expansion
$$C=1+C_1\alpha _\mathrm{S}+C_2\alpha _\mathrm{S}^2+\mathrm{}$$
(14)
and contains all the constant contributions in the limit $`L\mathrm{}`$ (the coefficients $`C_1,C_2,\mathrm{}`$ do not depend on the conjugate variable). The singular dependence on $`L`$ (more precisely, on the logarithm $`\stackrel{~}{L}`$ of the conjugate variable) is entirely exponentiated in the factor $`S`$:
$$S=\mathrm{exp}\left\{Lg_1(\alpha _\mathrm{S}L)+g_2(\alpha _\mathrm{S}L)+\alpha _\mathrm{S}g_3(\alpha _\mathrm{S}L)+\mathrm{}\right\}.$$
(15)
In the exponent, the function $`Lg_1`$ resums all the leading logarithmic (LL) contributions $`\alpha _\mathrm{S}^nL^{n+1}`$, while $`g_2`$ contains the next-to-leading logarithmic (NLL) terms $`\alpha _\mathrm{S}^nL^n`$ and so forth<sup>10</sup><sup>10</sup>10To compare this notation with that of Ref. , we can notice that our functions $`g_i`$ are obtained by the straightforward integration over $`\overline{\mu }`$ of the functions $`A(\alpha _\mathrm{S}(\overline{\mu }))`$ and $`B(\alpha _\mathrm{S}(\overline{\mu }))`$ of Ref. . In particular, our terms $`g_1,g_2,g_3`$ are not to be confused with the non-perturbative parameters of the same name used in Ref. . (all the functions $`g_i`$ are normalized as $`g_i(\lambda =0)=0`$). Note that the LL terms are formally suppressed by a power of $`\alpha _\mathrm{S}`$ with respect to the NLL terms, and so forth for the successive classes of logarithmic terms. Thus, this logarithmic expansion is as systematic as the fixed-order expansion in Eq. (4). In particular, using a matched NLL+NLO calculation, we can consistently $`i)`$ introduce a precise definition (say $`\overline{\mathrm{MS}}`$) of $`\alpha _\mathrm{S}(\mu )`$ and $`ii)`$ investigate the theoretical accuracy of the calculation by studying its renormalization-scale dependence.
The structure of the exponentiated resummed calculations discussed so far has to be contrasted with that obtained by organizing the logarithmic expansion on the right-hand side of Eq. (10) in terms of towers as
$$\widehat{\sigma }\alpha _\mathrm{S}^k\widehat{\sigma }^{(LO)}\left\{t_1(\alpha _\mathrm{S}L^2)+\alpha _\mathrm{S}Lt_2(\alpha _\mathrm{S}L^2)+\alpha _\mathrm{S}^2L^2t_3(\alpha _\mathrm{S}L^2)+\mathrm{}\right\},$$
(16)
where the double-logarithmic function $`t_1(\alpha _\mathrm{S}L^2)`$ and the successive functions are normalized as $`t_i(0)=\mathrm{const}.`$ While the ratio of two successive terms in the exponent of Eq. (15) is formally of the order of $`\alpha _\mathrm{S}`$, the ratio of two successive towers in Eq. (16) is formally of the order of $`\alpha _\mathrm{S}L`$. In other words, the tower expansion sums the double-logarithmic terms $`(\alpha _\mathrm{S}L^2)^n`$, then the terms $`\alpha _\mathrm{S}^nL^{2n1}\alpha _\mathrm{S}L(\alpha _\mathrm{S}L^2)^{n1}`$, and so forth; it thus assumes that the resummation procedure is carried out with respect to the large parameter $`\alpha _\mathrm{S}L^2`$ ($`\alpha _\mathrm{S}L^2<1`$). On the contrary, in Eq. (15) the large parameter is $`\alpha _\mathrm{S}L<1`$. The tower expansion allows us to formally extend the applicability of perturbative QCD to the region $`L<1/\sqrt{\alpha }_\mathrm{S}`$, and exponentiation extends it to the wider region $`L<1/\alpha _\mathrm{S}`$. This fact can also be argued by comparing the amount of information on the logarithmic terms that is included in the truncation of Eqs. (15) and (16) at some logarithmic accuracy. The reader can easily check that, after matching to the NLO (LO) calculation as in Eq. (11), the NLL (LL) result of Eq. (15) contains all the logarithms of the first four (two) towers in Eq. (16) (and many more logarithmic terms).
In the case of $`Q_{}`$-distributions, full NLL resummation has been performed for lepton pairs, $`W`$ and $`Z`$ bosons produced by the DY mechanism and for Higgs bosons produced by gluon fusion . Corresponding resummed calculations are discussed in Refs. and references therein.
Threshold logarithms in hadron collisions have been resummed to NLL accuracy for DIS and DY production and for Higgs boson production . Recent theoretical progress regards the extension of NLL resummation to processes produced by LO hard-scattering of more than two coloured partons, such as heavy-quark hadroproduction and leptoproduction , as well as prompt-photon , quarkonium and vector-boson production.
An important feature of threshold resummation is that the resummed soft-gluon contributions regard the partonic cross section rather than the hadronic cross section. This fact has two main consequences: $`i)`$ soft-gluon contributions can be sizeable long before the threshold region in the hadronic cross section is actually approached, and $`ii)`$ the resummation effects typically enhance the fixed-order perturbative calculations.
The first consequence follows from the fact that the evolution of the parton densities sizeably reduces the energy that is available in the partonic hard-scattering subprocess. Thus, the partonic cross section $`\widehat{\sigma }`$ in the factorization formula (2) is typically evaluated much closer to threshold than the hadronic cross section. In other words, the parton densities are strongly suppressed at large $`x`$ (typically, when $`x1`$, $`f(x,\mu ^2)(1x)^\eta `$ with $`\eta 3`$ and $`\eta 6`$ for valence quarks and sea-quarks or gluons, respectively); after integration over them, the dominant value of the square of the partonic centre-of-mass energy $`\widehat{s}=x_1x_2S`$ is therefore substantially smaller than the corresponding hadronic value $`S`$.
The second consequence, which depends on the actual definition of the parton densities, follows from the fact that the resummed contributions are those soft-gluon effects that are left at the partonic level after factorization of the parton densities. After having absorbed part of the full soft-gluon contributions in the customary definitions (for instance, those in the $`\overline{\mathrm{MS}}`$ or DIS factorization schemes) of the parton densities, it turns out that the residual effect in the partonic cross section is positive and tends to enhance the perturbative predictions.
A quantitative illustration of these consequences is given below by discussing top-quark and prompt-photon production. The discussion also shows another relevant feature of NLO+NLL calculations, namely, their increased stability with respect to scale variations.
The effects of soft-gluon resummation on the top-quark production cross sections at hadron colliders have been studied in Refs. . In the case of $`p\overline{p}`$ collisions, the comparison between QCD predictions at NLO and those after NLL resummation is shown in Fig. 10 . At the Tevatron the resummation effects are not very large and the NLO cross section is increased by about $`4\%`$. This had to be expected because the top quark is not produced very close to threshold ($`x=2m_t/\sqrt{S}0.2`$, at the Tevatron). Note, however, that the dependence on the factorization/renormalization scale of the theoretical cross section is reduced by a factor of almost 2 by including NLL resummation. More precisely, the scale dependence $`(\pm 5\%)`$ of the NLO+NLL calculation becomes comparable to that obtained by using different sets of parton densities . Combining linearly scale and parton density uncertainties, the NLO+NLL cross section is $`\sigma _{t\overline{t}}=5.0\pm 0.6`$, with $`m_t=175\mathrm{GeV}`$ and $`\sqrt{S}=1.8\mathrm{TeV}`$ .
At the LHC $`(x=2m_t/\sqrt{S}0.03)`$ the top quark is produced less close to the hadronic threshold than at the Tevatron. However this is compensated by the fact that the gluon channel<sup>11</sup><sup>11</sup>11Since $`f_g`$ is steeper than $`f_q`$ at large $`x`$, partonic cross sections in gluon subprocesses are typically closer to threshold than in quark subprocesses. Moreover, the intensity of soft-gluon radiation from gluons is larger than that from quarks by a factor of $`C_A/C_F2`$. is more important at the LHC. As a result, the effect of including soft-gluon resummation to NLL accuracy is very similar: the NLO cross section is enhanced by $`5\%`$ and its scale dependence is reduced from $`\pm 10\%`$ to $`\pm 5\%`$. Note, however, that the uncertainty $`(\pm 10\%)`$ coming from the parton (gluon) densities is larger than at the Tevatron .
Similar qualitative results are obtained when NLL resummation is applied to prompt-photon production at fixed-target experiments. The scale dependence of the theoretical calculation is highly reduced and the resummed NLL contributions lead to large corrections at high $`x_T=2E_T/\sqrt{S}`$ (and smaller corrections at lower $`x_T`$). Of course, the impact of soft-gluon resummation is quantitatively more sizeable in prompt-photon production than in top-quark production, because $`x_T`$ can be as large as 0.6, the hard scale $`E_T`$ is much smaller than $`m_t`$ (thus, $`\alpha _\mathrm{S}(E_T)>\alpha _\mathrm{S}(m_t)`$) and the gluon channel is always important. The scale dependence of the theoretical cross section for the E706 kinematics is shown in Fig. 11. Fixing $`\mu _R=\mu _F=\mu `$ and varying $`\mu `$ in the range $`E_T/2<\mu <2E_T`$ with $`E_T=10\mathrm{GeV}`$, the cross section varies by a factor of $`6`$ at LO, by a factor of $`4`$ at NLO and by a factor of $`1.3`$ after NLL resummation. The highly reduced scale dependence of the NLO+NLL cross section is also visible in Fig. 12, which, in particular, shows that when $`E_T=10\mathrm{GeV}`$ and $`E_{\mathrm{beam}}=530\mathrm{GeV}`$ the central value (i.e. with $`\mu =E_T`$) of the NLO cross section increases by a factor of $`2.5`$ after NLL resummation. As expected, the size of these effects is reduced by increasing $`\sqrt{S}`$ at fixed $`E_T`$ (see Fig. 12) or by decreasing $`E_T`$ at fixed $`\sqrt{S}`$ (see Fig. 8).
The comparison with the E706 data shown in Fig. 13 suggests that the NLO+NLL calculation can help to better understand prompt-photon production at large $`x_T`$. Note, however, that this comparison has to be regarded as preliminary in several respects . In particular, the parton densities used in Fig. 13 are those extracted from NLO fits. Owing to the soft-gluon enhancement at large $`x_T`$, refitting the parton densities may lead to a smaller $`f_g`$ at large $`x`$ and, consequently (because of the momentum sum rule), a larger $`f_g`$ at intermediate $`x`$. As a result, this procedure could somehow increase the theoretical cross section also at smaller values of $`x_T`$.
Soft-gluon resummation at NLL accuracy is now available for all the processes (namely, DIS, DY and prompt-photon production) that are typically used to perform global fits to parton densities. A detailed extraction/evolution of parton densities by consistently using NLL resummed calculations is thus nowadays feasible.
## 6 Other topics
The activity of the QCD Working Group at this Workshop has also been devoted to other topics, such as automatic computation of matrix elements and LO cross sections for multiparticle processes at high-energy colliders, definition and properties of jets algorithms, definition of isolated photons and related NLO calculations. Corresponding contributions are included in these Proceedings.
Other studies performed during this Workshop have a large overlap with the activity of the related Workshops at FERMILAB and CERN and can be found in those Proceedings .
Acknowledgements. This work was supported in part by the EU Fourth Framework Programme “Training and Mobility of Researchers”, Network “Quantum Chromodynamics and the Deep Structure of Elementary Particles”, contract FMRX–CT98–0194 (DG 12 – MIHT). I would like to thank the members of the Local Organizing Committee for the excellent Workshop.
## References
1. M.L. Mangano, plenary talk at the 1999 International Europhysics Conference on High Energy Physics Tampere, Finland, 15–21 July 1999, preprint CERN-TH/99-337 (hep-ph/9911256), to appear in the Proceedings, and references therein.
2. B.R. Webber, plenary talk at the 19th International Symposium on Lepton-Photon Interactions at High-Energies, LP 99, Stanford, California, 9–14 Aug 1999, preprint CERN-TH/99-387 (hep-ph/9912292), to appear in the Proceedings, and references therein.
3. J. Huston, in Proc. of the 29th International Conference on High-Energy Physics, ICHEP 98, eds. A. Astbury, D. Axen and J. Robinson (World Scientific, Singapore, 1999), vol. 1, p. 283, and references therein.
4. S. Catani, hep-ph/9712442, in Proc. of the XVIII International Symposium on Lepton-Photon Interactions, LP97, eds. A. De Roeck and A. Wagner (World Scientific, Singapore, 1998), p. 147, and references therein.
5. G. Jarlskog and D. Rein eds., Proc. of ECFA LHC Workshop, report CERN 90-10 (December 1990); ATLAS Coll., ATLAS TDR 15, report CERN/LHCC/99-15 (May 1999).
6. Proceedings of the Workshop on Physics at the Tevatron in Run II, Fermilab, 2000 (to appear). See: http://www-theory.fnal.gov/people/ellis/QCDWB/QCDWB.html
7. Proceedings of the Workshop on Standard Model Physics (and more) at the LHC, CERN 1999 (to appear). See: http://home.cern.ch/~mlm/lhc99/lhcworkshop.html
8. See J.C. Collins, D.E. Soper and G. Sterman, in Perturbative Quantum Chromodynamics, ed. A.H. Mueller (World Scientific, Singapore, 1989), p. 1, and references therein.
9. V.N. Gribov and L.N. Lipatov, Sov. J. Nucl. Phys. 15 (1972) 438, 675; G. Altarelli and G. Parisi, Nucl. Phys. B126 (1977) 298; Yu.L. Dokshitzer, Sov. Phys. JETP 46 (1977) 641.
10. E. G. Floratos, D. A. Ross and C. T. Sachrajda, Nucl. Phys. B129 (1977) 66, E ibid. B139 (1978) 545, Nucl. Phys. B152 (1979) 493; A. Gonzalez-Arroyo, C. Lopez and F. J. Yndurain, Nucl. Phys. B153 (1979) 161; A. Gonzalez-Arroyo and C. Lopez, Nucl. Phys. B166 (1980) 429; G. Curci, W. Furmanski and R. Petronzio, Nucl. Phys. B175 (1980) 27; W. Furmanski and R. Petronzio, Phys. Lett. B97 (1980) 437; E. G. Floratos, C. Kounnas and R. Lacaze, Nucl. Phys. B192 (1981) 417.
11. A.D. Martin, R.G. Roberts, W.J. Stirling and R.S. Thorne, Eur. Phys. J. C4 (1998) 463, Phys. Lett. 443B (1998) 301, preprint DTP-99-64 (hep-ph/9907231); M. Glück, E. Reya and A. Vogt, Eur. Phys. J. C5 (1998) 461; H.L. Lai et al., Eur. Phys. J. C12 (2000) 375.
12. LHC Guide to Parton Distribution Functions and Cross Sections, ATLAS note ATL-PHYS-99-008, http://www.pa.msu.edu/~huston/lhc/lhc<sub>-</sub>pdfnote.ps .
13. R.D. Ball and J. Huston, in these Proceedings, and references therein.
14. U.K. Yang and A. Bodek, Phys. Rev. Lett. 82 (1999) 2467, preprint UR-1581 (hep-ex/9908058).
15. W. Melnitchouk, I.R. Afnan, F. Bissey and A.W. Thomas, preprint ADP-99-47-T384 (hep-ex/9912001).
16. S. Kuhlmann et al., hep-ph/9912283.
17. J. Huston et al., Phys. Rev. D 58 (1998) 114034.
18. R. Hamberg, W.L. van Neerven and T. Matsuura, Nucl. Phys. B359 (1991) 343; W.L. van Neerven and E.B. Zijlstra, Nucl. Phys. B382 (1992) 11.
19. R. Bonciani, S. Catani, M.L. Mangano and P. Nason, Nucl. Phys. B529 (1998) 424.
20. L.N. Lipatov, Sov. J. Nucl. Phys. 23 (1976) 338; E.A. Kuraev, L.N. Lipatov and V.S. Fadin, Sov. Phys. JETP 45 (1977) 199; Ya. Balitskii and L.N. Lipatov, Sov. J. Nucl. Phys. 28 (1978) 822.
21. T. Jaroszewicz, Phys. Lett. 116B (1982) 291.
22. S. Catani and F. Hautmann, Phys. Lett. 315B (1993) 157, Nucl. Phys. B427 (1994) 475.
23. V.S. Fadin and L.N. Lipatov, Sov. J. Nucl. Phys. 50 (1989) 712, Phys. Lett. 429B (1998) 127.
24. M. Ciafaloni and G. Camici, Nucl. Phys. B496 (1997) 305, Phys. Lett. 430B (1998) 349.
25. S. Catani, Z. Phys. C70 (1996) 263, Z. Phys. C75 (1997) 665.
26. P. Nason, S. Dawson and R.K. Ellis, Nucl. Phys. B303 (1988) 607; W. Beenakker, H. Kuijf, W.L. van Neerven and J. Smith, Phys. Rev. D 40 (1989) 54; M.L. Mangano, P. Nason and G. Ridolfi, Nucl. Phys. B373 (1992) 295.
27. M. Baarmand, talk presented at the Workshop on Standard Model Physics (and more) at the LHC, CERN, January 1999.
See: http://home.cern.ch/n/nason/www/lhc99/14-01-99/program-1-14-99
28. S. Frixione, M.L. Mangano, P. Nason and G. Ridolfi, in Heavy Flavours II, eds. A.J. Buras and M. Lindner (World Scientific, Singapore, 1998), p. 609, and references therein.
29. M. Cacciari, M. Greco and P. Nason, JHEP 05 (1998) 007; F.I. Olness, R.J. Scalise and Wu-Ki Tung, Phys. Rev. D 59 (1999) 014506.
30. S. Catani, M. Ciafaloni and F. Hautmann, Phys. Lett. 242B (1990) 97, Nucl. Phys. B366 (1991) 135; J.C. Collins and R.K. Ellis, Nucl. Phys. B360 (1991) 3; E.M. Levin, M.G. Ryskin, Yu.M. Shabel’skii and A.G. Shuvaev, Sov. J. Nucl. Phys. 53 (1991) 657.
31. P. Aurenche, R. Baier, M. Fontannaz and D. Schiff, Nucl. Phys. B297 (1988) 661; H. Baer, J. Ohnemus and J.F. Owens, Phys. Rev. D 42 (1990) 61; P. Aurenche, R. Baier and M. Fontannaz, Phys. Rev. D 42 (1990) 1440; L.E. Gordon and W. Vogelsang, Phys. Rev. D 48 (1993) 3136.
32. P. Aurenche, M. Fontannaz, J.Ph. Guillet, B. Kniehl, E. Pilon and M. Werlen, Eur. Phys. J. C9 (1999) 107.
33. E706 Coll., L. Apanasevich et al., Phys. Rev. Lett. 81 (1998) 2642.
34. L. Apanasevich et al., Phys. Rev. D 59 (1999) 074007.
35. M.A. Kimber, A.D. Martin and M.G. Ryskin, Eur. Phys. J. C12 (2000) 655.
36. E. Laenen, G. Sterman and W. Vogelsang, preprint YITP-99-69 (hep-ph/0002078).
37. P. Aurenche, M. Fontannaz, J.P. Guillet, B.A. Kniehl and M. Werlen, preprint LAPTH-751-99 (hep-ph/9910252).
38. J. Huston et al., paper in preparation.
39. V. Del Duca and G. Heinrich, in these Proceedings, and references therein.
40. R.K. Ellis, W.J. Stirling and B.R. Webber, QCD and Collider Physics (Cambridge University Press, Cambridge, 1996) and references therein.
41. Yu.L. Dokshitzer, V.A. Khoze, A.H. Mueller and S.I. Troian, Basics of Perturbative QCD (Editions Frontières, Gif-sur-Yvette, 1991) and references therein.
42. A. Bassetto, M. Ciafaloni and G. Marchesini, Phys. Rep. 100 (1983) 201, and references therein.
43. S. Catani, G. Marchesini and B.R. Webber, Nucl. Phys. B349 (1991) 635.
44. M.H. Seymour, Comput. Phys. Commun. 90 (1995) 95.
45. J. Andre and T. Sjostrand, Phys. Rev. D57 (1998) 5767.
46. C. Friberg and T. Sjostrand, in Proc. of the Workshop Monte Carlo Generators for HERA Physics, eds. T.A. Doyle, G. Grindhammer, G. Ingelman and H. Jung, (DESY, Hamburg, 1999), p. 181.
47. J.C. Collins, hep-ph/0001040 and in these Proceedings.
48. G. Corcella and M.H. Seymour, Phys. Lett. 442B (1998) 417.
49. G. Miu and T. Sjostrand, Phys. Lett. 449B (1999) 313.
50. S. Mrenna, preprint UCD-99-4 (hep-ph/9902471).
51. G. Corcella and M.H. Seymour, preprint RAL-TR-1999-051 (hep-ph/9908388).
52. C. Balazs, J. Huston and I. Puljak, preprint FERMILAB-PUB-00-032-T
(hep-ph/0002032), and in these Proceedings.
53. G. Sterman, in Proc. 10th Topical Workshop on Proton-Antiproton Collider Physics, eds. R. Raja and J. Yoh (AIP Press, New York, 1996), p. 608; S. Catani, Nucl. Phys. Proc. Suppl. 54A (1997) 107, and in Proc of the 32nd Rencontres de Moriond: QCD and High-Energy Hadronic Interactions, ed. J. Tran Than Van (Editions Frontières, Paris, 1997), p. 331.
54. S. Catani, G. Turnock, B.R. Webber and L. Trentadue, Nucl. Phys. B407 (1993) 3.
55. M.H. Seymour, Nucl. Phys. B513 (1998) 269; J.R. Forshaw and M.H. Seymour, JHEP 09 (1999) 009.
56. Yu.L. Dokshitzer, D.I. Diakonov and S.I. Troian, Phys. Rep. 58 (1980) 269, and references therein.
57. N. Brown and W.J. Stirling, Phys. Lett. 252B (1990) 657.
58. S. Catani, Yu.L. Dokshitzer, M. Olsson, G. Turnock and B.R. Webber, Phys. Lett. 269B (1991) 432.
59. S. Catani, Yu.L. Dokshitzer and B.R. Webber, Phys. Lett. 285B (1992) 291; S. Catani, Yu.L. Dokshitzer, M.H. Seymour and B.R. Webber, Nucl. Phys. B406 (1993) 187; S.D. Ellis and D.E. Soper, Phys. Rev. D 48 (1993) 3160.
60. G. Parisi and R. Petronzio, Nucl. Phys. B154 (1979) 427.
61. J.C. Collins, D.E. Soper and G. Sterman, Nucl. Phys. B250 (1985) 199.
62. G. Sterman, Nucl. Phys. B281 (1987) 310.
63. S. Catani and L. Trentadue, Nucl. Phys. B327 (1989) 323, Nucl. Phys. B353 (1991) 183.
64. J.C. Collins and D.E. Soper, Nucl. Phys. B197 (1982) 446.
65. S. Catani, M.L. Mangano, P. Nason and L. Trentadue, Nucl. Phys. B478 (1996) 273.
66. C. Balázs, J.C. Collins and D.E. Soper, in these Proceedings.
67. J. Kodaira and L. Trentadue, Phys. Lett. 112B (1982) 66, Phys. Lett. 123B (1983) 335; C.T.H. Davies, B.R. Webber and W.J. Stirling, Nucl. Phys. B256 (1985) 413.
68. S. Catani, E. D’Emilio and L. Trentadue, Phys. Lett. 211B (1988) 335; R.P. Kauffman, Phys. Rev. D 45 (1992) 1512.
69. H. Contopanagos, E. Laenen and G. Sterman, Nucl. Phys. B484 (1997) 303.
70. M. Kramer, E. Laenen and M. Spira, Nucl. Phys. B511 (1998) 523.
71. N. Kidonakis and G. Sterman, Phys. Lett. 387B (1996) 867, Nucl. Phys. B505 (1997) 321.
72. N. Kidonakis, G. Oderda and G. Sterman, Nucl. Phys. B531 (1998) 365.
73. E. Laenen and S. Moch, Phys. Rev. D 59 (1999) 034027.
74. E. Laenen, G. Oderda and G. Sterman, Phys. Lett. 438B (1998) 173.
75. S. Catani, M.L. Mangano and P. Nason, JHEP 07 (1998) 024.
76. S. Catani, M.L. Mangano, P. Nason, C. Oleari and W. Vogelsang, JHEP 03 (1999) 025.
77. N. Kidonakis and J.F. Owens, preprint FSU-HEP-991216 (hep-ph/9912388).
78. M. Cacciari, preprint CERN-TH/99-312 (hep-ph/9910412).
79. N. Kidonakis and V. Del Duca, preprint FSU-HEP-991123 (hep-ph/9911460).
80. E. Laenen, J. Smith and W.L. van Neerven, Nucl. Phys. B369 (1992) 543, Phys. Lett. 321B (1994) 254.
81. E. Berger and H. Contopanagos, Phys. Rev. D 54 (1996) 3085, Phys. Rev. D 57 (1998) 253.
82. S. Catani, M.L. Mangano, P. Nason and L. Trentadue, Phys. Lett. 378B (1996) 329.
83. N. Kidonakis, preprint EDINBURGH-99-4 (hep-ph/9904507).
Partons for the LHC
R.D. Ball <sup>12</sup><sup>12</sup>12Royal Society University Research Fellow. and J. Huston
## Abstract
We discuss some of the experimental, theoretical and methodological issues in the determination of parton distributions with meaningful error estimates, and their impact on physical cross sections to be measured at the Tevatron and LHC.
## Abstract
In this section we summarize the formalism which extends the usual hadronic factorization theorem to the low transverse momentum region for the inclusive production of colorless final states, while resumming logarithms with the ratio of the invariant mass and transverse momentum. Among the various recent applications the calculation of the $`Z^0`$ and Higgs boson transverse momentum distributions are highlighted.
## Abstract
Monte Carlo event generators are being increasingly relied upon for predictions of experimental observables at colliders. In this section, the parton shower formalism for Monte Carlos is compared to that of analytic resummation calculations. Predictions for the transverse momentum distribution of $`Z^0`$ bosons, photon pairs, and the Higgs boson are compared for the Tevatron and the LHC. <sup>1</sup><sup>1</sup>1A more complete treatment of this subject can be found in Ref. .
## Abstract
I am concerned here with QCD calculations for processes with a hard scattering — production of heavy particles, jets, etc. The most accurate calculations are by the “analytic” methods. But the most useful calculations for direct comparison with data are done by Monte-Carlo event generators, and these are limited in accuracy. In particular, there is as yet no known method of systematically improving the Monte-Carlo calculations by incorporating the non-logarithmic parts of higher order perturbative corrections. This creates limitations on the analysis of future data. Therefore I summarize some ideas for remedying the situation.
## Abstract
In order to determine more accurately the energy contribution in a jet cone due to the underlying event, and in order to understand better the ambient event environment at both the Tevatron and the LHC, we have studied the energy distribution in a cone of radius 0.7 in both jet and in minimum bias events. We have compared the results from CDF data from Run 1b with results from HERWIG passed through the detector simulation program QFL .
## Abstract
Direct <sup>1</sup><sup>1</sup>1“Direct” or “prompt” mean that these photons do not result from the decay of $`\pi ^0`$ and $`\eta `$. photon pairs with large invariant mass are the so-called irreducible background in the search for Higgs bosons at the LHC in the channel $`h\gamma \gamma `$, in the mass range $`80140`$ GeV$`/c^2`$. This huge background requires an understanding and quantitative evaluation. Photon pair production at the Tevatron offers the opportunity to already test our understanding of this process. In the same mass range, the production of a Higgs boson ($`hrightarrow\gamma \gamma `$) in association with a hard jet at the LHC is a promising channel, as the corresponding $`\gamma \gamma +`$ jet background may be under better control.
## 1 Introduction
The calculation of production cross sections at the Tevatron and LHC, for both interesting physics processes and their backgrounds, relies upon a knowledge of the distribution of the momentum fraction $`x`$ of the partons in a proton at the relevant scale. These parton distribution functions (pdfs) are at present determined by global fits to data from deep inelastic scattering (DIS), Drell-Yan (DY), and jet and direct photon production at current energy ranges. Two groups, CTEQ and MRS, provide semi-regular updates to their best-fit parton distributions when new data and/or theoretical developments become available. The newest pdfs, in most cases, currently provide the single most accurate overall description of the world’s data, and should be utlilized in preference to older pdf sets. The most recent sets from the two groups are CTEQ5 and MRST .
In this contribution we will discuss the data sets used in the fits, the way in which the fits are performed in practice (in particular, issues such as the parametrization of initial distributions, the solution of the evolution equations, and scheme dependence), and the main uncertainties in the fitted pdfs due to uncertain or incomplete experimental data. In particular, we will concentrate on the difficulties involved in determining the gluon distribution through direct photons or jets. We then move on to discuss more general issues which may affect future pdf determinations: the inclusion of correlated systematics and the difficulties involved in combining these for different experiments, purely theoretical uncertainties arising from the limitations of NLO perturbative QCD, and finally, methodological uncertainties such as the dependence on the form of the parametrization and the assumption of Gaussian error propagation. We conclude with a summary of the progress that might be made before the LHC turns on, and the role of LHC data in determining pdfs.
## 2 Processes Involved in Global Analysis Fits
Lepton-lepton, lepton-hadron and hadron-hadron interactions probe complementary aspects of perturbative QCD (pQCD). Lepton-lepton processes provide clean measurements of $`\alpha _s(Q^2)`$ and of the fragmentation functions of partons into hadrons. Measurements of deep-inelastic scattering (DIS) structure functions ($`F_2,F_3`$) in lepton-hadron scattering and of lepton pair production cross sections in hadron-hadron collisions are the main source of information on the quark distributions $`q^a(x,Q)`$ inside hadrons. Scaling violations in deep inelastic processes give some information about the gluon distribution $`g(x,Q)`$. Furthermore the gluon distribution function enters directly (i.e. at leading order) in hadron-hadron scattering processes with direct photon and jet final states. Modern global parton distribution fits are carried out to next-to-leading (NLO) order which allows $`\alpha _s(Q^2),q^a(x,Q)`$ and $`g(x,Q)`$ to all mix and contribute in the theoretical formulae for all processes. Nevertheless, the broad picture described above still holds to some degree in global pdf analyses.
In pQCD, the gluon distribution is always accompanied by a factor of $`\alpha _s`$, in both the hard scattering cross sections and in the evolution equations for parton distributions. Thus, determination of $`\alpha _s`$ and the gluon distribution is, in general, a strongly coupled problem. One can determine $`\alpha _s`$ separately from $`e^+e^{}`$ or determine $`\alpha _s`$ and $`g(x,Q)`$ jointly in a global pdf analysis. In the latter case, though, the coupling of $`\alpha _s`$ and the gluon distribution may not lead to a unique solution for either (see for example the discussion in the CTEQ4 paper where good fits were obtained to a global analysis data set, including the inclusive jet data, for a wide range of $`\alpha _s`$ values .)
Currently, the world average value of $`\alpha _s(M_Z)`$ is $`0.119\pm 0.004`$ . This is in agreement with the average value from LEP, while the DIS experiments prefer a slightly smaller value (of the order of $`0.1160.118`$)). Since global pdf analyses are dominated by the high statistics DIS data, they would tend to favor the values of $`\alpha _s`$ closer to the lower DIS values. The more logical approach is to adopt the world average value of $`\alpha _s(M_Z)`$ and concentrate on the determination of the pdfs. This is what both CTEQ and MRS currently do. <sup>13</sup><sup>13</sup>13One can either quote a value of $`\alpha _s(M_Z)`$ or the value of $`\mathrm{\Lambda }^{\overline{MS}}`$. In the latter case, however, the number of flavors has to be clearly specified, since the value of $`\alpha _s`$ (and not $`\mathrm{\Lambda }^{\overline{MS}}`$) has to be continuous across flavor thresholds.
The data from DIS, DY, direct photon and jet processes utilized in pdf fits cover a wide range in $`x`$ and $`Q\sqrt{Q^2}`$. The kinematic ‘map’ in the $`(x,Q)`$ plane of the data points used in a recent parton distribution function analyses is shown in Figure 1. The HERA data (H1+ZEUS) are predominantly at low $`x`$, while the fixed target DIS and DY data are at higher $`x`$. There is considerable overlap, however, with the degree of overlap increasing with time as the statistics of the HERA experiments increases. DGLAP-based NLO pQCD provides an accurate description of the data (and of the evolution of the parton distributions) over the entire kinematic range shown. At very low $`x`$ and $`Q^2`$, DGLAP evolution is believed to be no longer applicable due to unresummed small $`x`$ logarithms. Similarly at very large $`x`$ there are significant contributions from unresummed soft logarithms (logarithms of $`1x`$). However, no evidence for such corrections is seen in the current range of data; thus all global analyses use conventional DGLAP evolution of pdfs.
There is a remarkable consistency between the data in the pdf fits and the NLO QCD theory fit to them. Over $`1300`$ data points are shown in Figure 1 and the $`\chi ^2`$/d.o.f. for the fit of theory to data is on the order of one.
Parton distributions determined at a given $`x`$ and $`Q^2`$ propagate down to lower $`x`$ values at higher $`Q^2`$ values. The accuracy of the extrapolation to higher $`Q^2`$ depends both on the accuracy of the original measurement and any uncertainty on $`\alpha _s(Q^2)`$. For the structure function $`F_2`$, the typical measurement uncertainty at medium to large $`x`$ is on the order of $`\pm 3\%`$. At large $`x`$, the DGLAP equation for $`F_2`$ can be approximated as $`\frac{F_2}{\mathrm{log}Q^2}=\alpha _s(Q^2)P^{qq}F_2`$. There is an extrapolation uncertainty of around $`\pm 5\%`$ in $`F_2`$ from low to high $`Q^2`$ ($`10^5`$ $`GeV^2`$) from the uncertainty in $`\alpha _s`$. Evolved distributions are also susceptible to uncertainties from an anomalously large contribution to $`F_2`$ near $`x`$ values of 1. Such a contribution may not be evident in fixed target measurements at low $`x`$ and low $`Q^2`$, but may influence higher $`Q^2`$ measurements .
For comparison, the kinematics appropriate for the production of a state of mass $`M`$ and rapidity $`y`$ at the LHC is shown in Figure 2 . For example, to produce a state of mass $`100`$ GeV and rapidity $`2`$ requires partons with $`x`$ values between $`0.05`$ and $`0.001`$ at a $`Q^2`$ value of $`10^4`$ $`GeV^2`$. Also shown in the figure is another view of the kinematic coverage of the fixed target and HERA experiments used in pdf fits. It can be seen that parton distributions determined from these experiments are sufficient to predict most LHC cross-sections of interest, provided that DGLAP evolution at small and large $`x`$ is sufficiently reliable.
## 3 Evolution, Schemes and Parametrizations
### 3.1 Evolution Codes
In order to fit the initial pdfs to experimental data they need to be evolved up to the correct scale by solving the DGLAP equations either to LO or NLO. The evolution can be carried out in either moment space or configuration space: both MRS and CTEQ use configuration space codes. Improvements have been made in the CTEQ and MRST evolution programs so that both now agree with the ‘DESY standard’ evolution prescription . The CTEQ and MRST packages should be able to carry out the evolution using NLO DGLAP to an accuracy of a few percent over the LHC kinematic range, except perhaps at very large and very small $`x`$. Note that the theoretical predictions for the W and Z total cross sections at the LHC may have uncertainties of less than 5%. This puts a great demand for the pdf evolution to have accuracies of better than a few percent, since any error on a pdf gets doubled in the cross section calculation. Mellin space codes might be the answer here.
A global pdf analysis carried out at next-to-leading order needs to be performed in a specific renormalization and factorization scheme. The evolution kernels are in a specific scheme and to maintain consistency, any hard scattering cross section calculations used for the input processes or utilizing the resulting pdfs need to also have been implemented in that same renormalization scheme. Almost universally, the $`\overline{MS}`$ scheme is used: pdfs are also available in the DIS scheme, a fixed flavor scheme (as in ref.) and several schemes that differ in their specific treatment of the charm quark mass .
It is also possible to use only leading-order matrix element calculations in the global fits which results in leading-order parton distribution functions. Such pdfs are preferred when leading order matrix element calculations (such as Monte Carlo programs like HERWIG and PYTHIA ) are used. The differences between LO and NLO pdfs, though, are formally NLO; thus, the additional error introduced by using a NLO pdf with HERWIG rather than a LO pdf, for example, should not be significant, in principle, and NLO pdfs can be used when no LO alternatives are available. The accuracy of current DIS/DY data is such that the $`\chi ^2`$ values for LO fits are noticeably worse than those from the NLO fits: the data are sensitive to the differences between LO and NLO partonic cross-sections and evolution kernels.
### 3.2 Parametrization of Initial Distributions
All current global analyses use a generic form for the parametrization of both the quark and gluon distributions at some reference value $`Q_0`$:
$$f(x,Q_0)=a_0x^{a_1}(1x)^{a_2}P(x;a_3,\mathrm{}).$$
(1)
The reference value $`Q_0`$ is usually chosen in the range of $`12`$ GeV. The parameter $`a_1`$ is associated with small-$`x`$ behaviour while $`a_2`$ is associated with large-$`x`$ valence counting rules. In some pdf fits, $`a_1^{\mathrm{gluon}}`$ has been tied to $`a_1^{\mathrm{seaquark}}`$; in more recent fits like CTEQ4, CTEQ5 and MRST, the two small $`x`$ exponents are allowed to vary independently. The current statistical power of the low $`x`$ and $`Q^2`$ DIS data from HERA warrants this separation.
The first two factors, in general, are not sufficient to describe either quark or gluon distributions. The term $`P(x;a_3,\mathrm{})`$ is a suitably chosen smooth function, depending on one or more parameters, that adds more flexibility to the pdf parametrization. In general, both the number of free parameters and the functional form can have an influence on the global fit. For example, the MRS group traditionally uses $`P_{MRS}(x;a_3,a_4)=1+a_3\sqrt{x}+a_4x`$. The CTEQ3 pdf used $`P_{CTEQ3}=1+a_3x`$ while CTEQ2, CTEQ4 and CTEQ5 all use the more general form $`P_{CTEQ2,4,5}=1+a_3x^{a_4}`$. The flexibility in the latter form, for example, makes possible the larger gluon at high $`x`$ observed in the CTEQ4HJ pdf.
Although the pdfs determined from global analyses should, in principle, be universal, in practice they could depend on the choice of data sets, and in particular on the choice of $`Q_{cut}`$ values that specify the minimum hard physical scale $`(Q,p_T,..)`$ required for data points to be included in the fit.
The parton distributions from the recent CTEQ pdf release are plotted in Figure 3 at a $`Q`$ value of 5 $`GeV`$. The gluon distribution is largest at small $`x`$ values while the valence quark distributions dominate at higher $`x`$.
### 3.3 Evolution in time and $`Q^2`$
As discussed in the introduction, the MRS and CTEQ groups provide semi-regular updates to their parton distributions as new data and/or theory becomes available. The latest parton distributions are the most accurate and should be used in preference to previous pdfs. However, in some cases calculations using older pdfs are necessary: for example, until recently <sup>14</sup><sup>14</sup>14In the most recent version of PYTHIA (6.1), the CTEQ5 pdf’s are available. none of the more recent pdfs were implemented in PYTHIA, and most comparisons in the ATLAS TDR have been made with the CTEQ2L pdf (the default pdf in PYTHIA version 5.7).
A comparison of the CTEQ1M , CTEQ2M , CTEQ3M and CTEQ4M parton distributions (in particular the up sea quark and gluon distributions) are shown in Figure 4, at a $`Q^2`$ value of $`5`$ GeV<sup>2</sup>. The CTEQ2-4 up quark sea distributions are substantially steeper than that of CTEQ1, reflecting the influence of the HERA data. A similar effect is seen with the gluon distribution. There is little change in the valence distributions.
The up sea quark and gluon distributions are shown in Figure 5 at a larger $`Q^2`$ value of $`10^4`$ GeV<sup>2</sup>. Evolution has evened out many of the differences observed at lower $`Q^2`$ values. A $`Q^2`$ value of $`10^4`$ GeV<sup>2</sup> corresponds to a mass scale at the LHC of about $`100`$ GeV.
The effects of evolution are examined in more detail in Figure 6 where the up sea quark and gluon distributions are plotted at $`Q^2`$ values of $`2`$, $`10`$, $`50`$, $`10^4`$ and $`10^6`$ GeV<sup>2</sup>. There are two interesting features that can be noted. Most of the evolution takes place at low $`Q^2`$ and there is little evolution for $`x`$ values in the vicinity of $`0.1`$. In contrast, at large $`x`$ value the distributions decrease by an order of magnitude from the lowest to the highest $`Q^2`$ value, while at small $`x`$ they increase by an order of magnitude.
## 4 Estimating Uncertainties
In addition to having the best estimates for the values of the pdfs in a given kinematic range, it is also important to understand the allowed range of variation of the pdfs, i.e. their uncertainties. The crudest method of estimating parton distribution uncertainties is to compare different published parton distributions. This is unreliable since most published sets of parton distributions (for example from CTEQ and MRS) adopt similar assumptions and the differences between the sets do not fully explore the full range uncertainties that actually exist. Here and in the next section we concentrate on estimating the uncertainties due to to the limitations of available data sets.
The sum of the quark distributions $`\mathrm{\Sigma }(q(x)+\overline{q}(x))`$ is, in general, well-determined over a wide range of $`x`$ and $`Q^2`$. As stated above, the quark distributions are predominantly determined by the DIS and DY data sets which have large statistics, and systematic errors in the few percent range ($`\pm 3\%`$ for $`10^4<x<0.75`$). Thus the sum of the quark distributions is basically known to a similar accuracy. The individual quark flavors, though, may have a greater uncertainty than the sum. This can be important, for example, in predicting distributions that depend on specific quark flavors, like the W asymmetry distribution and the W rapidity distribution.
Information on the $`\overline{d}`$ and $`\overline{u}`$ distributions comes, at small $`x`$, from HERA and at medium $`x`$ from fixed target DY production on $`H_2`$ and $`D_2`$ targets. It is now well-established that the $`\overline{d}`$ and $`\overline{u}`$ distributions are not the same. The difference in these distributions between the CTEQ4M and CTEQ5M pdfs is due primarily to the influence of the data from the E866 experiment. It is worth noting that our detailed knowledge of $`\overline{d}/\overline{u}`$ is limited primarily to the $`x`$ region (.03-.35) covered by E866.
The strange quark sea is determined from dimuon production in $`\nu `$ DIS (CCFR), with the strange quark distribution ($`s+\overline{s}`$) being approximately $`\frac{1}{2}`$($`\overline{u}+\overline{d}`$). The charm and bottom quark distributions are calculated perturbatively from gluon splitting for given masses of $`m_c`$ and $`m_b`$. (See also the previous discussion on schemes.)
Current information on $`d/u`$ at large $`x`$ comes from fixed target DY production on $`H_2`$ and $`D_2`$ and the lepton asymmetry in W production at the Tevatron. In the CTEQ5 and MRST fits, the NMC $`D_2/H_2`$ data are used to constrain the large $`x`$ $`d`$ quark distribution in this way. Bodek and Yang have argued that the $`D_2`$ data need to be corrected for nuclear binding effects, which would lead to a larger $`d/u`$ ratio at large $`x`$ (and thus a larger $`d`$ quark distribution as the $`u`$ quark distribution is well-determined from DIS) . The need for the nuclear binding corrections is still an open question . The larger $`d`$ quark distribution would lead to an increase in the high $`E_T`$ Tevatron jet cross section of about 10%. A similar excess would be expected for high $`E_T`$ jet production at the LHC.
The parton distribution with the greatest uncertainty is the gluon distribution, simply because it does not couple directly to an external probe. The LHC is essentially a gluon-gluon collider and many hadron-collider signatures of physics both within and beyond that Standard Model involve gluons in the initial state. Thus, it is very important to estimate the theoretical uncertainty due to the uncertainty in the gluon distribution.
The gluon distribution can be determined indirectly at low $`x`$ by measuring the scaling violations in the quark distributions ($`F_2/\mathrm{log}Q^2`$), but a direct measurement is necessary at moderate to high $`x`$. Direct photon production has long been regarded as potentially the most useful source of information on the gluon distribution with fixed target direct photon data, especially from the experiment WA70 , being used in a number of global analyses. However, as will be discussed in the next section, there are a number of theoretical complications with the use of direct photon data.
The momentum fraction of the proton carried by quarks is determined very well from DIS data; at a $`Q_0`$ value of 1.6 GeV, in the CTEQ4 analysis for example, the momentum fraction carried by quarks is 58% with an uncertainty of $`\pm 2\%`$. Thus, the momentum fraction carried by gluons is 42% with a similar uncertainty. This constraint is important; if the gluon distribution increases in one $`x`$ range, momentum conservation forces it to decrease in another $`x`$ range. Thus, if the gluon flux in the $`x`$ range from $`0.01`$ to $`0.3`$ were to decrease by $`20\%`$, the gluon flux would have to increase by a fairly dramatic amount in the other $`x`$ ranges to compensate. For example, if this compensation were to come in the high $`x`$ region, the gluon distribution would have to double.
A simple way of estimating the uncertainty in the gluon distribution is to systematically vary the gluon parameters in a global analysis and then look for incompatibilities with the data sets that make up the global analysis database. This study has been carried out by CTEQ using only DIS and Drell-Yan data where the theoretical and experimental systematic errors are under good control . Except at larger values of $`x(x>0.20.3)`$, the variation in the gluon distributions is less than 15% at low values of $`Q^2`$, decreasing to less than 10% at high values: as noted earlier, evolution is the great equalizer for parton distributions. Note that the DIS and DY datasets used in this analysis do not provide any strong constraints on the gluon distribution at high values of $`x`$. This study used the CTEQ4 value of $`\alpha _s`$ (i.e. $`0.116`$). If $`\alpha _s`$ is varied in the range from 0.113 to 0.122, the gluon distribution varies by 3% for $`x<0.15`$.
In order to assess the range of predictions for hadronic cross sections, it is more important to know the uncertainties in the gluon-gluon and gluon-quark luminosity functions at the appropriate kinematic region (in $`\tau =x_1x_2=\widehat{s}/s)`$ rather than the uncertainties in the parton distributions themselves. Therefore it is useful to define the relevant integrated parton-parton luminosity functions: for example the gluon-gluon luminosity function can be defined as:
$$\tau \frac{dL}{d\tau }=_\tau ^1\frac{dx}{x}g(x,Q^2)g(\tau /x,Q^2).$$
(2)
This quantity is directly proportional to the cross section for s-channel production of a single particle and it also gives a good estimate for more complicated production mechanisms. In Figure 7 is shown the range of allowed gluon-gluon luminosities (normalized to the CTEQ4M values) for the variations discussed above. Here, $`Q^2`$ is taken to be $`\tau s`$, which naturally takes the $`Q^2`$ dependence of the gluon distribution into account as one changes $`\sqrt{\tau }`$. The top region is for the LHC and the bottom is for the Tevatron. Above a $`\sqrt{\tau }`$ value of $`0.1`$, the allowed variation grows dramatically; this indicates the need for more information about the gluon distribution at large $`x`$ than provided by the DIS and DY data sets used in this analysis.
In analogy with the discussion of gluon-gluon luminosities, one can also study the gluon-quark luminosity (again normalized to the CTEQ4M result). The uncertainties on the parton-parton luminosities, as a function of $`\sqrt{\tau }`$, are summarized in Table 1. Note that the region of production of a $`100140`$ GeV Higgs at the LHC lies in the region where the range of variation in the gg luminosity is $`\pm 10\%`$.
## 5 Direct Photons and Jets in Global Fits
### 5.1 Direct Photons
As mentioned previously in this section and in Reference , direct photon production has long been viewed as an ideal vehicle for measuring the gluon distribution in the proton. The quark-gluon Compton scattering subprocess $`(gq\gamma q)`$ dominates photon production in all kinematic regions of $`pp`$ scattering, as well as for low to moderate values of parton momentum fraction $`x`$ in $`\overline{p}p`$ scattering. As described previously, the gluon distribution is relatively well constrained at low $`x(x<0.1)`$ by DIS and DY data, but less so at higher $`x`$. Consequently, fixed target direct photon data have been incorporated in several modern global parton distribution function analyses with the hope of providing a major constraint on the gluon distribution at moderate to high $`x`$.
A pattern of systematic deviations of direct photon data from NLO predictions has been observed , however, these being particularly striking for the E706 experiment. The origin of the deviations is still quite controversial. One possibility that has been suggested is that the deviations are due to the effects of soft gluon radiation, or $`k_T`$ . This view, however, is not universally held; see, for example, the discussion in Reference and in Reference . The $`k_T`$ values needed to describe the data are too large to be viewed as purely ‘intrinsic’ or non-perturbative in origin. But, as discussed in Reference , in the standard formalism for direct photon production there are no double-logs to be resummed. This is in contrast to double-arm observables such as Drell-Yan or diphoton production; since direct photon production is, by definition, a single-arm observable, there is no restriction of phase space for gluon emission, and thus no double logarithmic enhancement to the $`p_T`$ distribution. The only enhancement effects that survive arise from the purely ‘intrinsic’ $`k_T`$ present in the colliding hadrons.
Nonetheless, there is generally a substantial amount of $`k_T`$ that results from the emission of soft gluons in hard scattering processes. Direct evidence of this $`k_T`$ has long been evident in Drell-Yan, diphoton and heavy quark measurements. The values of $`<k_T>`$/parton for these processes vary from $`1`$ GeV at fixed target energies to $`34`$ GeV at the Tevatron Collider. The growth is approximately logarithmic with center of mass energy. (The value expected at the LHC for relatively low mass states ($`3040`$ GeV) is in the range of $`6.57.0`$ GeV.)
Perturbative QCD corrections are insufficient to explain the size of the observed $`k_T`$ and fully resummed calculations are required to explain Drell-Yan, W/Z and diphoton distributions . These resummed calculations qualitatively describe the growth of the $`<k_T>`$ with center-of-mass energy. Currently there is no rigorous $`k_T`$-type resummation calculation available for single photon production, for the reasons cited above. In addition, this calculation is quite challenging in that the final state parton takes part in soft gluon emission and in color exchange with initial state partons, in contrast with the Drell-Yan and diphoton cases. Also, the calculation is complicated by the fact that several overlapping power-suppressed corrections can contribute and, at high $`x`$, threshold effects are important.
Nevertheless, there has been recent theoretical progress in single photon resummation . In particular, in Reference , a technique has been presented for simultaneously treating recoil and threshold corrections in single photon inclusive cross sections, working within the formalism of collinear factorization. In the preliminary results, substantial enhancements have been observed, at moderate $`p_T`$ and $`x`$, from higher order perturbative and power-law non-perturbative corrections. This approach is still quite new and the efficacy of the formalism still has to be evaluated.
There is an intuitive picture that describes the effects of this soft gluon radiation, both perturbative and non-perturbative, on the direct photon cross section. The presence of soft gluon radiation, or $`k_T`$, can give a ‘kick’ in the photon direction. Due to the steeply falling cross sections, the $`k_T`$ kick can lead to the promotion of photons from lower $`p_T`$ to higher values of $`p_T`$. The more steeply falling the cross section, the larger the resulting enhancement. Using this intuitive picture, the effects of soft gluon radiation can be approximated by a convolution of the NLO cross section with a Gaussian $`k_T`$ smearing function. The value of $`<k_T>`$ to be used for each kinematic regime should be taken directly from relevant experimental observables, given the lack of a rigorous formalism, rather than from a theoretical prediction. The behaviour of the $`k_T`$ smearing correction is quite different for the Tevatron collider and for fixed target experiments. For the Tevatron, there are two points to note: (1) the agreement with the data is improved if the $`k_T`$ correction is taken into account and (2) the $`k_T`$ smearing effects fall off roughly as $`1/p_T^2`$ . The latter behaviour is the expectation for such a power-suppressed type of effect and is the behaviour expected at the LHC, where the effects of the $`k_T`$ smearing should not be important beyond $`p_T`$ values of $`30`$ GeV <sup>15</sup><sup>15</sup>15Similar $`k_T`$ smearing effects should be present in all hard scattering cross sections, for example jet production at the Tevatron. The size of the experimental and theoretical systematic errors in the low $`E_T`$ region make such a confirmation difficult..
The $`k_T`$ correction obtained for E706 at a center-of- mass energy of $`31.6`$ GeV is shown in Figure 8. The value of $`<k_T>`$ of $`1.2`$ GeV was obtained from measurements of several kinematic observables in the experiment . The $`k_T`$ smearing effect is much larger here then observed at the collider and does not have the $`1/p_T^2`$ falloff. Also shown are the $`k_T`$ corrrections using values of $`<k_T>`$ of $`1.0`$ and $`1.4`$ GeV (a reasonable estimate of the range of experimental uncertainty in the $`<k_T>`$ determination). In addition, the $`k_T`$ correction for the E706 data used in the recent MRST pdfs is shown. The MRST $`k_T`$ correction, utilizing a different model, is larger leading to a smaller gluon distribution in the relevant $`x`$ range. (Both the CTEQ4 and MRST pdfs, with their respective $`k_T`$ corrections, lead to good agreement with the E706 direct photon cross sections.) The differences between the $`k_T`$ correction from Reference and that from the MRST pdfs can be taken as an indication of the uncertainty in the value of this correction. Good agreement with the E706 direct photon and cross section at $`\sqrt{s}=31.6`$ GeV is observed when the nominal $`k_T`$ correction of $`1.2`$ GeV is used; however, the allowed range of variation of $`k_T`$ ($`1.01.4`$ GeV) makes quantitative comparisons, and thus an extraction of the gluon distribution, difficult <sup>16</sup><sup>16</sup>16NLO QCD predictions for fixed-target direct photon production (as is also true for other fixed target processes) also contain a non-negligible renormalization and factorization scale dependence, as discussed in Reference . Since the high $`p_T`$ E706 data agrees well with CTEQ4M, it would thus disfavor the CTEQ4HJ pdf. As stated before, however, a definitive conclusion must await a more rigorous theoretical treatment.
Other related fixed target processes, such as $`\pi ^0`$ production, in the same $`p_T`$ range as the measured direct photon cross section, may perhaps shed some light on the puzzle. It has been noted that essentially all of the fixed target $`\pi ^0`$ cross sections disagree with NLO predictions, by essentially a constant factor. Thus, there may be a common problem causing the deviations, such as uncertainties in the high $`z`$ quark and gluon fragmentation functions and possible sizeable higher order corrections. In addition, the importance of the high $`z`$ fragmentation region implies the need for threshold resummation techniques to be applied, in processes with non-trivial color flow.
It is worthwhile pointing out, though, that the same $`k_T`$ model used for for single photon production was shown to also provide an adequate description of the experimental $`\pi ^o`$ cross sections . As in the case of direct photon production, the controversy regarding the theory/data discrepancies is still open. The $`\pi ^0`$ cross sections may form a crucial role in the ultimate understanding for a number of reasons: if $`k_T`$ are important for photon production, they should also have a measureable impact on the $`\pi ^0`$ cross sections as well. In addition, $`\pi ^0`$’s form the primary experimental background to direct photon production.
Finally, it is not clear if any theoretical treatment for photon production is capable of describing all of the current fixed target direct photon data. There are discrepancies between the different experiments which may imply experimental difficulties, which are in addition to any of the theoretical problems discussed above.
### 5.2 Influence of Jets
An important process that is sensitive to the gluon distribution is jet production in hadron-hadron collisions. Processes responsible for jet production include gluon-gluon, gluon-quark and quark-quark(or anti-quark) scattering. Precise data on jet production at the Fermilab Tevatron are now available over a wide range of transverse energy, and the theoretical uncertainties in most of this range are well-understood. Thus, it is to be expected that jet production can provide a good constraint on the gluon distribution.
The jet data that has been utilized in global pdf fits has been from the CDF and D0 collaborations <sup>17</sup><sup>17</sup>17The experimental and theoretical errors associated with the UA2 jet cross section make its use in pdf fits difficult.. The data cover a wide kinematic range ($`E_T`$ values from $`15`$ to $`450`$ GeV corresponding to an $`x`$ range of $`0.02`$ to $`0.5`$). The CDF jet data from Run IA were utilized in the CTEQ4HJ pdf fit . Here, a large emphasis was given to the high $`E_T`$ data points which show a deviation from NLO QCD predictions with “conventional” pdfs. Given the lack of constraints on the high $`x`$ gluon distribution discussed in Section VI, the extra emphasis on the high $`E_T`$ region was enough to cause a significant increase in the gluon distribution; for example, the gluon distribution at an $`x`$ value of $`0.5`$ ($`Q=100`$ GeV) increases by a factor of two. Since the dominant jet subprocess in this region is $`\overline{q}q`$ scattering the increase in the gluon distribution of a factor of two causes only a 20% increase in the jet cross section. This is sufficient to pass through the bottom of the CDF high $`E_T`$ jet error bars. The preliminary jet cross sections from Run 1B (90 $`pb^1`$) from both the CDF and D0 experiments were used in the CTEQ4M fits, but with statistical errors only and only for $`E_T`$ in the range $`50200`$ GeV. The points with $`E_T`$ lower than $`50`$ GeV have substantial systematic errors on both the theoretical and experimental sides while the points with $`E_T`$ higher than $`200`$ GeV contain the CDF excess. The inclusion of the jet data serves to considerably constrain the gluon distribution over the $`x`$ range of $`0.1`$ to $`0.2`$. The resulting gluon (CTEQ4M) does not decrease the excess observed by CDF at high $`E_T`$.
The published D0 jet cross section along with the (soon-to-be published) CDF jet cross section from Run 1B were used in the recently released CTEQ5 parton distributions. The fits use the full $`E_T`$ range for the cross sections and use the correlation information on the systematic errors as contained in the covariance matrices for both experiments. The two experiments are in agreement with each other except for a slight normalization shift <sup>18</sup><sup>18</sup>18 A shift on the order of 3% is expected since the two experiments use values for the total inelastic cross section that differ by that amount.; the two highest $`E_T`$ data points for CDF are above those for D0, but both experiments have large statistical errors in this region. As can be seen in Figure 9 the NLO QCD prediction with the CTEQ5M pdf is in good agreement with the CDF data. The conclusions are exactly the same for the D0 jet data. The CTEQ5M gluon is very similar to CTEQ4M, except perhaps at very high x. The CTEQ4HJ pdf has been updated to complement the new CTEQ5M pdf. The CTEQ5HJ pdf gives almost as good a global fit as CTEQ5M to the full set of data on DIS and DY processes, and has the feature that the gluon distribution is significantly enhanced in the high $`x`$ region, resulting in improved agreement with the observed trend of jet data at high $`E_T`$ in both the CDF and D0 experiments.
## 6 Systematic Uncertainties
There is currently an increasing awareness of the need and possibility of propagating errors in the data into error estimates on parton distribution functions . Ideally, one might hope to perform a full error analysis and provide correlated errors for all the parton distributions determined in a global fit. This goal is difficult to carry out for several reasons. Firstly, there is no established way of quantifying the theoretical uncertainties for the diverse physical processes that are used. More pragmatically, only a subset of the experiments usually involved in global analyses provide correlation information on their data sets in a way suitable for the analysis. In these circumstances, comparing data from different experiments becomes very difficult. Furthermore the standard fitting procedure introduces methodological uncertainties due in particular to the necessity of choosing specific choices of parametrization. All of these uncertainties are of course all highly correlated. We discuss each in turn.
### 6.1 Theoretical Uncertainties
The most important theoretical uncertainty in the determination of parton densities is the truncation of the resummed perturbation series at NLO. Consistent NNLO determinations will require NNLO splitting functions: there has recently been some progress in this direction , and it is hoped that NNLO calculations might be available before the LHC is turned on. Meanwhile there are some ‘approximate NNLO’ calculations , which attempt to reconstruct the NNLO splitting functions from their known integer moments and behaviour at large and small $`x`$: these analyses suggest that NNLO corrections might reduce theoretical uncertainties due to truncation of the perturbative expansion by at least a factor of two.
One of the most important consequences of the theoretical uncertainty from unknown NNLO corrections is that it currently limits the accuracy of most of the experimentally more reliable determinations of $`\alpha _s`$. This in turn inevitably limits the accuracy of all extrapolations from low to high $`Q^2`$: for example one of the largest uncertainties in the prediction of the $`W`$ and $`Z`$ cross-sections is that due to the uncertainty in $`\alpha _s`$ .
Uncertainties at low $`Q^2`$ due to higher twist may be estimated from phenomenological fits: recent studies have shown that there are important correlations between empirical higher twist and the value of $`\alpha _s`$. It has also been shown that the fitted higher twist contribution drops when estimates of NNLO corrections are included . The empirical higher twist is qualitatively consistent with renormalon estimates. Taken together, these observations suggest that it is difficult to disentangle genuine higher twist from higher order perturbative corrections: the true higher twist contribution may be much smaller than is suggested by the fits.
The correct treatment of heavy quarks close to threshold was developed some time ago ; more recently it was proven that this procedure works to all orders in perturbation theory . This treatment is now included in some of the CTEQ fits ; a closely related but not identical procedure is used by MRS . A simpler version of ACOT, which nonetheless accurately reproduces its essential features, has also been developed .
An accurate treatment of heavy quark production, and indeed $`W`$ and Higgs production, requires the resummation of threshold logarithms. Recently it has been suggested that resummation of soft gluons may solve some of the problems with prompt photons . A fully consistent treatment will require the inclusion of soft gluon resummations in parton determinations, but as yet this has not been attempted. Renormalon studies suggest that such resummations may substantially improve the reliability of perturbation theory at large $`x`$. Again there will be strong correlations with higher twist. It would be particularly interesting to see the effect of such resummations on the predictions for the parton-parton luminosities eq.2 in the region relevant for Higgs production at the LHC.
The resummation of high energy (small $`x`$) logarithms is more problematic. Present data suggest that their effect on inclusive cross sections must be very small, at least at HERA and the Tevatron if not at the LHC. Furthermore, conventional theoretical approaches based on summations of LLx and NLLx corrections have been shown to break down: the NLLx corrections are overwhelmingly large and negative . Various suggestions for the resummation of these large corrections have been put forward . Hopefully a detailed phenomenological analysis based on one or other of these procedures will eventually provide a reliable estimate of the error due to uncertainties in small $`x`$ evolution when using parton distributions measured at HERA to predict those to be used at the LHC.
### 6.2 Combining Different Experiments
On the experimental side, one of the major problems with combining results from different experiments lies in the degree of ‘rigour’ in the interpretation of the experimental errors. Experimental results may be conveniently expressed as probabilities $`P(\mathrm{data}|\mathrm{theory})`$, i.e. the probabilities of obtaining the given set of data given a certain theoretical prediction . Often these probabilities are expressed in terms of predictions and (Gaussian) errors: for a given experiment, $`P(d|t)=\mathrm{exp}(\frac{1}{2}\chi ^2(d|t))`$, where $`d`$ are the data, $`t`$ the theoretical predictions and
$$\chi ^2(d|t)=\underset{\mathrm{data}}{}(dt)\mathrm{\Sigma }^1(dt)$$
(3)
where $`\mathrm{\Sigma }`$ is the matrix of correlated errors. Maximizing the probability, and thus obtaining the most likely ‘prediction’, then corresponds to minimizing the $`\chi ^2`$. It should be emphasized that it is not necessary to present experimental results in this way, and in particular some systematics may be completely non-Gaussian; however if the experiment is to be useful it must always provide a (clear) estimate of $`P(d|t)`$, otherwise the error analysis is at best incomplete and at worst useless.
In the present situation, the predictions will be constrained functionals of the input pdfs (the constraints being the result of perturbative evolution and cross-sections). If the errors have been estimated correctly, and the theory which constrains the predictions is sufficiently accurate, then there should be pdfs for which the $`\chi ^2`$ per degree of freedom is of order unity. Unfortunately for many important datasets this is not the case, and thus if one were to insist on the rigour of the statistical method, then many important experiments would not be included in the analysis . Such a strict criterion is probably unrealistic: rather the emphasis should be placed on using the maximal experimental constraints from experimental data . In this case the standard statistical techniques may not apply, but must be supplemented by physical considerations, taking into account experimental and theoretical limitations .
As an example of how this works in practice, we consider a recent CTEQ error analysis of the $`W`$-production cross-section . This uses the standard CTEQ5 analysis as the starting point: there are fifteen experimental data sets, with a total of $`1300`$ data points, and experimental errors are generally treated by ignoring correlations and combining statistical and systematic errors in quadrature (so $`\mathrm{\Sigma }`$ in eq.(3) is taken to be diagonal, with each diagonal entry set to $`\sigma _{\mathrm{stat}}^2+\sigma _{\mathrm{syst}}^2`$ of the corresponding data point). The initial pdfs are parameterised by $`18`$ parameters $`a_i,i=1,\mathrm{},18`$: each theoretical prediction is then a function of these parameters. The ‘best-fit’ distribution (CTEQ5M1 in this case) is then given by the set of parameters $`a`$ which minimise $`_{\mathrm{expts}}_{\mathrm{data}}\chi ^2(d|t[f(a)])`$, where $`t[f(a)]`$ are the theoretical predictions for each data point given the pdf $`f(a)`$ for the fifteen base experimental data sets.
A natural way to find the limits of a physical observable which depends on the pdfs, call it $`𝒪[f(a)]`$, such as the $`W`$-production cross-section $`\sigma _W`$ at $`\sqrt{s}=1.8`$ TeV, is then to study the dependence of the total $`\chi ^2`$ on $`𝒪`$. An efficient way of doing this is to use Lagrange’s method of undetermined multipliers: one minimizes
$$F(\lambda )=\underset{\mathrm{expts}}{}\underset{\mathrm{data}}{}\chi ^2(d|t[f(a)])+\lambda 𝒪[f(a)]$$
(4)
for fixed $`\lambda `$, and then varies $`\lambda `$ in order to map out the $`\chi ^2`$ as a function of $`𝒪`$.
Figs. 10a,b show the $`\chi ^2`$ for the fifteen base experimental data sets as a function of $`\sigma _W`$ at the Tevatron and LHC energies respectively . Two curves with points corresponding to specific global fits are included in each plot<sup>19</sup><sup>19</sup>19The third line in Figs. 10a refers to an alternative technique based on the assumption of Gaussian errors in the parameters $`a_i`$.: one obtained with all experimental normalizations fixed; the other with these included as fitting parameters (with the appropriate experimental errors). We see that the $`\chi ^2`$’s for the best fits corresponding to various values of the W cross-section are close to being parabolic, as expected. Indicated on the plots are 3% and 5% ranges for $`\sigma _W`$. The two curves for the Tevatron case are farther apart than for LHC, reflecting the fact that the W-production cross-section is more sensitive to the quark/anti-quark distributions and these are tightly constrained by existing DIS data.
The important question is: how large an increase in $`\chi ^2`$ should be taken to define the likely range of uncertainty in $`𝒪`$? The elementary statistical theorem that $`\mathrm{\Delta }\chi ^2=1`$ corresponds to one standard deviation of the measured quantity $`𝒪`$ relies on assuming that the errors are gaussian, uncorrelated, and with their magnitudes correctly estimated. Because these conditions do not hold here, this theorem cannot be naively applied quantitatively: rather one must examine in detail how well the fits along the parabolas shown in Fig.10 compare with the individual precision experiments included in the global analysis, in order to arrive at reasonable quantitative estimates on the uncertainty range for the W cross-section. In the meantime, based on past (admittedly subjective) experience with global fits, it seems that a $`\chi ^2`$ difference of $`4050`$ points represents a ‘reasonable’ estimate of current uncertainty of parton distributions. This implies that the uncertainty of $`\sigma _W`$ is about 3% at the Tevatron, and 5% at the LHC.
### 6.3 Correlated Experimental Systematics
There is now an increasing awareness of the necessity and possibility of carrying out a careful treatment of correlated systematic errors when attempting to determine errors on pdfs. For example a systematic study of the uncertainties in the parton distribution in the small $`x`$ region has been made recently by experimentalists at H1 and ZEUS . These studies include a proper treatment of correlated systematic errors, and some attempt is made to quantify parametrization uncertainties. Similar studies of the errors in polarized parton densities have been made by the SMC . Besides showing that careful estimates of parton uncertainties are useful and necessary, these studies also show that it is possible to include correlated systematics and combine data sets from different (albeit similar) experiments in a meaningful way. However they also show that doing something similar for a global parton determination would be very difficult and extremely tedious, unless new techniques are developed.
The importance of correlations in experimental systematic errors has been underlined by a recent reanalysis of the $`F_2`$ BCDMS data. A more careful treatment of the correlations between data taken at different beam energies, and the correlations between the fitted parton distributions and higher twist, results in a significant increase in the value of $`\alpha _s(M_Z^2)`$ extracted from the data: Alekhin quotes a value of $`0.118\pm 0.002`$. This is consistent with the current world average and the value $`0.119\pm 0.002`$ recently extracted from the reanalysed CCFR data (though after a more careful treatment of correlated higher twist this rises to $`0.122\pm 0.005`$).
In this context it should be noted that in the usual global analyses, in which correlations between systematic errors are ignored, and higher twist effects are not included, neither the BCDMS or the CCFR $`F_2`$ data show a minimum in their $`\chi ^2`$ as $`\alpha _s`$ is varied , despite the fact that when treated separately each is capable of yielding an excellent determination of $`\alpha _s`$. Only the minima in the H1 and ZEUS datasets are strong enough to survive this treatment: this may be helped by the fact that empirical higher twists are very small at small $`x`$ . It will be interesting to repeat the preliminary determination using the 95-97 HERA datasets when these finally become available.
### 6.4 Methodological Issues
While the issues addressed in the previous three sections are no doubt all important, there are also some methodological issues which need to be considered if we are to achieve our aim of a reliable determination of the errors in a global determination of parton distributions. In particular, we need a technique which can give parton distribution functions and their errors, such that:
(i) there is no inbuilt methodological bias (for example dependence on a particular parametrization of the input distributions)
(ii) it is easy to propagate the effects of correlated systematic errors in the data to correlated uncertainties in the parton distributions
(iii) it is easy to add new data sets or estimate theoretical errors or test models of new physics without redoing the whole of the analysis.
All of these criteria can be met if we ‘quantise’ our parton distributions: instead of trying to determine a single ‘best fit’ set of parameterised parton distributions with an associated error matrix, we construct an ensemble of sets of partons, distributed according to how well they fit the data . The expected result for a parton dependent observable, call it $`𝒪[f]`$, would then be given by an ensemble average:
$$𝒪[f]=𝒵^1[𝒟f]𝒪[f]J[f]s[f]\underset{\mathrm{expts}}{}P(d|t[f]),$$
(5)
where $`[𝒟f]`$ means functional integration over all possible input distributions $`f`$ (subject to basic constraints such as sum rules and positivity) and $`𝒵=1`$ is a normalization factor. The measure of integration is given essentially by the probability distributions $`P(d|t[f])`$ for each of the experiments used as input. These probabilities are, as explained above, the essential input of the experimental data used in the fit: they support distributions which fit the data well, and suppress the contribution of distributions which fit badly. If the errors on the data were assumed Gaussian, these probabilities would come in the form of a $`\chi ^2`$, as in eq.(3), though the technique does not depend on such an assumption, and non Gaussian errors could also be incorporated. There is also a Jacobian factor $`J[f]`$, which turns the integration measure from an integration over theoretical predictions $`t[f]`$ to one over the pdfs themselves, and enforces the theoretical constraint that the theoretical predictions are related through pQCD. It is also necessary to introduce a ‘smoothness’ factor $`s[f]`$ into the measure, to enforce the natural theoretical prejudice that the initial pdfs should be smooth functions of $`x`$, without wiggles or jumps: a suitable form for such a factor would be $`\mathrm{exp}\frac{1}{2}\epsilon _x(_xf)^2`$, where $`\epsilon `$ is a small parameter which quantifies the extent of this prejudice. Final results should be independent of the form of this term, and in particular the parameter $`\epsilon `$, provided that it is varied in a suitable range.
The way in which this procedure works should now be clear, since it is similar to the quantum mechanics (or more precisely statistical mechanics) of a particle in a (highly nonlocal) potential : the parton distributions may be thought of as quantum fields, with, in the case of Gaussian experimental errors, the action
$$𝒮[f]=\frac{1}{2}\underset{\mathrm{expts}}{}\underset{\mathrm{data}}{}(dt[f])\mathrm{\Sigma }^1(dt[f])+\frac{1}{2}\epsilon \underset{x}{}(_xf)^2.$$
(6)
The best fit parton distribution is then the solution of the classical equations of motion (since it minimises the action), while the error bands are given by the ‘quantum’ fluctuations around the classical field. Since the determination of the classical field is itself nontrivial, the system is best solved numerically: we discretise the field by introducing a parametrization with a finite number of parameters $`a_i`$, $`i=1,\mathrm{},N`$, so that $`[𝒟f]J[f]_ida_iJ(a_i)`$, rather as we would for a lattice field theory. Here the best discretization would not necessarily be a naive discretization in $`x_{\mathrm{Bj}}`$ with spline interpolation: rather it might involve expansion of each pdf in sets of orthogonal polynomials, or other sets of (orthogonal) functions, for example eq.(1) and its obvious generalizations. The integration over the parameters $`a`$ would then be done by Monte Carlo, using an algorithm such as Metropolis or HMC to generate an ensemble of configurations distributed according to the measure of integration, and thus according to its likelihood given the input datasets. Finding each such configuration will involve a similar computational effort to that of finding a best fit. Finally, we would like to increase the number of parameters $`N`$ (taking the ‘continuum limit’) until we are sufficiently close to a truly parametrization independent ensemble, at which stage we can readily compute expectation values of observables and their associated errors as averages over the ensemble of pdfs.
This procedure has several advantages:
(i) it is intrinsically parametrization independent as the number of parameters increases, because of the universality of the continuum limit. Flat directions are no longer the problem that they are in a best fit procedure: the total number of parameters is now limited only by computational resources. Indeed the flat directions are now interesting, since they give the most important uncertainties in the parton distribution functions.
(ii) the propagation of correlated systematics is automatically taken care of by the procedure. The only limitation is the reliability of the probabilities $`P(d|t)`$ produced by experimentalists. This should give added impetus to the determination of meaningful (and thus comparable) estimates of systematic errors by different experimental collaborations, and their presentation in such a way that they can be readily input into such an analysis. Preliminary explorations of the technique indicate that the errors in the pdf parameters are not only highly correlated, but also in many cases significantly non-Gaussian, even when the errors in the data are assumed to be Gaussian.
(iii) Data from new experiments can be added using the old configurations, since different experiments are (in principle!) statistically independent, so $`𝒮_{\mathrm{tot}}[f]=_{\mathrm{expts}}𝒮_{\mathrm{exp}}[f]`$. Similarly we could estimate theoretical errors due, for example, to NLO truncation, by using the standard configurations reweighted by varying renormalization and factorization scales. Similarly, we could test the effect of resummations by reweighting the configurations generated using standard NLO evolution, or indeed test models for new physics by reweighting the configurations generated using the Standard Model .
The main problems to be faced in actually implementing the procedure are computational: we need a fast evolution code, and high performance computing. The advantages of parallelization should be obvious. In fact the computational requirements are very similar to those of the lattice gauge theorists: calculating the ‘action’ is more difficult, but the ‘continuum limit’ should be reached much more quickly.
## 7 From here to the LHC and Beyond
### 7.1 Progress Before the LHC Turns on
Perturbative QCD has been extremely successful in describing data in DIS, DY and jet production, as well as describing the evolution of parton distributions over a wide range of $`x`$ and $`Q^2`$. From the point of view of pdf determination, the primary problem lies in the calculation of the direct photon cross sections which could serve as a primary probe of the gluon distribution at high $`x`$. However, a rigorous theoretical treatment of soft gluon effects (perhaps requiring both $`k_T`$ and Sudakov resummation) will be required before the data can be used with confidence in pdf fits. On the experimental side, it will also be necessary to resolve the inconsistency between the WA70 and E706 data.
D0 has recently presented a new result for the measurement of the inclusive jet cross section as a function of the jet rapidity (up to values of three) . Such a measurement probes a greater kinematic range than the central inclusive jet cross sections. In addition, the differential dijet data from the Tevatron explore a wider kinematic range than the inclusive jet cross section. Both CDF and D0 have dijet cross section measurements from Run I which may also serve probe the high $`x`$ gluon distribution, in regions where new physics is not expected but where any parton distribution shifts should be observable. The ability to perform such cross-checks is essential.
CDF and D0 will accumulate on the order of 2-4 $`fb^1`$ of data in Run II (from 2000-2003), a factor of 20-40 greater than the current sample. This sample should allow for more detailed information on parton distributions to be extracted from direct photon and DY data, as well as from jet production. Run III (2003-2007) offers a data sample potentially as large as 30 $`fb^1`$.
H1 and ZEUS will continue the analysis of the data taken with positrons in 1991-97. HERA switched to electron running in 1998 and plans to deliver approximately 60 In 2000, the HERA machine will be upgraded for high luminosity running, with yearly rates of 150 integrated luminosity of about 1 $`fb^1`$ by 2005. This will allow for an error of a few percent on the structure function $`F_2`$ for $`Q^2`$ scales up to $`10^4GeV^2`$. The gluon density, derived from scaling violations of $`F_2`$, should be known to an accuracy of less than 3% in the kinematic range $`10^4<x<10^1`$.
It is also hoped that over the next five years the Monte Carlo outlined in the previous section will begin to bear fruit, perhaps to the point where they can make a serious contribution to global pdf error analysis.
### 7.2 Physics cross sections at the LHC and the role of LHC data in pdf determination
ATLAS measurements of DY (including W and Z), direct photon, jet and top production will be extremely useful in determining pdfs relevant for the LHC. The data can be input to the global fitting programs, where it will serve to confirm/constrain the pdfs in the LHC range. Again, DY production will provide information on the quark (and anti-quark) distributions while direct photon, jet and top production will provide, in addition, information on the gluon distribution.
Other processes might also prove useful. For example diphoton production might be useful for determining the gluon distribution, and this in turn would lead to an improved knowledge of the relevant parton pdfs and parton-parton luminosity functions for the production of the Higgs (which is largely due to $`gg`$ scattering for low to moderate Higgs’ masses).
Another possibility that has been suggested is to directly determine parton-parton luminosities (and not the parton distributions per se) by measuring well-known processes such as W/Z production . This technique would not only determine the product of parton distributions in the relevant kinematic range but would also eliminate the difficult measurement of the proton-proton luminosity. It may be more pragmatic, though, to continue to separate out the measurements of parton pdfs (through global analyses which may contain LHC data) and of the proton-proton luminosity. The measurement of the latter quantity can be pegged to well-known cross sections, such as that of the W/Z, as has been suggested for the Tevatron.
## 8 Conclusions
The determination of parton distributions and uncertainties is an important ingredient of our preparations for physics at the LHC. The global fitting techniques used for the past fifteen years may soon be superseded by more sophisticated methods. Developing and exploiting these techniques will be a great challenge to theorists and experimentalists alike.
## 9 Acknowledgements
RDB would like to thank Sergey Alekhin, John Collins, Tony Doyle, Stefano Forte, Stefane Keller, Tony Kennedy, David Kosower, Brian Pendleton, Dave Soper, James Stirling, Wu-Ki Tung and Andreas Vogt for various stimulating and useful discussions. JH would like to thank James Stirling, Steve Mrenna and his CTEQ colleagues for useful comments. We would also like to thank James Stirling and Lenny Apanasevich for providing many of the figures. This work was supported in part by an EU TMR contract FMRX-CT98-0194 (DG 12 - MIHT) and by the NSF under grant PHY-9901946.
Generalized factorization and resummation
C. Balázs, J.C. Collins, D.E. Soper
## 1 The Collins-Soper-Sterman formalism
The standard factorization formula fails near kinematic boundaries. We discuss the case of low transverse momentum in $`Z_0`$ production, etc; this is an important case because the cross section peaks there. The failure of the factorization formula is symptomized by large corrections involving a factor of $`\mathrm{ln}^2Q/Q_T`$ for each power of $`\alpha _s`$.
Although the solutions to the problem are all commonly referred to as “resummations”, there are in fact two very different approaches. One is resummation in its strict sense: One performs a selective and approximation summation of the largest parts of the perturbative series for the hard scattering in the standard factorization formalism.
The second approach is that of Collins, Soper and Sterman (CSS) . These authors observed that the conventional factorization formalism is in fact wrong at low transverse momentum and they derive a correct factorization for this region. In an intermediate region of transverse momentum, the standard factorization with resummation is applicable with somewhat reduced accuracy, and there is an overlap between the two approaches, which we will discuss later.
In any case, it is essential to improve on the standard fixed-order factorization formalism, and the reward is an improved method that
* includes large, logarithmic QCD corrections up to all orders in the strong coupling,
* improves the renormalization scale dependence of the prediction,
* enables prediction of certain quantities reliably, which cannot be done in a fixed order calculation,
* provides an independent, analytic check for parton shower Monte Carlo’s.
### 1.1 $`k_T`$-dependent parton densities
CSS realized that the failure of the standard factorization when $`Q_TQ`$ occurs because it neglects the transverse motion of the incoming partons in the hard scattering. (Here $`Q`$ can be the invariant mass of a colorless particle, or set of particles, created in a hard partonic collision, and $`Q_T`$ is the related transverse momentum.) The approximation of neglecting parton transverse momentum is only valid when the cross section is integrated over a large range of $`Q_T`$. But if, for example, $`Q_T`$ is of order 1 GeV, then we are outside of the domain in which the factorization is applicable.
A fully satisfactory approach must use a factorization theorem that is valid for any $`Q_T`$ that is small compared to $`Q`$. CSS’s theorem gives the cross section as a convolution of transverse momentum distributions
$$\frac{d\sigma }{d^4Q}d^2k_TP(x_1,k_T)P(x_2,Q_Tk_T),$$
(1)
where $`P`$ is a partonic density distribution that is a function of both longitudinal ($`x`$) and transverse ($`k_T`$) momenta. The partonic recoil against soft gluons as well as the intrinsic partonic transverse momentum are included in $`P`$.
Such a treatment completely formalizes the intuitive notion that partons must have transverse momentum and that this transverse momentum gives rise to transverse momentum of the Drell-Yan pair. There is then no need to convolute a calculated cross section with “intrinsic transverse momentum” for the quarks; this manoeuvre is only necessary as an ad hoc correction to a formalism that is incomplete.
In QCD, complications arise from soft-gluon effects, because these effects do not cancel, in contrast to the case of the cross section integrated over $`Q_T`$. A consequence, proved by CSS, is a particular form of the evolution equations for the $`k_T`$-dependent parton densities. These equations are *not* the normal DGLAP equations<sup>1</sup><sup>1</sup>1 Although all the physics associated with the DGLAP equations is present. . The kernel of the evolution contains a perturbatively calculable part and non-perturbative part. The non-perturbative part can be summarized by saying that there is a fixed amount of gluon radiation per unit rapidity, so that the transverse momentum distribution of the partons broadens in a characteristic way with energy. The non-perturbative part of this energy-dependent radiation is fitted by the $`g_2`$ term of Eq.(1.3) below.
This feature may be the dominant reason why transverse momentum distributions are so broad at high energies, as in $`Z^0`$ production: the transverse momentum of the $`Z^0`$ has a component due to the recoil against non-perturbative glue emitted into many units of rapidity.
The CSS formalism clearly entails a phenomenological fitting of the non-perturbative part of the $`k_T`$-dependent parton densities and the evolution kernel. In principle, this can be done at fairly low energy, and then the evolution equations predict the results for higher energies with no further adjustable parameters. The more conventional resummation formalism is compatible with the CSS formalism, but it is not as complete.
Because the CSS formalism is designed to treat correctly the $`Q_TQ`$ region, it also provides an appropriate resummation of the large logarithms, $`\mathrm{ln}(Q/Q_T)`$ in the standard factorization formula.
We can gauge how important these logarithms are in practice by examining the cross section for $`Z`$ production at the Tevatron. The bulk of the cross section is in the low $`Q_T`$ region, and, as can be seen from Fig. 1, there is a peak at around $`Q_T=2.7`$ GeV, which is much smaller than the invariant mass $`Q=m_Z=91.187`$ GeV. This implies that for the bulk of the events $`\mathrm{ln}(Q/Q_T)`$ is large enough that $`\alpha _s(m_Z)\mathrm{ln}^2(Q/Q_T)>1`$. Since we have a double logarithm for each radiated gluon, higher orders in the perturbative series are not suppressed.
### 1.2 From fixed order to resummed
In this section we show how the results of the standard factorization theorem are related to a resummation in terms of leading logarithms, etc.
When the $`Z^0`$ is produced in a hadron-hadron collision its transverse momentum is balanced by some hadronic activity which stems from partons emitted by the initial state partons. (In the first order in the strong coupling a $`Z^0`$ and a gluon is produced.) The $`Q_T`$ distribution given by the usual factorization in the low $`Q_T`$ region is written as
$$\underset{Q_T0}{lim}\frac{d\sigma }{dQ_T^2}=\underset{n=1}{\overset{\mathrm{}}{}}\underset{m=0}{\overset{2n1}{}}\alpha _s^n\frac{{}_{n}{}^{}v_{m}^{}}{Q_T^2}\mathrm{ln}^m\left(\frac{Q^2}{Q_T^2}\right)+𝒪\left(\frac{1}{Q_T}\right),$$
(2)
where the coefficients $`{}_{n}{}^{}v_{m}^{}`$are perturbatively calculable apart from some factors of parton densities. When the two scales $`Q`$ and $`Q_T`$ are very different, the logarithmic terms $`\mathrm{ln}^m(Q^2/Q_T^2)`$ are large, and for $`Q_TQ`$ the perturbative series is dominated by these terms. For $`Q_TQ`$ truncation of the perturbative series, i.e. any fixed order calculation, gives an answer which neglects these important all order logarithmic contributions. At the lowest order, $`𝒪(\alpha _s^0)`$, the $`Z^0`$ boson is produced alone, that is with a $`Q_T`$ distribution of $`\delta (Q_T)`$. The singularity at $`Q_T=0`$ prevails at any fixed order in $`\alpha _s`$, as Eq. (2) shows.
One way of reorganizing the perturbation series is to make the expansion one in terms of $`\alpha _s\mathrm{ln}^2(Q/Q_T)`$ instead of $`\alpha _s`$ itself. In this simplified picture, calculating fixed order QCD corrections means calculating the perturbative series
$`\underset{Q_T0}{lim}{\displaystyle \frac{d\sigma }{dQ_T^2}}=`$
$`Q_T^2\{\alpha _s({}_{1}{}^{}v_{1}^{}L+{}_{1}{}^{}v_{0}^{})+\alpha _s^2({}_{2}{}^{}v_{3}^{}L^3+{}_{2}{}^{}v_{2}^{}L^2)+\alpha _s^3({}_{3}{}^{}v_{5}^{}L^5+{}_{3}{}^{}v_{4}^{}L^4)+\mathrm{}`$
$`+\alpha _s^2({}_{2}{}^{}v_{1}^{}L_2+{}_{2}{}^{}v_{0}^{}L^0)+\alpha _s^3({}_{3}{}^{}v_{3}^{}L^3+{}_{3}{}^{}v_{2}^{}L^2)+\mathrm{}`$
$`+\mathrm{}\mathrm{}\},`$
column by column. In the leading logarithm approach, on the other hand, we calculate the above series line by line . While in the fixed order (column by column) calculation the convergence for low $`Q_T`$ is spoiled by the higher order uncalculated logs ($`L=\mathrm{ln}(Q/Q_T)`$), in the resummed (line by line) calculation convergence is preserved in each “order” (by each line), and higher order corrections are systematically included.
### 1.3 The CSS formula
The improved factorization theorem of CSS together with their evolution equation for the $`k_T`$ dependent parton distributions, leads to a useful formula<sup>2</sup><sup>2</sup>2 While solving the RGE, an integro-differential equation, specific choices of integration constants were made (c.f. Ref. ): $`C_1=C_3=2e^{\gamma _E}C_0`$ and $`C_2=C_4=1`$, to optimize logarithmic contributions. This is similar to the $`\mu =Q`$ choice in case of the ultraviolet renormalization, to make terms like $`\mathrm{ln}(\mu /Q)`$ vanish. for the cross section. For $`Z^0`$ production it can be written as
$$\frac{d\sigma (h_1h_2Z^0X)}{dQ^2dQ_T^2dy}=\underset{j}{}\sigma _{0,j}W_{j\overline{ȷ}}(Q,Q_T,x_1,x_2)+Y(Q,Q_T,x_1,x_2),$$
(3)
where the “resummed” part, $`W(Q,Q_T,x_1,x_2)`$, is defined as
$`W_{j\overline{ȷ}}(Q,Q_T,x_1,x_2)`$ $`=`$
$`{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle d^2be^{i\stackrel{}{Q}_T\stackrel{}{b}}𝒞_{j/h_1}(Q,b_,x_1,\mu )e^{𝒮(Q,b_{})}𝒞_{\overline{ȷ}/h_2}(Q,b,x_2,\mu )}`$
for a given partonic initial state with flavor $`j`$.<sup>3</sup><sup>3</sup>3The lowest order partonic total cross section is $`\sigma _{0,j}=\pi ^2g^2((14Q_js_w^2)^21)/(48Q^2c_w^2)`$, where $`g`$ is the weak coupling constant, $`s_w^2`$ ($`c_w^2`$) is the sine (cosine) of the weak mixing angle squared, and $`Q_j`$ is the charge of the quark flavor $`j`$. The Fourier integral is introduced because transverse momentum conservation is explicit in the impact parameter, $`b`$, space .
All the dangerous logarithms are included in the perturbative Sudakov exponent
$$𝒮(Q,b_{})=_{C_0^2/b_{}^2}^{Q^2}\frac{d\overline{\mu }^2}{\overline{\mu }^2}\left[A\left(\alpha _s(\overline{\mu })\right)\mathrm{ln}\left(\frac{Q^2}{\overline{\mu }^2}\right)+B\left(\alpha _s(\overline{\mu })\right)\right].$$
(5)
Here $`C_0`$ is an arbitrary parameter which cuts off the perturbative low $`Q_T`$ region.<sup>4</sup><sup>4</sup>4In practice $`C_0=2e^{\gamma _E}`$ is used, which is related to the values of the integration constants of the RGE for the $`k_T`$ dependent PDF’s. To prevent perturbative calculations from being done in region where perturbation theory is inapplicable, the “impact parameter” $`b`$ in the Sudakov exponent was replaced by
$$b_{}=\frac{b}{\sqrt{1+(b/b_{\mathrm{max}})^2}}.$$
(6)
The errors caused by this replacement are of the same form as the non-perturbative contributions to be discussed below, and are therefore correctly treated by being absorbed into the non-perturbative part of the formula.
The $`A`$ and $`B`$ functions are free of large logarithms and can be reliably calculated perturbatively for a given process as
$$A\left(\alpha _s(\overline{\mu })\right)=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{\alpha _s(\overline{\mu })}{\pi }\right)^nA^{(n)},B\left(\alpha _s(\overline{\mu })\right)=\underset{n=1}{\overset{\mathrm{}}{}}\left(\frac{\alpha _s(\overline{\mu })}{\pi }\right)^nB^{(n)}.$$
(7)
The distributions
$$𝒞_{j/h}(Q,b,x,\mu )=\underset{a}{}_x^1\frac{d\xi }{\xi }C_{ja}(b_{},\frac{x}{\xi },\mu )f_{a/h}(\xi ,\mu )_{a/h}(b,x)e^{r(b)\mathrm{ln}Q}$$
(8)
depend on virtual and real emission contributions for a given process, via the Wilson coefficients $`C_{ja}`$. Just as the $`A`$ and $`B`$ functions the Wilson coefficients are expanded in terms of the strong coupling $`\alpha _s`$
$$C_{ij}(b_{},z,\mu )=\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{\alpha _s(\mu )}{\pi }\right)^nC_{ij}^{(n)}(z,b_{}).$$
(9)
Since the Sudakov exponent integrates to unity, the $`C_{ij}`$ function sets the normalization of the resummed distribution. In particular, if coefficients up to $`C_{ij}^{(n)}`$ are included in the calculation then the resummed rate will equal the rate calculated in fixed order at $`𝒪(\alpha _s^n)`$ . The function $`f_{a/h}(x,\mu )`$ is the usual renormalized momentum fraction ($`x`$) distribution of parton $`a`$ in hadron $`h`$ at the energy scale $`\mu `$. Observe that the impact parameter dependence of the perturbative coefficient functions is cut off by the use of $`b_{}`$ instead of $`b`$.
Included in Eq. (8) are two non-perturbative factors, $`_{a/h}(b,x)`$ and $`e^{r(b)\mathrm{ln}Q}`$. These implement the parts of the CSS factorization and evolution equation that cannot be implemented as a resummation of the standard factorization theorem. They also compensate for the errors in the resummation at large $`b`$. The overall effect is that (8) define $`k_T`$-dependent parton densities. The $``$ factor can be interpreted as allowing for intrinsic transverse momentum, and the $`e^{r(b)\mathrm{ln}Q}`$ factor allows for the recoil against soft gluon radiation. The $`\mathrm{ln}Q`$ in the exponent of the soft-gluon factor comes from the solution of the CSS evolution equation and can be interpreted by saying that soft gluons are emitted uniformly in rapidity.
The perturbative part of the formula uses $`b_{}`$ instead of $`b`$, as defined by Eq. (6). The parameter $`b_{\mathrm{max}}`$ provides an infra-red cutoff on the perturbative part of the formula. In practice the empirically optimal value, $`b_{\mathrm{max}}=1/2`$ GeV<sup>-1</sup>, is used. This arbitrary cutoff of the $`b`$ integration is compensated by the parameterization of the non-perturbative part of the formula, which is
$`W_{ij}^{\mathrm{NP}}(Q,b,x_1,x_2)`$ $`=`$ $`_{i/h_1}(Q,b,x_1)_{j/h_2}(Q,b,x_2)e^{r(b)\mathrm{ln}Q}`$
$`=`$ $`\mathrm{exp}\left[g_1b^2g_2b^2\mathrm{ln}\left({\displaystyle \frac{Q}{2Q_0}}\right)g_1g_3b\mathrm{ln}(100x_1x_2)\right],`$
where $`Q_0`$ is chosen to be the initial scale of the parton evolution<sup>5</sup><sup>5</sup>5For recent CTEQ PDF’s $`Q_0=1.6`$ GeV. and the $`g_i`$ parameters have to be determined using experimental data.<sup>6</sup><sup>6</sup>6The $`\mathrm{ln}\left(Q^2/Q_0^2\right)`$ term is introduced to match the logarithmic term of the Sudakov exponent and its coefficient is expected to be process independent, depending only on the initial partonic state.
### 1.4 Matching
The resummed term, defined by Eq. (1.3), was derived in the context of a generalized factorization, under the assumption that $`Q_TQ`$. This assumption will break down within and beyond the intermediate $`Q_T\stackrel{<}{}Q`$ region. In the high $`Q_T`$ region (where $`Q_T\stackrel{>}{}Q`$) the conventional perturbative factorization formalism is reliable. To obtain sufficiently accurate results for all $`Q_T`$, it is necessary to combine the formalisms.
The $`Y`$ term in Eq. (3) was introduced by CSS to correct the behavior of the resummed piece in the intermediate and high $`Q_T`$ regions.<sup>7</sup><sup>7</sup>7The exact definition of the $`Y`$ piece for $`Z^0`$ production can be found in Refs. . It is defined as the difference of the cross section calculated from the standard factorization formula at a fixed order $`n`$ of perturbation theory and the $`Q_TQ`$ asymptote of this cross section:
$$Y(Q,Q_T,x_1,x_2)=\left(\frac{d\sigma }{dQ^2dQ_T^2dy}\right)_n\left(\frac{d\sigma }{dQ^2dQ_T^2dy}\right)_{n,Q_TQ}.$$
(11)
Thus, the full CSS formula can be written as
$$\frac{d\sigma }{dQ^2dQ_T^2dy}=\left(\frac{d\sigma }{dQ^2dQ_T^2dy}\right)_{\mathrm{res}}+\left(\frac{d\sigma }{dQ^2dQ_T^2dy}\right)_n\left(\frac{d\sigma }{dQ^2dQ_T^2dy}\right)_{n,Q_TQ}.$$
(12)
This method of matching the resummed and fixed order pieces is valid because the low $`Q_T`$ asymptote used in Eq. (11) is the same as the large $`Q_T`$ asymptote of the resummed term $`W`$. At low $`Q_T`$ the asymptotic part dominates the $`Q_T`$ distribution (the logs are large), and the last two terms cancel in Eq.(12), while the resummed term is significant near $`Q_T=0`$. At high $`Q_T`$ the logs are small, and the expansion of the resummed term cancels the $`Q_T`$ singular terms up to higher orders in $`\alpha _s`$.<sup>8</sup><sup>8</sup>8The cancellation is higher order than the order at which the singular pieces were calculated. In this situation the first and third terms cancel and CSS formula reduces to the fixed order perturbative result. After matching the resummed and fixed order cross sections in such a “smooth” manner, it is expected that the normalization of the CSS cross section reproduces the fixed order total rate, since when expanded and integrated over $`Q_T`$ it deviates from the fixed order result only in small higher order terms in $`\alpha _s`$ .
Unfortunately the above argument does not completely work in practice. The problem arises because at large $`Q_T`$ the $`W`$ term in Eq. (3) is an extrapolation of the cross section from small $`Q_T`$. So it has a $`1/Q_T^2`$ behavior, modified by logarithms. This falls less steeply than the true cross section, which is subject to kinematic limits. The errors in the CSS formula at large $`Q_T`$ are indeed suppressed by a power of $`\alpha _s`$. But the coefficient of this power is the $`1/Q_T^2`$ part of the formula, and so the error can be easily larger than the true cross section. A symptom of the problem is that the cross section calculated from Eq. (3) is typically negative at large enough $`Q_T`$.
One possible remedy is to abandon the CSS formalism. But we regard this as undesirable, because it also abandons the important physical result of CSS that goes beyond mere resummation: their proper treatment of non-perturbative transverse momentum.
A second, commonly used remedy, is to utilize the fact that in the high $`Q_T`$ region the fixed order result is a good description of the distribution. So when calculating the $`Q_T`$ distribution one can simply switch from the CSS to the fixed order distribution whenever they cross for high $`Q_T`$’s. Since the mismatch between the resummed and the asymptotic terms in Eq.(12) decreases as the perturbative order of the calculation ($`n`$) increases, it is expected that the crossing point shifts toward $`Q_T=Q`$, and the slope of the resummed and fixed order curves approaches each other as $`n`$ increases (cf. Ref. ). Indeed, calculations at $`𝒪(\alpha _s^2)`$ blend closer to $`m_Z`$, and smoother than at $`𝒪(\alpha _s)`$, as shown in Fig. 2.
When this prescription for the switching is followed at the fully differential, $`d\sigma /dQ_TdQdy`$, level the result is a smooth and differentiable $`Q_T`$ distribution, after the invariant mass and rapidity is integrated out. This is illustrated in Fig. 3. It was shown in Ref. that the integral of the $`Q_T`$ distribution calculated using this prescription recovers the fixed order total rate within an error which is the size of the higher orders, as it is expected.
### 1.5 Improved matching
Since the calculations of $`W`$ and $`Y`$ are done using truncations of perturbation theory, the switching between calculational methods introduces an artificial discontinuity in the slope of the cross section. This practical problem arises in the matching because of a mismatch of the orders of perturbation theory at which $`W`$ and $`Y`$ are calculated. From the point of view of a standard factorization calculation, $`W`$ contains a selective summation of arbitrarily high orders of perturbation theory. The possibility of getting such a resummation relies on performing certain approximations that are only valid at small $`Q_T`$. The difficulty of performing complete higher-order calculation means that $`Y`$ can only be calculated at fixed order.
At large transverse momentum, $`|W|`$ is much larger than the actual cross section, and so the cross section Eq. (12) is obtained by the cancellation of two almost equal terms. This is clearly a recipe for bad numerical work.
Examination of the lowest-order calculation of $`Y`$, for the $`q\overline{q}`$ annihilation term shows some of the sources of the problems:
$`Y`$ $`=`$ $`{\displaystyle \frac{C}{Q_T^2}}{\displaystyle \frac{d\xi _1}{\xi _1}\frac{d\xi _2}{\xi _2}f_q(\xi _1)f_{\overline{q}}(\xi _2)}`$
$`\{{\displaystyle \frac{(Q^2\widehat{t})^2+(Q^2\widehat{u})^2}{\widehat{s}}}\delta (\widehat{s}+\widehat{t}+\widehat{u}Q^2)`$
$`2\delta (1x_1/\xi _1)\delta (1x_2/\xi _2)\left[\mathrm{ln}(Q^2/Q_T^2){\displaystyle \frac{3}{2}}\right]`$
$`\delta (1x_1/\xi _1)\left[{\displaystyle \frac{1+x_2^2/\xi _2^2}{1x_2/\xi _2}}\right]_+\left[{\displaystyle \frac{1+x_1^2/\xi _1^2}{1x_1/\xi _1}}\right]_+\delta (1x_2/\xi _2)\}.`$
Here $`x_1`$ and $`x_2`$ are the longitudinal momentum fractions of the Drell-Yan pair. The first term contains the usual perturbative calculation of the differential cross section, and the other 3 terms give the negative of its low $`Q_T`$ asymptote. The intrinsic rate of fall off of the cross section with $`Q_T`$ is given by the explicit $`1/Q_T^2`$ factor which is present in the parton cross section. But an extra fall off is caused by the fact that the parton densities are probed at larger fractional momenta when $`Q_T`$ is increased.
Some symptoms of the problems can already be seen. One is that the first subtraction term, on the second line of Eq. (1.5), changes sign at large $`Q_T`$: the extrapolation of a positive cross section becomes negative. The second is the plus distribution in the third line; if the parton distributions are steeply falling, the plus distributions give a misleading size for the integrand. This last effect really indicates that there is an additional scale in the process, so that the relevant scales are:
* The transverse momentum $`Q_T`$ of the Drell-Yan pair.
* The invariant mass $`Q`$ of the pair.
* The increase $`\mathrm{\Delta }Q`$ of $`Q`$ that is necessary to make the typical parton densities in the factorization formula decrease by a factor 2.
We believe the overall approach of a subtraction method is correct: $`W`$ correctly represents the physics at low $`Q_T`$, and we do not wish to give up a method that uses the intuitive notion of $`k_T`$-dependent parton densities. We therefore cannot expect to obtain a perfect estimate of the large $`Q_T`$ cross section from $`W`$ alone. The idea of adding a correction term $`Y`$ is a good way of combining the information in standard fixed order calculations with the resummed calculations.
But improvements in its implementation are needed. We suggest the following strategies that could be tried, individually or even in combination:
* Multiply $`W`$ by an ad hoc factor $`F(Q_T/M)`$. Correspondingly the formula for the subtraction term in $`Y`$ will also have the same factor. The parameter $`M`$ is in principle arbitrary, and it should be chosen so that the fall off in the modified $`W`$ term mimics that of the actual cross section. The cut-off function obeys $`F(0)=1`$, so that the small $`Q_T`$ behavior is unchanged, and the function should be zero for large $`Q_T`$.
* Change the argument of $`W`$ from $`Q_T`$ to some other function of $`Q_T`$. One possible choice would be $`Q_T^{}=Q_T/(1Q_T/M)`$, where $`M`$ is again a parameter to be chosen. One would replace $`W`$ by zero if $`Q_T>M`$. The effect of the variable change is to leave $`W`$ unaltered at small $`Q_T`$ and to give a more rapid fall off at large $`Q_T`$. Again one would make an identical redefinition in the subtraction term in $`Y`$.
* Redefine the $`+`$ distributions such as those in Eq. (1.5), by:
$$_0^1𝑑zf(z)\left[\frac{1}{z}\right]_{+,z_0}=_0^1𝑑z\frac{1}{z}\left[f(z)f(0)\theta (z_0z)\right].$$
(14)
(The usual definition has $`z_0=1`$.)
In each case we have a generalized renormalization-group invariance of the exact cross section under changes of the parameter $`M`$ or $`z_0`$. But approximations obtained by truncation of a perturbation series are invariant only up to a term of order the first uncalculated correction. The aim is to choose the parameters on physical grounds to be such as to keep these higher order terms small, to eliminate their reason(s) for being large.
### 1.6 Applications
Beyond $`Z^0`$ production, in its present form, the CSS formalism can be applied in hadron-hadron collisions whenever the final state is colorless. The phenomenological significance of this ”transverse momentum resummation” ranges from Drell-Yan pair production, through lepton pair production via $`W^\pm `$ and $`Z^0`$ bosons , di-gauge boson (e.g. photon or $`Z^0`$ boson pair) production , to Higgs production . In recent years it was tested in hadronic processes taking place at fixed target (e.g. in DY photon and diphoton production) and collider energies (e.g. in DY, $`W^\pm `$, $`Z^0`$, and diphoton production). It was applied for different hadronic initial states in pion–nucleon, proton–nucleon, and proton–anti-proton collisions. It was also modified and tested for DIS processes . Finally, since was first devised for the calculation of the energy correlation of jets in $`e^+e^{}`$ collisions , it can be used in jet production at lepton colliders. Such a wide variety of applicability, and good agreement with existing experimental results for different processes, colliders, center of mass energies, and initial states gives us a confidence in the resummed predictions for the LHC.
## 2 Higgs production
At the LHC the SM Higgs boson will be mainly produced through the gluon fusion subprocess via a top quark loop: $`gg`$ (top quark loop) $`HX`$ . The Higgs boson can be detected in its $`H\gamma \gamma `$ decay mode, if its mass is in the 100-150 GeV range . If the Higgs mass is higher than about 130 GeV then its $`HZ^0Z^0`$ decay mode is the cleanest and most significant . To distinguish these signals from the substantial QCD background, besides the sharp peak in the invariant mass distribution, the most straightforward measurable to use is the transverse momentum. According to earlier studies, a statistical significance on the order of 5-10 can be reached for the inclusive $`H\gamma \gamma `$ signal, actual values depending on luminosity and background estimates. Once their transverse momentum distribution is reliably predicted, the difference in the $`Q_T`$ of the signal and background can be utilized to devise kinematic cuts to enhance the statistical significance of the signal. After the discovery, when determining the properties of the Higgs boson, besides the total cross section and the invariant mass distribution, the simplest and most fundamental measurable to use is the transverse momentum. For a recently proposed new detection mode, $`H\gamma \gamma \text{jet}`$, in Ref. it was also found that in order to optimize the significance it is necessary to impose a 30 GeV cut on the transverse momentum of the jet, or equivalently (at NLO precision), on the $`Q_T`$ of the photon pair. With this cut in place extraction of the signal in the Higgs plus jet mode requires the precise knowledge of both the signal and background distributions in the mid- to high-$`Q_T`$ region.
To reliably predict the $`Q_T`$ distribution of Higgs bosons at the LHC, especially in the low to mid $`Q_T`$ region where the bulk of the rate is, the effects of the multiple soft–gluon emission have to be included. In practice, performing soft gluon resummation within the CSS formalism is equivalent to the determination of the $`A^{(n)}`$, $`B^{(n)}`$, and $`C^{(n)}`$ coefficients and the $`Y`$ part at some order in $`\alpha _s`$. One way to calculate the coefficients is to expand the resummed part in terms of the strong coupling (expanding the exponent an the Wilson coefficients), and compare the expansion with a fixed order calculation. Luckily, because of its significance, there was much work done on fixed order QCD corrections to Higgs production in the $`ggHX`$ channel. These fixed order QCD corrections are known to substantially increase the rate: by about 70 to 100 percent, depending on the Higgs mass, at $`𝒪(\alpha _s^3)`$, and by an additional 50 to 70 percent at $`𝒪(\alpha _s^4)`$ . It is expected that the calculation of even higher order corrections is important to reliably predict the cross section. In Ref. it was shown that multiple soft–gluon emission dominates the higher order corrections.
### 2.1 Soft gluon resummation for the $`ggHX`$ channel
Resummed calculations, taking into account the soft–gluon effect, attempted to estimate the size of the non-calculated higher order corrections , as well as provide a reliable shape of the Higgs transverse momentum distribution . Our present approach surpasses these by calculating the $`Q_T`$ distribution while including $`𝒪(\alpha _s^4)`$ terms in the Sudakov exponent, using the state of the art matching to the latest fixed order distributions, using a QCD improved gluon-Higgs effective coupling , and using an improved non-perturbative function. We utilize the approximation that the object which couples the gluons to the Higgs (the top quark in the SM), is much heavier than the Higgs itself. This approximation is not essential to our calculation and can be released by the calculation of the further Wilson coefficients keeping the relevant masses. The heavy quark approximation in the SM was shown to be reliable within 5 percent for $`m_H<2m_t`$ , and still reasonable even in the range of $`m_H\stackrel{>}{}2m_t`$ . It has also been shown that the approximation remains valid for the $`Q_T`$ distribution in the large $`Q_T`$ region, provided that $`m_H<m_t`$ and $`Q_T<m_t`$ . In this work we assume that the approximation is valid in the whole $`Q_T`$ region. Unlike the authors of Ref. we do not assume that the QCD corrections to the $`ggHX`$ cross section can be factorized into a multiplicative term in the heavy quark limit in all orders of $`\alpha _s`$. We can release this approximation because the CSS formalism, by definition, systematically incorporates higher order fixed order corrections via the definition of the Sudakov exponent and the Wilson coefficients as perturbative expansions .
Multiple soft–gluon emission affects the $`ggHX`$ cross section when the transverse momentum of the Higgs is low, while for high transverse momenta the hard gluon radiation is dominant. Thus, using the CSS formalism we resum large logs of the type $`\mathrm{ln}(Q/Q_T)`$ in the low $`Q_T`$ region, and we match the resummed result to the fixed order calculation which is valid for high $`Q_T`$ . We also include the $`qg`$ and $`q\overline{q}`$ subprocesses which, depending on the Higgs mass, together constitute 0 to 10 percent of the total rate .
The resummed differential cross section of the Higgs boson production in hadronic collisions is written as
$`{\displaystyle \frac{d\sigma (h_1h_2H^0X)}{dQ^2dydQ_T^2}}=\sigma _0{\displaystyle \frac{Q^2}{S}}\pi \delta (Q^2m_H^2)`$
$`\times \{{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle }d^2be^{i\stackrel{}{Q}_T\stackrel{}{b}}\stackrel{~}{W}_{gg}(b_{},Q,x_1,x_2,C_{1,2,3})`$
$`\times \stackrel{~}{W}_{gg}^{\mathrm{NP}}(b,Q,x_1,x_2)+Y(Q_T,Q,x_1,x_2,C_4)\}.`$ (15)
The kinematic variables $`Q`$, $`y`$, and $`Q_T`$ are the invariant mass, rapidity, and transverse momentum of the Higgs boson in the laboratory frame. The parton momentum fractions are defined as $`x_1=e^yQ/\sqrt{S}`$, and $`x_2=e^yQ/\sqrt{S}`$, with $`\sqrt{S}`$ being the center–of–mass (CM) energy of the hadrons $`h_1`$ and $`h_2`$. The lowest order cross section, with the QCD corrected effective coupling of the Higgs boson to gluons is
$$\sigma _0=\kappa _\varphi (Q)\frac{\sqrt{2}G_F\alpha _s^2(Q^2)}{576\pi },$$
(16)
where $`G_F`$ is the Fermi constant and $`\kappa _\varphi `$ is defined in Ref. . The renormalization group invariant kernel of the Fourier integral $`\stackrel{~}{W}_{gg}`$ and the regular terms $`Y(Q_T,Q,x_1,x_2,C_4)`$ (together with the variables $`b_{}`$ and $`C_1`$ to $`C_4`$) are defined in Ref. . In addition to Ref. we use the process independent coefficient
$$A^{(2)}=C_A\left[\left(\frac{67}{36}\frac{\pi ^2}{12}\right)N_C\frac{5}{18}N_f\right],$$
(17)
in the expansion of the $`A`$ function ($`N_C=3`$ the number of colors and $`N_f=5`$ the number of active quark flavors).
### 2.2 Some numerical results
The resummation formula is coded in the ResBos Monte Carlo event generator , which uses the following electroweak input parameters : $`G_F=1.16639\times 10^5\mathrm{GeV}^2,m_Z=91.187\mathrm{GeV},m_W=80.36\mathrm{GeV}`$. The NLO expressions for the running electromagnetic and strong couplings $`\alpha (\mu )`$ and $`\alpha _S(\mu )`$ are used, as well as the NLO parton distribution function set CTEQ4M (defined in the modified minimal subtraction, $`\overline{MS}`$, scheme). The renormalization and factorization scales are set equal to the Higgs invariant mass. In the choice of the non-perturbative parameters we follow Ref. . Since we are not concerned with the decays of Higgs bosons in this work, we do not impose any kinematic cuts.
Fig. 4 displays production cross sections at the LHC, calculated in the SM as the function of the Higgs mass. Our $`𝒪(\alpha _s^3)`$ curve agrees well with the result in Ref. . The ratio of the fixed order $`𝒪(\alpha _s^3)`$ (dashed) and the lowest order $`𝒪(\alpha _s^2)`$ (dotted) curves varies between 2.35 and 2.00. We note that less than 2 percent of the $`𝒪(\alpha _s^3)`$ corrections come from the $`qg`$ and $`q\overline{q}`$ initial states for Higgs masses below 200 GeV. The resummed curve is slightly (5 to 6 percent) higher than the $`𝒪(\alpha _s^3)`$ one, as expected based on the findings that the CSS formalism preserves the fixed order rate within the error of the matching (which is expected to be higher order) . The resummed rate is close to the $`𝒪(\alpha _s^3)`$, because we used the $`𝒪(\alpha _s^3)`$ fixed order results to derive the Wilson coefficients which are utilized in our calculation. In Ref. the $`𝒪(\alpha _s^4)`$ corrections were utilized to show that in the high $`Q_T`$ region the $`𝒪(\alpha _s^4)`$ to $`𝒪(\alpha _s^3)`$ $`K`$-factor is nearly constant and is about 1.5 (for CTEQ4M parton distributions). Based on this finding we also plot the $`𝒪(\alpha _s^3)`$ curve rescaled by 1.5, to illustrate the size of the $`𝒪(\alpha _s^4)`$ corrections and to establish the normalization of our resummed calculation among the fixed order results.
Fig. 5 illustrates the effect of the various contributions of the CSS formalism on the Higgs boson transverse momentum distribution. The lower peaking curves, drawn by the same type line, contain the coefficients $`A^{(1,2)}`$. The others lack the $`A^{(2)}`$ coefficient. Comparison of pairs of curves shows that the log multiplied by the $`A^{(2)}`$ coefficient increases the rate by about 10% around the peak, and decreases it in the mid-$`Q_T`$ region. The figure also shows that exclusion of the $`B^{(1)}`$ term leads to about 40% decrease around the peak, and an increase away from it. Finally, the exclusion of the $`C^{(1)}`$ coefficient decreases the overall rate by about a factor of 2, coupled with some shape change similar to the $`B^{(1)}`$ case.
Fig. 6 displays transverse momentum distributions of Higgs bosons produced at the LHC. The $`Q_T`$ distribution is calculated under several different assumptions for the non-perturbative sector of the CSS formalism, in order to span the range of scatter of these different predictions. In Fig. 6a the (solid) curve using the result of the latest 3-parameter fit for the non-perturbative function is shown. (The actual values of the parameters used are: $`g_1=0.15`$ GeV<sup>2</sup>, $`g_2=(C_F/C_A)0.48`$ GeV<sup>2</sup>, and $`g_3=0.58`$ GeV<sup>-1</sup>.) Also shown the (dashed) curve using the result of the latest 3-parameter fit of Ref. . (The values were used are: $`g_1=0.24`$ GeV<sup>2</sup>, and $`g_2=(C_F/C_A)0.34`$ GeV<sup>2</sup>.) We plotted the (dotted) curve using the previous 3-parameter fit of Ref. , as well. In the lower portion of the figure we show the ratios of the different curves to the solid curve. From this we conclude that the three different parameterizations differ by about 5 percent, at most, in the relevant $`Q_T`$ region. At $`Q_T=10`$ GeV, in the region of the peak of the distribution, the difference is about 2 percent.
In Fig. 6b the solid curve is the same as in Fig. 6a. In this figure results using $`g_2=(C_F/C_A)0.33`$ GeV<sup>2</sup>, and $`g_2=(C_F/C_A)0.69`$ GeV<sup>2</sup> values are plotted (dashed). These $`g_2`$ values are 3 $`\sigma `$ deviations from the central value $`g_2=(C_F/C_A)0.48`$ GeV<sup>2</sup> of the new 3-parameter fit. Also shown a curve with $`g_2=0.48`$ GeV<sup>2</sup>, where the assumption that the non-perturbative parameter $`g_2`$ scales by $`C_A/C_F`$ for the gluonic initial state was not utilized. The lower portion of the figure shows that the ratios of the various curves to the solid curve do not deviate from 1 significantly except in the very low $`Q_T`$ ($`<`$ 5 GeV) region.
## Acknowledgments
We thank the organizers of the les Houches workshop for their hospitality. We are indebted for the CTEQ Collaboration for many invaluable discussions and W. Sakumoto for the CDF results. C.B. thanks M. Spira, and C.-P. Yuan for discussions. This work was supported in part by the DOE under grant DE-FG-03-94ER40833.
A Comparison of the Predictions from Monte Carlo Programs and Transverse Momentum Resummation
C. Balázs, J. Huston, I. Puljak
## 1 Introduction
Parton shower Monte Carlo programs such as PYTHIA, HERWIG and ISAJET are commonly used by experimentalists, both as a way of comparing experimental data to theoretical predictions, and also as a means of simulating experimental signatures in kinematic regimes for which there is not yet experimental data (such as the LHC). The final output of the Monte Carlo programs consists of the 4-vectors of a set of final state hadrons; this output can either be compared to reconstructed experimental quantities or, when coupled with a simulation of a detector response, can be directly compared to raw data taken by the experiment, and/or passed through the same reconstruction procedures as the raw data. In this way, the parton shower programs can be more useful to experimentalists than analytic calculations. Indeed, almost all of the physics plots in the ATLAS physics TDR involve comparisons to PYTHIA (version 5.7).
For many physical quantities, the predictions from parton shower Monte Carlo programs should be nearly as precise as those from analytic theoretical calculations. This is expected, among others, for calculations which resum logs with the transverse momentum of partons initiating the hard scattering. In the recent literature, most calculations of this kind are either based on or originate from the formalism developed by J. Collins, D. Soper, and G. Sterman (CSS) <sup>2</sup><sup>2</sup>2See, for example, the discussion in the previous section., which we choose as the analytic ‘benchmark’ of this section. In this case, both the Monte Carlo and analytic calculations should accurately describe the effects of the emission of multiple soft gluons from the incoming partons, an all orders problem in QCD. The initial state soft gluon emission can affect the kinematics of the final state partons. This may have an impact on the signatures of physics processes at both the trigger and analysis levels and thus it is important to understand the reliability of such predictions. The best method for testing the reliability is the direct comparison of the predictions to experimental data. If no experimental data is available for certain predictions, then some understanding of the reliability may be gained from the comparison of the predictions from the two different methods.
## 2 Parton Showering and Resummation
For technical reasons, the initial state parton shower proceeds by a $`\mathrm{𝑏𝑎𝑐𝑘𝑤𝑎𝑟𝑑𝑠}`$ evolution, starting at the large (negative) $`Q^2`$ scale of the hard scatter and then considering emissions at lower and lower (negative) virtualities, corresponding to earlier points on the cascade (and earlier points in time), until a scale corresponding to the factorization scale is reached. The transverse momentum of the initial state is built up from the whole series of splittings (and boosts). The showering process is independent of the hard scattering process being considered (as long as one does not introduce any matrix element corrections), and depends only on the initial state partons and the hard scale of the process.
In the case of parton showering, the leading order collinear singularities factorize for cross sections in the collinear limit
$`\underset{p_gp_b}{lim}|_{n+1}|^2=g_s^2(p_b.p_g)^1P_{ga}(z)|_n|^2,`$ (1)
where $`_{n+1}`$ is the invariant amplitude for the process producing $`n`$ partons and a gluon, $`g_s`$ is the strong coupling constant, $`p_b`$ and $`p_g`$ are the 4-momenta of the daughters of the n’th parton $`a`$ (i.e. $`a`$ splits into $`b`$ and $`g`$, and when they are collinear then $`p_b.p_g0`$). Finally $`P_{ga}(z)`$ is the DGLAP splitting kernel belonging to the $`ag`$ splitting. The leading order collinear singularities can be factorized into a Sudakov form factor: $`S=1P(noemission)=exp(𝑑p^2/p^2𝑑zP(z))`$. The distribution $`1S`$ can be used to generate the $`Q^2`$ for the first emission and hence for the whole cascade. The formalism can be extended to soft singularities as well by using angular ordering. In this approach, the choice of the hard scattering is based on the use of evolved parton distributions, which means that the inclusive effects of initial-state radiation are already included. What remains is therefore to construct the exclusive showers.
Parton showering resums primarily the leading logs, which are universal, i.e. process independent, and depend only on the given initial state. In this lies one of the strengths of Monte Carlos, since parton showering can be incorporated into a wide variety of physical processes. An analytic calculation, in comparison, can resum all logs. For example, the CSS formalism sums all of the logarithms with $`Q^2/p_T^2`$ in their arguments, where (for Higgs boson production) $`Q`$ is the four momentum of the Higgs and $`p_T`$ is its transverse momentum. As discussed in the previous section on resummation, all of the ‘dangerous logs’ are included in the Sudakov exponent, which can be written in the impact parameter ($`b`$) space as:
$`𝒮(p,b)={\displaystyle _{1/b^2}^{Q^2}}{\displaystyle \frac{d\overline{\mu }^2}{\overline{\mu }^2}}\left[A\left(\alpha _s(\overline{\mu })\right)\mathrm{ln}\left({\displaystyle \frac{Q^2}{\overline{\mu }^2}}\right)+B\left(\alpha _s(\overline{\mu })\right)\right],`$
with the $`A`$ and $`B`$ functions being free of large logs and perturbatively calculable:
$`A\left(\alpha _s(\overline{\mu })\right)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\alpha _s(\overline{\mu })}{\pi }}\right)^nA^{(n)},B\left(\alpha _s(\overline{\mu })\right)={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\alpha _s(\overline{\mu })}{\pi }}\right)^nB^{(n)}.`$
These functions contain an infinite number of coefficients, with the $`A^{(n)}`$ coefficients being universal while the $`B^{(n)}`$ are process dependent. In practice, the number of towers of logarithms included in the Sudakov exponent depends on the level to which a fixed order calculation was performed for a given process. For example, if only a next-to-leading order calculation is available, only the coefficients $`A^{(1)}`$ and $`B^{(1)}`$ can be included. If a NNLO calculation is available, then $`A^{(2)}`$ and $`B^{(2)}`$ can be extracted and incorporated into a resummation calculation, and so on. This is the case, for example, for $`Z^0`$ boson production. So far, only the $`A^{(1)}`$, $`A^{(2)}`$ and $`B^{(1)}`$ coefficients are known for Higgs production but the calculation of $`B^{(2)}`$ is in progress. If we try to interpret parton showering in the same language, which is admittedly risky, then we can say that the Monte Carlo Sudakov exponent always contains a term analogous to $`A^{(1)}`$. It was shown in Reference that a suitable modification of the Altarelli-Parisi splitting function, or equivalently the strong coupling constant $`\alpha _s`$, also effectively approximates the $`A^{(2)}`$ coefficient. <sup>3</sup><sup>3</sup>3This is rigorously true only for the high x or $`\sqrt{\tau }`$ region.
In contrast with the shower Monte Carlos, analytic resummation calculations integrate over the kinematics of the soft gluon emision, with the result that they are limited in their predictive power for inclusive final states. While the Monte Carlo maintains an exact treatment of the branching kinematics, in the original CSS formalism no kinematic penalty is paid for the emission of the soft gluons, although an approximate treatment of this can be incorporated into its numerical implementations, like ResBos . Neither the parton showering process nor the analytic resummation translate smoothly into kinematic configurations where one hard parton is emitted (at large $`p_T`$). In the Monte Carlo matrix element corrections, and in the analytic resummation calculation matching is necessary. This matching is standard procedure for resummation calculations and matrix element corrections are becoming increasingly common in Monte Carlos .
With the appropriate input from higher order cross sections, a resummation calculation has the corresponding higher order normalization and scale dependence. The normalization and scale dependence for the Monte Carlo, though, remains that of a leading order calculation. The parton showering process redistributes the event particles in phase space, but does not change the total cross section (for example, for the production of a Higgs boson). <sup>4</sup><sup>4</sup>4Technically, one could add the branching for $`qq`$+Higgs in the shower, which would have the capability of increasing somewhat the Higgs cross section; however, the main contribution to the higher order $`K`$-factor comes from the virtual corrections and the ‘Higgs Bremsstrahlung’ contribution is neglible.
In particular, one quantity which should be well-described by both calculations is the transverse momentum ($`p_T`$) of the final state electroweak boson in a subprocess such as $`q\overline{q}WX`$, $`ZX`$ or $`ggHX`$, where most of the $`p_T`$ is provided by initial state parton showering. The parton showering supplies the same sort of transverse kick as the soft gluon radiation in a resummation calculation. Indeed, very similar Sudakov form factors appear in both approaches, with the caveats about the $`A^{(n)}`$ and $`B^{(n)}`$ terms mentioned previously. This correspondence between the Sudakov form factors in resummation and Monte Carlo approaches may seem trivial, but there are many subtleties between the two approaches relating to both the arguments of the Sudakov factors as well as the impact of subleading logs .
At a point in its evolution corresponding to (typically) the virtuality of a few GeV<sup>2</sup>, the parton shower is cut off and the effects of gluon emission at softer scales must be parameterized and inserted by hand. This is similar to the (somewhat arbitrary) division between perturbative and non-perturbative regions in a resummation calculation. The parameterization is typically done with a Gaussian formalism similar to that used for the non-perturbative $`k_T`$ in a resummation program. In general, the value for the non-perturbative $`k_T`$ needed in a Monte Carlo program will depend on the particular kinematics being investigated. In the case of the resummation calculation the non-perturbative physics is determined from fits to fixed target data and then automatically evolved to the kinematic regime of interest.
A value for the average non-perturbative $`k_T`$ of greater than 1 GeV does not imply that there is an anomalous intrinsic $`k_T`$ associated with the parton size; rather this amount of $`k_T`$ needs to be supplied to provide what is missing in the truncated parton shower. If the shower is cut off at a higher virtuality, more of the ‘non-perturbative’ $`k_T`$ will be needed.
## 3 $`Z^0`$ Boson Production at the Tevatron
The 4-vector of a $`Z^0`$ boson, and thus its transverse momentum, can be measured with great precision in the $`e^+e^{}`$ decay mode. Resolution effects are relatively minor and are easily corrected. Thus, the $`Z^0`$ $`p_T`$ distribution is a great testing ground for both the resummation and Monte Carlo formalisms for soft gluon emission. The (resolution corrected) $`p_T`$ distribution for $`Z^0`$ bosons (in the low $`p_T`$ region) for the CDF experiment<sup>5</sup><sup>5</sup>5We thank Willis Sakumoto for providing the figures for CDF $`Z^0`$ production is shown in Figure 1, compared to both the resummed prediction from ResBos, and to two predictions from PYTHIA (version 6.125). One PYTHIA prediction uses the default (rms)<sup>6</sup><sup>6</sup>6For a Gaussian distribution, $`k_T^{rms}=1.13k_T`$. value of intrinsic $`k_T`$ of 0.44 GeV and the second a value of 2.15 GeV (per incoming parton). <sup>7</sup><sup>7</sup>7A previous publication indicated the need for a substantially larger non-perturbative $`k_T`$, of the order of 4 GeV for the case of $`W`$ production at the Tevatron. The data used in the comparison, however, were not corrected for resolution smearing, a fairly large effect for the case of $`We\nu `$ production and decay. The latter value was found to give the best agreement for PYTHIA with the data.<sup>8</sup><sup>8</sup>8A similar conclusion has been reached for comparisons of the CDF $`Z^0`$ $`p_T`$ data with HERWIG . All of the predictions use the CTEQ4M parton distributions . The shift between the two PYTHIA predictions at low $`p_T`$ is clearly evident. As might have been expected, the high $`p_T`$ region (above 10 GeV) is unaffected by the value of the non-perturbative $`k_T`$. Note that much of the $`k_T`$ ‘given’ to the incoming partons at their lowest virtuality, $`Q_0`$, is reduced at the hard scatter due to the number of gluon branchings preceding the collision. The emitted gluons carry off a sizeable fraction of the original non-perturbative $`k_T`$. This point will be investigated in more detail later for the case of Higgs production.
As an exercise, one can transform the resummation formula in order to bring it to a form where the non-perturbative function acts as a Gaussian type smearing term. Using the Ladinsky-Yuan parameterization of the non-perturbative function in ResBos leads to an rms value for the effective $`k_T`$ smearing parameter, for $`Z^0`$ production at the Tevatron, of 2.5 GeV. This is similar to that needed for PYTHIA and HERWIG to describe the $`Z^0`$ production data at the Tevatron.
In Figure 1, the normalization of the resummed prediction has been rescaled upwards by 8.4%. The PYTHIA prediction was rescaled by a factor of 1.3-1.4 (remember that this is only a leading order comparison) for the shape comparison.
As stated previously, the resummed prediction correctly describes the shape of the $`Z^0`$ $`p_T`$ distribution at low $`p_T`$, even with the optimal non-perturbative $`k_T`$, although there is still a noticeable difference in shape between the Monte Carlo and the resummed prediction. It is interesting to note that if the process dependent coefficients ($`B^{(1)}`$ and $`B^{(2)}`$) were not incorporated into the resummation prediction, the result would be an increase in the height of the peak and a decrease in the rate between 10 and 20 GeV, leading to a better agreement with the PYTHIA prediction .
The $`Z^0`$ $`p_T`$ distribution is shown over a wide $`p_T`$ range in Figure 2. The PYTHIA and ResBos predictions both describe the data well. Note especially the agreement of PYTHIA with the data at high $`p_T`$, made possible by explict matrix element corrections (from the subprocesses $`q\overline{q}Z^0g`$ and $`gqZ^0q`$) to the $`Z^0`$ production process.<sup>9</sup><sup>9</sup>9Slightly different techniques are used for the matrix element corrections by PYTHIA and by HERWIG . In PYTHIA, the parton shower probability distribution is applied over the whole phase space and the exact matrix element corrections are applied only to the branching closest to the hard scatter. In HERWIG, the corrections are generated separately for the regions of phase space unpopulated by HERWIG (the ‘dead zone’) and the populated region. In the dead zone, the radiation is generated according to a distribution using the first order matrix element calculation, while the algorithm for the already populated region applies matrix element corrections whenever a branching is capable of being ‘the hardest so far’.
## 4 Diphoton Production
Most of the experience that we have for comparisons of data to resummation calculations/Monte Carlos deals with Drell-Yan production, i.e. $`q\overline{q}`$ initial states. It is important then to examine diphoton production at the Tevatron, where a large fraction of the contribution at low mass is due to $`gg`$ scattering. The prediction for the diphoton $`k_T`$ distribution at the Tevatron, from PYTHIA (version 6.122), is shown in Figure 3, using the experimental cuts applied in the CDF analysis . It is interesting to note that about half of the diphoton cross section at the Tevatron is due to the $`gg`$ subprocess, and that the diphoton $`p_T`$ distribution is noticeably broader for the $`gg`$ subprocess than the $`q\overline{q}`$ subprocess.
A comparison of the $`p_T`$ distributions for the two diphoton subprocesses $`(q\overline{q},gg)`$ in PYTHIA versions 5.7 and 6.1 is shown in Figure 4. There seems to be little difference in the $`p_T`$ distributions between the two versions for both subprocesses.
In Figure 5 are shown the ResBos predictions for diphoton production at the Tevatron from $`q\overline{q}`$ and $`gg`$ scattering compared to the PYTHIA predictions (using the same experimental cuts). The $`gg`$ subprocess predictions in ResBos agree well with those from PYTHIA while the $`q\overline{q}`$ $`p_T`$ distribution is noticebly broader in ResBos. The latter behavior is due to the presence of the $`Y`$ piece in ResBos at moderate $`p_T`$, and the matching of the $`q\overline{q}`$ cross section to the fixed order $`q\overline{q}\gamma \gamma g`$ at high $`p_T`$. The corresponding matrix element correction is not in PYTHIA. It is interesting to note that the PYTHIA and ResBos predictions for $`gg\gamma \gamma `$ agree in the moderate $`p_T`$ region, even though the ResBos prediction has the $`Y`$ piece present and is matched to the matrix element piece $`gg\gamma \gamma g`$ at high $`p_T`$, while there is no such matrix element correction for PYTHIA. This shows the smallness of the $`Y`$ piece for the $`gg`$ subprocess, which is the same conclusion that was reached in Ref. . One way to understand this is recalling that the $`gg`$ parton-parton luminosity falls very steeply with increasing partonic center of mass energy, $`\sqrt{\widehat{s}}`$. This falloff tends to suppress the size of the $`Y`$ piece since the production of the diphoton pair at higher $`p_T`$ requires larger $`x_1`$, $`x_2`$ values. In the default CSS formalism, there is no such kinematic penalty in the resummed piece since the soft gluon radiation comes for “free”. (Larger $`x_1`$ and $`x_2`$ values are not required.)
A comparison of the CDF diphoton data to NLO and resummed (ResBos) QCD predictions is shown in Figure 6. Plotted are the diphoton mass, the angle $`\mathrm{\Delta }\varphi `$ between the two photons and the transverse momentum $`k_T`$ of the diphoton pair. The transverse momentum distribution, in particular, is sensitive to the effects of the soft gluon radiation and better agreement can be observed with the ResBos prediction than with the NLO one. The data shown in this figure is from an integrated luminosity of 87 $`pb^1`$. A much more precise comparison with the effects of soft gluon radiation will be possible with the 2 $`fb^1`$ or greater data sample that is expected for both CDF and D0 in Run 2.
The prediction for the diphoton production cross section, as a function of the diphoton $`p_T`$ and using cuts appropriate to ATLAS and CMS, is shown in Figure 7. Note that, as at the Tevatron, about half of the cross section is due to $`gg`$ scattering and the diphoton $`p_T`$ distribution from $`gg`$ scattering is noticeably broader than that from $`q\overline{q}`$ production.
In Figure 8 is shown a comparison of the diphoton $`p_T`$ distribution for two different versions of PYTHIA, for the two different subprocesses. Note that the $`p_T`$ distribution in PYTHIA version 5.7 is somewhat broader than that in version 6.122 for the case of $`gg`$ scattering. The effective diphoton mass range being considered here is lower than the 150 GeV Higgs mass that will be considered in the next section. As will be seen, the differences in soft gluon emission between the two versions of PYTHIA are larger in that case.
In Figure 9 are shown the ResBos predictions for diphoton production at the LHC from $`q\overline{q}`$ and $`gg`$ scattering compared to the PYTHIA predictions (using the same experimental cuts). Again, the $`gg`$ subprocess predictions in ResBos agree well with those from while the $`q\overline{q}`$ $`p_T`$ distribution is noticebly broader in ResBos, for the reasons cited previously.
## 5 Higgs Boson Production
A comparison of the Higgs $`p_T`$ distribution at the LHC, for a Higgs mass of 150 GeV, is shown in Figure 10, for ResBos and the two recent versions of PYTHIA. As before, PYTHIA has been rescaled to agree with the normalization of ResBos to allow for a better shape comparison. Note that the peak of the resummed distribution has moved to $`p_T`$ 11 GeV (compared to about 3 GeV for $`Z^0`$ production at the Tevatron). This is partially due to the larger mass (150 GeV compared to 90 GeV), but is primarily because of the larger color factors associated with initial state gluons ($`C_A=3`$) rather than quarks ($`C_F=4/3`$), and also because of the larger phase space for initial state gluon emission at the LHC.
The newer version of PYTHIA agrees well with ResBos at low to moderate $`p_T`$, but falls below the resummed prediction at high $`p_T`$. This is easily understood: ResBos switches to the NLO Higgs + jet matrix element at high $`p_T`$ while the default PYTHIA can generate the Higgs $`p_T`$ distribution only by initial state gluon radiation, using as maximum virtuality the Higgs mass squared. High $`p_T`$ Higgs production is another example where a $`21`$ Monte Carlo calculation with parton showering can not completely reproduce the exact matrix element calculation, without the use of matrix element corrections. The high $`p_T`$ region is better reproduced if the maximum virtuality $`Q_{max}^2`$ is set equal to the squared partonic center of mass energy, $`s`$, rather than $`m_H^2`$. This is equivalent to applying the parton shower to all of phase space. However, this has the consequence of depleting the low $`p_T`$ region as ‘too much’ showering causes events to migrate out of the peak. The appropriate scale to use in PYTHIA (or any Monte Carlo) depends on the $`p_T`$ range to be probed. If matrix element information is used to constrain the behavior, the correct high $`p_T`$ cross section can be obtained while still using the lower scale for showering. The incorporation of matrix element corrections to Higgs production (involving the processes $`gqqH`$,$`q\overline{q}gH`$, $`gggH`$) is the next logical project for the Monte Carlo experts, in order to accurately describe the high $`p_T`$ region.
A comparison of the two versions of PYTHIA and of ResBos is also shown in Figure 11 for the case of Higgs production (at a Higgs mass of 100 GeV) at the Tevatron with center-of-mass energy of 2.0 TeV. The same qualititative features are observed as at the LHC: the newer version of PYTHIA agrees better with ResBos in describing the low $`p_T`$ shape, and there is a falloff at high $`p_T`$ unless the larger virtuality is used for the for the parton showers. The default (rms) value of the non-perturbative $`k_T`$ (0.44 GeV) was used for the PYTHIA predictions for Higgs production.
The older version of PYTHIA produces too many Higgs events at moderate $`p_T`$ (in comparison to ResBos) at both the Tevatron and the LHC. Two changes have been implemented in the newer version. The first change is that a cut is placed on the combination of $`z`$ and $`Q^2`$ values in a branching: $`\widehat{u}=Q^2\widehat{s}(1z)<0`$, where $`\widehat{s}`$ refers to the subsystem of the hard scattering plus the shower partons considered to that point. The association with $`\widehat{u}`$ is relevant if the branching is interpreted in terms of a $`22`$ hard scattering. The corner of emissions that do not respect this requirement occurs when the $`Q^2`$ value of the spacelike emitting parton is little changed and the $`z`$ value of the branching is close to unity. This effect is mainly for the hardest emission (largest $`Q^2`$). The net result of this requirement is a substantial reduction in the total amount of gluon radiation <sup>10</sup><sup>10</sup>10Such branchings are kinematically allowed, but since matrix element corrections would assume initial state partons to have $`Q^2=0`$, a non-physical $`\widehat{u}`$ results (and thus no possibility to impose matrix element corrections). The correct behavior is beyond the predictive power of leading log Monte Carlos. In the second change, the parameter for the minimum gluon energy emitted in spacelike showers is modified by an extra factor roughly corresponding to the $`1/\gamma `$ factor for the boost to the hard subprocess frame . The effect of this change is to increase the amount of gluon radiation. Thus, the two effects are in opposite directions but with the first effect being dominant.
This difference in the $`p_T`$ distribution between the two versions of PYTHIA could have an impact on the analysis strategies for Higgs searches at the LHC. For example, for the CMS detector, the higher $`p_T`$ activity associated with Higgs production in version 5.7 would have allowed for a more precise determination of the event vertex from which the Higgs (decaying into two photons) originated. Vertex pointing with the photons is not possible in the CMS barrel region, and the large number of interactions occuring with high intensity running will mean a substantial probability that at least one of the interactions will produce jets at low to moderate $`E_T`$. This could lead to the wrong vertex being chosen for the Higgs, leading to a significant degradation in the $`\gamma \gamma `$ effective mass resolution. In principle, this problem could affect the $`p_T`$ distribution for all PYTHIA processes. In practice, it affects only $`gg`$ initial states, due to the enhanced probability for branching with such an initial state.
As an exercise, an 80 GeV $`W`$ and an 80 GeV Higgs were generated at the Tevatron using PYTHIA5.7 . A comparison of the distribution of values of $`\widehat{u}`$ and the virtuality $`Q`$ for the two processes indicates a greater tendency for the Higgs virtuality to be near the maximum value and for there to be a larger number of Higgs events with positive $`\widehat{u}`$ (than W events).
## 6 Comparison with HERWIG
The variation between versions 5.7 and 6.1 of PYTHIA gives an indication of the uncertainties due to the types of choices that can be made in Monte Carlos. The requirement that $`\widehat{u}`$ be negative for all branchings is a choice rather than an absolute requirement. Perhaps the better agreement of version 6.1 with ResBos is an indication that the adoption of the $`\widehat{u}`$ restrictions was correct. Of course, there may be other changes to PYTHIA which would also lead to better agreement with ResBos for this variable.
Since there are a variey of choices that can be made in Monte Carlo implementations, it is instructive to compare the predictions for the $`p_T`$ distribution for Higgs production from ResBos and PYTHIA with that from HERWIG (version 5.6, also using the CTEQ4M parton distribution functions). The HERWIG prediction is shown in Figure 12 along with the PYTHIA and ResBos predictions, all normalized to the ResBos prediction. <sup>11</sup><sup>11</sup>11The normalization factors (ResBos/Monte Carlo) are PYTHIA (both versions)(1.68) and HERWIG (1.84). (In all cases, the CTEQ4M parton distribution was used.) The predictions from HERWIG and PYTHIA 6.1 are very similar, with the HERWIG prediction matching the ResBos shape somewhat better at low $`p_T`$. For reference, the absolutely normalized predictions from ResBos, PYTHIA and HERWIG for the $`p_T`$ distribution of a 150 GeV Higgs at the LHC are shown in Figure 13.
## 7 Non-perturbative $`k_T`$
A question still remains as to the appropriate value of non-perturbative $`k_T`$ to input in the Monte Carlos to achieve a better agreement in shape, both at the Tevatron and at the LHC. In Figures 14 and 15 are shown comparisons of ResBos and PYTHIA predictions for the Higgs $`p_T`$ distribution at the Tevatron and LHC. The PYTHIA prediction (now version 6.1 alone) is shown with several values of non-perturbative $`k_T`$. Suprisingly, no difference is observed between the predictions with the different values of $`k_T`$, with the peak in PYTHIA always being somewhat below that of ResBos. This insensitivity can be understood from the plots at the bottom of the two figures which show the sum of the non-perturbative initial state $`k_T`$ ($`k_{T1}`$+$`k_{T2}`$) at $`Q_0`$ and at the hard scatter scale $`Q`$. Most of the $`k_T`$ is radiated away. with this effect being larger (as expected) at the LHC. The large gluon radiation probability from a gluon-gluon initial state (and the greater phase space available at the LHC) lead to a stronger degradation of the non-perturbative $`k_T`$ than was observed with $`Z^0`$ production at the Tevatron.
For completeness, a comparison of PYTHIA and ResBos is shown in Figure 16 for $`Z^0`$ boson production at the LHC. There are two points that are somewhat surprising. There is still a very strong sensitivity to the value of the non-perturbative $`k_T`$ used in the smearing, and the best agreement with ResBos is obtained with the default value (0.44 GeV), in contrast to the 2 GeV needed at the Tevatron. Note again the agreement of PYTHIA with ResBos at the highest values of $`Z^0`$ $`p_T`$ due to the explicit matrix element corrections applied.
The sum of the incoming parton $`k_T`$ distributions, both at the scale $`Q_0`$ and at the hard scattering scale, are shown in Figure 17 for several different starting (rms) values of primordial $`k_T`$ (per parton). There is substantially less radiation for a $`q\overline{q}`$ initial state than for a gg initial state (as in the case of the Higgs), leading to a noticeable dependence of the $`Z^0`$ $`p_T`$ distribution on the primordial $`k_T`$ distribution.
## 8 Conclusions
An understanding of the signature for Higgs boson production at either the Tevatron or LHC depends upon the understanding of the details of soft gluon emission from the initial state partons. This soft gluon emission can be modelled either in a Monte Carlo or in a $`k_T`$ resummation program, with various choices possible in both implementations. A comparison of the two approaches is useful to understand the strengths and weaknesses of each. The data from the Tevatron that either exists now, or will exist in Run 2, will be extremely useful to test both approaches.
## 9 Acknowledgements
We would like to thank Claude Charlot, Gennaro Corcella, Steve Mrenna, Willis Sakumoto, Torbjorn Sjostrand and Valeria Tano for useful conversations and plots. This work was supported in part by the NSF under grant PHY-9901946 and by the DOE under grant DE-FG-03-94ER40833.
Automatic Computation of LHC Processes
E. Boos, V. Ilyin, K. Kato, A. Pukhov, A. Semenov, A. Skatchkova
Automatic computation is a new approach to HEP computing. The first such systems, GRACE , FeynArt/FeynCalc and CompHEP , were reported at the 1st International Workshop AIHENP held on March, 1990 in Lyon-Villeurbanne (France). Under this terminology, automatic computation system (ACS), we assume, as a distinguishing feature, the generation of the computing code for a specific collision process with the aid of another code.
ACS’s are now used widely by phenomenologists for the calculation of many collision processes. For example, the GRACE and CompHEP systems were used in the LEP2 Workshop , and for evaluation of processes at TeV linear colliders . With ACS one can calculate all collision processes within a given physical model, where by physical model we mean the set of Feynman rules. Recent developments with the LanHEP package have opened a possibility to derive Feynman rules in the form of the ACS intrinsic physical model in a fully automatic way, starting from the Lagrangian. Now, not only the Standard Model but a number of its extensions, like SUSY models, are implemented in ACS. A general review of this new approach is given in these Proceedings by K.Kato together with discussion of main directions of the ACS development. Here, we discuss in more detail specifics of the ACS applications in LHC phenomenology, and in particular to the evaluation of QCD processes.
To close this preview we list below the main ACS options in order to provide an idea for users of what is available:
1. selection of physical model (Lagrangian) and hard subprocess,
2. Feynman diagram generation,
3. generation of the code for matrix element,
4. convolution with parton distributions,
5. generation of kinematics (phase space parameterization) with regularization of kinematical peaks,
6. integration over the phase space (evaluation of cross section),
7. generation of events at partonic level, including the interface to hadronization tools.
## 1 The problem of multiparticle final states: why automatic computations?
We start from the problem of the accurate evaluation of hard subprocesses in the case of multiparticle final states.
When physicists simulate HEP processes with such generators as PYTHIA , ISAJET or HERWIG they use a data base of hard subprocesses implemented in these packages. It means that a) matrix elements are stored as formulas, and b) a knowledge about the behaviour of matrix elements as phase space integrands are coded in the form of modelling functions in order to get a fast generation of the partonic events. One can note that these data bases include a rather simple variety of subprocesses, mainly of the $`22`$ type.
If one tries to include a hard subprocesses with 3, 4 and more particles in the final state, large problems appear. Indeed, the size of matrix elements increases very fast. For example, in the $`24`$ case, the size of the code for evaluation of helicity amplitudes for one subprocess is at the 100’s of Kbyte level. However, the main problem lies elsewhere; it is impossible to construct an analytical formula for matching peaks and other structures of the rather singular behaviour of matrix elements. In the $`23`$ case, phase space has 4 dimensions plus two for convolution with the PDF’s; in the $`24`$ case 7+2 dimensions are present, and so on. As a result, the set of kinematical singularities has, as a rule, a very complicated positioning in the multidimensional phase space. This particular problem was not solved accurately, e.g., when the $`Zbb`$ final state was implemented in PYTHIA 5.7.
Let us discuss further this somewhat delicate point. It is necessary to integrate the squared matrix element over the phase space in order to obtain the cross section. Precise information about the behaviour of the integrand then is necessary for further event generation. This information can be obtained at the step of the phase space integration. The problem is that the integrand, as a rule, has a singular behaviour with sharp kinematical peaks connected with different denominators (propagators) of Feynman diagrams. This problem is caused, in particular, by the circumstance that one has to take into account nonzero masses of particles in many important cases, especially if accurate calculations are needed. The masses of elementary particles can have extremely small values, e.g. the masses of the 1st generation quarks (few MeV), and can also be zero (for the photon and gluon). At the same time, other parameters are of the order of a hundred GeV, e.g. masses of $`W`$ and $`Z`$ bosons and $`t`$-quark. Moreover the collision energy can also have a very large value, e.g. 14000 GeV for LHC processes, and some other important variables, like the transverse energy of jets, are at the hundred GeV scale or even greater. This huge scale interval for different parameters causes serious computational problems which result in the appearance of sharp peaks for the integrand. So, at the step of phase space parameterization, one has to include a regularization of the integration measure in order to smooth the singularities of the integrand.
LHC phenomenology requires the computation of a wide spectrum of hard subprocesses with 3,4 and even more particles (partons) in the final state. This is a common need for all of the physics working groups: QCD, EW, Higgs, SUSY etc. These requirements are especially common for new physics searches. Furthermore, a major challenge results from background analyses, where QCD subprocesses play a major role with, in many cases, multiparton final states. As a rule, for each LHC discovery reaction, one should calculate several QCD processes giving both irreducible and reducible backgrounds. The parton-shower generation of multiparticle final states is usually utilized in this situation. However, this can be too crude an approximation for many important studies leading to sometimes grossly unreliable predictions.
We emphasize that ACS can give the possibility to compute accurately a variety of LHC processes (and in particular QCD processes) with 3, 4 and more bodies in the final partonic state. Indeed, the first problem (the size of the matrix element computing code and the difficulty to obtain the exact matrix elements) is solved in ACS by the automatic generation of the corresponding code. This step is fast and pain-free from the viewpoint of the user. The second problem (the accurate integration over the multidimensional phase space) is solved in ACS by the generation of kinematics where the necessary regularizations are included. For example, in CompHEP the user has to list a set of singular propagators using the menu system. After that, the code for kinematics (with regularizations) is generated automatically. In GRACE, a library of kinematics (with regularizations) is used and the user has to make the necessary choices. Thus, the high art (mathematics and programming), needed to elaborate the sharp peaks, is enclosed in a form hidden from the user, giving him a possibility to compute complicated processes.
At the step of integration over the phase space, ACS uses adaptive Monte Carlo integrators (VEGAS in CompHEP, and BASES in GRACE). To match the complete set of singularities, the multichannel MC approach is utilized. As a result, the phase space grid is created with an accurate mapping of the singular behaviour of the matrix element. This complex body of information (let us call it MEgrid) has a rather large size that rapidly increases with the number of phase space dimensions. One can consider MEgrid as a multidimensional analog of the modelling function used in PYTHIA and other similar packages for the effective generation of partonic events. Of course, this information can not be expressed in analytical form. It is necessary to point out also that the convolution with parton distributions should be made at the same stage as the integration over the phase space. Indeed, the contributions of different subspaces (in particular different kinematical peaks) can depend largely on the partonic collison energy, $`\widehat{s}`$, resulting from the information stored in MEgrid.
ACS can be considered as a tool for the automatic generation of the data base of hard subprocesses for physical generators like PYTHIA, ISAJET and HERWIG. However, it is difficult to imagine that the data base created can be implemented in the code of these generators. This is due, first of all, to the size of the generated codes. Thus, we propose a two stage approach. At the first stage, ACS is used resulting in a cross section and MEgrid for the subprocess under evaluation. This can be stored in a special LHC data base. This data base can be used for the effective generation of partonic events. In GRACE, it is available with the SPRING generator, and in CompHEP by a relatively straightforward procedure and an effective generator is under construction). The output is a partonic event flow that can be used as an input for physical generators like PYTHIA, ISAJET and HERWIG; this is second stage of the full simulations. At this stage partons (quarks and gluons) should be hadronized and unstable particles decayed. We note that in PYTHIA there exists a rather flexible interface for such a two stage approach, the option for inclusion of external processes through the routine PYUPEV.
This is a general view on the way in which ACS (GRACE and CompHEP in particular) can be used for the simulation of LHC processes. Below we discuss some specific features of this technology with special attention to QCD aspects.
## 2 General Considerations about GRACE and CompHEP
With CompHEP and GRACE the user can evaluate hard subprocesses at the tree level, i.e. Feynman diagrams are generated without loops. This corresponds to the basic request for LHC phenomenology. However, it is well known that QCD next-to-leading corrections are large, as a rule, for LHC processes. In many cases these corections can be accounted for in the form of so-called K-factors and one can include them easily in tree level calculations. Nevertheless, in many important cases an explicit evaluation of higher order corrections is necessary. At this moment it is not clear how to automate calculations of LHC processes at NLO level. The problem is connected, in particular, with the circumstance that different resummations of large logarithms should be included in order to get reliable NLO predictions. The interface between resummation techniques and event generators is under intensive discussion now, and at the present Workshop also. We note in this respect, that the GRACE package includes the code for the generation and evaluation of one-loop diagrams.
The user interface should provide the possibility to calculate complicated processes for users not experienced in programming. CompHEP has a (graphical) menu driven system where the user proceeds through all steps of the calculation without any programming. In GRACE, the user needs to write a few simple interface routines.
The information on the GRACE system and its products can be found at
http://www-sc.kek.jp/minami/
The code of CompHEP is free for users and one can take it from the following Web page
http://theory.npi.msu.su/comphep
where the user’s manual is available in PS format (see also hep-ph/9908288). The CompHEP package, adapted for LHC processes (see next section) is installed on the SUN platform
/afs/cern.ch/cms/physics/COMPHEP/v33-SUN
and on the PC/Linux platform
/afs/cern.ch/cms/physics/COMPHEP/v33-Linux
The interface between CompHEP and PYTHIA has beencreated with the corresponding code available from the address:
/afs/cern.ch/cms/physics/comp-pyth
where one can find a short description in the file README. With this interface, the partonic event flow for any processes calculated with CompHEP can be sent to PYTHIA to generate physical events.
## 3 QCD aspects in automatic computations
In this section we discuss the treatment of QCD effects in the case of automatic computations, and consider CompHEP options as an example. As has been discussed above, CompHEP calculates only at tree level, and so at leading order (LO). Thus, the main problems concerning an accurate accounting of QCD effects are outside the discussion. Nevertheless, some important QCD dependencies can not be avoided even at tree level and the corresponding options are available for users. These aspects are: a) parton distributions, b) QCD scale, and c) running strong coupling constant.
Parton distributions.
In CompHEP the specification of initial states in the collision process under evaluation can include the convolution with structure function. So, in the case of hadron collisions, the cross section is evaluated as an integral
$$\sigma (s)=_0^1𝑑x_1𝑑x_2f_i(x_1,Q)f_j(x_2,Q)\widehat{\sigma }_{ij}(x_1x_2s)$$
where $`f_i`$ are the corresponding parton distributions, $`\widehat{\sigma }`$ is the partonic cross section and $`Q`$ is the QCD scale.
In CompHEP v.33, installed at CERN (see address above), parton distributions from two pdf families are implemented, MRS and CTEQ, and in particular the following versions: 1) MRS(A’) and MRS(G) , 2) CTEQ4l and CTEQ4m . Note that CTEQ4l is a LO parametrization, while in all others the evolution of parton distributions is realized in the next-to-leading (NLO) approximation.
In addition, a special interface is available to include a user’s defined parton distribution. By this way one can implement the most recent parametrizations (at this moment CTEQ5 and MSRT). See the CompHEP user’s manual for the corresponding procedure (section 3.6.2).
Choice of QCD scale.
The factorization theorem states that parton distributions depend not only on the Bjorken variable $`x`$ but also on some parameter $`Q`$ which characterizes the energy (or momentum) scale at which the QCD effects give the main contribution to the hard subprocess. This parameter is set by the user for each specific QCD process. It is possible to set a fixed scale or a running scale. In the later case, $`Q^2`$ can be a squared linear combination of any set of initial and outgoing particles momenta, for example, $`(p_1p_3)^2`$, $`(p_1p_3p_4)^2`$, $`(p_3+p_4)^2`$ and so on (initial and outgoing momenta enter with opposite signs). The corresponding settings are made through the option QCD SCALE of the numerical menu.
Running $`\alpha _s`$.
It is the nature of strong interactions that there is no absolute normalization of the corresponding coupling constant. This is in contrast to the value $`1/137`$ for the electromagnetic constant known with high accuracy from classical electrodynamical experiments. Instead, we have a function for $`\alpha _s`$ rather than a constant. Even in the leading order approximation, $`\alpha _s^{LO}=6\pi /[(332n_f)\mathrm{log}Q/\mathrm{\Lambda }^{(n_f)}]`$, where $`Q`$ is the QCD scale of the hard subprocess under the evaluation with $`\mathrm{\Lambda }`$ the so-called QCD fundamental parameter. Then, $`n_f`$ is the number of parton flavours with masses lower than $`Q`$. The $`n_f`$ dependence in QCD parameter $`\mathrm{\Lambda }`$ matches the quark mass threshold effects.
In the version of CompHEP installed at CERN, the running $`\alpha _s`$ is realized in LO, NLO and NNLO. All of the corresponding formulas are based on the choice of $`\mathrm{\Lambda }^{(6)}`$ (see Review of Particle Physics p.81. The user can find the corresponding switch in the option QCD SCALE in the numerical menu.
Therefore, to evaluate QCD processes with CompHEP, one has, first of all, to fix the normalization of $`\alpha _s`$. The popular normalization point is the mass of $`Z`$ boson, $`Q=M_Z`$. By changing the parameter $`\mathrm{\Lambda }^{(6)}`$, the user should set the strong coupling at the appropriate value, say $`\alpha _s(M_Z)=0.118`$. Then, the user has to choose the order for the running $`\alpha _s`$ (LO, NLO or NLO). Finally, the user has to define the QCD scale $`Q`$, which will be used both for the evaluation of $`\alpha _s`$ at this scale and in the parton distributions.
Thus, the complete LO calculations of LHC processes are available, with the matrix element, parton distributions and running strong coupling constant calculated in the lowest order of the perturbation theory. This is a self-consistent starting point in the phenomenological analysis; when/where higher order corrections are available, all elements of the calculation can be calculated at the higher order and then compared to the leading order result.
However, it is also common for phenomenologists use a mixed approach, with the matrix element evaluated at LO but the parton distributions and running $`\alpha _s`$ taken in NLO approximation. Surely, only a part of the NLO corrections is accounted for in this case. We note that this option is also available for users in CompHEP calculations.
## 4 Partonic Subprocesses
When hadronic collision processes are evaluated, especially in the case of a large number of final state particles, one serious problem is the large number of contributing partonic subprocesses. This occurs because of the quark and gluon content of the initial hadrons and CKM quark mixing. For example, at LHC energies, 180 subprocesses contribute to the $`W+2jets`$ and 292 subprocesses to the $`W+3jets`$ production (taking into account only quarks of the first two generations ). During this workshop a new method has been proposed to avoid a multiplication of channels due to the mixture of quark states . The method leads to a simple modification of the rules for the evaluation of the cross sections and distributions. It is based on the unitary rotation of down quarks, thus providing the transportation of mixing matrix elements from vertices of Feynman diagrams to the parton distribution functions. As a result, one can calculate cross sections with a significantly smaller number of subprocesses contributing. For the examples mentioned above, one needs to evaluate (with the new rules) only 21 and 33 subprocesses, respectively, in order to compute the cross sections for the $`W+2jets`$ and $`W+3jets`$ processes. The matrix elements of the subprocesses are calculated without quark mixing, but with a modified PDF convolution which now depends on the quark mixing angle and the topologies of the gauge invariant classes of diagrams contributing to the subprocesses. The method proposed has been incorporated into the CompHEP program and checked with many examples.
## 5 PEVLIB - library of LHC processes
Now the library of CompHEP based event generators for LHC processes has been started at the address:
/afs/cern.ch/cms/physics/PEVLIB
The following QCD processes are stored already in this library: $`Zb\overline{b}`$, $`Wb\overline{b}`$, $`t\overline{t}b\overline{b}`$ and some others. In the corresponding directories (with the names literally corresponding to the final states) unweighted events are stored (see the files README in these directories for details about evaluation of the corresponding samples of events).
Together with the CompHEP-PYTHIA interface code (see discussion above) these event files can be used for full LHC simulations with the help of PYTHIA package and detector simulation software in the standard way.
Let us discuss the process $`Zb\overline{b}`$ in order to supply more details. In the directory
/afs/cern.ch/cms/physics/PEVLIB/Z\_b\_b
the file \__pevZbb includes about 200000 unweighted events with the final state $`Zb\overline{b}`$. Each event includes the Lorentz momenta of all particles in the initial and final states. In the present version of this library, there is no information about the color flow in the event. Thus, only the Independent Fragmentation Model can be used for the hadronization. Of course, the user can use the Lund model; for this one has to define the corresponding color flows by hand in the routine PYUPEV. The same remark is valid also for FSR (final state radiation), what is switched off by default in CompHEP-PYTHIA interface. In the same time ISR (initial state radiation) is switched on automatically.
Note that the user can generate more events than stored in the library. In the corresponding subdirectories (indicated in the file README) the generators are stored in the form of the executable code (at this moment for SUN platform only). These generators are the corresponding CompHEP codes for the process with the proper set of kinematical regularizations.
The library PEVLIB is under construction now. New processes will be added. The structure and user’s interface will be developed.
## 6 Acknowledgements
The work of V.I., A.P., A.S, E.B. and A.S. was partially supported by the CERN-INTAS grant 377 and RFBR-DFG grant 99-02-04011.
Monte Carlo Event Generators at NLO
J. Collins
##### Factorization for inclusive processes
* Normal proofs of factorization are for inclusive cross sections.
* To get the simplifications in the factorization formula, as compared with the exact cross section, one makes suitable approximations.
* The approximation is to the hard scattering part $`H`$ of a cross section, as in:
which is a graph for the DIS cross section. The kinematics of the external lines of $`H`$ are changed to massless on-shell partons with zero transverse momentum. Also the internal (light parton) lines of $`H`$ are made massless.
* The approximation is correct to the leading power of $`m^2/Q^2`$, where $`m`$ is a typical hadronic scale and $`Q`$ is the scale of the hard scattering (e.g., $`Q^2`$ is the virtuality of the virtual photon in DIS).
* Subtractions are applied to the hard scattering to cancel double counting. A high-order graph for the hard scattering has subtractions that correspond to smaller hard subgraphs (and hence smaller regions of momentum space).
* The exact form of the subtractions corresponds to the approximations made for the smaller regions.
##### Monte-Carlo event generators at NLO (and beyond)
Current event generators essentially use the leading order for the hard-scattering coupled to an algorithm that approximates the exclusive structure of the low virtuality parts of graphs for the cross section. The algorithm is in an improved leading logarithm approximation.
With the exception of the recent paper by Friberg and Sjöstrand , previous attempts, e.g., , at incorporating NLO corrections have tended to implement them by a reweighting of the events generated by showering from the LO matrix elements. Normal NLO “analytic” calculations for the hard-scattering coefficients for inclusive scattering are not usable as they stand in a Monte-Carlo event generator, because they involve singular distributions.
Here I summarize some ideas to remedy this situation. I have applied them to one specific case in DIS in Ref. , but I think they can be generalized:
* Separate groups of events are generated with LO and NLO hard scattering coefficients.
* To obtain the NLO coefficients, the same general ideas are used as in inclusive hard scattering.
* The methods are applied both to the hard scattering itself and to the showering kernels.
* However, because we are now working with exclusive processes, the form of the approximations is different. The approximated graphs must satisfy the following requirements:
+ Exact 4-momentum conservation must be obeyed for each subprocess (hard scattering or one stage of the showering). I.e., $`p_i^\mu =p_f^\mu `$, where $`p_i^\mu `$ and $`p_f^\mu `$ are the total initial and final momenta the the subprocess.
+ The approximation may change the momenta of the internal lines but it must preserve the momenta of the external lines. This avoids the problem with having singular distributions.
* A cut-off is applied, for otherwise the approximated graphs give ultra-violet divergences when integrated to large transverse momentum. The kinematics associated with exact momentum conservation do provide a cutoff, but such a cutoff tends to violate factorization of the momentum-space integrals. So a separate artificial cutoff is better, and probably makes for a better implementation of the algorithm.
* Conventionally, in a Monte-Carlo a sharp cutoff is used. But a smooth cut-off will probably be better for numerical work. It will also make it easier to get positive cross sections.
* The dependence on the cutoff is a generalized renormalization-group transformation, and the exact cross section is independent of the form of the cutoff.
* Separate explicit soft factors are needed, as in the factorization theorem for the $`q_T`$ distribution for the Drell-Yan process; unlike the case of an inclusive cross section the cancellation of the soft region is not complete. This issue is not treated in Refs. , but will need further work, which is in progress.
##### Example of application to DIS
The parton model graph of Fig. 2 is combined with showering to give a LO cross section that can be summarized as
$$\sigma _{\mathrm{LO}}=\text{Born graph}\times \text{initial-state showering}\times \text{final-state showering}.$$
(1)
The NLO cross section is obtained from subtracted one-loop graphs, and the hard-scattering coefficient is of the form
$`\text{}\text{}\times \text{Jacobian}\times \text{cutoff function}.`$ (2)
The first term is an unsubtracted NLO graph. The subtraction corresponds to the approximation made in LO. Above the horizontal line, the following replacement for the internal momentum $`l^\mu `$ is made:
$`l`$ $``$ massless on-shell, zero transverse momentum (3)
$``$ $`\text{Lorentz-transformed momentum with correct final state}.`$ (4)
The first step is the standard approximation. The second step is needed to obtain the correct kinematics with conservation of 4-momentum. It is somewhat non-trivial to implement consistently since the subtraction term is needed when it is far from the collinear region. The bulk of my work in Refs. is about constructing a definite correct and consistent implementation. Correctness here means that the subtraction term is the order $`\alpha _s`$ approximation to the showering in the LO Monte Carlo algorithm.
##### Differences between conventional “analytic” method and Monte-Carlo method
These differences can be illustrated by the following mathematical example. Warning: In a number of respects this example is over-simplified. For example, it does not take account of the parton densities. As explained in Ref. , the parton densities in the Monte Carlo are not in the usual $`\overline{\mathrm{MS}}`$ scheme. However, they are related to them by definite formulae.
Let the unapproximated unsubtracted integrand at NLO be
$$\frac{d\widehat{\sigma }}{d^2k_T}=\frac{Q^2}{(Q^2+k_T^2+m^2)(k_T^2+m^2)}.$$
(5)
The conventional approach obtains the hard-scattering coefficient by setting the mass $`m`$ to zero. The integral is then infinite, and a subtraction is inserted which consists of a delta function at $`k_T=0`$ with an infinite coefficient. The result is a $`+`$ distribution, which has a finite integral:
$$\frac{d\widehat{\sigma }}{d^2k_T}=\frac{Q^2}{Q^2+k_T^2}\left(\frac{1}{k_T^2}\right)_+=\frac{Q^2}{(Q^2+k_T^2)k_T^2}C\delta ^{(2)}(k_T),$$
(6)
where $`C`$ is an infinite constant, defined with the aid, for example, of dimensional regularization.
In my new approach the subtracted integrand is
$$\frac{d\widehat{\sigma }}{d^2k_T}=\frac{Q^2}{(Q^2+k_T^2)k_T^2}\frac{f(k_T/\mu )J(k_T)}{k_T^2},$$
(7)
where $`f(k_T/\mu )`$ is the previously mentioned cutoff function, which is unity for small $`k_T`$ and zero for large $`k_T`$. As usual $`\mu `$ is the factorization scale. The factor $`J(k_T)`$ symbolizes the Jacobian that is necessary in the transformation from the variables appropriate for generation of events and the variables for the measured particles.
NLO and NNLO Calculations
V. Del Duca and G. Heinrich
## 1 The NLO and NNLO program
QCD calculations of multijet rates beyond the leading order (LO) in the strong coupling constant $`\alpha _s`$ are usually quite involved. Nowadays we know (see Section 1.2) how to perform in general calculations of the next-to-leading order (NLO) corrections to multijet rates, and almost every process of interest has been computed to that accuracy. Instead, the calculation of the next-to-next-to-leading order (NNLO) corrections is still at an organizational stage and represents a main challenge. Why should we perform calculations which are technically so complicated ?
The general motivation is that the calculation of the NLO corrections allows us to estimate reliably a given production rate, while the NNLO corrections allow us to estimate the theoretical uncertainty on the production rate. This is achieved by reducing the dependence of the cross section on the renormalization scale, $`\mu _R`$, and for processes with strongly-interacting incoming particles the dependence on the factorization scale, $`\mu _F`$, as well.
An example is the determination of $`\alpha _s`$ from event shape variables in $`e^+e^{}3`$ jets . Although the NLO contributions to $`e^+e^{}3jets`$ have been computed for some time now , the NNLO contributions have yet to be obtained. A calculation of these NNLO contributions would be needed to further reduce the theoretical uncertainty in the determination of $`\alpha _s`$.
We present in Section 1.1 an additional motivation for performing QCD calculations at NNLO, which is specific to the LHC program, and we outline in Section 1.2 how QCD calculations at NLO are implemented and in Section 1.3 how QCD calculations at NNLO could be performed.
### 1.1 Higgs production
The main goal of the LHC physics program is the investigation of the mechanism of the electroweak symmetry breaking, and namely the search and detection of the Higgs boson. If the Higgs boson is light (100 GeV $`m_H`$ 140 GeV), the rare decay channel in two photons, H$`\gamma \gamma `$, provides the best signature . Since the signal-to-background ratio is quite low ($`7\%`$), the analysis of this channel promises to be demanding. Our theoretical understanding of signal and background is still preliminary: the NLO QCD corrections to the signal are known to be quite large $`(𝒪(100\%))`$ . Also the QCD background $`pp\gamma \gamma `$, given at LO by the parton subprocess $`q\overline{q}\gamma \gamma `$, is known to NLO , with the full NLO fragmentation contributions having just been evaluated . However, $`pp\gamma \gamma `$ receives a sizeble contribution from NNLO corrections because of the large gluon luminosity of the subprocess $`gg\gamma \gamma `$ appearing first at NNLO . Thus in order to have a reliable theoretical estimate both the signal and the background need to be determined at least to NNLO accuracy.
In order to improve the signal-to-background ratio, Higgs production in association with a high transverse energy ($`E_T`$) jet, $`ppHjet\gamma \gamma jet`$, has been considered . This production rate offers the advantage of being more flexible in choosing suitable acceptance cuts to curb the background. $`ppHjet`$ is known to LO exactly , while the NLO corrections have been computed in the infinite top-mass limit. The NLO corrections to the signal are large. However, it is believed that the background, $`pp\gamma \gamma jet`$, can be more reliably calculated because LO production is dominated by the parton subprocess $`qgq\gamma \gamma `$, which benefits from the large gluon luminosity, while the subprocess $`ggg\gamma \gamma `$, which is believed to dominate the NNLO contribution, yields a comparatively smaller contribution . Thus, even though the signal, $`ppHjet`$, likely needs be computed at NNLO accuracy, it should suffice to evaluate the background, $`pp\gamma \gamma jet`$, at NLO. The NLO corrections to the background, though, have yet to be computed, with the appropriate QCD amplitudes having just been evaluated .
### 1.2 NLO algorithms and one-loop amplitudes
In recent years it has become clear how to construct general-purpose algorithms for the calculation of multijet rates at NLO accuracy. The crucial point is to organise the cancellation of the infrared (i.e. collinear and soft) singularities of the QCD amplitudes in a universal, i.e. process-independent, way. The universal terms in a NLO calculation are given by the tree-level splitting and eikonal functions, and by the universal structure of the poles of the one-loop amplitudes . The universal NLO terms and the process-dependent amplitudes are combined into effective matrix elements, which are devoid of singularities. The various NLO algorithms (phase-space slicing , subtraction method , dipole formalism and subtraction-improved slicing ) provide different methods to construct the effective matrix elements. These can be integrated, analytically or otherwise numerically, in four dimensions. The integration can be performed with arbitrary experimental acceptance cuts.
Then the remaining work to be performed to calculate a production rate at NLO is to compute the appropriate tree and one-loop amplitudes. To compute $`n`$-jet production at NLO, two sets of amplitudes are required: a) $`n`$-particle production amplitudes at tree level and one loop; b) $`(n+1)`$-particle production amplitudes at tree level. If the one-loop amplitudes are regularised through dimensional regularization (DR) by evaluating them in $`d=42ϵ`$ dimensions, it suffices at NLO to compute them to $`𝒪(ϵ^0)`$. As an example, in Fig. 1 we show the squared matrix elements which are required to calculate the NLO corrections to $`e^+e^{}3jets`$.
Efficient methods based on the color decomposition of an amplitude in color-ordered subamplitudes, which are then projected onto the helicity states of the external partons, have largely enhanced the ability of computing tree and one-loop amplitudes. Accordingly, tree amplitudes with up to seven massless partons and with a vector boson and up to five massless partons have been computed analytically. In addition, efficient techniques to evaluate numerically tree multi-parton amplitudes have been introduced , and have been used to compute tree amplitudes with up to eleven massless partons . The calculation of one-loop amplitudes can be reduced to the calculation of one-loop $`n`$-point scalar integrals . The reduction method allowed the computation of one-loop amplitudes with four massless partons and with a vector boson and three massless partons . However, one-loop scalar integrals present infrared divergences, induced by the massless external legs. For one-loop multi-parton amplitudes, the infrared divergences hinder the reduction methods of ref. . This problem has been overcome in ref. . Accordingly, one-loop amplitudes with five massless partons and with a vector boson and four massless partons have been computed analytically. The reduction procedure of ref. has been generalised in ref. , where it has been shown that any one-loop $`n`$-point scalar integral, with $`n>4`$, can be reduced to box scalar integrals, and that in the reduction of $`n`$-point tensor integrals, all higher dimensional ($`d>42ϵ`$) $`n`$-point integrals with $`n>4`$ drop out. The calculation of one-loop multi-parton amplitudes thus can be pushed a step further in the near future.
### 1.3 NNLO calculations
Eventually, a procedure similar to the one followed at NLO will permit the construction of general-purpose algorithms at NNLO accuracy. It is mandatory then to fully investigate the infrared structure of the phase space at NNLO. The universal pieces needed to organise the cancellation of the infrared singularities are given by the tree-level double-splitting , double-eikonal and splitting-eikonal functions, by the one-loop splitting and eikonal functions, and by the universal structure of the poles of the two-loop amplitudes . These universal pieces have yet to be assembled together, to show the cancellation of the infrared divergences at NNLO.
Then to compute $`n`$-jet production at NNLO, three sets of amplitudes are required: a) $`n`$-particle production amplitudes at tree level, one loop and two loops; b) $`(n+1)`$-particle production amplitudes at tree level and one loop; c) $`(n+2)`$-particle production amplitudes at tree level. In Fig. 2 we show the squared matrix elements which are required to calculate the NNLO corrections to $`e^+e^{}3jets`$. In DR at NNLO, the two-loop amplitudes need be computed to $`𝒪(ϵ^0)`$, while the one-loop amplitudes must be evaluated to $`𝒪(ϵ^2)`$ . The main challenge is the calculation of the two-loop amplitudes. At present, the only amplitude known at two loops is the one for $`Vq\overline{q}`$ , with $`V`$ a massive vector boson, which depends only on one kinematic variable. It has been used to evaluate the NNLO corrections to Drell-Yan production and to deeply inelastic scattering (DIS) . No two-loop computations exist for configurations involving more than one kinematic variable, except in the case of maximal supersymmetry . One of the main obstacles for configurations involving two kinematic variables is the analytic computation of the two-loop four-point functions with massless external legs, where significant progress has just been achieved. These consist of planar double-box integrals , non-planar double-box integrals , single-box integrals with a bubble insertion on one of the propagators and single-box integrals with a vertex correction . The two-loop four-point functions with massless external legs are needed for the computation of two-loop amplitudes in parton-parton scattering. Finally, the topical processes considered above, i.e. $`e^+e^{}3jets`$ and $`ppHjet`$ sport configurations involving three kinematic variables and require the analytic computation of two-loop four-point functions with a massive external leg. The two-loop four-point functions of this kind with up to five different denominators have been derived recently , while those with six and seven different propagators are still missing. Another obstacle is the color decomposition of two-loop amplitudes, which is not known yet. Substantial progress is expected in the next future on all the issues outlined above, which should make the present note soon outdated.
Finally, we mention that in the factorization of collinear singularities for strongly-interacting incoming particles, the evolution of the parton distribution functions ($`pdf`$’s) in the jet cross section should be determined to an accuracy matching the one of the parton cross section. For hadroproduction of jets computed at NLO, one needs the NLO, or two-loop, evolution of the $`pdf`$’s . Accordingly for hadroproduction at NNLO the evolution of the $`pdf`$’s should be computed to NNLO, or three-loop, accuracy. Except for the lowest five (four) even-integer moments of the three-loop non-singlet (singlet) splitting functions , no calculation of the NNLO evolution of the $`pdf`$’s exists yet. However, NNLO analysis based on the finite set of known moments have been performed for $`xF_3`$ and $`F_2`$ (non-singlet and singlet ). Furthermore, in ref. a quantitative assessment of the importance of the yet unknown higher-order terms has been performed, with the conclusion that they should be numerically significant only for Bjorken-scaling $`x<10^2`$.
The computation of the evolution of the $`pdf`$’s at NNLO accuracy is a main challenge in QCD. The NLO computation was performed with two different methods, one using the operator product expansion (OPE) in a covariant gauge , the other using the light-cone axial (LCA) gauge with principal value prescription . However, the prescription used in ref. has certain shortcomings. Accordingly, the calculation has been repeated in the LCA gauge using a generally correct prescription , which makes it amenable to extensions beyond NLO. On the other hand, using the OPE method, there had been a problem with operator mixing in the singlet sector, which has been fixed only recently, and the result finally coincides with the one obtained in the LCA gauge in ref. . Thus the calculation of the $`pdf`$ evolution at NLO accuracy is fully under control. Recent proposals for a calculation beyond NLO include extensions of the OPE technique, which have been used to recompute the NNLO corrections to DIS , and a computation of the $`pdf`$ evolution by combining the universal gauge-invariant collinear pieces . For the two-loop $`pdf`$ evolution, e.g., they are the collinear pieces mentioned at the beginning of this section.
Jet Algorithms
S. Catani and D. Zeppenfeld
## 1 Jet algorithms
Jet algorithms have the task to assign streams of hadrons in hard scattering processes to a jet, who’s energy, mass and momentum can then be related to a collection of partons in a perturbative QCD calculation. Although, at the experimental level, jets can be defined by using rather general and intuitive procedures, if we would like to compute jet cross sections and properties by using QCD perturbation theory, the definition of jets should fulfil stronger constraints to guarantee its perturbative safety. Perturbative safety means that the definition has to be infrared safe (jet properties cannot depend on the presence of arbitrarily soft partons), collinear safe (jet properties cannot change by replacing a parton with a set of collinear partons carrying the same total momentum) and collinearly factorizable (jet properties should be insensitive to partons radiated collinearly to the beam direction). If the jet definition is not perturbative safe, we cannot perform calculations order-by-order in perturbation theory because they are affected by uncancelled infrared divergences. Of course, in the full QCD theory (i.e. beyond perturbation theory) the perturbative divergences are regularized by small physical cutoffs related to hadron masses and the finite experimental resolution (size of calorimeter cells, energy thresholds, etc.). The physical cutoffs are always present, independently of the jet definition. However, in the case of a perturbative safe definition, their effects are suppressed by some inverse power of the jet transverse energy $`E_T`$, and thus they can be made small by sufficiently increasing $`E_T`$. This power suppression is not at work in perturbative unsafe jet definitions, where the effects of the small physical cutoffs can amount to large corrections (of order unity) to the perturbative results. Thus, perturbative safe definitions are preferred.
Jet algorithms start from a list of “particles” which we would like to freely associate with calorimeter cells or hadrons at the experimental level, and with partons in a QCD calculation. Each particle $`i`$ carries a 4-momentum $`p_i^\mu `$, which we take to be massless. The task is to select a set of particles which are emitted close to each other in angle and combine their momenta to form the momentum of a jet. The selection process is called the “jet algorithm”, the momentum addition rule is called the “recombination scheme”.
Let us start with a discussion of recombination schemes. In a hadron collider environment the arbitrary boost of the hard scattering system along the beam axis needs to be taken into account in the definition of angles to ensure collinear factorizability. This is achieved by using transverse momentum, $`p_T=\sqrt{p_x^2+p_y^2}`$, rapidity $`y=1/2\mathrm{log}(E+p_z)/(Ep_z)`$ and azimuthal angle $`\varphi `$ of the massless particles as the kinematic variables. When adding the massless 4-vectors of particles we obtain massive objects which only approximately correspond to the massless partons which we would like to associate with jets at tree level.
One popular choice, the Snowmass convention , leaves the question of jet mass open, by only defining total transverse energy, rapidity and azimuthal angle of a set of parton momenta, as the $`E_T`$ weighted sums of the individual particle variables. For the original massless particles $`E_{Ti}=p_{Ti}`$ and $`\eta _i=y_i`$. The corresponding recombined variables for a cluster of particles are then given by the total transverse energy
$$E_T=\underset{i}{}E_{Ti},$$
(1.1)
the cluster pseudorapidity
$$\eta =\underset{i}{}\frac{E_{Ti}}{E_T}\eta _i,$$
(1.2)
and the azimuthal angle of the cluster
$$\varphi =\underset{i}{}\frac{E_{Ti}}{E_T}\varphi _i.$$
(1.3)
Note that the designation of $`\eta `$ as pseudorapidity is purely conventional. It corresponds to neither the pseudorapidity nor the rapidity of the massive cluster and is approximately equal to either only in the limit of small cluster mass ($`<<E_T`$). The concomitant loss of Lorenz invariance is a serious disadvantage of the Snowmass convention. Another serious problem appears in resummation calculations (see Sect. 2): the kinematic boundary of jet $`E_T`$ shifts (from $`\sqrt{\widehat{s}}/2`$ in e.g. dijet kinematics) when including additional final state partons.
Because of these shortcomings we formulate all jet algorithms in the 4-momentum recombination scheme (also called E-scheme) in the following, i.e. the kinematic variables of a cluster of particles is given by direct addition of the 4-momenta of the individual massless particles:
$$p^\mu =(E,p_x,p_y,p_z)=\underset{i}{}p_i^\mu .$$
(1.4)
Since the resulting clusters have a clearly defined mass, we must distinguish transverse energy $`E_T`$ from transverse momentum $`p_T`$, and pseudorapidity $`\eta `$ from rapidity $`y`$. We define
$$p_T=\sqrt{p_x^2+p_y^2},\varphi =\mathrm{tan}^1\frac{p_x}{p_y},y=\frac{1}{2}\mathrm{log}\frac{E+p_z}{Ep_z},$$
(1.5)
The rapidity $`y`$ and azimuthal $`\varphi `$ should be used as the legoplot position of the jet when calculating its separation from other particles or jets. Auxiliary quantities are
$$\theta =\mathrm{cos}^1\frac{p_z}{\sqrt{p_x^2+p_y^2+p_z^2}},E_T=E\mathrm{sin}\theta ,\eta =\mathrm{log}\mathrm{tan}\frac{\theta }{2}.$$
(1.6)
### 1.1 The $`k_T`$ algorithm
The $`k_T`$ algorithm is a successive recombination algorithm. The idea is to recombine particles with nearly parallel momenta, beginning with the softest particles in the sample. This recombination stops once all clusters of particles are separated by a distance larger than $`D`$ in the legoplot. The $`k_T`$ algorithm starts from a list of protojets and their momenta, $`p_i^\mu `$, which in the beginning consists of the list of all particles:
* For each protojet $`i`$ define
$$d_i=E_{Ti}^2$$
(1.7)
and for each pair of protojets $`i,j`$ define a distance
$$d_{ij}=\mathrm{min}(E_{Ti}^2,E_{Tj}^2)\frac{(y_iy_j)^2+(\varphi _i\varphi _j)^2}{D^2}.$$
(1.8)
* Find the smallest of all $`d_i`$ and $`d_{ij}`$ and call it $`d_{min}`$.
* If $`d_{min}`$ is $`d_{ij}`$ then merge protojets $`i`$ and $`j`$ to form a new protojet $`k`$ of momentum
$$p_k^\mu =p_i^\mu +p_j^\mu $$
(1.9)
* If $`d_{min}`$ is $`d_i`$ then remove protojet $`i`$ from the protojet list and move it to the list of completed jets.
* Continue with step 1 until the list of protojets is empty.
The algorithm is infrared safe because it renders all soft partons harmless: it either takes them off the protojet list or it combines them with nearby harder partons. The only effect of the soft parton then is a shift in the momentum of the recombined cluster. However, this shift is small, disappearing in the infrared limit, which guarantees infrared safety. Also for collinear emission the algorithm is safe, because in the limit of zero angle between two partons, these two will have the smallest $`d_{ij}`$ and will thus be combined early on in the recombination process, thus restoring the momentum of the almost on-shell parton from which they originated by splitting. Finally, every original particle is assigned to exactly one jet, i.e. there are no splitting/merging issues to be resolved for the $`k_T`$ algorithm.
### 1.2 ILCA: an infrared safe cone algorithm
Cone algorithms are intended to cluster all energy within a given radius, $`R`$, around a point in the legoplot, to form jets. Naively, this procedure is both infrared and collinear safe: the effect of infrared radiation on the cluster momentum vanishes in the infrared limit, and the energy measured for the jet is the same whether a single particle is at the core of the cone or whether there has been collinear splitting. This naive expectation can easily be violated, however, by the prescription for selecting cones. Two examples illustrate this point .
Assume that cones are constructed around actual energy depositions only. In Fig. 1(a) two particles are emitted at a distance greater than the cone radius $`R`$ but smaller than $`2R`$ and therefore are assigned to separate cones, which are then identified as two distinct jets. The only difference in Fig. 1(b) is the emission of a third soft particle (a “soft gluon”) between the original two particles. Now the additional cone around the soft energy deposition encompasses all three particles and they will be classified as a single jet. The presence of a soft particle changes the classification of a hard event: this is an example for an infrared unsafe algorithm. In perturbation theory, at sufficiently high order, an arbitrary number of soft gluons will be radiated, hence, cones should be allowed anywhere in phase space to anticipate this feature of higher order corrections. An arbitrary restriction on allowed cone positions may lead to an infrared unsafe algorithm.
Similarly, an infinite number of collinear splittings occurs at higher order in perturbation theory. A possible collinear problem, resulting from $`E_T`$ ordering of particles, is illustrated in Fig. 2. The difference between the two situations is that the central (hardest) parton may split into two almost collinear partons. On the left-hand-side the distance between the lateral partons is larger than $`R`$ but the three hard partons all fall within a cone of radius $`R`$ around the central parton, which happens to have the largest $`E_T`$. As a result, all three partons are recombined to a single jet. Collinear splitting renders the right hand parton to be the one with the largest $`E_T`$. Drawing the first cone around the highest $`E_T`$ parton will recombine it with the two central partons and a separate jet is likely to be assigned to the remaining fourth parton. A differing jet number which depends on the presence or absence of collinear splitting must be avoided because the incomplete cancellation of the logarithmic divergences of real emission and virtual contributions will lead to a collinear unsafe jet algorithm. One can eliminate such ambiguities by making the selection or ordering of jet definition cones independent of the $`E_T`$ of individual particles. Also, allowing trial cones anywhere in phase space would have made the two situations in Fig. 2 more similar: allowing cones centered between the central and the outside partons from the start would lead to a more similar jet identification in the two cases.
The above considerations lead us to consider a cone algorithm which allows trial cones to be positioned anywhere in phase space, irrespective of the transverse momentum carried by individual particles or calorimeter cells. We start by formulating this seed-less algorithm at the calorimeter level, where the basic entities are calorimeter towers.
* Make a list of all calorimeter towers.
* Select the next tower on the list as the center of a trial cone of radius $`R`$.
Goto (4) if the list of towers is exhausted.
* Add the momenta of all towers inside the trial cone and determine the legoplot position $`(y,\varphi )`$ corresponding to this momentum.
* If this position is outside the selected tower, discard the trial cone and go to (1).
If $`(y,\varphi )`$ is inside the selected tower, add the set of towers inside the trial cone as a new entry to the list of protojets.
At this stage we have a list of protojets, and we need to split/merge them to make jets.
* Select the highest $`E_T`$ protojet remaining on the list. (If the list is exhausted jet identification for the event is complete.)
* Does the selected protojet share any towers with other protojets?
+ No: Move protojet to list of jets and continue with (4).
+ Yes: Find the highest $`E_T`$ protojet that shares towers with the selected protojet. (Call this the neighbor protojet.) Decide whether the $`E_T`$ in the shared cells is greater than a fraction $`f`$ of the $`E_T`$ in the neighbor protojet.
- No: Split the shared towers.
* Allocate shared towers to either the selected or the neighbor protojet depending on which jet center is closer.
* Calculate the new momenta for the modified protojets, i.e. their $`E_T`$, and legoplot positions $`(y,\varphi )`$. Continue with step (4).
- Yes: Merge the selected and neighbor protojets to form a new protojet.
Add the momenta of both protojets and determine the total $`E_T`$ and and the legoplot position $`(y,\varphi )`$. Continue with step (4).
The procedure that defines the list of protojets is infrared and collinear safe. The additional steps completely define how to solve the problem of overlapping cones. The critical overlap fraction $`f`$ is a free parameter of the algorithm and may be chosen as 50%, similar to the D0 choice in run I of the Tevatron (CDF uses 75%). The $`E_T`$-ordering of protojets in this split/merge step does not introduce collinear problems provided the cone size $`R`$ is chosen sufficiently large.
The definition of calorimeter towers, i.e. a discretization of $`(y,\varphi )`$ space, would be cumbersome in a theoretical calculation, and is indeed not necessary. In a perturbative calculation at fixed order, the maximal number, $`n`$, of partons is fixed. The only possible positions of stable cones are then given by the partitions of the $`n`$ parton momenta, i.e. there are at most $`2^n1`$ possible locations of protojets. They are given by the legoplot positions of individual partons, all pairs of partons, all combinations of three partons etc. In a perturbative calculation, e.g. via a NLO Monte Carlo program, the protojet selection of the seedless algorithm (steps (0) to (3) above) can then be replaced as follows:
* Make a list of all possible cone centers. These are the legoplot coordinates of all parton momenta $`p_i`$, of all pairs of parton momenta $`p_i+p_j`$, of all triplets of parton momenta $`p_i+p_j+p_k`$, etc. For each cone center record which set of partons defines it.
* Select the next cone center on the list as the center of a trial cone of radius $`R`$.
Goto (4) if the list of cone centers is exhausted.
* Add the momenta of all partons inside the trial cone and determine the legoplot position $`(y,\varphi )`$ corresponding to this momentum.
* If this position is different from the trial cone center, i.e. if the cone center record and the list of partons inside the trial cone disagree, discard the trial cone and go to (1).
If $`(y,\varphi )`$ is the trial cone center, add the set of partons inside the trial cone as a new entry to the list of protojets.
As before, different protojets may share partons, i.e. they may overlap. The required split/merge step is then identical to the calorimeter level steps (4) and (5), with towers replaced by partons as elements of protojets.
In an actual experiment the number of calorimeter towers may be very large (order 6000 for tower sizes of $`\mathrm{\Delta }\eta \times \mathrm{\Delta }\varphi =0.1\times 0.1`$ and an $`\eta `$ coverage of $`\pm 5`$ units of pseudorapidity). The calorimeter level algorithm may then be rather slow computationally. The question arises whether an acceptable approximation of the seedless algorithm can be constructed, analogous to the parton level short-cut, by considering only those towers which have energy depositions above a minimal seed threshold. One would like to replace the list of parton momenta above by the list of tower momenta with
$$p_{Ti}>E_{T,seed}.$$
(1.10)
Since the algorithm is infrared and collinear safe when $`E_{T,seed}=0`$, it is always possible to chose the seed threshold $`E_{T,seed}`$ low enough so that variations of $`E_{T,seed}`$ lead to negligeable variations in any observable under consideration.
One would like to include in the determination of jet momenta all towers, of course, which lie inside the cone of radius $`R`$ around the protojet axis. This requires an additional iteration of the cone axis in the parton level algorithm when a seed threshold is imposed. The steps leading to the definition of protojets can then be modified as follows:
* Make a list of all possible cone centers. These are the legoplot coordinates of all parton/tower momenta $`p_i`$ with $`p_{Ti}>E_{T,seed}`$, of all pairs of such parton/tower momenta $`p_i+p_j`$, of all triplets $`p_i+p_j+p_k`$, etc.
* Select the next cone center on the list as the center of a trial cone of radius $`R`$.
Goto (3) if the list of cone centers is exhausted.
* Add the momenta of all partons/towers inside the trial cone (also those with $`p_{Ti}<E_{T,seed}`$) and determine the legoplot position $`(y,\varphi )`$ corresponding to this cone momentum. Use $`(y,\varphi )`$ as the new center of the trial cone and iterate this step until the position is stable. The set of all towers/partons inside the final trial cone constitutes a new protojet. Continue with step (1).
* Eliminate all duplicate protojets, i.e. protojets with an identical set of towers/partons.
With these changes, the resulting algorithm (named Improved Legacy Cone Algorithm or ILCA) is quite close to those used in run I of the Tevatron. The main change is the inclusion of midpoints of seeds (the $`p_i+p_j`$ pairs) and of centers of larger numbers of seeds as additional seed locations for trial cones. Including these additional midpoints is absolutely crucial in perturbative calculations in order to achieve infrared safety (see discussion on Fig. 1). When dealing with data, these effects are somewhat diminished, because with sufficiently low seed thresholds $`E_{T,seed}`$, a large number of trial cones will be generated from actual soft energy depositions in the calorimeter. However, because these soft energy depositions will decide how many jets are reconstructed, one potentially introduces a high sensitivity of jet observables to soft hadrons, Monte Carlo modelling of soft particles etc. The inclusion of the extra midpoints eliminates these soft effects because observables no longer depend on whether soft emission actually took place.
## 2 Resummed calculations
The ILCA and $`k_T`$-algorithm eventually lead to jets whose topology is not extremely different from that expected on the basis of a naive definition in terms of cones in azimuth-rapidity space. This is obviously true for the ILCA, where jets can contain particles whose distance is smaller than $`2R`$ and have a shape that differs from a cone-shape only because of the merging/splitting procedure. In the $`k_T`$-algorithm, jets have no sharp boundaries, but opening angles of particles within each jet are, typically, smaller than $`D`$ and all opening angles between jets are larger than $`D`$. The detailed jet structure is, however, different in the two algorithms. Although both algorithms are perturbative safe, the differences show up in higher-order perturbative calculations.
Higher-order perturbative computations and, in particular, resummed calculations can be necessary in special kinematics configurations that lead to large logarithmically-enhanced contributions at any fixed order in perturbation theory. Typical examples are calculations of jet cross sections near the phase-space boundary and of the fine internal structure of jets (shape variables, subjets, etc.). These quantities can be strongly dependent on the jet definition. The corresponding perturbative calculations do strongly depend on the jet definition, because they are the result of the integration of the QCD matrix elements (which do not depend on jets) over phase-space regions whose boundaries depend on the fine details of jet kinematics.
As an example of this strong sensitivity, we can consider the one-jet inclusive cross section as a function of the transverse momentum $`p_T`$ of the jet. If the jet variable $`p_T`$ is defined by using the 4-momentum recombination scheme (see Eqs. (1.4)–(1.6)), the kinematical boundary is $`x_T1`$, where $`x_T=2p_T/\sqrt{S}`$. Close to the boundary $`x_T1`$, the perturbative contributions are enhanced by large logarithmic corrections $`(\alpha _S\mathrm{log}^2(1x_T))^n`$ that need to be resummed to all orders in $`\alpha _S`$. Techniques to perform this resummation can be developed (see below). However, if the jet variable $`p_T`$ is defined by using a recombination scheme that does not conserve the 4-momentum (e.g. the true $`p_T`$ in Eq. (1.5) is replaced by the variable $`E_T`$ in Eq. (1.1)), the kinematical boundary for $`x_T`$, although close to $`x_T=1`$, is not fixed: the corresponding large logarithms cannot be resummed because the $`x_T`$-boundary shifts in a complicated manner depending on the number of final-state partons in the calculation .
The feasibility of resummed calculations depends not only on the recombination scheme but also on the jet algorithm, as is well known for jets in $`e^+e^{}`$ annihilation . The jet definition of the $`k_T`$-algorithm is inspired by the parton shower picture of jet fragmentation . Thus resummed calculations can be carried out by using the analytic version of the recurrence techniques used to generate multiparton final states in Monte Carlo parton showers. This is demonstrated by explicit calculations of subjet multiplicity and rates in hadron collisions . The ILCA has still to be investigated in this respect. The two algorithms differ only slightly at highly inclusive level . Thus, in these cases (such as the large-$`x_T`$ behaviour of the one-jet inclusive cross section), resummed calculations in the ILCA should be feasible as in the $`k_T`$-algorithm. Studies of the internal structure of ILCA jets may instead be more difficult since it depends on the procedure to merge/split overlapping jets.
## 3 Conclusions
During the Les Houches workshop discussions were centered on the general properties of jet algorithms, in particular their infrared and collinear safety at the perturbative level. The $`k_T`$-algorithm and the seedless cone algorithm described in Section 1.2 fulfil these requirements. Beyond these theoretical concerns there are many experimental issues which need to be addressed to obtain a practical algorithm. Among these are ease of energy calibration, effects of underlying event and overlapping events in a high luminosity hadron collider environment, high jet reconstruction efficiency, and efficient use of computer resources in reconstructing jets. These issues have been addressed in a parallel study, during the Run II QCD Workshop at Fermilab . In particular it has been shown that the ILCA produces small corrections only, when compared with the jet algorithms used in run I of the Tevatron. We refer the reader to the Proceedings of the Run II Workshop for a detailed study of these effects.
A Study of the Underlying Event in Jet and Minimum Bias Events
J. Huston and V. Tano
## 1 Introduction
Due to the importance of the inclusive jet cross section as a test of perturbative QCD over a wide range of $`Q^2`$ values, it is necessary to carefully consider all systematic effects that influence its measurement. In addition to the hard interaction that produces the jets in the final state there is also an underlying event, originating mostly from $`soft`$ spectator parton interactions. Because of the softness of the scale, their contribution cannot be perturbatively calculated. There may also be a contribution due to $`semihard`$ interactions between spectator partons, which create $`minijets`$ at transverse momenta almost large enough for perturbative calculations, but much smaller than that of the primary interaction responsible for the highest $`E_T`$ jets in the event. This process is known as double parton scattering. Both of the above processes, as well as higher order radiation from the $`22`$ hard subprocess , contribute to the underlying event.
The experimental cross sections are most commonly compared to theoretical calculations at next-to-leading order (NLO) in the coupling constant $`\alpha _s`$, such as JETRAD or EKS . At NLO, there can be at most 3 partons in the final state, leading to the presence of either 2 or 3 jets, depending on whether the third parton is present in the final state and whether it ends up being clustered with one of the other two partons. As the jet clustering is based on a fixed cone algorithm, the contribution due to the underlying event must be subtracted from the jet cone, in order to compare the results with NLO QCD calculations. As will be seen below, one of the largest sources of systematic error for the inclusive jet cross section at low $`E_T`$ is due to the uncertainty on the subtraction of this underlying event.
The current ‘paradigm’ <sup>1</sup><sup>1</sup>1Dave Soper claims that this term is vastly overused but we choose to employ it anyway. is that the underlying event in jet events is similar to the average energy level found in ‘active’ minimum bias events. Thus, this energy level needs to be determined and subtracted from the energy in a jet cone before the jet data is compared to NLO theory predictions. CDF assumes an uncertainty of 30$`\%`$ on this underlying event subtraction which makes it the dominant error at low $`E_T`$.
The flip side of the above ‘paradigm’ is that the underlying event energy in a jet event (once the two leading jets have been subtracted) should be the sum of the minimum bias level contribution and the third parton in the NLO calculation. The preliminary results of a study in CDF designed to test the accuracy of these assumptions are described in this section.
In this section, the experimental results will be compared to the HERWIG Monte Carlo which has the $`22`$ matrix elements for jet production, parton showering in the initial and final state, and a model for the underlying event. The ultimate result from HERWIG consists of the 4-vectors of the final state hadrons. In HERWIG the soft underlying event in the hadron-hadron collision is assumed to be a soft collision between the two ‘beam clusters’, which contain the spectators from the incoming hadrons. The model for the simulation of underlying event uses the $`p\overline{p}`$ event generator from the UA5 Collaboration, which is modified to make use of the HERWIG fragmentation algorithm.
The results from HERWIG can be quoted at the parton level (excluding the soft underlying event portion), the hadron level and/or comparisons can be made at the detector level after the Monte Carlo hadrons are passed through the CDF detector simulation program QFL . For most of the results that we will be reporting, the QFL comparisons will be crucial.
To summarize, the purpose of this analysis is to examine in detail the underlying event in jet events and to understand whether the amount subtracted from the jet cones is correct and whether the uncertainty assumed can be reduced. In addition, a test will be made as to how well HERWIG models the underlying event energy in jet events as well as in minimum bias events.
## 2 Underlying Energy at $`90^0`$ in jet events
Events generated with HERWIG were passed through QFL. The HERWIG code was adapted to produce the same information as found in the data samples. This information includes the energy, the position and the number of calorimeter towers of the jets in the event, together with the energy and the number of towers in two cones situated at $`\pm 90^0`$ in $`\varphi `$ and at the same $`\eta `$ as the leading jet. $`\varphi `$ and $`\eta `$ are respectively the azimuthal and polar angle.
For each jet event, two cones of radius $`0.7`$ at $`\eta =\eta _{LeadJet}`$ and $`\varphi =\varphi _{LeadJet}\pm \frac{\pi }{2}`$ were examined. The energy in each cone was determined for two different calorimeter tower thresholds: 50 and 100 MeV. A cut of 100 MeV on tower energies is typically used for jet analyses. For most of the comparisons to follow, a 50 MeV cut was used, though, since we are interested in possible contributions to the tower energies from a number of different sources. The two cones were used to study the underlying event energy because they are supposed to be in a semi-quiet region, far away from the two leading jets, but still in the central rapidity region. Given the non-uniform response of the CDF detector as a function of rapidity, the latter criterion is essential. The leading jet was required to be in the central region, $`|\eta |<0.7`$, the same as in the inclusive jet analysis. No requirement was made on the location of the second jet. In Fig 1 the calorimeter central region is shown as ‘unrolled’; $`\eta `$ ranges are between -1 and +1, while $`\varphi `$ goes from $`0^0`$ to $`360^0`$. The leading jet cone and the two cones under study are shown.
The $`E_T`$ distributions inside the two cones provide an idea of the contribution of the underlying event in the jet cone. For each event the cone which has the maximum energy and the cone with the minimum energy were labelled. This is useful because NLO perturbative corrections to the $`22`$ hard scattering can contribute only to one of these two regions . The difference between the maximum and the minimum cone provides information on this contribution, while the minimum cone gives an indication of the amount of underlying event. The (roughly constant) underlying event contribution should be suppressed in the difference.
The data were required to have one and only one vertex in order to insure that there is only one interaction per event. A similar cut was made in the simulation. In Figure 2 the transverse energy inside the two cones (max and min) is plotted as a function of the $`E_T`$ of the leading jet. It can be clearly observed that HERWIG and the data have a similar behaviour for the max and min cone; the min cone stays flat while the max cone increases with the $`E_T`$ of the leading jet. The increase of the max cone energy with increasing jet $`E_T`$ is easily understandable. What may be surprising is the flatness of the min cone energy as the lead jet transverse energy increases. Contrary to the pronouncements of some of the politicians of our day, a rising tide does not raise all boats (or cones), but instead favors the cone in the highest tax bracket. Of course, the division into a max and min cone partially encourages this effect through selection. However, the level of flatness is still somewhat surprising.
It is evident that there is an offset between data and the HERWIG+QFL simulation of about $`800`$ MeV for the max cone and $`500`$ MeV for the min cone. If the tower threshold is increased from $`50`$ to 100 MeV, the transverse energy decreases by about $`180`$ MeV in the data (both cones), while in HERWIG the transverse energy decreases by about $`70`$ MeV in the max cone and $`40`$ MeV in the min cone.
The difference between the transverse energy in the max and in the min cones has a similar trend in both data and simulation(Figure 3). There is still an offset but the offset decreases to about $`300`$ MeV. It appears that the max-min distribution starts going down again at very high $`E_T`$ (perhaps due to kinematic suppression), although the statistics become poor.
In Figure 4 the $`E_T`$ frequency distributions for data and HERWIG+QFL are compared for four different jet sub-samples. In this plot, the $`E_T`$ values for max-min are plotted, for both data and HERWIG+QFL. The number of entries is scaled to easily allow a direct comparison. The $`E_T`$ distribution of the max-min cone for HERWIG+QFL looks very similar to that of the data. Here the contribution of the underlying event as minimum bias data is presumably removed.
### 2.1 Parton-Hadron-Detector level
With HERWIG (unlike the data), we have the advantage of being able to examine the energy distributions not only at the detector level, but also at the hadron and parton levels. The HERWIG model for the soft underlying event, though, does not show any effect at the parton level because the energy contribution is calculated directly at the hadron level. In the following discussion, in order to examine the differences between hadron, detector and parton level, the underlying event in HERWIG has been switched off.
Figures 5 and 6 show the transverse energy inside the max and min cones at $`\eta =\eta _{LeadJet}`$ and $`\varphi =\varphi _{LeadJet}\pm \frac{\pi }{2}`$ as a function of the leading jet transverse energy at the parton, hadron and detector level. The lead jet is always in the central region. Because of the degradation due to the detector response, the amount of energy is higher at the hadron level than at the detector level.
It is also interesting to note that the hadron level energy is larger than the parton level energy, by the order of several hundred MeV. This is due to hadronization effects of the partons produced in or near the lead and second jet cones. Most of the hadronization effects come from resonance production ($`\rho ,A_1,A_2,\mathrm{}`$) and their subsequent decays. The hadronization effects from the partons inside the jet cone have previously been termed “splashout”. It is important to note that this splashout is not currently taken into account in either the CDF or D0 jet analyses. Both experiments implicitly assume that the hadron and parton levels produce the same energy in the jet cone. This is especially relevant for low $`E_T`$ jet production.
In order to evaluate to what level resonance decays influence the energy inside the two cones at $`90^0`$ from the leading jet, all resonance decays are switched off and the energy in the cones were examined at both the parton and the hadron level. The difference of $`E_T`$ inside the min cone (between hadron and parton level) decreases from an average of 300 MeV to 100 MeV, while the difference in the max cone goes from 500 MeV to 100 MeV.
## 3 Underlying Energy in minimum bias events
The model used in HERWIG to simulate minimum bias events is the same as used for the soft underlying energy in hard scattering events. Minimum bias events were studied in order to see if the reason for the offset observed between the data and simulation results from the HERWIG description of the soft underlying event. The minimum bias events generated with HERWIG were passed through the detector simulation program QFL and the information on the energy released in the calorimeter towers stored.
The amount of transverse energy in the calorimeter, in a random cone of radius 0.7 that is required to be in the central region ($`|\eta |<0.7`$), was determined. The transverse energy distributions for the two different tower thresholds are summarized in Table 1 where a comparison with data also can be found.
The offset of about 650 MeV between data and HERWIG is slightly higher than the one found comparing the min cones in the jet events.
## 4 $`E_T`$ summed in the central region (Swiss Cheese)
For these comparisons, the transverse energy in every calorimeter tower in the central region ($`|\eta |<1`$) is summed, excluding the towers in a radius 0.7 from the center of the two (or three) most energetic jets in the event:
$$SumofE_T=\underset{towers}{}E_T^{towers}\underset{2/3jets}{}\left[\underset{towers}{}E_T^{towers_{jet}}\right]$$
where $`E_T^{towers_{jet}}`$ are all the towers in a radius 0.7 from the center of the jet. We require $`E_{TJet}>5`$ GeV. This configuration has been labelled ‘Swiss cheese’ <sup>2</sup><sup>2</sup>2Or specifically Emmental.
There are an average of between 2 and 2.5 jets in the central rapidity region, with this average having a slight slope as a function of the lead jet transverse energy. The Swiss cheese energy in the central region is plotted in Fig 7 at the hadron, parton and detector level. The approximate minimum bias level for HERWIG and data is shown with a flat line on the picture. In the simple picture presented earlier, and on which the CDF and D0 jet analyses are based, the difference between the Swiss cheese energy with two jets subtracted and the minimum bias level should be proportional to the NLO (third parton) contribution. The Swiss cheese level with three jets subtracted should have little or no NLO contribution and can be directly compared to the minimum bias data level. The 3-jet subtracted Swiss cheese energy is larger than the minimum bias level and there is a small slope as a function of the lead jet $`E_T`$ (the offset varies from 6-8 GeV over the $`E_T`$ range). This indicates perhaps that there is more complexity here than in the simple picture. Other possible contributions to the Swiss cheese energy include hadronization from the jets (“splashout”), double parton scattering and higher order radiation effects.
As was done for the min and max cone studies, the underlying event in HERWIG can be switched off and the hadron/detector level in the Swiss Cheese plots compared when the resonance decay is not allowed. At the hadron level, when the resonance decay is not allowed, we find about 1.5 GeV energy less then when allowing the resonance decay. This implies a 600-700 MeV contribution of splashout per jet (again at the detector level) to the Swiss cheese energy.
The comparison of the Swiss cheese results for the data and HERWIG+QFL is complex and its interpretation is continuing.
## 5 Conclusions and where do we go from here
The energy from the underlying event is not perturbatively calculable and must be subtracted from a jet cone in order for comparisons to be made to NLO calculations. Because of the ambiguous definition of what constitutes this underlying event, a relatively large uncertainty has been assigned to the value of this subtraction. In order to study the underlying event, we have considered two cones in the calorimeter far away from the leading jet and we examined the energy in the cones, both in the CDF data and with the HERWIG simulation. We discovered that both the data and HERWIG exhibited a similar behaviour for the max and the min cone; the min cone stays flat, while the max cone increases as a function of the leading jet $`E_T`$. There is an offset, however, of about 500 MeV for the min and of 800 MeV for the max cone between data and HERWIG. If we examine the difference between the max and min cones, where the underlying event energy contribution should be minimized, we find very similar distributions for data and HERWIG. In minimum bias, the HERWIG model predicts a level of energy substantially below the one found in minimum bias data (400 MeV compared to 1 GeV). Part of this difference is due to the lack of any kind of hard interaction in the minimum bias model.
With HERWIG we investigated max/min cone distributions at the parton, hadron and detector level and we found out that the energy inside the cones is higher at the hadron then at the parton level. This is mainly due to resonance decay.
An improved understanding of the underlying event is desired for a number of reasons:
* The underlying event subtraction is the largest uncertainty for the jet cross section at low transverse energy (below 60 GeV). In order to have a good comparison of the data with theory, a better understanding of the proper level of this subtraction must be obtained. This uncertainty is especially important for the measurement of the jet cross section at 630 GeV, since most of the data points are below 60 GeV, and similar considerations to those at 1800 GeV also apply.
* This analysis probes the interface between perturbative and non-perturbative QCD, an arena where a great deal of work still needs to be done.
* The authors of the Monte Carlo programs are trying to predict the environments for physics measurements at the LHC. This can be difficult/uncertain without the proper understanding of what is happening at the Tevatron.
This analysis will be extended to the jet and minimum bias data taken at 630 GeV by CDF. It will be especially interesting to observe the level of agreement of the HERWIG minimum bias predictions with the CDF data, given that the HERWIG model parameters were determined from the UA5 taken at a similar energy. It may be that there is an increase in the semi-hard component of minimum bias energy when going from 630 to 1800 GeV. After the comparisons at 630 and 1800 GeV are complete, extrapolations will made made to LHC energies for the underlying event in both jet and minimum bias events.
## 6 Acknowledgements
This work was performed in conjunction with our colleagues on CDF, Anwar Bhatti and Eve Kovacs.
Isolated Photon Production
S. Frixione and W. Vogelsang
## 1 Isolated-photon production
### 1.1 General features of photon production at colliders
When mentioning the photon in the framework of high-energy collider physics, one is immediately led to think – with good reasons – to Higgs searches through the gold-plated channel $`H\gamma \gamma `$. However, the production of photons also deserves attention on its own. Firstly, a detailed understanding of the continuum two-photon production is crucial in order to clearly disentangle any Higgs signals from the background. Secondly, in hadronic collisions, where a very large number of strong-interacting particles is produced, photon signals are relatively clean, since the photon directly couples only to quarks. Therefore, prompt-photon data can be used to study the underlying parton dynamics, in a complementary way with respect to analogous studies performed with hadrons or jets. For the same reason, these data represent a very important tool in the determination of the gluon density in the proton, $`g(x)`$. Indeed, in recent years almost all the direct information (that is, not obtained through scaling violations as predicted by Altarelli-Parisi equations) on the intermediate- and high-$`x`$ behaviour of $`g(x)`$ came from prompt-photon production, $`pp\gamma X`$ and $`pN\gamma X`$, in fixed-target experiments. The main reason for this is that, at leading order, a photon in the final state is produced in the reactions $`qg\gamma q`$ and $`q\overline{q}\gamma g`$, with the contribution of the former subprocess being obviously sensitive to the gluon and usually dominant over that of the latter. It is the ‘point-like’ coupling of the photon to the quark in these subprocesses that is responsible for a much cleaner signal than, say, for the inclusive production of a $`\pi ^0`$, which proceeds necessarily through a fragmentation process.
There is, however, a big flaw in the arguments given above. In fact, photons can also be produced through a fragmentation process, in which a parton, scattered or produced in a QCD reaction, fragments into a photon plus a number of hadrons. The problem with the fragmentation component in the prompt-photon reaction is twofold: first, it introduces in the cross section a dependence upon non-perturbative fragmentation functions, similar to those relevant in the case of single-hadron production, which are not calculable in perturbative QCD and are, at present, very poorly determined by the sparse LEP data available. Secondly, all QCD partonic reactions contribute to the fragmentation component; thus, when addressing the problem of the determination of the gluon density, the advantage of having a priori only one partonic reaction ($`q\overline{q}\gamma g`$) competing with the signal ($`qg\gamma q`$) is lost, even though some of the subprocesses relevant to the fragmentation part at the same time result from a gluon in the initial state.
The relative contribution of the fragmentation component with respect to the direct component (where the photon participates in the short-distance, hard-scattering process) is larger the larger the center-of-mass energy and the smaller the final-state transverse momentum <sup>1</sup><sup>1</sup>1Actually, in the fixed-target $`pp\gamma X`$ reaction, one can see the fragmentation component increasing relatively to the direct one also at very large $`p_{T\gamma }`$, because of the direct cross section dying out very quickly at such momenta. This effect is of no phenomenological relevance at the LHC.: at the LHC, for transverse momenta of the order of few tens of GeV, it can become dominant. However, here the situation is saved by the so-called ‘isolation’ cut, which is imposed on the photon signal in experiments. Isolation is an experimental necessity: in a hadronic environment the study of photons in the final state is complicated by the abundance of $`\pi ^0`$’s, eventually decaying into pairs of $`\gamma `$’s. The isolation cut simply serves to improve the signal-to-noise ratio: if a given neighbourhood of the photon is free of energetic hadron tracks, the event is kept; it is rejected otherwise. Fortunately, by requiring the photon to be isolated, one also severely reduces the contribution of the fragmentation part to the cross section. This is because fragmentation is an essentially collinear process: therefore, photons resulting from parton fragmentation are usually accompanied by hadrons, and are therefore bound to be rejected after the imposition of an isolation cut.
Thus, the fragmentation contribution, that threatened to spoil the cleanliness of the photon signals at colliders, is relatively well under control in the case of isolated-photon cross sections. There is of course a price to pay for this gain: the isolation condition poses additional problems in the theoretical computations, which are not present in the case of fully-inclusive photon cross sections. This topic will be the argument of the next subsection.
### 1.2 Isolation prescriptions
Consistently with what written above, we write the cross section for the production of an isolated-photon in hadronic collisions as follows:
$`d\sigma _{AB}(K_A,K_B;K_\gamma )=`$
$`{\displaystyle 𝑑x_1𝑑x_2f_a^{(A)}(x_1,\mu _F)f_b^{(B)}(x_2,\mu _F)𝑑\widehat{\sigma }_{ab,\gamma }^{isol}(x_1K_A,x_2K_B;K_\gamma ;\mu _R,\mu _F,\mu _\gamma )}`$
$`+{\displaystyle 𝑑x_1𝑑x_2𝑑zf_a^{(A)}(x_1,\mu _F^{})f_b^{(B)}(x_2,\mu _F^{})𝑑\widehat{\sigma }_{ab,c}^{isol}(x_1K_A,x_2K_B;K_\gamma /z;\mu _R^{},\mu _F^{},\mu _\gamma )D_\gamma ^{(c)}(z,\mu _\gamma )},`$ (1.1)
where $`A`$ and $`B`$ are the incoming hadrons, with momenta $`K_A`$ and $`K_B`$ respectively, and a sum over the parton indices $`a`$, $`b`$ and $`c`$ is understood. In the first term on the RHS of eq. (1.1) (the direct component) the subtracted partonic cross sections $`d\widehat{\sigma }_{ab,\gamma }^{isol}`$ get contributions from all the diagrams with a photon leg. On the other hand, the subtracted partonic cross sections $`d\widehat{\sigma }_{ab,c}^{isol}`$ appearing in the second term on the RHS of eq. (1.1) (the fragmentation component), get contribution from the pure QCD diagrams, with one of the partons eventually fragmenting in a photon, in a way described by the parton-to-photon fragmentation function $`D_\gamma ^{(c)}`$. As the notation in eq. (1.1) indicates, the isolation condition is embedded into the partonic cross sections.
It is a well-known fact that, in perturbative QCD beyond leading order, and for all the isolation prescriptions known at present, with the exception of that of ref. , neither the direct nor the fragmentation components are separately well defined at any fixed order in perturbation theory: only their sum is physically meaningful. In fact, the direct component is affected by quark-to-photon collinear divergences, which are subtracted by the bare fragmentation function that appears in the unsubtracted fragmentation component. Of course, this subtraction is arbitrary as far as finite terms are concerned. This is formally expressed in eq. (1.1) by the presence of the same scale $`\mu _\gamma `$ in both the direct and fragmentation components: a finite piece may be either included in the former or in the latter, without affecting the physical predictions. The need for introducing a fragmentation contribution is physically better motivated from the fact that a QCD hard scattering process may produce, again through a fragmentation process, a $`\rho `$ meson that has the same quantum numbers as the photon and can thus convert into a photon, leading to the same signal.
As far as the isolation prescriptions are concerned, here we will restrict to those belonging to the class that can be denoted as ‘cone isolations’ . In the framework of hadronic collisions, where the need for invariance under longitudinal boosts suggests not to define physical quantities in terms of angles, the cone is drawn in the pseudorapidity–azimuthal angle plane, and corresponds to the set of points
$$𝒞_R=\{(\eta ,\varphi )\sqrt{(\eta \eta _\gamma )^2+(\varphi \varphi _\gamma )^2}R\},$$
(1.2)
where $`\eta _\gamma `$ and $`\varphi _\gamma `$ are the pseudorapidity and azimuthal angle of the photon, respectively, and $`R`$ is the aperture (or half-angle) of the cone. After having drawn the cone, one has to actually impose the isolation condition. We consider here two sub-classes of cone isolation, whose difference lies mainly in the behaviour of the fragmentation component. Prior to that, we need to define the total amount of hadronic transverse energy deposited in a cone of half-angle $`R`$ as
$$E_{T,had}(R)=\underset{i=1}{\overset{n}{}}E_{Ti}\theta (RR_{\gamma i}),$$
(1.3)
where
$$R_{\gamma i}=\sqrt{(\eta _i\eta _\gamma )^2+(\varphi _i\varphi _\gamma )^2},$$
(1.4)
and the sum runs over all the hadrons in the event (or, alternatively, $`i`$ can be interpreted as an index running over the towers of a hadronic calorimeter). For both the isolation prescriptions we are going to define below, the first step is to draw a cone of fixed half-angle $`R_0`$ around the photon axis, as given in eq. (1.2). We will denote this cone as the isolation cone.
The photon is isolated if the total amount of hadronic transverse energy in the isolation cone fulfils the following condition:
$$E_{T,had}(R_0)ϵ_cp_{T\gamma },$$
(1.5)
where $`ϵ_c`$ is a small number, and $`p_{T\gamma }`$ is the transverse momentum of the photon.
The photon is isolated if the following inequality is satisfied:
$$E_{T,had}(R)ϵ_\gamma p_{T\gamma }𝒴(R),$$
(1.6)
for all the cones lying inside the isolation cone, that is for $`RR_0`$. The function $`𝒴`$ is arbitrary to a large extent, but must at least have the following property:
$$\underset{R0}{lim}𝒴(R)=0,$$
(1.7)
and being different from zero everywhere except for $`R=0`$.
Definition A was proven to lead to an infrared-safe cross section at all orders of perturbation theory in ref. . The smaller $`ϵ_c`$, the tighter the isolation. Loosely speaking, for vanishing $`ϵ_c`$ the direct component behaves like $`\mathrm{log}ϵ_c`$, while the fragmentation component behaves like $`ϵ_c\mathrm{log}ϵ_c`$. Thus, for $`ϵ_c0`$ eq. (1.1) diverges. This is obvious since the limit $`ϵ_c0`$ corresponds to a fully-isolated-photon cross section, which cannot be a meaningful quantity, whether experimentally (because of limited energy resolution) or theoretically (because there is no possibility for soft particles to be emitted into the cone).
Definition B was proposed and proven to lead to an infrared-safe cross section at all orders of perturbation theory in ref. . Eq. (1.7) implies that the energy of a parton falling into the isolation cone $`𝒞_{R_0}`$ is correlated to its distance (in the $`\eta `$$`\varphi `$ plane) from the photon. In particular, a parton becoming collinear to the photon is also becoming soft. When a quark is collinear to the photon, there is a collinear divergence; however, if the quark is also soft, this divergence is damped by the quark vanishing energy. When a gluon is collinear to the photon, then either it is emitted from a quark, which is itself collinear to the photon – in which case, what was said previously applies – or the matrix element is finite. Finally, it is clear that the isolation condition given above does not destroy the cancellation of soft singularities, since a gluon with small enough energy can be emitted anywhere inside the isolation cone. The fact that this prescription is free of final-state QED collinear singularities implies that the direct part of the cross section is finite. As far as the fragmentation contribution is concerned, in QCD the fragmentation mechanism is purely collinear. Therefore, by imposing eq. (1.6), one forces the hadronic remnants collinear to the photon to have zero energy. This is equivalent to saying that the fragmentation variable $`z`$ is restricted to the range $`z=1`$. Since the parton-to-photon fragmentation functions do not contain any $`\delta (1z)`$, this means that the fragmentation contribution to the cross section is zero, because an integration over a zero-measure set is carried out. Therefore, only the first term on the RHS of eq. (1.1) is different from zero, and it does not contain any $`\mu _\gamma `$ dependence.
We stress again that the function $`𝒴`$ can be rather freely defined. Any sufficiently well-behaved function, fulfilling eq. (1.7), could do the job, the key point being the correlation between the distance of a parton from the photon and the parton energy, which must be strong enough to cancel the quark-to-photon collinear singularity. Throughout this paper, we will use
$$𝒴(R)=\left(\frac{1\mathrm{cos}R}{1\mathrm{cos}R_0}\right)^n,n=1.$$
(1.8)
We also remark that the traditional cone-isolation prescription, eq. (1.5), can be recovered from eq. (1.6) by setting $`𝒴=1`$ and $`ϵ_\gamma =ϵ_c`$.
### 1.3 Isolated photons at the LHC
In this section, we will present results for isolated-photon cross sections in $`pp`$ collisions at 14 TeV. These results have been obtained with the fully-exclusive NLO code of ref. , and are relevant to the isolation obtained with definition B; the actual parameters used in the computation are given in eq. (1.8), together with $`ϵ_\gamma =1`$. We let $`R_0=0.4`$. We will comment in the following on the outcome of definition A.
Any sensible perturbative computation should address the issue of the perturbative stability of its results. A rigorous estimate of the error affecting a cross section at a given order can be given if the next order result is also available. If this is not the case, it is customary to study the dependence of the physical observables upon the renormalization ($`\mu _R`$) and factorization ($`\mu _F`$) scales. It is important to stress that the resulting spread should not be taken as the ‘theoretical error’ affecting the cross section; to understand this, it is enough to say that the range in which $`\mu _R`$ and $`\mu _F`$ are varied is arbitrary. Rather, one should compare the spread obtained at the various perturbative orders; only if the scale dependence decreases when including higher orders the cross section can be regarded as perturbatively stable and sensibly compared to data.
Usually, $`\mu _R`$ and $`\mu _F`$ are imposed to have the same value, $`\mu `$, which is eventually varied. However, this procedure might hide some problems, because of a possible cancellation between the effects induced by the two scales. It is therefore desirable to vary $`\mu _R`$ and $`\mu _F`$ independently. Here, an additional problem arises at the NLO. The expression of any cross section in terms of $`\mu `$ (that is, when $`\mu _R=\mu _F`$) is not ambiguous, while it is ambiguous if $`\mu _R\mu _F`$. In fact, when $`\mu _R\mu _F`$, the cross section can be written as the sum of a term corresponding to the contribution relevant to the case $`\mu _R=\mu _F`$, plus a term of the kind:
$$\alpha _\mathrm{S}(\mu _A)(\alpha _\mathrm{S}(\mu _R))\mathrm{log}\frac{\mu _R}{\mu _F},$$
(1.9)
where $``$ has the same power of $`\alpha _\mathrm{S}`$ as the LO contribution, say $`\alpha _\mathrm{S}^k`$. The argument of the $`\alpha _\mathrm{S}`$ in front of eq. (1.9), $`\mu _A`$, can be chosen either equal to $`\mu _R`$ or equal to $`\mu _F`$, since the difference between these two choices is of NNLO. Thus, it follows that the dependence upon $`\mu _R`$ or $`\mu _F`$ of a NLO cross section reflects the arbitrariness of the choice made in eq. (1.9), which is negligible only if the NNLO ($`\alpha _\mathrm{S}^{k+2}`$) corrections are much smaller than the NLO ones ($`\alpha _\mathrm{S}^{k+1}`$). This leads to the conclusion that a study of the dependence upon $`\mu _R`$ or $`\mu _F`$ only can be misleading. In other words: $``$ in eq. (1.9) is determined through RG equations in order to cancel the scale dependence of the cross section up to terms of order $`\alpha _\mathrm{S}^{k+2}`$. This happens regardless of the choice made for $`\mu _A`$ in eq. (1.9). However, here we are not discussing the cancellation to a given perturbative order of the effects due to scale variations; we are concerned about the coefficient in front of the $`𝒪(\alpha _\mathrm{S}^{k+2})`$ term induced by such variations, whose size is dependent upon the choice made for $`\mu _A`$ and therefore, to some extent, arbitrary. We have to live with this arbitrariness, if we decide to vary $`\mu _R`$ or $`\mu _F`$ only. However, we can still vary $`\mu _R`$ and $`\mu _F`$ independently, but eventually putting together the results in some sensible way, that reduces the impact of the choice made for $`\mu _A`$. In this section, we will consider the quantities defined as follows:
$`\left({\displaystyle \frac{\delta \sigma }{\sigma }}\right)_\pm `$ $`=`$ $`\pm \{\left[{\displaystyle \frac{\sigma (\mu _R=\mu _0,\mu _F=\mu _0)\sigma (\mu _R=a_\pm \mu _0,\mu _F=\mu _0)}{\sigma (\mu _R=\mu _0,\mu _F=\mu _0)+\sigma (\mu _R=a_\pm \mu _0,\mu _F=\mu _0)}}\right]^2`$ (1.10)
$`+\left[{\displaystyle \frac{\sigma (\mu _R=\mu _0,\mu _F=\mu _0)\sigma (\mu _R=\mu _0,\mu _F=a_\pm \mu _0)}{\sigma (\mu _R=\mu _0,\mu _F=\mu _0)+\sigma (\mu _R=\mu _0,\mu _F=a_\pm \mu _0)}}\right]^2\}^{\frac{1}{2}},`$
where $`a_+`$ and $`a_{}=1/a_+`$ are two numbers of order one, which we will take equal to 1/2 and 2 respectively; the $`\pm `$ sign in front of the RHS of eq. (1.10) is purely conventional. We can evaluate $`(\delta \sigma /\sigma )_\pm `$ by using $`\mu _A=\mu _R`$ or $`\mu _A=\mu _F`$ in eq. (1.9). The reader can convince himself, with the help of the definition of the QCD $`\beta `$ function, that the difference between these two choices is of order $`\alpha _\mathrm{S}^4`$ in the expansion of the contribution to $`(\delta \sigma /\sigma )_\pm ^2`$ due to eq. (1.9); on the other hand, this difference is only of order $`\alpha _\mathrm{S}^3`$ in each of the two terms under the square root in the RHS of eq. (1.10). This is exactly what we wanted to achieve: a suitable combination of the cross sections resulting from independent $`\mu _R`$ and $`\mu _F`$ variations is less sensitive to the choice for $`\mu _A`$ made in eq. (1.9) with respect to the results obtained by varying $`\mu _R`$ or $`\mu _F`$ only.
In table 1 we present the results for the total isolated-photon rates, both at NLO and at LO. The latter cross sections have been obtained by retaining only the lowest order terms ($`𝒪(\alpha _{\mathrm{em}}\alpha _\mathrm{S})`$) in the short-distance cross section, and convoluting them with NLO-evolved parton densities. Also, a two-loop expression for $`\alpha _\mathrm{S}`$ has been used. There is of course a lot of freedom in the definition of a Born-level result. However, we believe that with this definition one has a better understanding of some issues related to the stability of the perturbative series. In order to obtain the rates entering table 1, we required the photon transverse momentum to be in the range $`40<p_{T\gamma }<400`$ GeV, and we considered the rapidity cuts $`\left|\eta _\gamma \right|<1.5`$ and $`\left|\eta _\gamma \right|<2.5`$, in order to simulate a realistic geometrical acceptance of the LHC detectors. We first consider the scale dependence of our results (last column), evaluated according to eq. (1.10). We see that the NLO results are clearly more stable than the LO ones; this is reassuring, and implies the possibility of a sensible comparison between NLO predictions and the data. Notice that the size of the radiative corrections ($`K`$ factor, defined as the ratio of the NLO result over the LO result) is quite large. From the table, we see that the cross sections obtained with different parton densities differ by 6% at the most (relative to the result obtained with MRST99-1 , which we take as the default set). MRST99 sets 2 and 3 are meant to give an estimate of the effects due to the current uncertainties affecting the gluon density, whereas sets 4 and 5 allow to study the sensitivity of our predictions to the value of $`\alpha _\mathrm{S}(M_Z)`$ (sets 1, 4 and 5 have $`\mathrm{\Lambda }_5^{\overline{\mathrm{MS}}}=`$220, 164 and 288 MeV respectively). On the other hand, the difference between MRST99-1 and CTEQ5M results is due to the inherent difference between these two density sets (CTEQ5M has $`\mathrm{\Lambda }_5^{\overline{\mathrm{MS}}}=`$226 MeV, and therefore the difference in the values of $`\alpha _\mathrm{S}(M_Z)`$ plays only a very minor role).
From inspection of table 1, we can conclude that isolated-photon cross section at the LHC is under control, both in the sense of perturbation theory and of the dependence upon non-calculable inputs, like $`\alpha _\mathrm{S}(M_Z)`$ and parton densities. The relatively weak dependence upon the parton densities, however, is not a good piece of news if one aims at using photon data to directly access the gluon density. On the other hand, the expected statistics is large enough to justify attempts of a direct measurement of such a quantity. In the remainder of this section, we will concentrate on this issue. We will consider
$$_x=\frac{d\sigma _0/dxd\sigma /dx}{d\sigma _0/dx+d\sigma /dx},$$
(1.11)
where $`x`$ is any observable constructed with the kinematical variables of the photon and, possibly, of the accompanying jets. $`\sigma `$ and $`\sigma _0`$ are the cross sections obtained with two different sets of parton densities, the latter of which is always the default one (MRST99-1). We can imagine a gedanken experiment, where it is possible to change at will the parton densities; in this way, we can assume the relative statistical errors affecting $`\sigma `$ and $`\sigma _0`$ to decrease as $`1/\sqrt{N}`$ and $`1/\sqrt{N_0}`$, $`N`$ and $`N_0`$ being the corresponding number of events. It is then straightforward to calculate the statistical error affecting $`_x`$; by imposing $`_x`$ to be larger than its statistical error, one gets
$$_x>\left(_x\right)_{min}\frac{1}{\sqrt{2ϵ\sigma (x,\mathrm{\Delta }x)}},$$
(1.12)
where $``$ is the integrated luminosity, $`ϵ1`$ collects all the experimental efficiencies, and
$$\sigma (x,\mathrm{\Delta }x)=_{x\mathrm{\Delta }x/2}^{x+\mathrm{\Delta }x/2}𝑑x\frac{d\sigma }{dx}$$
(1.13)
is the total cross section in a range of width $`\mathrm{\Delta }x`$ around $`x`$.
In fig. 1 we present our predictions for $`_x`$. In the left panel of the figure we have chosen $`x=p_{T\gamma }`$, while in the right panel we have $`x=x_{\gamma j}`$, where
$$x_{\gamma j}=\frac{p_{T\gamma }\mathrm{exp}(\eta _\gamma )+p_{Tj}\mathrm{exp}(\eta _j)}{\sqrt{S}}.$$
(1.14)
In this equation $`\sqrt{S}`$ is the center-of-mass energy of the colliding hadrons, and $`p_{Tj}`$ and $`\eta _j`$ are the transverse momentum and rapidity of the hardest jet recoiling against the photon. In order to reconstruct the jets, we adopted here a $`k_T`$-algorithm, namely that proposed in ref. , with $`D=1`$. Notice that $`x_{\gamma j}`$ exactly coincides at the leading order with the Bjorken-$`x`$ of the partons in one of the incoming hadrons; NLO corrections introduce only minor deviations.
For all the density sets considered, the dependence of $``$ upon $`p_{T\gamma }`$ is rather mild. The values in the low-$`p_{T\gamma }`$ region could also be inferred from table 1, since the cross section is dominated by small $`p_{T\gamma }`$’s. Analogously to what happens in the case of total rates, the sets MRST99-4 and MRST99-5 give rise to extreme results for $`_{p_{T\gamma }}`$, since the value of $`\mathrm{\Lambda }_{QCD}`$ is quite different from that of the default set. From the figure, it is apparent that, by studying the transverse momentum spectrum, it will not be easy to distinguish among the possible shapes of the gluon density. On the other hand, it seems that, as far as the statistics is concerned, a distinction between any two sets can be performed. Indeed, the symbols in the figure display the quantity defined in eq. (1.12), for $`=100`$ fb<sup>-1</sup>, $`\mathrm{\Delta }p_{T\gamma }=10`$ GeV and $`ϵ=1`$. Of course, the latter value is not realistic. However, a smaller value (leading to a larger $`()_{min}`$), can easily be compensated by enlarging $`\mathrm{\Delta }p_{T\gamma }`$ and by the fact that the total integrated luminosity is expected to be much larger than that adopted in fig. 1.
Turning to the right panel of fig. 1, we can see a much more interesting situation. Actually, it can be shown that the pattern displayed in the figure is rather faithfully reproduced by plotting the analogous quantity, where one uses the gluon densities instead of the cross sections. This does not come as a surprise. First, $`x_{\gamma j}`$ is in an almost one-to-one correspondence with the $`x`$ entering the densities. Secondly, photon production is dominated by the gluon-quark channel, and therefore the cross section has a linear dependence upon $`g(x)`$, which can be easily spotted. It does seem, therefore, to be rather advantageous to look at more exclusive variables, like photon-jet correlations (this is especially true if one considers the procedure of unfolding the gluon density from the data: in the case of single-inclusive variables, the unfolding requires a de-convolution, which is not needed in the case of correlations). Of course, there is a price to pay: the efficiency $`ϵ`$ will be smaller in the case of photon-jet correlations, with respect to the case of single-inclusive photon observables, mainly because of the jet-tagging. However, from the figure it appears that there should be no problem with statistics, except in the very large $`x_{\gamma j}`$ region.
Finally, we would like to comment on the fact that, for the case of single-inclusive photon observables, we also computed the cross section by isolating the photon according to definition A, using $`ϵ_c=2`$ GeV$`/p_{T\gamma }`$. The two definitions return a $`p_{T\gamma }`$ spectrum almost identical in shape, with definition B higher by a factor of about 9%. It is only at the smallest $`p_{T\gamma }`$ values that we considered, that definition B returns a slightly steeper spectrum. The fact that such different definitions produce very similar cross sections may be surprising. This happens because, prior to applying the isolation condition, partons tend to be radiated close to the photon; therefore, most of them are rejected when applying the isolation, no matter of which type. This situation has already been encountered in the production of photons at much smaller energies. The reader can find a detailed discussion on this point in ref. .
Direct photon pair production at colliders,
an irreducible background to Higgs boson searches at the LHC
T. Binoth, J.P. Guillet, V.A. Ilyin, E. Pilon and M. Werlen
## 1 Role and relevance of higher order corrections
Our theoretical understanding of direct photon pair production (as any hard hadronic process, cf. ) is based on the QCD improved parton model, according to which long and short distance effects factorize from each other. Short distance subprocesses are safely computed in perturbative QCD. Long distance effects cannot be completely calculated from QCD at present, although their scaling violations can be. Instead, they are extracted from experimental data and encoded into non-perturbative quantities, such as the parton distribution functions in incoming hadrons and, if necessary, inclusive fragmentation functions of partons into observed outgoing particles, e.g. photons. Yet these quantities are universal, i.e. independent of the hard subprocess: schematically they can be measured in one process, then transported to predict another one.
However, the border between short and long distance scales is arbitrary. The separation requires the introduction of unphysical parameters; e.g. the factorization scale $`M^2`$, and similarly the fragmentation scale $`M_f^2`$ in the case of a fragmentation process. In an ideal exact calculation, the dependence on these spurious parameters (as well as on the arbitrary renormalization scale $`\mu ^2`$) would cancel between the short and long distance parts. In an actual expansion in powers of $`\alpha _s`$ truncated at some finite order, this cancellation is only partial; it holds up to a term of the lowest uncalculated order. As the order of truncation increases, theoretical estimates become flatter and flatter over broader and broader ranges of these spurious scales. The uncertainty induced by this actual dependence restricts the accuracy and predictive character of QCD calculations. In particular, the result of a lowest order calculation is plagued by a large monotonic dependence with respect to $`M^2`$, $`M_f^2`$ and $`\mu ^2`$. It changes by a large factor (two or more) when these scales are varied around the typical hard scale of the process; such a lowest order estimate is not at all quantitative. This is the first reason why any tentatively accurate QCD calculation has to be carried out to at least next-to-leading order (abbreviated below as NLO).
Another important motivation is that higher order corrections to some given process may reveal new mechanisms, whose rates may be not necessarily negligible compared to the leading order contribution. The production of photons is a typical example of this phenomenon, as will be explained below. This amounts to large higher order corrections, which affect substantially both the magnitude and the shape of the distributions, not due to poor apparent convergence of the first terms in the perturbative expansion, but for physically understood reasons.
Finally, finite order calculations may not be accurate enough, as in the case of infrared sensitive observables, i.e. observables controlled by multiple soft gluon emission <sup>2</sup><sup>2</sup>2See for example . Yet, in some less well-known cases of infrared sensitivity, they may reveal perturbative instabilities or even divergences plaguing the calculation at any further order inside the physical spectrum . This is, for example, the case for the transverse momentum distribution of pairs of isolated photons, as will be discussed below. The calculation of higher order corrections is therefore the first step towards a deeper understanding of what happens in such cases.
## 2 Mechanisms of production.
Schematically, three possible mechanisms may produce prompt photon pairs with large invariant mass: one (which may be called “two direct”) in which both photons take part directly in the hard subprocess, another one (“one fragmentation”) in which one of the photons undergoes the hard subprocess while the other results from the fragmentation of a hard parton (quark or gluon), itself produced at large tranverse momentum, and yet another mechanism (“two fragmentation”) in which both photons result from such a fragmentation. This schematical splitting into these three contributions emerges from a factorization procedure sketched in what follows. Although this splitting provides a convenient picture, one must however keep in mind that it is arbitrary; none of these contributions can be measured separately. Only their sum is physical.
### 2.1 Direct vs. fragmentation mechanisms
From a topological point of view, a photon produced from fragmentation is with a high probability accompanied by a jet of hadrons. From a technical point of view, the lowest order of the “one fragmentation” contribution emerges in the calculation of higher order perturbative corrections to the “two direct” contribution given by the Born process $`q\overline{q}\gamma \gamma `$. Some of these higher order corrections, such as $`qgq\gamma \gamma `$, are plagued by final state collinear singularities associated with the collinear splitting $`qq\gamma `$. The latter have to be factorized and absorbed into quark and gluon fragmentation functions to a photon, $`D_{\gamma /qorg}(z,M_f^2)`$ defined at some fragmentation scale <sup>3</sup><sup>3</sup>3and in some given factorization scheme. Here we use the $`\overline{MS}`$ scheme, . $`M_f^2`$. Analogously to the so-called anomalous component of the photon structure function, a collinear logarithmic enhancement occurs, induced by the pointlike quark-photon coupling. To all orders in $`\alpha _s`$, this phenomenon results in $`D_{\gamma /qorg}(z,M_f^2)`$ behaving asymptotically <sup>4</sup><sup>4</sup>4i.e. when the fragmentation scale $`M_f^2`$ (chosen of the order of the hard scale of the subprocess) is large compared to any typical hadronic scale $`1`$ GeV<sup>2</sup>. as $`\alpha /\alpha _s(M_f^2)`$. This compensates the one extra power of $`\alpha _s`$ involved in the short distance subprocess, so that fragmentation contributions are asymptotically of the same order as the Born term, by power counting in $`\alpha _s`$. What is more, given the high gluon density at LHC, the $`gq`$ (or $`\overline{q}`$) initiated process involving one photon from fragmentation even dominates the inclusive production rate in the range $`80140`$ GeV.
In turn, higher order corrections to “one fragmentation” reveal the “two fragmentation” mechanism. Similarly, a collinear enhancement associated with each photon fragmentation compensates two extra powers of $`\alpha _s`$ in the short distance subprocess, so that the power counting in $`\alpha _s`$ is here also asymptotically the same as for the Born and “one fragmentation” parts. Higher order corrections to both fragmentation contributions have in principle to be computed in order to provide a consistent NLO study. This has been done in . The actual quantitative significance of these contributions is discussed in sect. 4.
### 2.2 The box contribution
Beyond this, the gluon-gluon fusion contribution $`gg\gamma \gamma `$, of the “two direct” type, cannot be neglected. Indeed, although it is an $`𝒪(\alpha ^2\alpha _s^2)`$ i.e. next-to-next-to-leading order contribution, the suppression due to higher powers of $`\alpha _s`$ is compensated by the large gluon luminosity at colliders. This is especially true at LHC in the relevant range for Higgs search, where the so-called box contribution has the same magnitude as the Born term, roughly 50 to 80 %.
Moreover it is the lowest order of a new mechanism, whose spurious scale dependences are thus monotonic, and only higher corrections to it would reduce the sensitivity with respect to spurious scale dependences. Finally, this lowest order is a $`22`$ process which yields only back to back photons in the direction transverse to the beam axis; it gives no contribution to the tail of the transverse momentum distribution of photon pairs. An evaluation of the distortion of the transverse momentum distribution of photon pairs due to the process of gluon-gluon fusion requires the computation of at least the next order correction . The sum “two direct + box” will be refered to as the “direct” contribution.
## 3 Isolation
At TeV colliders, the fragmentation contributions are far from negligible. In particular, the “one fragmentation” component dominates the inclusive production of photon pairs in the lower range of the invariant mass spectrum; in the range of interest for Higgs search at LHC it happens to be 2 to 5 times larger than the “two direct” contribution, depending on choice of scales <sup>5</sup><sup>5</sup>5As already mentioned in the beginning of sect. 2, this statement is strongly fragmentation scale dependent. . A NLO evaluation of the fragmentation contribution is thus necessary to have a tentatively reliable prediction.
Actually, collider experiments do not measure inclusive photon pairs. Isolation cuts <sup>6</sup><sup>6</sup>6More about isolation issues in processes involving direct photons is discussed in . are imposed experimentally to drastically reduce the gigantic background coming from decays of $`\pi ^0`$ and $`\eta `$ mesons. Schematically, a candidate-photon is considered isolated if, in some given cone in azimuthal angle and rapidity about the photon defined by
$$(\varphi \varphi _\gamma )^2+(yy_\gamma )^2R^2,$$
(3.1)
the deposited hadronic transverse energy $`E_T^{had}`$ is smaller than some maximal amount $`E_T^{max}`$,
$$E_T^{had}E_T^{max}$$
(3.2)
$`R`$ and $`E_T^{max}`$ being fixed by the experiments. Such isolation cuts affect also the production rate of direct photon pairs, especially the “single-” and “two fragmentation” contributions, whose topologies are similar to the one of the background. Those are severely reduced when $`E_T^{max}`$ is chosen to be very small compared to the transverse momenta of the photons.
Yet a NLO evaluation of fragmentation contributions is still relevant for various reasons. First, the actual isolation cuts used by collider experiments may be quite more complicated than the schematical criterion given by eqns. (3.1,3.2). Higher order partonic calculations are not designed to account for such criteria accurately, contrary to Monte-Carlo event generators such as PYTHIA or HERWIG . Since these Monte-Carlos and NLO partonic calculations are based on different QCD approximations, it is worthwhile to compare these two approaches whenever possible, as for the inclusive production rate, as well as with rather simple isolation cuts such as the one of eqns. (3.1,3.2). Secondly, the above cone-type isolation criterion induces infrared sensitivity inside the physical spectrum for observables such as the $`q_T`$ spectrum of photon pairs. This effect appears at the NLO - and every higher order - in the “one fragmentation” component, as will be shown in the next section.
## 4 Phenomenology
In earlier works on di-photon production , only the “two direct” contribution was calculated at NLO, while the fragmentation contribution included only the lowest order “one fragmentation” part <sup>7</sup><sup>7</sup>7The “box” contribution was included too.. Moreover, these works were not implemented in a form suited to compute observables such as the invariant mass relevant for Higgs search, nor flexible enough to accomodate experimental selection cuts. A further refinement has implemented the same approximation in a more flexible approach combining analytical and Monte-Carlo integration techniques, thus allowing the computation of several observables within the same calculation, and the possibility to account for selection/isolation cuts. Two recent developments are presented in this workshop (see also C. Balazs’ contribution).
We have implemented the “two direct”, “single-” and “double fragmentation” contributions at NLO accuracy, together with the box gluon-gluon contribution, into a general purpose computer program of “partonic event generator” type (DIPHOX) presented in detail in . The results which we present here are derived from this analysis.
### 4.1 Comparison with Tevatron data
The NLO results agree with the preliminary D0 data reasonably, as seen in Figs. (1-4), except for the tails of each photon’s transverse momentum distribution $`d\sigma /dE_T`$ and the invariant mass distribution of pairs $`d\sigma /dM_{\gamma \gamma }`$ at large $`E_T`$ and $`M_{\gamma \gamma }`$ respectively, where the three highest data points are affected by correlated systematic uncertainties due to background evaluation in both cases. However, more instructive conclusions will be drawn after a finalized understanding of the systematics, and even more so after the Tevatron Run II with the statistics improved by a factor of 20.
### 4.2 Estimates for LHC
We now give some theoretical estimates in the domain relevant for Higgs search at the LHC, for the invariant mass distribution cf. Fig. (5) with and without isolation. One has to keep in mind that the theoretical uncertainties are still large. Firstly these results are still plagued by rather large scale uncertainties, as discussed below. Secondly, for a given scale choice, they may still underestimate the actual background to Higgs search.
### 4.3 Critical examination of various theoretical issues
#### 4.3.1 Scale uncertainties
As mentioned above all results depend on three unphysical scales. Varying these between $`M_{\gamma \gamma }^2/4`$ and $`4M_{\gamma \gamma }^2`$ along the first diagonal $`\mu ^2=M^2=M_f^2`$, the NLO results for the invariant mass distribution appear surprisingly stable, since they change by about 5% only. Alternatively, anti-diagonal variations of $`\mu ^2`$ and $`M^2=M_f^2`$ in the same interval about the central value $`M_{\gamma \gamma }^2`$ lead to a variation still rather large (up to 20 %). This is because variations with respect to $`\mu ^2`$ and $`M^2`$ act in opposite ways. When $`\mu ^2`$ is increased, $`\alpha _s(\mu ^2)`$ and hence the NLO corrections decrease; on the other hand the relevant values of the momentum fraction of incoming partons are small, $`𝒪(10^3`$ to $`10^2)`$, so that the gluon and sea quark distribution functions increase when $`M^2`$ is increased. Scale changes with respect to $`\mu ^2`$ and $`M^2`$ thus nearly cancel again each other along the first diagonal but add up in the other case. Actually, the stability along the first diagonal is accidental at this order <sup>8</sup><sup>8</sup>8In processes for which the lowest order involves some power of $`\alpha _s`$, an explicit $`\mu ^2`$ dependence appears in next-to-leading order correction, which partially compensates the $`\mu ^2`$ dependence in $`\alpha _s(\mu ^2)`$. Unlike this, in the two direct component which dominates the cross section when a drastic isolation is required, the lowest order involves no $`\alpha _s`$. The explicit $`\mu ^2`$ dependence would thus appear only at $`𝒪(\alpha _s^2)`$, i.e. at next-to-next-to leading order. At next to leading order, the $`\mu ^2`$ dependence occurs only through the monotonous decrease of the $`\alpha _s(\mu ^2)`$ weighting the first correction: there is no partial cancellation of $`\mu ^2`$ dependence. The mechanism is more complicated in the fragmentation components, and the situation becomes mixed up between all components when the severity of isolation is reduced.. These observations hold separately for the box contribution.
In conclusion, the $`\mu ^2`$, $`M^2`$ dependences are thus not completely under control yet. On the other hand, accounting for the NLO corrections to the fragmentation components provides stability with respect to $`M_f^2`$ variations about orthodox choices of the fragmentation scale.
#### 4.3.2 Quantitative importance of fragmentation contributions
For orthodox choices of the fragmentation scale, $`M_f^2`$ of order $`M_{\gamma \gamma }^2`$, the “single fragmentation” contribution is small at Tevatron, given the stringent isolation cuts used, and the “two fragmentation” one is even smaller. For example, as can be seen in Figs. (1,2) the “one fragmentation” contribution is about one order of magnitude less than the “two direct” one. It may still have a small visible effect, as in the tail of the azimutal angle distribution $`d\sigma /d\varphi _{\gamma \gamma }`$ ($`\varphi _{\gamma \gamma }`$ being the azimutal angle between the two photons of a pair) in the low $`\varphi _{\gamma \gamma }`$ range, cf. Fig. (4).
The situation is the same for LHC predictions, see Fig. (6). Contributions from fragmentation are drastically reduced when very stringent cuts are imposed, e.g. $`E_T^{max}=2.5`$ GeV in $`R=0.4`$. However, in practice such isolation cuts will be nearly saturated by underlying events: their veto on the hard event itself is thus even more severe, allowing almost no transverse energy leakage from the hard process inside the cone. This may be experimentally most suitable. However, requiring that no transverse energy be deposited in a cone of fixed size about a photon is not infrared safe order by order in perturbation. It means that finite but very stringent isolation cuts imposed in fixed order partonic calculations would lead to unreliable results. When less severe isolation cuts are used, the “one-”, and to a lesser extend, “two fragmentation” components are subdominant but not negligible.
#### 4.3.3 Infrared sensitive distributions
Being based on a finite order calculation, our computer code is not suited for the study of observables controled by multiple soft gluon emission . Among those, one may distinguish the following examples, most of which would require an improved account of soft gluon effects.
The transverse momentum distribution of pairs $`d\sigma /dq_T`$ near $`q_T=0`$
The problematics of the “two direct” contribution is similar to the well-known Drell-Yan process, see . On the other hand, the fragmentation contributions do not diverge order by order when $`q_T0`$. Indeed, in the “one fragmentation” case,
$`\text{parton 1}+\text{parton 2}`$ $``$ $`\gamma _1+\text{parton 3}`$ (4.1)
parton 3 $``$ $`\gamma _2+X`$ (4.2)
the NLO contribution to the hard subprocess (4.1) yields a double logarithm
$$\alpha _s\mathrm{ln}^2\text{p}_T(\gamma _1)+\text{p}_T(\text{parton}\mathrm{\hspace{0.17em}3})$$
(4.3)
when $`\text{p}_T(\gamma _1)+\text{p}_T(\text{parton}\mathrm{\hspace{0.17em}3})0`$. However the extra convolution associated with the fragmentation (4.2) involves an integral over $`z_2=p_T(\gamma _2)/p_T(\text{parton}\mathrm{\hspace{0.17em}3})`$ which smears out this integrable singularity. The “two fragmentation” contribution involves two such convolutions, and hence one more smearing.
The azimuthal angle distribution $`d\sigma /d\varphi _{\gamma \gamma }`$ near $`\varphi _{\gamma \gamma }=\pi `$
This case differs from the previous one for two reasons. Firstly, not only the “two direct” contribution diverges order by order when $`\varphi _{\gamma \gamma }\pi `$, but also both “single-” and “double-fragmentation” contributions do, as can be seen in Fig. (4). Moreover, in both fragmentation cases, soft gluons may couple to both initial- and final-state hard emitters. Indeed, consider the example of the “one fragmentation case”, cf. eqns. (4.1). Selecting $`\varphi _{\gamma \gamma }\pi `$ emphasizes $`\varphi (\text{parton}\mathrm{\hspace{0.17em}3})\varphi (\gamma _1)\pi `$, so that all the emitted partons besides parton 3 have to be collinear to either of the incoming or outgoing particles, and/or soft, which yields double logarithms
$$\alpha _s\mathrm{ln}^2\left[\pi \left(\varphi (\text{parton}\mathrm{\hspace{0.17em}3})\varphi (\gamma _1)\right)\right]$$
(4.4)
associated with each of the hard partons 1,2,3 - plus single logarithms as well. For the observable $`d\sigma /d\varphi _{\gamma \gamma }`$ near $`\varphi _{\gamma \gamma }=\pi `$, the integral involved in the convolution of the hard subprocess with the fragmentation functions does not smear these logarithmic divergences, since the fragmentation variable $`z_2=p_T(\gamma _2)/p_T(\text{parton}\mathrm{\hspace{0.17em}3})`$ is decoupled from the azimutal variable $`\varphi (\text{parton}\mathrm{\hspace{0.17em}3})`$ equal to $`\varphi (\gamma _2)`$ in the soft and collinear limits. A similar explanation holds for the “double fragmentation component”. An analagous problem affects the $`q_T`$ distribution of a pair photon $`+`$ jet at low $`q_T`$.
The azimuthal angle distribution $`d\sigma /d\varphi _{\gamma \gamma }`$ near $`\varphi _{\gamma \gamma }=0`$
Both fragmentation contributions to $`d\sigma /d\varphi _{\gamma \gamma }`$ diverge also order by order when $`\varphi _{\gamma \gamma }0`$. Here also soft gluons may couple to both initial- and final-state hard emitters. Actually, given the large invariant mass of the pairs, the vicinity of $`\varphi _{\gamma \gamma }=0`$ is never probed, so that an improved treatment of soft gluon effects is not needed in this case. Yet, as a consequence, an increase of the “single-fragmentation” contribution can be seen <sup>9</sup><sup>9</sup>9A similar behaviour occurs also for the “two fragmentation” contribution, which is however too tiny to have any significant effect. in the lower range of the $`\varphi _{\gamma \gamma }`$ spectrum, cf. Fig. (4).
An infrared divergence for $`d\sigma /dq_T`$ inside the physical spectrum
Besides the well-known issue at $`q_T=0`$, another infrared sensitive point appears in the $`q_T`$ spectrum due to isolation, at the critical value $`q_T=E_T^{max}`$. Indeed, at lowest order the “one fragmentation” component is not smooth, it instead behaves as a step function ,
$$\left(\frac{d\sigma }{dq_T}\right)^{\mathrm{`}\mathrm{`}singlefragm\mathrm{"},(LO)}\mathrm{\Theta }\left(E_T^{max}q_T\right)$$
(4.5)
Consequently, in agreement with the general study of , the NLO - and every higher order - correction has a double logarithmic divergence at the critical point. Such singularities are very sensitive to the kinematical constraints and the observable considered. The present case has a double logarithm below the critical point,
$$\left(\frac{d\sigma }{dq_T}\right)^{\mathrm{`}\mathrm{`}singlefragm\mathrm{"},(NLO)}\alpha _s\mathrm{ln}^2\left(E_T^{max}q_T\right)\mathrm{\Theta }\left(E_T^{max}q_T\right)+\mathrm{}$$
(4.6)
The infrared sensitive behaviour at this critical point can be infered on Fig. (7), where a rather large value for $`E_T^{max}`$ is used in order to split this critical point from the small $`q_T`$ region.
This effect is not visible on the theoretical prediction of Fig. (3) because of the low value $`E_T^{max}=4`$ GeV used by the D0 collaboration. The two infrared sensitive regions ($`q_T0`$ and $`q_TE_T^{max}`$) are not separated enough, and the bin smearing averages over both effects. The phenomenon is thus camouflaged. An all order summation of this soft gluon effect has to be carried out to restore a sensible shape to the $`q_T`$ distribution. A similar observation would be made about the $`q_T`$ distribution of a pair {photon $`+`$ jet}.
## 5 Perspectives
Future improvements of the theoretical understanding of di-photon production will require various inputs. Next-to-next to leading order corrections would hopefully stabilize the scale dependences and correct the normalisations. Another important quantitative improvement would be the higher order corrections to the box contribution, which is within reach thanks to some recent achievements. Progresses in this direction are reported on elsewhere in these proceedings .
This accounting of multiple soft gluon effects, already implemented in for the “two direct” contribution in absence of isolation cuts, is needed in the DIPHOX-based study to provide correct distributions in the infrared sensitive regions. This affects the $`q_T`$ distribution at low $`q_T`$. It also concerns the $`\varphi _{\gamma \gamma }`$ distribution when $`\varphi _{\gamma \gamma }\pi `$, which has been less studied, and for which the soft gluon effects in the fragmentation cases are more involved than in the “two direct” case . It is also important for the $`q_T`$ distribution in the vicinity of the critical point $`q_T=E_T^{max}`$, induced by the fixed cone type isolation criterion.
The understanding of the background to the Higgs search is quantitatively not yet on the same footing as for the signal. Hence, accounting for the higher order correction to the signal in numerical simulations might be instructive, but the results should be considered with care in order to avoid statements that are too optimistic.
## 6 Higgs search in association with a hard jet
In order to overcome the present insufficient control of higher order corrections in inclusive production, it has been suggested to study the associated production of $`h(\gamma \gamma )+`$ jet, for which both signal $`S`$ and background $`B`$ are lower (but still at the level of hundred signal events at low luminosity). The lowest order estimate has shown that the $`S/B`$ ratio is improved critically (up to $`1/21/3`$) with the same level of the significance $`S/\sqrt{B}`$. Furthermore, higher order corrections to the background have been shown recently to be under better control than in the inclusive (i.e. unassociated) case.
### 6.1 Background: associated vs. inclusive
Indeed, in the inclusive case, the magnitude of the NNLO box contribution is comparable to the LO cross section essentially because the latter is initiated by $`q\overline{q}`$, whereas the former involves $`gg`$. The $`gg`$ luminosity, much larger than the $`q\overline{q}`$ one, compensates numerically the extra $`\alpha _s^2`$ factor of the box. On the contrary, in the channel $`\gamma \gamma `$ \+ jet, the LO cross-section is dominated by a $`qg`$ initiated subprocess. The $`qg`$ luminosity is sizeably larger than the $`q\overline{q}`$ one, so that the corresponding NNLO remains small compared to the LO result . Thus, expecting that the subprocess $`gg\gamma \gamma g`$ gives the main NNLO correction, a quantitative description of the background with an accuracy better than 20% could be achieved already at NLO in the $`\gamma \gamma `$\+ jet channel for a high $`p_T`$ ($`30`$ GeV) jet. All the helicity amplitudes needed for the implementation of the (“direct” contribution to the) background to NLO accuracy <sup>10</sup><sup>10</sup>10We remind the reader unfamiliar with the LO, NLO, NNLO, etc, terminology that this terminology does not refer to the absolute power of $`\alpha _s`$ involved, but instead to the relative power with respect to the Born term of the process considered. Hence, NLO corrections to $`\gamma \gamma +`$ jet are also part of the NNLO corrections to $`\gamma \gamma `$ inclusive (converse not true, cf. box).are now available .
### 6.2 Signal vs. background
The origin of this improvement of the $`S/B`$ ratio is the following. At LO, the signal of associated $`h(\gamma \gamma )+jet`$ production has basically a 2-body kinematics, due to the extremely small Higgs width (a few MeV). On the contrary, the LO background contributions $`q\overline{q}\gamma \gamma g`$ and $`gq\gamma +\gamma +q`$ (as well as the NNLO $`gg\gamma \gamma g`$ one) have 3-body kinematics. This is in contrast to the inclusive channel where both signal and background LO subprocesses have 2-body kinematics. This circumstance opens the room for more refined cuts to suppress the background more efficiently . The contributions $`q\overline{q}\gamma \gamma g`$ and $`gg\gamma \gamma g`$ are suppressed down to 40% of the signal (to be compared with $`S/B1/7`$ for the inclusive channel). Unfortunately, as found in , the contribution of the other subprocess $`gq\gamma \gamma q`$, which dominates, yields the overall ratio $`S/B1/21/3`$.
The cuts which allow this efficient suppression of the irreducible background at LO are based on the differences in the shapes of the angular distributions in the partonic c.m.s.. Due to helicity and total angular momentum conservation, the S-wave does not contribute to the dominant signal subprocess $`ggHg`$. On the contrary, all angular momentum states contribute to the subprocesses $`gq\gamma \gamma q`$ and $`q\overline{q}\gamma \gamma g`$. Therefore, the signal has a more suppressed threshold behaviour compared to the background. The $`S/B`$ ratio can thus be improved by increasing the partonic c.m.s. energy $`\sqrt{\widehat{s}}`$ far beyond threshold, and a cut $`\sqrt{\widehat{s}}>300`$ GeV has been found to give the best S/B ratio for the LHC. Actually, the effect can not be fully explained by the threshold behavior only, since that would result in a uniform suppression factor. It was shown in (see Figs. 5 and 6 there) that the dependences of the background and the signal on the c.m.s. angular variables are quite different; therefore, the strong $`\widehat{s}`$ cut affects them with different suppression factors (see for more details). This effect can be exploited to enhance the significance $`S/\sqrt{B}`$ at the same level as $`S/B`$. If the cut $`\mathrm{cos}(\vartheta ^{})(j\gamma )<0.87`$ on the jet-photon in the partonic c.m.s. is applied for $`\sqrt{\widehat{s}}<300`$ GeV and combined with the cut $`\sqrt{\widehat{s}}>300`$ GeV, the change on $`S/B`$ is rather small, while the significance is improved by a factor $``$ 1.3. The same effect can be observed with the cut on the jet angle in the partonic c.m.s. $`(\vartheta ^{}(j)`$, cf. Fig. 5 of , but one should notice that the two variables, $`\vartheta ^{}(j\gamma )`$ and $`\vartheta ^{}(j)`$, are correlated. Therefore, it is desirable to perform a multi-variable optimization of the event selection. Notice that the present discussion is based on a LO analysis, and concerns only what was defined above as the “direct” component of the irreducible background. One now has to understand how this works at NLO.
Other, reducible, sources of background are potentially dangerous. The above-defined “one fragmentation” component to the so-called irreducible background, and the reducible background coming from misidentification of jet events were treated on a similar footing in the LO analysis of as a de facto reducible background. In , a rough analysis found that this reducible background is less than 20% of the irreducible one after cuts are imposed. The misidentification rate is given mainly by the subprocesses $`gq\gamma gq`$, $`gg\gamma q\overline{q}`$ and $`qq^{}\gamma q(g)q^{}(g)`$, when the final state parton produces an energetic isolated photon but other products of the hadronization escape the detection as a jet. There, a $`\gamma (\pi ^0)/jet`$ rejection factor equal to 2500 for a jet misidentified as a photon and 5000 for a well separated $`\gamma (\pi ^0)`$ production by a jet were used. No additional $`\pi ^0`$ rejection algorithms were applied. Furthermore, this reducible background is expected to be suppressed even more strongly than the irreducible background of “direct” type when a cut on $`\sqrt{\widehat{s}}`$ is applied.
In summary, the associated channel $`H(\gamma \gamma )+`$ jet with jet transverse energy $`E_T>30`$ GeV and rapidity $`|\eta |<4.5`$ (thus involving forward hadronic calorimeters) opens a promising possibility for discovering the Higgs boson with a mass of 100-140 GeV at LHC even at low luminosity. However, to perform a quantitative analysis, the NLO calculations, namely of the background, have to be completed, and included in a more realistic final state analysis.
## 7 Acknowledgements
T. B. is a EU fellow supported by the EU Fourth Programme “Training and Mobility of Researchers”, Network “Quantum Chromodynamics and the Deep Structure of Elementary Particles”, contract FMRX-CT98-0194 (DG12 - MIHT). V.A. I. acknowledges support by the CERN-INTAS grant 377 and RFBR-DFG grant 99-02-04011. LAPTH is a Unité Mixte de Recherche (UMR 5108) associée au CNRS et à l’Université de Savoie.
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# New solutions to covariant nonequilibrium dynamics
## I Introduction and Conclusions
A theoretical framework to study nonequilibrium dynamics is provided by Boltzmann transport theory. The dynamical variables of this theory are the Lorentz-covariant, one-particle phase space distributions $`f_i(x,p)`$; while the dynamics is governed by transition probabilities $`W_{cc^{}}`$, which are Lorentz-covariant functions of the particle momenta. The theory, while not exact, is rather general. First, it is not restricted to particular particle types. The particles could be partons, hadrons, or molecules. Second, the reaction rates that specify the dynamics are also unrestricted in their origin. For example, the rates could emerge from an effective quantum field theory or Newtonian mechanics. The primary limitations of the theory are the neglect of dynamical correlations (one-body truncation of the formal BBGKY hierarchy) and the inability, without additional classical fields, to model phase transition dynamics.
We consider here the simplest form of Lorentz-covariant Boltzmann transport theory in which the on-shell phase space density $`f(x,𝐩)`$, evolves with an elastic $`22`$ rate as:
$`p_1^\mu _\mu f_1`$ $`=`$ $`{\displaystyle \underset{2}{}}{\displaystyle \underset{3}{}}{\displaystyle \underset{4}{}}\left(f_3f_4f_1f_2\right)W_{1234}\delta ^4(p_1+p_2p_3p_4)`$ (2)
$`+S(x,𝐩_1).`$
Here $`W`$ is the square of the scattering matrix element, the integrals are shorthands for $`\underset{i}{}\frac{gd^3p_i}{(2\pi )^3E_i}`$, where $`g`$ is the number of internal degrees of freedom, while $`f_jf(x,𝐩_j)`$. The initial conditions are specified by the source function $`S(x,𝐩)`$, which we discuss later in Section II. For our applications below, we interpret $`f(x,𝐩)`$ as describing an ultrarelativistic massless gluon gas with $`g=16`$ (8 colors, 2 helicities).
Recall several important properties of Eq. (2). First, the particle number current and the energy-momentum tensor are given by
$$N^\mu (x)\frac{d^3p}{(2\pi )^3E}p^\mu f(x,𝐩)$$
(3)
and
$$T^{\mu \nu }(x)\frac{d^3p}{(2\pi )^3E}p^\mu p^\nu f(x,𝐩).$$
(4)
With these definitions, particle number and energy-momentum conservation follow from Eq. (2) (when the source term $`S(x,𝐩)=0`$). Second, there is a class of fixed points, called global equilibria, which are phase space densities of the form
$$f(x,𝐩)=\frac{g}{(2\pi )^3}\mathrm{exp}\left[\frac{\mu p_\mu u^\mu }{T}\right],$$
(5)
where $`u_\mu `$ is a constant four-vector that specifies a global flow velocity, while $`T`$ and $`\mu `$ correspond to the constant temperature and chemical potential. Furthermore, the H-theorem states that the Boltzmann transport equation drives the system towards global equilibrium.
Despite its relatively simple form, the Boltzmann equation is nonlinear with very few known analytic solutions. Until recently, progress to obtain even numerical solutions has been hampered by its numerical complexity. The rapid increase in computational power has finally made it possible to break through this barrier. For nuclear collision applications, new numerical algorithms are being developed, tested, and made available via the World Wide Web under a new Open Standard for Codes and Routines (OSCAR). The present work is a further step in that development.
For nuclear collisions at SPS energies ($`\sqrt{s}20`$ $`A`$GeV), numerical solutions of hadronic transport models have been available for some time. However, for higher collider energies, the emergence of massless partonic degrees of freedom creates the technical challenge of how to retain Lorentz covariance. In this paper we present results based on a new numerical technique, MPC 0.1.2 , that provides reliable solutions in this ultrarelativistic regime.
The difficulty of the analytic treatment of the Boltzmann transport equation has forced workers in the past to make strong simplifying assumptions. A common simplification has been to ignore the general nonequilibrium problem and to postulate that local equilibrium is maintained at all times. In the framework of the Boltzmann transport theory, this choice corresponds to substituting a local equilibrium ansatz in place of the fixed point global equilibria. Allowing $`(u_\mu ,T,\mu )`$ to vary with the coordinate $`x^\mu `$, this ansatz corresponds to
$$f(x,𝐩)=\frac{g}{(2\pi )^3}\mathrm{exp}\left[\frac{\mu (x)p_\mu u^\mu (x)}{T(x)}\right].$$
(6)
It is however well known that local equilibrium is not a solution of Eq. (2). The nonlinear collision term vanishes in local equilibrium, but $`p^\mu _\mu f0`$ in general. Only in the limit when all the rates go to infinity, i.e., when the mean free path goes to zero, can the solution approach local equilibrium.
A covariant, dynamical theory can nevertheless be postulated based on the assumption of local equilibrium. That is relativistic hydrodynamics, which is widely used in heavy-ion physics to calculate observables. Assuming an equilibrium initial condition specified on a hypersurface $`\sigma _{in}^\mu (x)`$, the local energy momentum and baryon number conservation laws
$$^\mu T_{\mu \nu }=0,^\mu N_{B,\mu }=0$$
(7)
reduce to (Euler) hydrodynamical equations under the assumptions that local chemical and thermal equilibrium are maintained and that dissipation (viscosity and thermal conductivity) can be neglected. In that case, the energy-momentum tensor and baryon current can be expressed as $`T_{\mu \nu }=u_\mu u_\nu (e+p)g_{\mu \nu }p`$ and $`N_{B,\mu }=u_\mu n(x)`$ in terms of the local flow velocity $`u_\mu (x)`$, local pressure $`p(x)`$, local energy density $`e(x)`$, and local proper density $`n(x)`$. The equations form a closed system if, in addition, the equation of state, $`e(p,n_B)`$, is specified.
It is clear that the idealization of local equilibrium may apply, if at all, only in the interior of the reaction volume, where the local mean free path $`\lambda (x)=1/(\sigma n(x))`$ may remain small for a while as compared to the characteristic dimensions and gradients of the system, $`L_\mu (x)|_\mu \mathrm{log}e(x)|^1`$. However, these assumptions are marginal for conditions encountered in heavy ion collisions and certainly break down near the surface region and throughout the freeze-out phase. Due to longitudinal expansion the density decreases as $`1/\tau `$ until $`\tau \sqrt{3}R`$ when 3-dimensional expansion rapidly increases $`\lambda `$ beyond $`L`$.
For small departures from local equilibrium, corrections to the ideal Euler hydrodynamic evolution can be calculated by taking the $`T^{\mu \nu }`$ and $`N_B^\mu `$ moments of an underlying kinetic theory. To first order in $`\lambda /L`$, the equations reduce to the Navier-Stokes equations. The solutions to Navier-Stokes depend therefore not only on the equation of state, but also on the bulk and shear viscosity and thermal conductivity transport coefficients of the medium. While such an approach has proven useful in nonrelativistic problems and in special relativistic geometries, severe problems of instability and acausality appear when extended into the ultrarelativistic domain.
The newly formulated, covariant, parton kinetic theory technique, MPC, allows us to compute the highly dissipative evolution during the densest partonic phase of the reaction in a covariant manner as well as investigate the final freeze-out dynamics. MPC is an extension of Zhang’s covariant parton cascade algorithm, ZPC. Both MPC and ZPC have been extensively tested and compared to analytic transport solutions and covariant Euler and Navier-Stokes dynamics in 1+1D geometry. A critical new element of both these algorithms is the parton subdivision technique proposed by Pang. As shown in detail in Section III, rather high subdivision $`100`$ is needed to preserve Lorentz covariance numerically for massless parton evolution.
Extensions of MPC to include inelastic $`23`$ partonic processes are under development, but in this paper we use MPC in the pure elastic parton interactions mode as in ZPC.
The aim of this work is to calculate the sensitivity of the evolution and freeze-out of an ultrarelativistic (massless) parton system to the transport rates and initial conditions expected in Au+Au reactions at RHIC energies ($`\sqrt{s}200`$ AGeV). We compare the results to the (OSCAR compliant) relativistic Euler hydrodynamic code developed by Rischke and Dumitru . An ideal gas ($`e=3p`$) equation of state is used in that analysis. This work extends Refs. , focusing on the freeze-out problem in 3+1D Bjorken expansion. It is less ambitious than for example Ref. by limiting the study to massless partons. This avoids introducing yet further complications due to dissipative hadronization and hadronic transport effects. The initial conditions are taken from the HIJING multiple mini jet generator.
One of the main results of this work is shown in Fig. 1. To test whether ideal hydrodynamics is an adequate approximation of the parton transport equation (2), we have followed the evolution both for the Boltzmann equation (2) and Euler hydrodynamics from the same RHIC initial condition from . (See Section V for details of the simulations.) Figure 1 shows the evolution of the transverse energy $`dE_t/dy`$ at midrapidity. There is a large difference between the transport and hydrodynamic solutions in both the 1+1 and 3+1 dimensional case, even for a physically extreme, 15mb, cross section. Note, that these 15mb curves are equivalent<sup>*</sup><sup>*</sup>* The equivalence is due to the scaling property of Eq. (2) explained in Section III. to a solution for 3 mb cross section with five times higher initial density than expected with HIJING. Therefore, large deviations from ideal hydrodynamics are to be expected for possible initial conditions at RHIC. In addition, this conclusion is independent of the initial system size (see explanation in Subsection V A).
This conclusion reinforces the results of Ref. , where it was shown that large deviations even from the Navier-Stokes evolution in 1+1D Bjorken expansion are expected for initial densities up to four times higher than predicted by the HIJING model.
A second main result of this work is shown in Fig. 2. We tested whether the widely used Cooper-Frye freeze-out prescription could “correct” final observables for the neglect of the early breakdown of ideal hydrodynamics. Figure 2 shows the final, experimentally observable $`p_{}`$-distributions divided by the thermal initial distribution. We find that there is an order of magnitude difference between transport and hydrodynamic solutions at both low and at high $`p_{}`$ for realistic ($``$ few mb) gluonic cross sections. The difference is still a factor of two to three even for a physically extreme, 15 mb, cross section. In addition, increasing the radius of the initial Bjorken cylinder from 2 to 6 fm does not reduce the discrepancy between the covariant transport solutions and those of Cooper-Frye frozen hydrodynamics. To get closer to the transport $`p_{}`$-distributions, one would have to choose a freeze-out temperature much above the commonly assumed $`100150`$ MeV range.
The last main result of this work is illustrated in Fig. 3. This shows that high hydrodynamic freeze-out temperatures that would be needed to “fit” the transport solutions are, however, inconsistent with the space-time freeze-out distributions of covariant transport theory. Unlike the “sharp” space-time freeze-out particle distributions commonly assumed using the Cooper-Frye freeze-out prescription, the transport theory freeze-out volume is four-dimensional. Particles freeze out over a large four-volume that forms a wedgelike freeze-out region in the $`\tau R`$ plane. This freeze-out distribution depends strongly on the microscopic reaction rates (higher rates lead to a later freeze-out). It is not possible to tune the Cooper-Frye freeze-out temperature to reproduce the cascade freeze-out distributions. Though one can arrange that the Cooper-Frye freeze-out curve follows more-or-less the ridge, the transport distribution along that ridge is not correctly reproduced by Cooper-Frye frozen hydrodynamics. In particular, hydrodynamic freeze-out surfaces with a timelike section result in unphysical spikes in the $`dN/d\tau `$ distribution which are not present in the transport theory calculations. These spikes arise when the freeze-out temperature is such that the interior of the system freezes out due to longitudinal Bjorken expansion (see Subsection V B 1 for further discussion).
In summary, our results show that for a rather wide range of initial conditions at RHIC energies, the evolution of the system deviates strongly from Eulerian hydrodynamics throughout the 3+1D evolution. It is not possible to mimic the observables from the nonequilibrium evolution by simply applying the isotherm Cooper-Frye freeze-out prescription to ideal hydrodynamics. The space-time four-volume of freeze-out, even for the largest ($`R6`$ fm) nuclei, does not resemble a timelike surface. In addition, the observable transverse momentum spectra are very sensitive to the microscopic reaction rates. These results indicate that while ideal hydrodynamics is a useful model to explore possible collective dynamics in nuclear collisions, the interpretation of experimental observables must take into account the finite transition probabilities $`W_{\{i\}\{j\}}`$ that govern the nonequilibrium evolution. Fortunately, numerical techniques such as MPC and ZPC are now readily available. Experimentally, the $`A`$, $`E_{cm}`$, and multiplicity dependence of the observables provides the best way to measure these effective reaction rates in the ultradense matter formed in nuclear collisions.
## II Covariant Parton Transport Theory
Equation (2) is the simplest form of classical Lorentz-covariant Boltzmann transport theory. In principle, the transport equation could be extended for bosons with the substitution $`f_1f_2f_1f_2(1+f_3)(1+f_4)`$ and a similar one for $`f_3f_4`$ (where we used the short-hand $`f_if(x,𝐩_i)`$). In practice, no covariant algorithm yet exists to handle such nonlinearities. We therefore limit our study to quadratic dependence of the collision integral on $`f`$.
The elastic gluon scattering matrix elements in dense parton systems are typically of the Debye-screened form: $`d\sigma /dq^2(9\pi \alpha _s^2/2)/(q^2+\mu ^2)^2`$, which favors small angle scattering. However, the relevant transport cross section is $`\sigma _t=𝑑\sigma \mathrm{sin}^2\theta _{cm}(9\pi \alpha ^2/2s)\mathrm{log}(s/4\mu ^2)`$, where $`s17T^2`$. In order to maximize the equilibration rate for a fixed cross section, we take here an isotropic differential cross section in the center-of-mass frame instead. We further assume an energy-independent cross section with a threshold specified by $`\mu ^2`$, i.e., our solutions therefore correspond to the microscopic dynamics specified by the following idealized model
$$d\sigma =\mathrm{\Theta }(s^2\mu ^2)\frac{\sigma _0}{4\pi }d\mathrm{\Omega }.$$
(8)
The transport cross section is $`2\sigma _0/3`$ in this case.
It is important to emphasize that while the cross section suggests a geometrical picture of action over finite distances, we use Eq. (8) only as a convenient parametrization to describe the effective local transition probability, $`W`$. In the present study this is simply modeled as $`dW/d\mathrm{\Omega }=sd\sigma /d\mathrm{\Omega }`$. The particle subdivision technique (see next Section) needed to recover covariance removes all notion of nonlocality in this approach, just like in hydrodynamics. Thus, the cross sections, e.g., 60 mb, used in the present study to simulate rapid local changes of the phase space density in no way imply that distances bigger than $`1`$ fm play any role.
With the above cutoff $`\mu `$, freeze-out of a test particle can arise in two different ways: either the system becomes too dilute, i.e., $`1/n\sigma _0>>L`$, or the system cools down and the threshold suppresses further interactions. By construction, the possibility for the latter case occurs along an isotherm, $`T_f\mu /\sqrt{17}`$. With Eq. (8), we can therefore study the influence of dissipative phenomena by varying the two scales $`\sigma _0`$ and $`\mu ^2`$. The evolution was performed with $`\sigma _0=3`$, 15, 30, 60, 121 mb and $`\mu =0`$, 0.1, 0.5 GeV.
The initial condition was taken to be a longitudinally boost invariant Bjorken cylinder in local thermal and chemical equilibrium at temperature $`T(\tau _0)=500`$ MeV at proper time $`\tau _0=0.1`$ fm/$`c`$ as by fitting the gluon mini-jet transverse momentum spectrum predicted by HIJING. In order to compare to hydrodynamics, we assume that the transverse density distribution is uniform up to a radius $`R_0=2`$, 4, 6, or 8 fm. The pseudo-rapidity $`\eta 1/2\mathrm{log}((t+z)/(tz))`$ distribution was taken as uniform between $`|\eta |<5`$. Since we want to compare to chemically and thermally equilibrated hydrodynamics, Technically, MPC was run with an out-of-chemical-equilibrium initial gluon density $`n_{\eta ,0}=4`$ fm<sup>-2</sup> as obtained via HIJING including final state radiation and with cross sections $`\sigma _0=2`$, 10, 20, 40, and 80 mb. As explained in Section III, because of the scaling property of the solutions of the transport equation, the solutions for the chemical equilibrium initial condition are identical when the cross section is rescaled by a factor $`l=2.6505/41/1.509`$. the equilibrium initial gluon density was taken for this $`T(\tau _0)`$ to be
$$n_{\eta ,0}\frac{dN}{d\eta d^2x_{}}|_{\tau _0}=\frac{g}{\pi ^2}T^3\tau _02.65\mathrm{fm}^2.$$
(9)
Evolutions from different initial densities can be obtained by varying the cross section only and using the scaling property explained in the next Section.
## III Parton Subdivision and Scaling of Solutions
We utilize the parton cascade method to solve the Boltzmann transport equation (2). A critical drawback of all cascade algorithms is that they inevitably lead to numerical artifacts because they violate Lorentz covariance. This occurs because particle interactions are assumed to occur whenever the distance of closest approach (in the relative c.m.) is $`d<\sqrt{\sigma _0/\pi }`$, which corresponds to action at a distance. To recover the local character of equation (2) and hence Lorentz covariance, it is essential to use the parton subdivision technique. This is based on the covariance of Eq. (2) under the transformation
$$ff^{}lf,WW^{}W/l(\sigma \sigma ^{}=\sigma /l).$$
(10)
As shown in Ref. , the magnitude of numerical artifacts is governed by the diluteness of the system $`\sqrt{\sigma }/\lambda _{MFP}`$, that scales with $`1/\sqrt{l}`$. Lorentz violation therefore formally vanishes in the $`l\mathrm{}`$ limit.
### A Convergence with subdivision
Figure 4 illustrates the severeness of the cascade numerical artifacts in the case of insufficient particle subdivision. The top plot in Fig. 4 shows that the parton cascade solution for the evolution of the transverse energy per unit rapidity does not converge until the subdivision factor reaches $`l100`$. The lack of covariance can be seen in the difference between the solutions in frames separated by 3 units of rapidity. The very fact that the cascade evolution is different for different particle subdivisions means that the subdivision covariance (10) is itself violated by the cascade algorithm. Nevertheless, both Lorentz and subdivision covariance are recovered when $`l`$ is sufficiently large.
The large overshoot in the $`dE_t/dy`$ evolution is a result of the superluminal signal propagation speed inherent to the cascade algorithm. A cascade particle can influence almost instantaneously another cascade particle that is within the interaction range $`r_\sigma \sqrt{\sigma /\pi }`$. In a very dense system, a “chain” of almost instantaneous interactions can occur causing long range superluminal artifacts.
As a measure of the signal propagation speed in a nonlocal collision in the cascade we define
$$𝐯_s\frac{𝐱_{partner}(t_{collision})𝐱_{particle}(t_{lastcollision})}{t_{collision}t_{lastcollision}}.$$
(11)
Analytically, the deviation of the signal propagation speed from the speed up light can be roughly approximated by
$$\mathrm{\Delta }v_s=\frac{r}{t},$$
(12)
where $`t`$ is the time between the collision and the previous collision, while $`r`$ is the distance between the colliding particles at the time of the collision ($`r<r_\sigma `$). This is a pessimistic estimate that maximizes the deviations. Assuming that subsequent collisions are uncorrelated, $`t`$ follows a Poisson distribution The scaling (10) leaves the mean free path $`\lambda `$ invariant.
$$P(t)\frac{dn}{dt}=\frac{1}{\lambda }\mathrm{exp}(t/\lambda ).$$
(13)
Hence the distribution of $`\mathrm{\Delta }v_s`$ (with $`r`$ fixed) is
$$P(\mathrm{\Delta }v_s)=\frac{dn}{dt}\frac{dt}{d\mathrm{\Delta }v_s}=\frac{r}{(\mathrm{\Delta }v_s)^2\lambda }\mathrm{exp}\left(\frac{r}{|\mathrm{\Delta }v_s|\lambda }\right).$$
(14)
Particle subdivision reduces $`r_\sigma `$ as $`r_\sigma (l)=r_\sigma (1)/\sqrt{l}`$. Therefore, in the large subdivision limit, the subluminal and superluminal tails of the signal velocity distribution scale as a power law
$$P(\mathrm{\Delta }v_s)\left(\frac{v_0/\sqrt[4]{l}}{\mathrm{\Delta }v_s}\right)^2,v_0\sqrt{\frac{r}{\lambda }},$$
(15)
i.e., the distribution gets narrower as $`P_l(\mathrm{\Delta }v_s)P_1\left(\mathrm{\Delta }v_s\sqrt[4]{l}\right)`$.
The “measured” cascade distributions of the magnitude of the signal propagation speed, $`dn/dv_s`$, as defined via (11), and the magnitude of its transverse component, $`dn/dv_{s,}`$, are shown in Fig. 4. Though the distribution of $`v_s`$ is strongly peaked at $`v_s=c`$, both super- and subluminal propagation are present. While increasing particle subdivision decreases the deviations from the exact propagation speed $`c`$, convergence is slow. Even for a subdivision of 100, $`13\%`$ of the collisions correspond to a signal propagation velocity larger than $`1.5c`$. One must keep in mind that Fig. 4 shows the distribution of the signal propagation speed measured over the length and time scale of a single collision. On larger scales, the deviation is reduced because the large scale signal velocity is the sum of many small scale signal velocities.
In summary we demonstrated that the numerical artifacts due to Lorentz violation and acausality are reduced by subdivision and the cascade solution converges as $`l`$ increases. In the $`l\mathrm{}`$ limit the cascade technique gives the correct numerical solution of the transport equation (2). In practice, rather high subdivisions were found necessary to recover covariance. We could explore convergence up to $`l=800`$, 200, 150, and 100, for $`R_0=2`$, 4, 6, and 8 fm with the workstations available to us.
### B Scaling of the transport solutions
Subdivision covariance (10) actually implies that the transport equation has a broad dynamical range, and the solution for any given initial condition and transport property immediately provides the solution to a broad band of suitably scaled initial conditions and transport properties. This is because solutions for problems with $`l`$ times larger the initial density $`dN/d\eta d^2x_{}`$, but with one $`l`$-th the reaction rate can be mapped to the original ($`l=1`$) case for any $`l`$. We must use subdivision to eliminate numerical artifacts. However, once that is achieved, we have actually found the solution to a whole class of suitably rescaled problems.
The dynamical range of the transport equation (2) is further increased by its covariance under coordinate rescaling
$$f(x,𝐩)f^{}(x,𝐩)f(\frac{x}{l_x},𝐩),WW^{}\frac{W}{l_x}.$$
(16)
This is a simultaneous rescaling of space-time and the transition probability. In addition, there is also a covariance under rescaling of the momenta
$`f(x,𝐩)`$ $``$ $`f^{}(x,𝐩)l_p^3f(x,{\displaystyle \frac{𝐩}{l_p}}),`$ (17)
$`W(\{p_i\})`$ $``$ $`W^{}(\{p_i\})l_p^2W\left(\left\{{\displaystyle \frac{p_i}{l_p}}\right\}\right),`$ (18)
such that the particle density is again unchanged. This scaling also implies a rescaling of the mass $`mm^{}=m/l_p`$. Combining the three scaling transformations, we find covariance of the transport theory under
$`f(x,𝐩)`$ $``$ $`f^{}(x,𝐩)l_p^3lf({\displaystyle \frac{x}{l_x}},{\displaystyle \frac{𝐩}{l_p}}),`$ (19)
$`W(\{p_i\})`$ $``$ $`W^{}(\{p_i\}){\displaystyle \frac{l_p^2}{l_xl}}W\left(\left\{{\displaystyle \frac{p_i}{l_p}}\right\}\right).`$ (20)
In our calculation using MPC, we vary the physical parameters: $`\sigma `$, $`\mu `$, $`T_0`$, $`R_0`$, $`\tau _0`$, and $`n_{\eta ,0}dN/d\eta d^2x_{}|_{\tau _0}`$ (the rapidity interval $`\eta _{max}=5`$ was fixed). Keeping in mind Eq. (39) and that
$$n_\eta \frac{dN}{dyd_x_{}^2}|_\tau =d^2p_{}𝑑\eta m_t\mathrm{ch}(y\eta )\tau f(𝐱_{},\eta ,\tau ,𝐩_{},y),$$
covariance under the transformation (20) implies that once the solution for a particular choice of these parameters is known, then the solution is known for any other choice of the parameters which are related to the original via
$`\sigma ^{}=l_x^1l^1\sigma ,T_0^{}=l_pT_0,R_0^{}=l_xR_0,`$ (21)
$`n_{\eta ,0}^{}=l_xln_{\eta ,0},\mu ^{}=l_p\mu ,\tau _0^{}=l_x\tau _0.`$ (22)
Therefore, we can scale one solution to others provided that $`\mu /T_0`$, $`R_0/\tau _0`$, and $`\sigma n_{\eta ,0}\tau _0/\overline{\lambda }_{MFP}`$ remain the same. For example, three times the density with one-third the cross section leaves all three parameters the same, hence the results can be obtained via scaling without further computation. Table I shows sets of the three ratios that we mapped out via MPC.
## IV Freeze-Out
### A Hydrodynamic Freeze-Out Problem
In Section I we argued that hydrodynamics cannot be valid during the complete evolution in nuclear collisions because the assumption of local equilibrium breaks down. Thus, in spite of its appeal, hydrodynamics cannot be compared with measurements without additional model assumptions needed to specify when and how it breaks down. The problem of determining those extra model assumptions is the so-called freeze-out problem.
For application to nuclear collisions, freeze-out cannot be formulated as an expansion in $`\lambda /L`$ since by definition it occurs when that ratio exceeds unity. Hence, even the Navier-Stokes hydrodynamics is inadequate to solve the freeze-out problem.
A common freeze-out prescription, which we here name “Cooper-Frye frozen hydrodynamics”, is to assume the validity of ideal hydrodynamics up to a “sharp” 3D freeze-out hypersurface $`\sigma ^\mu (x)`$. Assuming that all interactions suddenly cease on that hypersurface, the final (frozen-out) invariant differential distribution of particles is then computed via the Cooper-Frye formula:
$$EdN=\frac{d^3p}{(2\pi )^3}d\sigma ^\mu (x)p_\mu f(x,𝐩).$$
(23)
Here $`d\sigma ^\mu `$(x) is the normal vector to the 3D freeze-out hypersurface at the point $`x`$, while $`f(x,𝐩)`$ is assumed to be in local equilibrium and hence, for classical particles, given by Eq. (6). While this prescription is covariant and appealingly simple, it suffers from several well known problems:
First, because the hydrodynamical solutions do not contain dynamical information needed to compute the freeze-out hypersurface, the assumed one is simply an ad hoc external constraint. It is usually parameterized in terms of a few physically “reasonable” parameters, the most common being a freeze-out isotherm $`T(\sigma ^\mu )=T_f`$ or freeze-out energy density $`e_f`$. It is not possible to estimate the errors introduced by such a prescription.
Second, the Cooper-Frye formula allows negative contributions to the measurable particle yields. This can be avoided by choosing a non-equilibrium post freeze-out distribution that does not have particles in the phase space domain where $`d\sigma ^\mu p_\mu <0`$. However, such a choice still relies on the existence of a sharp 3D freeze-out surface.
Finally, while an idealized sharp freeze-out surface may be adequate for applications to quasi-stationary macroscopic systems, it cannot be justified in expanding mesoscopic systems in which $`L/\lambda `$ is never large. The very fact that such systems do freeze out, i.e., $`\lambda (\sigma )>L`$, means that the solution to freeze-out problem must entail global information as the system becomes more and more dilute. Furthermore, there is no way to justify the neglect of final state interactions during freeze-out stage of the reactions while expansion and rarefaction are causing the system to depart from local equilibrium.
In Ref. a continuous emission hydrodynamical freeze-out model was proposed to overcome some of these problems. The global information relevant to freeze-out in that model is taken there as the Glauber escape probability
$$P(x,p)=\mathrm{exp}\left(_\tau ^{\tau _{out}}𝑑\tau ^{}\sigma v_{rel}n(x(\tau ))\right).$$
(24)
This formula reveals clearly the highly nonlocal character of the freeze-out problem. At any point in spacetime, $`x^\mu `$, the line integral runs over the future trajectory and therefore is exponentially sensitive to the future evolution of the system. This leads to a formidable self-consistency problem. For special geometries such as Bjorken boost invariant expansion, a rough estimate of $`P`$ can be made using the approximate scaling Bjorken hydrodynamic solution
$$n(x(\tau ))=\frac{\tau _0}{\tau }n(x(\tau _0))\mathrm{\Theta }(R^2(𝐱_{}+𝐯_{}\tau )^2).$$
(25)
Together with the Glauber straight line trajectory, this leads to the characteristic power law survival probability
$$P\left(\frac{\tau }{\tau _{out}}\right)^{\sigma \tau _0n(\tau _0)},$$
(26)
that also appears, e.g., in the $`J/\psi `$ suppression problem. While Eq. (24) captures essential global physics of freeze-out, it is not complete since before freeze-out the trajectories cannot be straight if local equilibrium is maintained via the assumed hydrodynamic equations. Also, in the surface region where $`P1/2`$ neither hydrodynamics nor eikonal dynamics applies.
The solution to the freeze-out problem in classical mechanics is given by microscopic transport theory. A hybrid approach that partially reaches that end was proposed in Ref. , which combines partonic hydrodynamics with hadronic transport theory. In that approach, hydrodynamics is assumed to hold up to only some intermediate critical temperature hypersurface, $`T(\sigma _{int}^\mu )=T_c>T_f`$, on which the fluid is converted to hadrons via the Cooper-Frye formula. Subsequent evolution of the hadronic system is then calculated by solving the hadronic transport theory as encoded in UrQMD. As noted in Ref. , the freeze-out surface is actually a diffuse four-volume, and in addition different hadronic species freeze out over different four-volume domains. This sequential freeze-out leads to strong observable correlations such as the mass dependence of final transverse spectra. The main limitation (and/or advantage) of the above hybrid model is the need to assume the validity of hydrodynamics in the dense partonic phase of the collision. It is advantageous in that possible collective effects due to the quark-gluon confinement transition can be explored with hydrodynamics using “realistic” equations of state. It is disadvantageous in that it is far from clear that local equilibrium is ever reached during the evolution. Recall that dissipative effects on even global observables such as the transverse energy per unit rapidity cannot be accurately calculated using the Navier-Stokes equations.
Despite these known complications of the freeze-out problem, ideal hydrodynamics and Cooper-Frye freeze-out are still commonly used to fit experimental data using isotherm freeze-out hypersurfaces and draw inferences about the underlying dynamics. The consistency and significance of interpretations based on such fits can only be assessed by comparing detailed dynamical transport calculations to the hydrodynamic limit (see Section V).
### B Formal Definition for Freeze-out
An important experimental observable aspect of the space-time evolution of kinetic theory is the freeze-out distribution. In the framework of discrete parton cascade dynamics, the definition of the freeze-out distribution, $`dN_{fo}`$, is the number of partons per $`d^4xd^4p`$ invariant phase space volume with momentum $`p^\mu `$ that have a collision at $`x^\mu `$ but suffer no more collisions. Given the trajectories, $`(𝐱_a(t),𝐩_a(t))`$ or the world lines $`x_a^\mu (\tau )`$ of all partons $`a`$, that distribution is given by the ensemble average of the space-time coordinates, $`x_{af}^\mu (t_{af},𝐱_a(t_{af}))`$, of the last collision together with the final outgoing momentum, $`𝐩(t_{af}+0^+)`$:
$`f_{fo}(x,p){\displaystyle \frac{dN_{fo}}{d^4xd^4p}}={\displaystyle \underset{a}{}}\delta (tt_{af})\delta ^3(𝐱𝐱_a(t_{af}))`$ (27)
$`\times \delta ^4\left(pp_a(t_{af}+0^+)\right).`$ (28)
Because the freeze-out times, $`t_{af}`$, are distributed over a broad time interval, $`f_{fo}`$ does not correspond to
$$f(x,p)=N𝑑\tau \delta ^4(xx(\tau ))\delta ^4(pp(\tau ))$$
(29)
at any time or on any fixed 3-D hypersurface. Note that $`f`$ measures the phase space density of the world lines $`x^\mu (\tau )`$ and their four-velocities at a single point $`x^\mu `$. On the other hand, $`f_{fo}`$ measures the phase space density of last scattering events, where the momentum $`p`$ of a particle was last changed. Pion interferometry measures the Fourier transform of $`f_{fo}`$. Note that even after integrating over the freeze-out points, the final observed momentum spectrum, $`d^4xf_{fo}(x,p)`$, is only equal to the Cooper-Frye formula if $`x_{af}^\mu `$ happen to lie on a sharp 3D hypersurface, $`\sigma ^\mu (\zeta _1,\zeta _2,\zeta _3)`$. We can write the Cooper-Frye freeze-out distribution then as
$$\frac{dN_{CF}}{d^4xd^4p}=Nd^3\zeta \delta ^4(x\sigma (\zeta ))\delta ^4(pp(\sigma (\zeta ))).$$
(30)
As discussed in the Appendix, we can write the (on-shell) freeze-out distribution in terms of the solution of the Boltzmann equation as
$`E_1{\displaystyle \frac{dF_{fo}(x,𝐩_1)}{d^4xd^3p_1}}P_0(x,𝐩_1)\times [S(x,𝐩_1)+`$ (31)
$`+\mathrm{\hspace{0.17em}\hspace{0.17em}2}{\displaystyle \underset{3}{}}{\displaystyle \underset{4}{}}{\displaystyle \underset{5}{}}W_{3415}\delta ^4(p_3+p_4p_1p_5)f_3f_4].`$ (32)
While neither normalized nor unique, this expression provides at least a formal definition of the freeze-out distribution for the Boltzmann equation solely in terms of the phase space distribution $`f(x,𝐩)`$, and the assumed transition probabilities $`W_{ijkl}`$.
## V Numerical Results
### A Kinetic versus Hydrodynamic Evolution
To test the ideal hydrodynamical assumptions against transport theory, it is essential to eliminate as many model differences as possible. For example, both the hydrodynamic model and the kinetic theory should have the same degrees of freedom, the same equation of state, and the same initial conditions. Equation (2) describes a gas that in thermal equilibrium has the equation of state $`e=3p`$, if the partons are massless. We therefore used this ideal gas equation of state in the hydrodynamical simulations. We also chose the transport initial condition to be in local equilibrium, since hydrodynamics is limited to such initial conditions.
The hydrodynamic algorithm used is furthermore designed for particles without a conserved charge, i.e., the particle number changes as dictated by chemical equilibrium. The algorithm solves the energy-momentum conservation equation to obtain the energy density, pressure and flow evolution. Then, instead of the charge conservation equation, it exploits the relation between density, particle mass, and temperature in chemical equilibrium to compute the freeze-out particle distribution. It is important to note therefore that Eq. (2) with elastic collisions has the same hydrodynamic limit as the hydrodynamic model only if the partons are massless. This is because ideal hydrodynamics conserves entropy and for massless particles in thermal and chemical equilibrium entropy conservation is equivalent to particle number conservation.<sup>§</sup><sup>§</sup>§ Because in this case $`s=4n`$. For massive particles, we would have to compare transport to hydrodynamics with particle conservation. Conversely, we would need to supplement Eq. (2) to include inelastic channels, such as $`23`$ in Ref. , to compare to chemically equilibrated hydrodynamics. In the infinite rate limit we recover the hydrodynamic model even though we have a fixed number of particles. However, when the solution is out of equilibrium (either thermal, chemical, or both), it does make a difference whether we include particle number changing processes or not.
To test whether ideal hydrodynamics is an adequate description of the parton transport theory (2), we compare the evolution of the transverse energy $`dE_t/dy`$ at midrapidity from the two models. This comparison is free from any hydrodynamic freeze-out prescription because the transverse energy is given directly by the phase space distribution as
$`{\displaystyle \frac{dE_t}{dy}}|_\tau `$ $`=`$ $`\tau {\displaystyle d^2p_{}𝑑\eta d^2x_{}m_t\mathrm{cosh}(y\eta )}`$ (34)
$`\times m_tf(y,𝐩_{},\eta ,𝐱_{},\tau ),`$
where, through the local equilibrium ansatz (6), the hydrodynamic phase space evolution is determined by the evolution of the flow velocity and local temperature as dictated by the equations of motion (7).
Figure 1 shows the transverse energy evolution from transport theory and hydrodynamics, for an initial Bjorken cylinder radius of 2 fm, with $`\tau _0=0.1`$ fm$`/c`$, $`T_0=\mu =0.5`$ GeV, $`n_{\eta ,0}=2.6505`$ fm<sup>-2</sup> (via scaling), $`\sigma =`$ 15, and 60 mb. (We chose $`\eta _m=5`$, subdivisions 800 for 3+1D, 256 for 1+1D, and a $`100`$ fm<sup>2</sup> transverse area for the 1+1D evolution.)
The transverse energy decreases much faster from ideal hydrodynamics than from kinetic theory, both in 1+1D and 3+1D, showing that hydrodynamics does more work than the cascade. This is due to the different phase space evolution in the two models. The early discrepancy, even for cross sections as extreme as 15 or 60mb, indicates that either the transport evolution gets very quickly out of equilibrium, or the initial evolution is close to equilibrium but the energy-momentum tensor is not ideal. Note that even if the latter is true, it does not necessarily mean that this initial, locally equilibrated, nonideal dynamics can be described by the Navier-Stokes equations.
The above conclusion holds for any initial system size larger than 2fm as well. Since the 1+1D curves correspond to the infinite transverse size limit, the hydrodynamic and transport evolutions for initial sizes larger than 2fm will lie between the 2fm and the 1+1D curves for hydrodynamics and for transport theory, respectively. Because these two regions do not overlap, the discrepancy between ideal hydrodynamics and transport theory will not disappear with increasing system size.
### B Kinetic vs Hydrodynamic Freeze-out Results
In the previous Subsection we showed that parton kinetic theory does not reduce to ideal hydrodynamics for initial conditions at RHIC. Thus, the final observables from the two models can be similar only if the hydrodynamic freeze-out prescription helps mimic the observables from the nonequilibrium transport evolution.
Here we test whether one can reproduce the transport observables by a suitable choice of the hydrodynamic freeze-out parameters. We chose the widely-used Cooper-Frye freeze-out prescription (23) with isotherm freeze-out surfaces, despite all known problems discussed in Section IV. Hence, our only adjustable parameter is the freeze-out temperature. Since Eq. (2) describes Boltzmann classical particles, we must use the classical distribution (6) in the Cooper-Frye formula.
#### 1 Coordinate space evolution
Freeze-out distributions in space-time from MPC are shown in Figs. 3 and 5. Due to the assumed cylindrical symmetry and longitudinal boost invariance, that distribution is only a function of $`\tau `$ and $`R`$.
Figures 3, 5 show the freeze-out distribution for initial radii 6 fm and 2 fm, respectively, with $`\tau _0=0.1\mathrm{fm}/c`$, $`T_0=\mu =0.5`$ GeV, $`n_{\eta ,0}=2.6505`$ fm<sup>-2</sup> (via scaling), $`\sigma =`$ 3, 15, and 60 mb. For comparison, three different freeze-out isotherms are also shown from solution of Cooper-Frye frozen ideal hydrodynamics. (We chose $`\eta _m=5`$, subdivisions 800 for 2 fm, and 150 for 6 fm.)
Unlike the sharp hydrodynamic freeze-out surface, the freeze-out distribution from the cascade is a broad wedge. Particles originate from a hypervolume in space-time, rather than from a hypersurface. In the top left plot (3 mb, 6 fm) in Fig. 3, the wedge moved down to $`\tau =\tau _0`$, which is a general feature for very low reaction rates. In the limit of a vanishing reaction rate, all particles freeze out from $`\tau =\tau _0`$.
Figures 3, 5 show that particles freeze out later with increasing microscopic rates as expected. The maximum of the wedge moves outward with increasing rates, hence no freeze-out temperature can be universal. If we tune the freeze-out temperature to get as close as possible to the cascade freeze-out distribution, the freeze-out temperature will depend on the reaction rate.
Thus, the remarkable agreement seen in the bottom figure of Fig. 5 between Cooper-Frye frozen ideal hydrodynamics with a 130 MeV freeze-out temperature and the cascade with $`\sigma =60`$ mb is a mere coincidence; higher rates would lead to a later freeze-out. For very high reaction rates, the 130-MeV hypersurface from the cascade would be very close to that from hydrodynamics because the hydrodynamic evolution is the infinite reaction rate limit of the cascade evolution. But that does not mean that the freeze-out distributions are the same. On the contrary, if hydrodynamics and the cascade are close to each other at $`T=130`$ MeV then we have no justification to stop the hydrodynamic evolution and freeze out with Eq. (23) because we are still in equilibrium and particles will certainly collide in the future, i.e., they have not yet frozen out.
It is not possible to tune the Cooper-Frye freeze-out temperature to reproduce the cascade freeze-out distribution. Though the contour plots in Figs. 3, 5 suggest that such a tuning can get the hydrodynamic freeze-out curve close to the ridge of the wedge of the cascade freeze-out distribution, that is not enough. As the $`dN/d\tau `$ distributions show, the resulting hydrodynamic freeze-out distribution is not close to the cascade distribution because one has to reproduce not only the curve given by the ridge of the wedge but also the exact distribution along this curve.
If the freeze-out temperature is high enough to yield a freeze-out surface with a timelike portion, we get unphysical spikes in the freeze-out distribution that are not present in the cascade calculations. This can be seen in Fig. 5 for $`T_f=200`$ MeV, and in Fig. 3 for $`T_f=130`$ and 200 MeV. For example, for $`R_0=6`$ fm with $`T_f=130`$ MeV, Cooper-Frye frozen hydrodynamics produces most particles at around $`\tau =5.6`$ fm$`/c`$. Cooper-Frye frozen hydrodynamics produces most particles at around $`\tau =5.6`$ fm$`/c`$. This is because the inside of the cylinder follows a 1D Bjorken evolution with $`T(\tau )=T_0(\tau _0/\tau )^{1/3}`$ until the rarefaction wave from the boundary arrives. The rarefaction wave travels with a speed $`c_s=1/\sqrt{3}`$. If the system is large enough, most of the system reaches the freeze-out temperature before the rarefaction wave arrives, i.e., during the 1+1D Bjorken evolution. With our parameters $`T_0=0.5`$ GeV, $`T_{fo}=130`$ MeV, and $`\tau _0=0.1\mathrm{fm}/c`$, this gives a freeze-out for the inside of the cylinder at $`\tau _{fo}=5.6\mathrm{fm}/c`$, which is in complete disagreement with our transport theory solutions. Furthermore, it does not correspond to the infinite reaction rate limit either because in that case particles freeze out very late.
Hence the peaks in $`dN/d\tau `$ at $`\tau =5.6`$ fm$`/c`$ ($`T_f=130`$ MeV) and $`\tau =1.6`$ fm$`/c`$ ($`T_f=200`$ MeV) are a clear consequence of the arbitrary freeze-out prescription using Eq. (23) with isotherm freeze-out hypersurfaces. Smearing the peaks out around their maxima does not help either because that does not change the location of the peaks, while the maximum from the cascade moves outward with increasing reaction rates.
#### 2 Momentum space
The freeze-out distribution in momentum space is shown in Figs. 2 and 6. Figure 6 shows the freeze-out $`p_{}`$-distribution for initial radii 2 fm and 6 fm, cascade cross sections 3, 15, and 60 mb compared to ideal hydrodynamics with a Cooper-Frye freeze-out at temperatures $`T_f=100`$, 130, and 200 MeV. As the reaction rate increases, the small $`p_{}`$-slopes rise as the system cools due to longitudinal work. The $`p_{}`$-distribution seems to approach that of Cooper-Frye frozen hydrodynamics. However, this is only an illusion on a low-resolution logarithmic plot. Figure 2, where we plotted the final $`p_{}`$-spectra divided by the initial $`T_0=500`$ MeV thermal one, shows that there is a large, up to a factor of ten difference at both low ($`<0.5`$ GeV) and high $`p_{}`$ ($`>2`$ GeV), depending on the microscopic rates. For all the cases studied, Cooper-Frye frozen hydrodynamics has more low-$`p_{}`$ particles but fewer high-$`p_{}`$ ones than the cascade. This is not necessarily a general feature because the assumed hydrodynamic freeze-out temperature is an arbitrary number. A later freeze-out (lower temperature) gives a larger slope, an earlier freeze-out (higher temperature) gives a smaller one.
It is also striking that one would need rather high, $`T_f300450`$ MeV freeze-out temperatures to get closer to the cascade $`p_{}`$-spectra. We conclude that it is not possible to reproduce both the space-time and the momentum space transport theory freeze-out distributions using ideal hydrodynamics with the isotherm Cooper-Frye freeze-out prescription. Either one needs to treat hydrodynamic freeze-out more accurately than the Cooper-Frye prescription, or one needs to use full-scale transport theory instead of ideal hydrodynamics. The present work is a step in the latter direction, while Refs. \[21-23\] are important steps in the former direction looking for a simplification of the full transport theoretical problem that, hopefully, will still be applicable to a wide class of situations.
## VI Outlook
There are many open problems in the development of covariant transport theory. The most urgent need is to develop practical convergent algorithms to incorporate inelastic $`23`$ processes to allow studies of chemical equilibration. Preliminary work in Ref. indicated a rather slow convergence towards Lorentz covariance with particle subdivision. Unlike the $`l^{1/2}`$ convergence in $`22`$, a much slower $`l^{1/5}`$ convergence is expected in $`23`$ processes even when nonlocal formation physics ($`\mathrm{\Delta }t>\mathrm{}/\mathrm{\Delta }E`$) is neglected.
Also, we note that all results in this paper pertain to homogeneous initial conditions. In Ref. , it was shown that jets induce large nonstatistical local fluctuations that may evolve in a turbulent manner. A transport study of the evolution from such inhomogeneous initial conditions would be useful to compare to the known hydrodynamic solutions.
## VII Acknowledgments
We are grateful to Yang Pang and Bin Zhang for extensive discussions on transport theory and numerical cascade algorithms and their contributions to the OSCAR effort. We are also grateful to Adrian Dumitru and Dirk Rischke for discussions on hydrodynamics and use of their codes. We acknowledge useful discussions with László P. Csernai and George Bertsch on hydrodynamics and freeze-out, and the Parallel Distributed Systems Facility at the National Energy Research Scientific Computing Center for providing computing resources.
This work was supported by the Director, Office of Energy Research, Division of Nuclear Physics of the Office of High Energy and Nuclear Physics of the U.S. Department of Energy under contract No. DE-FG-02-93ER-40764.
## Formal definition for freeze-out
Unlike in the cascade solution where the freeze-out distribution is trivially defined by Eq. (28), in the Boltzmann equation $`f`$ changes continuously and no discrete final collisions can be identified. In this appendix we propose a generalization of Eq. (24) which is independent of the discrete numerical cascade picture. We motivate here a formal definition, Eq. (32), of the freeze-out distribution using solely $`f(x,𝐩)`$ and $`W_{ijkl}`$.
Following the notion of the “last collision”, one can first compute the probability that a particle starting at a coordinate $`x_1^\mu `$ with momentum $`𝐩_1`$ does not collide any further. The collision rate is given by
$`\mathrm{\Gamma }_{coll}`$ $``$ $`{\displaystyle \frac{dN_{coll}}{d^4x}}(x,𝐩_1,𝐩_2)`$ (35)
$`=`$ $`f_1(x,𝐩_1)f_2(x,𝐩_2)\sigma (p_1,p_2)v_{12}d^3p_1d^3p_2,`$ (36)
where the relative velocity and the total cross section are given with the Lorentz scalar
$$t_{12}\sqrt{(p_1^\mu p_{2}^{}{}_{\mu }{}^{})^2m_1^2m_2^2}$$
(37)
as
$`v_{12}={\displaystyle \frac{t_{12}}{E_1E_2}},`$ (38)
$`\sigma (p_1,p_2)={\displaystyle \frac{1}{t_{12}}}{\displaystyle \underset{3}{}}{\displaystyle \underset{4}{}}W_{1234}\delta ^4(p_1+p_2p_3p_4).`$ (39)
A free pointlike particle has the phase space distribution
$$f_1(x,p)=_0^{\mathrm{}}𝑑\tau \delta ^4(xx_1u_1\tau )\delta ^4(pp_1),$$
(40)
where for an on-shell particle
$$\delta \left(p^0\sqrt{𝐩^2+m^2}\right)=\delta (p^2m^2)2\sqrt{m^2+𝐩^2}\mathrm{\Theta }(p_0),$$
i.e.,
$$\delta ^4(pp_1)=\delta ^3(𝐩𝐩_1)\delta (p^2m^2)2\sqrt{m^2+𝐩_1^2}\mathrm{\Theta }(p_0)$$
and thus Recall, the on-shell phase space distribution is defined via
$$f(x,p)2m\delta (p^2m^2)\mathrm{\Theta }(p_0)f(x,𝐩).$$
(41)
$`f_1(x,𝐩)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\tau \delta ^4(xx_1u_1\tau )`$ (43)
$`\times \delta ^3(𝐩𝐩_1){\displaystyle \frac{\sqrt{m^2+𝐩_1^2}}{m}}.`$
Plugging this result into Eq. (36), the probability that a free particle will not have any further collisions is This can be derived assuming that subsequent collisions are uncorrelated (just like the similar formula for the inhomogeneous Poisson distribution).
$`P_0(x_1,𝐩_1)=\mathrm{exp}\left({\displaystyle \mathrm{\Gamma }_{coll}d^3p_1d^3p_2d^4x}\right)`$ (44)
$`=\mathrm{exp}\left({\displaystyle \frac{d\tau d^3p_2}{E_2m}f_2(x_1+u_1\tau ,𝐩_2)\sigma (p_1,p_2)t_{12}}\right).`$ (45)
Now we can write the freeze-out distribution as the number of particles having a collision at $`x^\mu `$ with outgoing momentum $`𝐩_1`$ times the probability that these particles do not collide any further, i.e.,
$`E_1{\displaystyle \frac{dF_{fo}^{coll}(x,𝐩_1)}{d^4xd^3p_1}}P_0(x,𝐩_1)`$ (46)
$`\times \mathrm{\hspace{0.17em}\hspace{0.17em}2}{\displaystyle \underset{3}{}}{\displaystyle \underset{4}{}}{\displaystyle \underset{5}{}}W_{3415}\delta ^4(p_3+p_4p_1p_5)f_3f_4.`$ (47)
This definition does not include those particles that are formed but suffer no collisions afterward. Their contribution is
$$E_1\frac{dF_{fo}^{form}(x,𝐩_1)}{d^4xd^3p_1}S(x,𝐩_1)P_0(x,𝐩_1).$$
(48)
Hence the final freeze-out distribution is given by Eq. (32).
The definition (32) should be regarded only one measure of the freeze-out distribution because it has several a shortcomings. The probabilities summed are not probabilities for disjoint events. One should exclude the volume in space time given by all the linear paths of the already frozen-out particles. This requires knowledge of multiparticle correlations beyond the scope of the Bolztmann equation. As long as those excluded volume effects are small, Eq. (32) is adequate. A clear problem with the present formal definition is that particle number and momentum are not conserved by it as is automatic in (28). It is interesting to contrast on the other hand, the trivial way that the cascade solves this problem through Eq. (28). In cascade, the $`N`$-body correlations are automatically calculated and freeze-out is easily defined conserving number and total four-momentum. The continuum limit is thus subtle. Our numerical results define that continuum limit as the limit of infinite subdivisions using the cascade technique.
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# Summing Sudakov logarithms in 𝐵→𝑋_𝑠𝛾 in effective field theory
## I Introduction
Effective field theories (EFT’s) provide a simple and elegant method for calculating processes with several relevant energy scales. Part of the utility of EFT’s is that they dramatically simplify the summation of powers of logarithms of ratios of mass scales, which would otherwise make perturbation theory poorly behaved. For example, in a theory with a very heavy particle of mass $`M`$, one-loop corrections will typically be enhanced by $`\mathrm{log}(M/\lambda )`$, where $`\lambda `$ is a low scale in the problem. In the EFT in which the heavy particle has been removed from the theory, such logarithms are replaced by factors of $`\mathrm{log}(\mu /\lambda )`$ (where $`\mu `$ is the renormalization scale in dimensional regularization, or the cutoff in cutoff regularization), and the complete series of leading logarithms $`\alpha _s^n\mathrm{log}^n(\mu /\lambda )`$ is straightforward to sum via the renormalization group.
The situation is more complicated for processes with highly energetic light particles. In this case, there are both collinear and infrared divergences in the theory, which give rise to the familiar Sudakov double logarithms . For example, the perturbative expansion of the $`N`$’th moment of the photon spectrum in inclusive $`bX_s\gamma `$ decay is of the form
$$\underset{n}{}\underset{m2n}{}C_{n,m}\alpha _s^n\mathrm{log}^mN.$$
(1)
Although the arguments of these logarithms are not obviously the ratio of two scales, they arise because the typical energy and invariant mass of light particles are widely separated, and they may be summed via well-known techniques based on factorization theorems into the form
$$\mathrm{exp}\left[\underset{n}{}\left(a_n\alpha _s^n\mathrm{log}^{n+1}N+b_n\alpha _s^n\mathrm{log}^nN\right)+\mathrm{}\right].$$
(2)
The terms $`\alpha _s^n\mathrm{log}^{n+1}N`$ are referred to as the leading logarithmic contribution, the terms $`\alpha _s^n\mathrm{log}^nN`$ are referred to as the next-to-leading logarithmic contribution, and the remaining terms are called subdominant.
Recently there has been some discussion in the literature of summing Sudakov logarithms using effective field theory techniques . Such an approach could have several advantages over the conventional method; in particular, while factorization formulas are based on perturbation theory, EFTs, by construction, are valid beyond perturbation theory, and by including higher dimension operators it should be straightforward (if tedious) to go beyond the leading twist approximation. In the various versions of the EFT approach which have been suggested, the effective theory is the so-called “Large Energy Effective Theory” (LEET) , which describes light-like particles coupled to soft degrees of freedom. However, a difficulty with the approaches presented to date is that, as pointed out in Refs. , in the minimal subtraction (MS) scheme logarithms arising at one loop in LEET do not match logarithms arising at one loop in QCD for any choice of the matching scale $`\mu `$; hence these logarithms may not be summed using the RGE’s.
In this paper we consider this problem in the context of $`BX_s\gamma `$ decays.<sup>§</sup><sup>§</sup>§In fact, the authors of argued that the resummation of subleading Sudakov logarithms is not necessary for practical purposes for this decay. Nevertheless, it provides a simple example in which we may compare our results to those in the literature. We show that the problem of matching scales may be resolved by introducing a new intermediate effective theory containing both soft and collinear degrees of freedom, which is then matched onto LEET (effectively integrating out the collinear modes) at a lower scale. We show that the matching conditions onto both effective theories contain no large logarithms at one loop. We then calculate the RGE’s in the two theories summing the leading logarithms and a certain subset of the next-to-leading logarithms. To this order the expression obtained for the resummed Sudakov logarithms is identical to that derived in Refs .
## II Sudakov Logarithms in $`BX_S\gamma `$ and LEET
Inclusive decays of heavy quarks have been well understood for many years in the context of an operator product expansion (OPE) in the inverse mass of the heavy quark. At leading order in the $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$ expansion the $`B`$ meson decay rate is equal to the $`b`$ quark decay rate, and nonperturbative effects are suppressed by at least two powers of $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$. However, the OPE only converges for sufficiently inclusive observables. Unfortunately, experimental cuts on measurements of rare decays such as $`BX_s\gamma `$, $`BX_s\mathrm{}^+\mathrm{}^{}`$, and $`BX_u\mathrm{}\overline{\nu }`$ are required, restricting the available phase space considerably. Since all of these decays are of phenomenological interest, either in the determination of $`|V_{ub}|`$ or detection of new physics, understanding inclusive decays in restricted regions of phase space is important.
If the phase space is restricted such that the final hadronic state is dominated by only a few resonances, the breakdown of the OPE simply reflects the fact that an inclusive treatment based on local duality is no longer appropriate. This is the case for the dilepton invariant mass spectrum in inclusive $`BX_s\mathrm{}^+\mathrm{}^{}`$ and $`BX_u\mathrm{}\overline{\nu }`$ decays . However, when the kinematic cut is in a region of phase space dominated by highly energetic, low invariant mass final states, the OPE breaks down even for quantities smeared over a parametrically larger region of phase space, where the decay is not resonance dominated. This situation arises in the endpoint region of the electron energy spectrum and the low hadronic invariant mass region in semileptonic $`BX_u\mathrm{}\overline{\nu }`$ decay, as well as the endpoint region of the photon spectrum in $`BX_s\gamma `$ decay.
Consider the dominant contribution to the decay $`BX_s\gamma `$, which arises from the magnetic penguin operator
$$\widehat{O}_7=\frac{e}{16\pi ^2}m_b\overline{s}\sigma ^{\mu \nu }\frac{1}{2}(1+\gamma _5)bF_{\mu \nu },$$
(3)
where the strange quark mass has been set to zero.Throughout this work we will ignore the contribution of operators other than $`\widehat{O}_7`$ to the decay. The OPE for this decay is illustrated in Fig. 1. We write the momenta of the $`b`$ quark, photon, and light $`s`$ quark jet as
$$p_b^\mu =m_bv^\mu +k^\mu ,q^\mu =\frac{m_b}{2}x\overline{n}^\mu ,p_s^\mu =\frac{m_b}{2}n^\mu +l^\mu +k^\mu $$
(4)
where, in the rest frame of the $`B`$ meson,
$$v^\mu =(1,\stackrel{}{0}),n^\mu =(1,0,0,1),\overline{n}^\mu =(1,0,0,1).$$
(5)
Here $`k^\mu `$ is a residual momentum of order $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, and $`l^\mu =\frac{m_b}{2}(1x)\overline{n}^\mu `$, where $`x=2E_\gamma /m_b`$. The invariant mass of the light $`s`$-quark jet
$$p_s^2m_bn(l+k)=m_b^2(1x+\widehat{k}_+),$$
(6)
(where $`\widehat{k}_+=k_+/m_b`$) is $`O(m_b^2)`$ except near the endpoint of the photon energy spectrum where $`x1`$. Inclusive quantities are calculated via the OPE by taking the imaginary part of the graphs in Fig. 1 and expanding in powers of $`k^\mu /\sqrt{p_s^2}`$. As long as $`x`$ is not too close to the endpoint, this is an expansion in powers in $`k^\mu /m_b`$, which matches onto local operators.
This leads to an expansion for the photon energy spectrum as a function of $`x`$ in powers of $`\alpha _s`$ and $`1/m_b`$ :
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}`$ $`=`$ $`\mathrm{\Gamma }_0\{[1{\displaystyle \frac{\alpha _sC_F}{4\pi }}(2\mathrm{log}{\displaystyle \frac{\mu ^2}{m_b^2}}+5+{\displaystyle \frac{4}{3}}\pi ^2)]\delta (1x)`$ (10)
$`+{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left[7+x2x^22(1+x)\mathrm{log}(1x)\left(4{\displaystyle \frac{\mathrm{log}(1x)}{1x}}+{\displaystyle \frac{7}{1x}}\right)_+\right]`$
$`+{\displaystyle \frac{1}{2m_b^2}}[(\lambda _19\lambda _2)\delta (1x)(\lambda _1+3\lambda _2)\delta ^{}(1x){\displaystyle \frac{\lambda _1}{3}}\delta ^{^{\prime \prime }}(1x)]\}`$
$`+O(\alpha _s^2,1/m_b^3),`$
where
$$\mathrm{\Gamma }_0=\frac{G_F^2|V_{tb}V_{ts}^{}|^2\alpha |C_7(\mu )|^2}{32\pi ^4}m_b^5\left[\frac{m_b(\mu )}{m_b}\right]^2,$$
(11)
and the subscript “+” denotes the usual plus distribution,
$`{\displaystyle \frac{1}{(1x)_+}}`$ $``$ $`\underset{\beta 0}{lim}\left\{{\displaystyle \frac{1}{1x}}\theta (1x\beta )+\mathrm{log}(\beta )\delta (1x\beta )\right\}`$ (12)
$`\left({\displaystyle \frac{\mathrm{log}(1x)}{(1x)}}\right)_+`$ $``$ $`\underset{\beta 0}{lim}\left\{{\displaystyle \frac{\mathrm{log}(1x)}{1x}}\theta (1x\beta )+{\displaystyle \frac{1}{2}}\mathrm{log}^2(\beta )\delta (1x\beta )\right\}.`$ (13)
The parameters $`\lambda _1`$ and $`\lambda _2`$ are matrix elements of local dimension five operators.
Near the endpoint of the photon spectrum, $`x1`$, both the perturbative and nonperturbative corrections are singular and the OPE breaks down. The severity of the breakdown is most easily seen by integrating the spectrum over a region $`1\mathrm{\Delta }<x<1`$. When $`\mathrm{\Delta }\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$ the most singular terms in the $`1/m_b`$ expansion sum up into a nonperturbative shape function of characteristic width $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$. The perturbative series is of the form
$$\frac{1}{\mathrm{\Gamma }_0}_{1\mathrm{\Delta }}^1\frac{d\mathrm{\Gamma }}{dx}=1+\frac{\alpha _sC_F}{4\pi }\left(2\mathrm{log}^2\mathrm{\Delta }7\mathrm{log}\mathrm{\Delta }+\mathrm{}\right)+O\left(\alpha _s^2\right),$$
(14)
where the ellipses denote terms that are finite as $`\mathrm{\Delta }0`$. These Sudakov logarithms are large for $`\mathrm{\Delta }1`$, and can spoil the convergence of perturbation theory. The full series has been shown to exponentiate and the leading and next-to-leading logarithms must be resummed for $`\mathrm{\Delta }\mathrm{exp}\left(\sqrt{\pi /\alpha _s(m_b)}\right)`$, which is parametrically larger than $`\mathrm{\Lambda }_{\mathrm{QCD}}/m_b`$ in the $`m_b\mathrm{}`$ limit.
In general, “phase space” logarithms are to be expected whenever a decay depends on several distinct scales. For example, in $`bX_ce\overline{\nu }_e`$ decay the rate calculated with the OPE performed at $`\mu =m_b`$ contains large logarithms of $`m_c/m_b`$. In an EFT was used to run from $`m_b`$ to $`m_c`$, summing phase space logarithms of the ratio $`m_c/m_b`$. Similarly, in $`bX_s\gamma `$ near the endpoint of the photon energy spectrum the invariant mass of the light quark jet scales as $`m_b\sqrt{1x}`$, and is widely separated from the scale $`\mu =m_b`$ where the OPE is performed. In order to sum logarithms of $`\mathrm{\Delta }`$ (or the more complicated plus distributions in the differential spectrum, Eq. (10)) we would expect to have to switch to a new effective theory at $`\mu =m_b`$, use the renormalization group to run down to a scale of order $`m_b\sqrt{1x}`$, at which point the OPE is performed. (In fact, we will see that the situation is slightly more complicated than this).
We are then left with the question of the appropriate theory below the scale $`m_b`$. The simplest possibility is to expand the theory in powers of $`k^\mu /m_b`$ and $`l^\mu /m_b`$. The heavy quark is then treated in the heavy quark effective theory (HQET) , while the light quark propagator is treated in the large energy effective theory (LEET) proposed many years ago by Dugan and Grinstein. Expanding the $`s`$ quark propagator in powers of $`1/m_b`$, we find the LEET propagator
$$\frac{ip/_s}{p_s^2}=\frac{n/}{2}\frac{i}{n(l+k)}+O\left(\frac{l^\mu +k^\mu }{m_b}\right).$$
(15)
LEET is an effective theory of lightlike Wilson lines, much as HQET is an effective theory of timelike Wilson lines. The hope would then be to match QCD onto LEET and then use the renormalization group to sum the Sudakov logarithms. This is the approach taken in . However, a simple attempt at matching shows that this does not sum the appropriate logarithms.
Consider the one-loop matching of the operator $`\widehat{O}_7`$ from QCD to LEET. We regulate ultraviolet (UV) divergences with dimensional regularization ($`d=42ϵ`$). We introduce a small invariant mass $`p_s^2`$ for the $`s`$ quark which regulates all infrared (IR) divergences except that in the heavy-quark wave function diagram, Fig. 2(b). This IR divergence is regulated using dimensional regularization. The vertex diagram, Fig. 2(a) yields
$$A_{QCD}^{(a)}=C_7(\mu )\overline{s}\mathrm{\Gamma }^\mu b\frac{\alpha _sC_F}{4\pi }\left[\mathrm{log}^2\frac{p_s^2}{m_b^2}+2\mathrm{log}\frac{p_s^2}{m_b^2}+\mathrm{}\right],$$
(16)
where
$$\mathrm{\Gamma }^\mu =\frac{e}{8\pi ^2}m_b\sigma ^{\mu \nu }\frac{(1+\gamma _5)}{2}q_\nu .$$
(17)
$`C_7(\mu )`$ is the Wilson coefficient of $`\widehat{O}_7`$ and the dots denote (here and in the rest of the paper) finite terms which are not logarithmically enhanced. Including a factor of $`\sqrt{Z}`$ for each external field
$`Z_b`$ $`=`$ $`1{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left[{\displaystyle \frac{3}{ϵ}}+3\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{}\right]`$ (18)
$`Z_s`$ $`=`$ $`1{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left[{\displaystyle \frac{1}{ϵ}}\mathrm{log}{\displaystyle \frac{p_s^2}{m_b^2}}+\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{}\right],`$ (19)
where
$$\stackrel{~}{\mu }^24\pi \mu ^2e^{\gamma _E}$$
(20)
and adding the counterterm required to subtract off the UV divergence
$$Z_7=1+\frac{\alpha _sC_F}{4\pi }\frac{1}{ϵ}$$
(21)
we find
$$A_{QCD}=C_7(\mu )\overline{s}\mathrm{\Gamma }^\mu b\left[1\frac{\alpha _sC_F}{4\pi }\left(\mathrm{log}^2\frac{p_s^2}{m_b^2}+\frac{3}{2}\mathrm{log}\frac{p_s^2}{m_b^2}+\frac{1}{ϵ}+2\mathrm{log}\frac{\stackrel{~}{\mu }^2}{m_b^2}+\mathrm{}\right)\right].$$
(22)
The corresponding LEET diagram is shown in Fig. 3. Neither of the wave function graphs gives a contribution, since the light quark wave function in Feynman gaugeWe will work in Feynman gauge throughout this paper. is proportional to $`n^2=0`$, and the heavy quark wave function vanishes in dimensional regularization. Thus the only contribution is from the vertex graph. Denoting the coefficient of the corresponding operator in LEET as $`C^{(0)}(1+(\alpha _sC_F/4\pi )C^{(1)}+\mathrm{})`$, we find
$`A_{LEET}`$ $`=`$ $`C^{(0)}(\stackrel{~}{\mu })\overline{\xi }_n\mathrm{\Gamma }^\mu h\left[1{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left(\left(4\pi {\displaystyle \frac{\mu ^2m_b^2}{p_s^4}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }(1+2ϵ)\mathrm{\Gamma }(12ϵ)\mathrm{\Gamma }(1+ϵ)}{ϵ^2}}C^{(1)}(\stackrel{~}{\mu })\right)\right]`$ (23)
$`=`$ $`C^{(0)}(\stackrel{~}{\mu })\overline{\xi }_n\mathrm{\Gamma }^\mu h\left[1{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left({\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{2}{ϵ}}\mathrm{log}{\displaystyle \frac{p_s^2}{m_b\stackrel{~}{\mu }}}+2\mathrm{log}^2{\displaystyle \frac{p_s^2}{m_b\stackrel{~}{\mu }}}+\mathrm{}C^{(1)}(\stackrel{~}{\mu })\right)\right],`$ (24)
where $`p_s^2/m_b`$ is the soft scale. So
$`C^{(0)}(\stackrel{~}{\mu })`$ $`=`$ $`C_7(\stackrel{~}{\mu })`$ (25)
$`C^{(1)}(\stackrel{~}{\mu })`$ $`=`$ $`{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{1}{ϵ}}\left(2\mathrm{log}{\displaystyle \frac{p_s^2}{m_b\stackrel{~}{\mu }}}+1\right)+2\mathrm{log}^2{\displaystyle \frac{p_s^2}{m_b\stackrel{~}{\mu }}}\mathrm{log}^2{\displaystyle \frac{p_s^2}{m_b^2}}{\displaystyle \frac{3}{2}}\mathrm{log}{\displaystyle \frac{p_s^2}{m_b^2}}`$ (27)
$`2\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{}.`$
We immediately notice two problems<sup>\**</sup><sup>\**</sup>\**Note that $`C_7(\mu )`$ includes a factor of $`\alpha _sC_F/(4\pi )\mathrm{log}(m_W/\mu )`$, which converts one of the factors of $`\mathrm{log}(\mu /m_b)`$ in Eq. (16) to $`\mathrm{log}(m_W/m_b)`$. This is not important for our argument.:
1. There is no matching scale $`\stackrel{~}{\mu }`$ at which all the large single and double logarithms in $`C^{(1)}`$ vanish. Thus, there are logarithms in the rate which cannot be summed using the renormalization group in LEET.
2. $`C^{(1)}`$ contains a divergence proportional to $`\frac{1}{ϵ}\mathrm{log}p_s^2`$. Since $`p_s^2`$ is an infrared scale in the problem, it is not clear how to sensibly renormalize this term. In Ref. this divergence was cancelled by a nonlocal counterterm in the inclusive rate; however, this term indicates that LEET cannot be used for exclusive processes . Furthermore, the matching of the inclusive rate performed in still leaves large logarithms in the coefficient of the operator.
The problem is that LEET only describes the coupling of light-like particles to soft gluons, but does not describe the splitting of an energetic particle into two almost collinear particles. Thus, by matching onto LEET one is integrating out the collinear modes which also contribute to infrared physics. As we will show below, once collinear degrees of freedom are included, both of the above problems are resolved.
## III The Collinear-Soft Theory
### A Collinear and soft modes
It is convenient to work in light-cone coordinates $`p^\mu =(p^+,p^{},p_{}^i)`$, where $`p^+=np`$ and $`p^{}=\overline{n}p`$, and to define a power-counting parameter $`\lambda =\sqrt{1x}`$ that becomes small in the limit $`x1`$. The momentum of the light-quark jet then scales as
$$p_s^\mu m_b(\lambda ^2,1,\lambda ).$$
(28)
This scaling is unchanged by emission of either soft or collinear degrees of freedom, with momenta scaling as
$$p_{\mathrm{soft}}m_b(\lambda ^2,\lambda ^2,\lambda ^2),p_{\mathrm{collinear}}m_b(\lambda ^2,1,\lambda ),$$
(29)
and so emission of both modes is kinematically allowed. It is the presence of infrared sensitive graphs with collinear loop momentum that makes this EFT more complicated than other, more familiar, EFT’s, where infrared sensitivity comes purely from soft modes. This is similar to the situation in non-relativistic QCD (NRQCD), in which power counting is complicated by the fact that a given amplitude receives contributions from loop momenta which are small compared to the heavy quark mass, but which have parametrically different dependence on the heavy quark velocity $`v`$. In NRQCD, the relevant scales are known as soft, ultrasoft and potential, and must be treated separately in order to obtain consistent power counting.
We follow a similar approach here, and introduce separate fields for both soft and collinear degrees of freedom<sup>††</sup><sup>††</sup>††At two loops an additional gluon field scaling as $`(\lambda ,\lambda ,\lambda )`$ might have to be included .. Between the scales $`m_b`$ and $`m_b\lambda `$ the effective theory contains separate fields for both collinear and soft modes, while at scales below $`m_b\lambda `$ (the exact scale depends on the operator under consideration, as will be discussed in the next section), the collinear modes are integrated out of the theory and it is matched onto LEET. We will refer to this intermediate theory as the collinear-soft theory, and resist the urge to create another acronym.
There is an important difference between the approach taken here and the one taken in Refs. where logarithms of $`v`$ are summed in NRQCD and NRQED. In the latter case no intermediate theory is introduced; instead the running is performed in one step through the velocity RGE. In NRQED these two approaches differ at subleading order , and it may be that such one-step running is needed here at two loops.
The power counting rules in the collinear-soft theory may be obtained by a field rescaling, analogous to that performed in . The scaling of the fields is chosen such that the propagators are all $`O(1)`$, putting the $`\lambda `$ dependence into the interaction terms. For example, in the kinetic term for a soft gluon,
$$d^4x(_\mu A_\nu ^a_\nu A_\mu ^a)^2$$
(30)
the typical length scale associated with soft excitations scales as $`\lambda ^2p_{\mathrm{soft}}^1`$, so the factor of $`d^4x`$ scales as $`\lambda ^8`$. Each derivative scales as $`p_{\mathrm{soft}}\lambda ^2`$, so the soft gluon field must scale as $`\lambda ^2`$ for the kinetic term to be $`O(1)`$.
Since the various collinear momentum components scale differently with $`\lambda `$, power counting for collinear gluons is gauge dependent (this is easily seen from the propagator, since in a covariant gauge the components of the $`k^\mu k^\nu `$ term scale differently). In this paper we are working in Feynman gauge, in which case the different components of collinear gluons have the same scaling. Performing a similar analysis for the other fields, we obtain the power-counting rules given in Table I.
Rather than write down the effective Lagrangian for the various fields, which is quite lengthy, we will instead just give the Feynman rules, which are obtained by expanding the QCD amplitudes in powers of $`\lambda `$. The spinors in the collinear-soft theory are related, at leading order in $`\lambda `$, to the QCD spinors via
$$h_v=P_+u,\xi _n=P_nu,\xi _{\overline{n}}=P_{\overline{n}}u,$$
(31)
where we have defined the projection operators
$$P_+=\frac{v/+1}{2},P_n=\frac{n/\overline{n}/}{4},P_{\overline{n}}=\frac{\overline{n}/n/}{4},$$
(32)
which project out the heavy quark spinor, a massless spinor in the $`n`$ direction, and a massless spinor in the $`\overline{n}`$ direction respectively.
The propagators for the different fields are shown in Fig. 4.
The interactions leading in $`\lambda `$ which we will need in this paper are shown along with their Feynman rules and scaling in Fig. 5.
Note that the interaction of a soft particle with a collinear particle leaves the minus and perpendicular momenta of the collinear particle unchanged, since they are parametrically larger for the collinear particle. This is analogous to the multipole expansion which is performed in NRQCD. As a result at one loop, soft-collinear interactions in this theory are equivalent to LEET, since collinear propagators in soft loops reduce to LEET propagators:
$$\frac{n/}{2}\frac{\overline{n}(pk)}{(pk)^2}\frac{n/}{2}\frac{p^{}}{(pk)^+p^{}(p^{})^2}=\frac{n/}{2}\frac{1}{nk}$$
(33)
where $`p`$ is a collinear momentum, $`k`$ is a soft momentum, and $`p^2=0`$ from the equations of motion. Once again, this is analogous to NRQCD, where in ultrasoft loops the Feynman rules reduce to those for HQET. By the same token, in soft-collinear interactions, the appropriate volume element is the collinear volume element, scaling as $`\lambda ^4`$.
Because the leading purely collinear interaction, Fig. 5(a), scales as $`\lambda ^1`$, power counting for collinear loops is less simple than for soft loops. Terms which would scale as $`\lambda ^2`$, such as the purely collinear wave-function graph in Fig. 6, are proportional to $`n^2=0`$ and so vanish in the effective theory.
However, the $`1/\lambda `$ coupling enhances terms which would naïvely be suppressed. In fact, although the $`\lambda `$-counting looks complicated, graphs with only collinear lines are identical to the corresponding graphs in QCD. This is because in any graph in which all the lines have the same scaling (and there are no purely soft graphs, so this only refers to purely collinear graphs), expanding in powers of $`\lambda `$ does not change the propagators. Since the locations of poles in the propagators are unaffected, it is irrelevant whether one calculates the full graph in QCD and then expands in powers of $`\lambda `$, or calculates each order in $`\lambda `$ in the collinear-soft theory. Thus, for purely collinear graphs, such as the wave function graph in Fig. 6, we will not bother to write down the complete set of operators, but simply calculate the graph in QCD and expand.
There is one important subleading operator, shown in Fig. 7, which can be enhanced by the $`1/\lambda `$ piece of the purely collinear coupling. By momentum conservation, there is no vertex coupling two heavy quarks and a collinear quark, since a heavy quark cannot emit a collinear gluon and stay on its mass shell. However, expanding the diagram in Fig. 7 in powers of $`\lambda `$ gives the nonlocal $`O(\lambda )`$ interaction shown in the figure. (This is similar to the nonlocal operators found in ). Though it is formally subleading, in graphs such as Fig. 8(a) it gives an $`O(1)`$ effect.
### B Matching onto the Collinear-Soft Theory
We now proceed to compute the matching conditions for the operator $`\widehat{O}_7`$, and demonstrate that there are no large logarithms in the matching coefficients. At tree level, the matching is trivial. Defining the current in the effective theory by
$$V^\mu =\overline{\xi }_n\mathrm{\Gamma }^\mu h_v,$$
(34)
where $`\mathrm{\Gamma }^\mu `$ is given in (17), the Wilson coefficient $`C_V`$ at tree level is
$$C_V=1+O(\alpha _s).$$
(35)
To perform this matching at one-loop, we repeat the one-loop matching calculation discussed in Section II, but now using the collinear-soft theory instead of LEET, hence including collinear modes. The calculation is simplest if we set the invariant mass of the $`s`$ quark to zero; this introduces additional infrared divergences to the calculation which cancel in the matching conditions. The one loop matrix element of $`\widehat{O}_7`$ in full QCD can be calculated from the diagrams in Fig. 2, and we find the amplitude
$$A_{\mathrm{QCD}}=\overline{s}\mathrm{\Gamma }^\mu b\left[1\frac{\alpha _sC_F}{4\pi }\left(\frac{1}{ϵ^2}+\frac{\mathrm{log}(\stackrel{~}{\mu }^2/m_b^2)}{ϵ}+\frac{5}{2ϵ}+\frac{1}{2}\mathrm{log}^2\frac{\stackrel{~}{\mu }^2}{m_b^2}+\frac{7}{2}\mathrm{log}\frac{\stackrel{~}{\mu }^2}{m_b^2}+\mathrm{}\right)\right].$$
(36)
where all the $`1/ϵ`$ divergences are infrared in origin. The one loop correction in the collinear-soft theory can be calculated from the Feynman diagrams in Figs. 3 and 8. In pure dimensional regularization all graphs are zero, as there is no scale present in the loop integrals. Thus, we find the matching condition
$$C_VZ_V=1+\frac{\alpha _sC_F}{4\pi }\left(\frac{1}{ϵ^2}+\frac{\mathrm{log}(\stackrel{~}{\mu }^2/m_b^2)}{ϵ}+\frac{5}{2ϵ}+\mathrm{}\right),$$
(37)
where $`Z_V`$ is the counterterm required to subtract the UV divergences in the collinear-soft theory.
This derivation of course assumes that the collinear-soft theory reproduces the infrared behaviour of QCD. We can check this by instead introducing a small invariant mass for the $`s`$ quark, as in Section II, and explicitly verifying that the dependence on the invariant mass in the collinear-soft theory is identical to that in full QCD given in Eq. (22). The soft gluon contribution in the collinear-soft theory is identical to the LEET result, given in (24)
$`A_s`$ $`=`$ $`C_V\overline{\xi }_n\mathrm{\Gamma }^\mu h{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left[{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{2}{ϵ}}\mathrm{log}{\displaystyle \frac{p_s^2}{m_b\stackrel{~}{\mu }}}+2\mathrm{log}^2{\displaystyle \frac{p_s^2}{m_b\stackrel{~}{\mu }}}+\mathrm{}\right].`$ (38)
The collinear vertex diagram, Fig. 8(a), gives
$`A_c^{(v)}`$ $`=`$ $`C_V\overline{\xi }_n\mathrm{\Gamma }^\mu h{\displaystyle \frac{\alpha _sC_F}{2\pi }}\left(4\pi {\displaystyle \frac{\mu ^2}{p_s^2}}\right)^ϵ{\displaystyle \frac{\mathrm{\Gamma }(1+ϵ)\mathrm{\Gamma }(1ϵ)\mathrm{\Gamma }(2ϵ)}{\mathrm{\Gamma }(22ϵ)}}{\displaystyle \frac{1}{ϵ^2}}`$ (39)
$`=`$ $`C_V\overline{\xi }_n\mathrm{\Gamma }^\mu h{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left[{\displaystyle \frac{2}{ϵ^2}}{\displaystyle \frac{2}{ϵ}}+{\displaystyle \frac{2}{ϵ}}\mathrm{log}{\displaystyle \frac{p_s^2}{\stackrel{~}{\mu }^2}}\mathrm{log}^2{\displaystyle \frac{p_s^2}{\stackrel{~}{\mu }^2}}+2\mathrm{log}{\displaystyle \frac{p_s^2}{\stackrel{~}{\mu }^2}}+\mathrm{}\right].`$ (40)
As previously discussed, the leading piece of the wave function graph Fig. 8(b) is $`O(1/\lambda ^2)`$, but fortunately vanishes. Expanding to higher orders in $`\lambda `$, the graph gives the same result as in full QCD, (19). We therefore obtain for the contribution of the collinear gluons
$`A_c`$ $`=`$ $`C_V\overline{\xi }_n\mathrm{\Gamma }^\mu h{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left[{\displaystyle \frac{2}{ϵ^2}}{\displaystyle \frac{3}{2ϵ}}+{\displaystyle \frac{2}{ϵ}}\mathrm{log}{\displaystyle \frac{p_s^2}{\stackrel{~}{\mu }^2}}\mathrm{log}^2{\displaystyle \frac{p_s^2}{\stackrel{~}{\mu }^2}}+{\displaystyle \frac{3}{2}}\mathrm{log}{\displaystyle \frac{p_s^2}{\stackrel{~}{\mu }^2}}+\mathrm{}\right].`$ (41)
Adding the soft and collinear contributions, as well as the counterterm given in (37), we obtain
$$A_{cs}=C_V\overline{\xi }_n\mathrm{\Gamma }^\mu h\frac{\alpha _sC_F}{4\pi }\left[\mathrm{log}^2\frac{p_s^2}{m_b^2}+\frac{3}{2}\mathrm{log}\frac{p_s^2}{m_b^2}+\frac{1}{ϵ}\frac{1}{2}\mathrm{log}^2\frac{\stackrel{~}{\mu }^2}{m_b^2}\frac{3}{2}\mathrm{log}\frac{\stackrel{~}{\mu }^2}{m_b^2}+\mathrm{}\right].$$
(42)
Note that the troublesome divergence $`\frac{1}{ϵ}\mathrm{log}p_s^2`$ cancels once the two contributions (38) and (41) are added. Thus, both collinear and soft modes are required for the theory to be renormalized sensibly. Comparing to the full theory result (22), we see that the collinear-soft theory reproduces the IR physics of QCD, and that at the scale $`\stackrel{~}{\mu }=m_b`$ all nonanalytic terms vanish. This determines the matching scale to be $`m_b`$, confirming the result found by calculating in pure dimensional regularization (37).
### C Renormalization group equations
From the counterterm given in (37) it is simple to extract the anomalous dimension of the operator $`V^\mu `$ in the collinear-soft theory. From the definition,
$$\gamma _V=Z_V^1\left(\stackrel{~}{\mu }\frac{}{\stackrel{~}{\mu }}+\beta \frac{}{g}\right)Z_V$$
(43)
we have
$`\stackrel{~}{\mu }{\displaystyle \frac{}{\stackrel{~}{\mu }}}Z_V`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\stackrel{~}{\mu })C_F}{2\pi ϵ}}`$ (44)
$`\beta {\displaystyle \frac{}{g}}Z_V`$ $`=`$ $`{\displaystyle \frac{\alpha _s(\stackrel{~}{\mu })C_F}{2\pi }}\left({\displaystyle \frac{1}{ϵ}}+\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+{\displaystyle \frac{5}{2}}\right),`$ (45)
where we have used $`\beta =gϵ+O(g^3)`$. This give the anomalous dimension
$$\gamma _V=\frac{\alpha _s(\stackrel{~}{\mu })C_F}{2\pi }\left(\mathrm{log}\frac{\stackrel{~}{\mu }^2}{m_b^2}+\frac{5}{2}\right).$$
(46)
Note that the divergent piece of the anomalous dimension cancels between the two terms . The RGE for the coefficient of the operator $`V^\mu `$ is therefore
$$\stackrel{~}{\mu }\frac{d}{d\stackrel{~}{\mu }}C_V(\stackrel{~}{\mu })=\gamma _V(\stackrel{~}{\mu })C_V(\stackrel{~}{\mu }).$$
(47)
Solving this RGE we obtain
$$C_V(\stackrel{~}{\mu })=\left(\frac{\alpha _s(\stackrel{~}{\mu })}{\alpha _s}\right)^{\frac{C_F}{2\beta _0}\left(5\frac{8\pi }{\beta _0\alpha _s}\right)}\left(\frac{\stackrel{~}{\mu }^2}{m_b^2}\right)^{\frac{C_F}{\beta _0}}C_V(m_b),$$
(48)
where $`\alpha _s\alpha _s(m_b)`$, $`\beta _0=112/3n_f`$, and $`C_V(m_b)=1+O(\alpha _s(m_b))`$. Note that in deriving the anomalous dimension (46) we have assumed that the nonlocal vertex given in Fig. 7 has the same running as the QCD coupling. This assumption needs to be checked in subsequent work.
## IV The soft theory
### A Matching
The collinear-soft effective theory is valid down to $`\stackrel{~}{\mu }m_b\sqrt{1x}`$, the typical invariant mass of the light $`s`$-quark jet. At this scale we integrate out the collinear modes, and perform an OPE to calculate the inclusive $`b`$ decay rate. Diagrammatically, this is illustrated in Fig. (9). This results in a nonlocal OPE in which the two currents are separated along a light-like direction. As in Eq. (4), we write the momentum of the eikonal line as
$$p_s^\mu =\frac{m_b}{2}n^\mu +k^\mu +\frac{m_b}{2}(1y)\overline{n}^\mu $$
(49)
where $`k^\mu `$ is the residual momentum of the heavy quark (note that we distinguish $`y`$ from $`x`$, the rescaled photon momentum, since beyond tree level they will differ). The imaginary piece of the first graph is then proportional to $`\delta (1y+\widehat{k}_+)`$ (where, as usual, hatted variables are divided by $`m_b`$), so the OPE is in terms of an infinite number of nonlocal operators, labelled by $`y`$:
$$O(y)=\overline{h}_v\delta (1y+i\widehat{D}_+)h_v.$$
(50)
Feynman rules for nonlocal operators of this type were obtained in , by writing them as the Fourier transform of operators in position space, and expanding out the path-ordered exponential in powers of the gauge field. Equivalently, the Feynman rules may be obtained by taking the imaginary piece of the time-ordered product in LEET with additional gluons; the single gluon Feynman rule is given in Fig. 9.
The matrix element of $`O(y)`$ between heavy quark states with residual momentum $`k`$ is
$$b(k)|\overline{h}_v\delta (1y+i\widehat{D}_+)h_v|b(k)=\delta (1y+\widehat{k}_+)+O(\alpha _s).$$
(51)
while its matrix element between hadrons is the well known structure function
$$f(y)=\frac{B|\overline{h}_v\delta (1y+i\widehat{D}_+)h_v|B}{B|\overline{h}_vh_v|B}.$$
(52)
Thus, LEET consists of a continuous set of operators labeled by $`y`$. Each operator has a coefficient that depends on the kinematic variable $`x`$, and the differential rate for $`BX_s\gamma `$ is given by the integral
$$\frac{d\mathrm{\Gamma }}{dx}=\mathrm{\Gamma }_0𝑑yC(y,x;\mu )f(y;\mu ),$$
(53)
where the $`C(y,x;\mu )`$’s are the coefficients of the OPE.
To match onto LEET at one loop we compare the differential decay rate in the parton model, $`bX_s\gamma `$, which in LEET is
$$\frac{d\mathrm{\Gamma }}{dx}|_{k_+}=\mathrm{\Gamma }_0𝑑yC(y,x;\mu )b(k)|O(y;\mu )|b(k).$$
(54)
We therefore need the one-loop matrix element of $`O(y)`$ between quark states. This may be calculated from the diagrams shown in Fig. 10.
Again all divergences are regulated in dimensional regularization. As an example, Fig. 10(a) gives
$$b(k)|O^{(a)}(y)|b(k)=iC_Fg^2\left(\frac{\mu }{m_b}\right)^{4d}\frac{d^{d2}\widehat{q}_{}}{(2\pi )^{d2}}\frac{d\widehat{q}_{}}{2\pi }\frac{d\widehat{q}_+}{2\pi }\frac{\delta (\widehat{k}_++1y)\delta (\widehat{k}_++\widehat{q}_++1y)}{(\widehat{q}_+\widehat{q}_{}\widehat{q}_{}^2+iϵ)(\widehat{q}_++\widehat{q}_{}+iϵ)\widehat{q}_+}.$$
(55)
The first term is proportional to
$`{\displaystyle \frac{d^d\widehat{q}}{(2\pi )^d}\frac{1}{(\widehat{q}^2+iϵ)(\widehat{q}v+iϵ)(\widehat{q}n)}}=8{\displaystyle \frac{d^d\widehat{q}}{(2\pi )^d}_0^1𝑑x_0^{\mathrm{}}𝑑\lambda \frac{\lambda }{\left(\widehat{q}^2+2\lambda \widehat{q}(v(1x)+xn)\right)^3}}`$ (56)
$`={\displaystyle \frac{4i}{(4\pi )^{d/2}}}\mathrm{\Gamma }(3d/2){\displaystyle _0^1}𝑑x{\displaystyle _0^{\mathrm{}}}𝑑\lambda \lambda ^{d5}\left((1x)^2+2(1x)\right)^{d/23}.`$ (57)
The $`\lambda `$ integral vanishes in dimensional regularization, so this term vanishes. After performing the trivial $`\widehat{q}_+`$ integral in the second term, we are left with
$`b(k)|O^{(a)}(y)|b(k)`$ $`=`$ $`i{\displaystyle \frac{C_Fg^2}{2\pi }}\left({\displaystyle \frac{\mu }{m_b}}\right)^{4d}{\displaystyle \frac{1}{\widehat{k}_++1y}}{\displaystyle \frac{d^{d2}\widehat{q}_{}}{(2\pi )^{d2}}\frac{d\widehat{q}_{}}{2\pi }\frac{1}{\widehat{q}_{}(\widehat{k}_++1y)+\widehat{q}_{}^2iϵ}}`$ (59)
$`\times {\displaystyle \frac{1}{\widehat{q}_{}(\widehat{k}_++1y)+iϵ}}`$
$`=`$ $`{\displaystyle \frac{C_Fg^2}{2\pi }}\left({\displaystyle \frac{\mu }{m_b}}\right)^{4d}{\displaystyle \frac{\theta (\widehat{k}_++1y)}{\widehat{k}_++1y}}{\displaystyle \frac{d^{d2}\widehat{q}_{}}{(2\pi )^{d2}}\frac{1}{\widehat{q}_{}^2+(\widehat{k}_++1y)^2}}`$ (60)
$`=`$ $`{\displaystyle \frac{C_Fg^2}{8\pi ^2}}\left(4\pi {\displaystyle \frac{\mu ^2}{m_b^2}}\right)^ϵ\mathrm{\Gamma }(ϵ){\displaystyle \frac{\theta (\widehat{k}_++1y)}{(\widehat{k}_++1y)^{1+2ϵ}}}.`$ (61)
Using the identity
$$\frac{\theta (yx)}{(yx)^{1+2ϵ}}=\frac{1}{2ϵ}\delta (yx)+\theta (yx)\left[\frac{1}{(yx)_+}2ϵ\left(\frac{\mathrm{log}(yx)}{(yx)}\right)_++O(ϵ^2)\right]$$
(62)
we find
$`b(k)|O^{(a)}(y)|b(k)`$ $`=`$ $`{\displaystyle \frac{\alpha _sC_F}{4\pi }}\{({\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{1}{ϵ}}\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+{\displaystyle \frac{1}{2}}\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{})\delta (1y+\widehat{k}_+)`$ (65)
$`\theta (1y+\widehat{k}_+)[({\displaystyle \frac{2}{ϵ}}+2\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}){\displaystyle \frac{1}{(1y+\widehat{k}_+)_+}}`$
$`4\left({\displaystyle \frac{\mathrm{log}(1y+\widehat{k}_+)}{1y+\widehat{k}_+}}\right)_+]\}.`$
The diagram in Fig. 10(b) gives the same result as (a), while the diagram in Fig. 10(c) gives
$$b(k)|O^{(\mathrm{c})}(y)|b(k)=\frac{\alpha _sC_F}{4\pi }\left[\left(\frac{2}{ϵ}2\mathrm{log}\frac{\stackrel{~}{\mu }^2}{m_b^2}\right)\delta (1y+\widehat{k}_+)+4\frac{\theta (1y+\widehat{k}_+)}{(1y+\widehat{k}_+)_+}\right].$$
(66)
In dimensional regularization the wavefunction diagrams vanish. Since the decay rate is infrared finite, including the wavefunction graphs simply converts an infrared $`1/ϵ`$ divergence to an ultraviolet divergence. Therefore, we may neglect the wavefunction counterterm, and combining all graphs we find the bare matrix element
$`b(k)|O^{\mathrm{bare}}(y)|b(k)`$ $`=`$ $`[1{\displaystyle \frac{\alpha _sC_F}{4\pi }}({\displaystyle \frac{2}{ϵ^2}}{\displaystyle \frac{2}{ϵ}}+{\displaystyle \frac{2}{ϵ}}\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}2\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}`$ (70)
$`+\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{})]\delta (1y+\widehat{k}_+)`$
$`+{\displaystyle \frac{\alpha _sC_F}{4\pi }}\theta (1y+\widehat{k}_+)[({\displaystyle \frac{4}{ϵ}}4\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+4){\displaystyle \frac{1}{(1y+\widehat{k}_+)_+}}`$
$`+8\left({\displaystyle \frac{\mathrm{log}(1y+\widehat{k}_+)}{1y+\widehat{k}_+}}\right)_+],`$
where all divergences are ultraviolet. The renormalized operator $`O(y;\mu )`$ is related to the bare operator by
$$O^{\mathrm{bare}}(y)=𝑑y^{}Z(y^{},y;\stackrel{~}{\mu })O(y^{};\stackrel{~}{\mu }).$$
(71)
Renormalizing in MS (generalized in the obvious way to cancel the $`1/ϵ^2`$ divergences), we find
$`Z(y^{},y;\stackrel{~}{\mu })`$ $`=`$ $`\{[1{\displaystyle \frac{\alpha _s(\stackrel{~}{\mu })C_F}{2\pi }}({\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{1}{ϵ}}\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}{\displaystyle \frac{1}{ϵ}})]\delta (y^{}y)`$ (73)
$`+{\displaystyle \frac{\alpha _s(\stackrel{~}{\mu })C_F}{\pi }}{\displaystyle \frac{1}{ϵ}}{\displaystyle \frac{1}{(y^{}y)_+}}\theta (y^{}y)\}.`$
Note that the counterterm consists of a diagonal piece which is proportional to $`\delta (y^{}y)`$, and an off-diagonal piece proportional to $`\theta (y^{}y)`$. This latter terms mixes the operator $`O(y)`$ with all operators $`O(y^{})`$ with $`y^{}>y`$.
Inserting the one-loop matrix element of the renormalized operator into (54) we find the the differential decay rate in the parton model $`bX_s\gamma `$
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}|_{k_+}`$ $`=`$ $`\mathrm{\Gamma }_0{\displaystyle 𝑑yC(y,x;\stackrel{~}{\mu })O(y;\stackrel{~}{\mu })}`$ (74)
$`=`$ $`\mathrm{\Gamma }_0{\displaystyle }dyC(y,x;\stackrel{~}{\mu })\{[1{\displaystyle \frac{\alpha _sC_F}{4\pi }}(\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}2\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{})]\delta (1y+k_+)`$ (77)
$`{\displaystyle \frac{\alpha _sC_F}{4\pi }}\theta (1y+\widehat{k}_+)[(44\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}){\displaystyle \frac{1}{(1y+\widehat{k}_+)_+}}`$
$`+8\left({\displaystyle \frac{\mathrm{log}(1y+\widehat{k}_+)}{1y+\widehat{k}_+}}\right)_+]\}.`$
One might worry about the appearance in (74) of logarithmic terms that depend on $`m_b`$, since this scale has been integrated out and thus should not be present in the effective theory. These terms are due to our choice of factoring the heavy quark mass out of the soft scale $`m_b(1y+\widehat{k}_+)`$ by writing our expressions in terms of hatted quantities. The logarithms of $`m_b`$ cancel in the matching coefficient.
The Wilson coefficients $`C(y,x;\mu )`$ are determined by matching the collinear-soft theory onto LEET. In the collinear-soft theory, the Feynman diagrams for the forward scattering matrix element are shown in Fig. 11.
As with LEET, all divergences are regulated in dimensional regularization. Expanding the expression for the forward scattering amplitude obtained from these graphs in powers of $`(1x+\widehat{k}_+)`$, we find for the differential decay rate
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}|_{k_+}`$ $`=`$ $`C_V^2(\stackrel{~}{\mu })\mathrm{\Gamma }_0\{[1+{\displaystyle \frac{\alpha _sC_F}{4\pi }}(\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+5\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{})]\delta (1x+\widehat{k}_+)`$ (79)
$`{\displaystyle \frac{\alpha _sC_F}{4\pi }}\theta (1x+\widehat{k}_+)[4\left({\displaystyle \frac{\mathrm{log}(1x+\widehat{k}_+)}{1x+\widehat{k}_+}}\right)_++7\left({\displaystyle \frac{1}{1x+\widehat{k}_+}}\right)_+]\}.`$
Comparing Eqs. (79) and (74) gives the short-distance coefficient $`C(y,x;\mu )`$. At tree level, the matching is trivial, and we write
$$C(y,x;\stackrel{~}{\mu })=C_V^2(\stackrel{~}{\mu })\left[\delta (yx)+\frac{\alpha _sC_F}{4\pi }C^{(1)}(y,x;\stackrel{~}{\mu })\right]+O(\alpha _s^2),$$
(80)
where $`\mu `$ is the matching scale. At one loop, we find
$`C^{(1)}(y,x;\stackrel{~}{\mu })`$ $`=`$ $`\left(2\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+3\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}+\mathrm{}\right)\delta (yx)\left(3+4\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}\right){\displaystyle \frac{\theta (yx)}{(yx)_+}}`$ (82)
$`+4\theta (yx)\left({\displaystyle \frac{\mathrm{log}(yx)}{yx}}\right)_+`$
$`=`$ $`\left(2\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2(yx)}}+3\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2(yx)}}+\mathrm{}\right)\delta (yx)`$ (84)
$`4{\displaystyle \frac{\theta (yx)}{yx}}\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2(yx)}}3{\displaystyle \frac{\theta (yx)}{yx}}.`$
At the scale $`\stackrel{~}{\mu }=m_b\sqrt{yx}`$ the logarithmic terms vanish, and we find
$$C^{(1)}(y,x;m_b\sqrt{yx})=3\frac{\theta (yx)}{yx}+\mathrm{}.$$
(85)
The matching scale is therefore different for each operator $`O(y)`$.
### B Renormalization group
The differential decay rate in LEET given in (74) may be written as
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}`$ $`=`$ $`\mathrm{\Gamma }_0{\displaystyle }dyC(y,x;\stackrel{~}{\mu })\{[1{\displaystyle \frac{\alpha _sC_F}{4\pi }}(\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2(1y+\widehat{k}_+)^2}}2\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2(1y+\widehat{k}_+)^2}})]`$ (87)
$`\times \delta (1y+\widehat{k}_+)+{\displaystyle \frac{\alpha _sC_F}{4\pi }}({\displaystyle \frac{4}{1y+\widehat{k}_+}}\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2(1y+\widehat{k}_+)^2}}{\displaystyle \frac{4}{1y+\widehat{k}_+}})\},`$
and so the large logarithms in the matrix element of $`O(y;\stackrel{~}{\mu })`$ vanish at the scale $`\stackrel{~}{\mu }=m_b(1y+\widehat{k}_+)`$. (This expression looks highly singular, but as can be seen from (74), the delta functions combine with the other terms to form plus functions.) Thus, in order to sum all logarithms of $`\mu `$ we must continue to run the operator $`O(y)`$ in LEET. From (71) and (73) we obtain the renormalization group equation
$$\mu \frac{d}{d\mu }C(y,x;\stackrel{~}{\mu })=𝑑y^{}\gamma (y,y^{};\stackrel{~}{\mu })C(y^{},x;\stackrel{~}{\mu }),$$
(88)
where $`\gamma (y,y^{};\stackrel{~}{\mu })`$ is the continuous anomalous dimension matrix
$$\gamma (y,y^{};\stackrel{~}{\mu })=\frac{\alpha _s(\stackrel{~}{\mu })C_F}{\pi }\left[\left(\mathrm{log}\frac{\stackrel{~}{\mu }^2}{m_b^2}1\right)\delta (y^{}y)\frac{2}{(y^{}y)_+}\theta (y^{}y)\right].$$
(89)
Solving (88) analytically, however, is nontrivial and beyond the scope of this work . Instead, we may diagonalize the anomalous dimension matrix by taking high moments of the spectrum. This will allow us to compare our results to those of Refs. . Note that in Refs. both leading and next-to-leading logarithms were resummed. This requires the two loop contribution to the $`1/ϵ^2`$ counterterm, the full one loop matching condition, and the two loop running of $`\alpha _s`$, none of which have been included here. As a result our calculation only resums the leading logarithms and a class of the subleading logarithms. However, it is straightforward to extract from the literature a resummation of exactly the same set of logarithms.
To calculate the moments we set the residual momentum $`k`$ to zero. (This residual momentum can easily be incorporated by boosting from the rest frame of the $`b`$ quark, $`p_b=m_bv`$, to the frame $`p=m_bv+k`$). Taking moments unconvolutes the expression for the differential decay rate in LEET (54) and we obtain
$`\mathrm{\Gamma }(N)`$ $`=`$ $`{\displaystyle _0^1}𝑑xx^{N1}{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}`$ (90)
$`=`$ $`\mathrm{\Gamma }_0{\displaystyle _0^1}𝑑xx^{N1}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑yC(yx;\mu )O(y;\mu )`$ (91)
$`=`$ $`\mathrm{\Gamma }_0{\displaystyle _0^1}𝑑zz^{N1}C^{}(1z;\stackrel{~}{\mu }){\displaystyle _0^1}𝑑yy^{N1}O(y;\stackrel{~}{\mu })`$ (92)
$``$ $`\mathrm{\Gamma }_0C(N;\stackrel{~}{\mu })O(N;\stackrel{~}{\mu }),`$ (93)
where we have used
$$C(yx)=\frac{1}{y}C^{}\left(1\frac{x}{y}\right)\mathrm{\Theta }(yx)$$
(94)
since $`C(yx)`$ just contains delta functions and plus distributions. Thus, the operator $`O(N;\mu )`$ is just a linear combination of the set of operators $`O(y;\mu )`$. The matching from the collinear-soft theory onto LEET at tree level is trivial, and we find
$$C(N;\stackrel{~}{\mu })=C_V^2(\stackrel{~}{\mu })\left[1+\frac{\alpha _sC_F}{4\pi }C^{(1)}(N;\stackrel{~}{\mu })\right]+O(\alpha _s^2).$$
(95)
Determining $`C^{(1)}(N;\mu _0)`$ requires the one-loop expression of $`\mathrm{\Gamma }(N)`$ in the collinear-soft theory and the one-loop matrix element of $`O(N;\mu )`$ between partonic states. The one-loop expression for the differential decay rate in the collinear-soft theory is given in (79). Setting $`k_+`$ to zero and taking moments we obtain
$`\mathrm{\Gamma }(N)`$ $`=`$ $`{\displaystyle _0^1}x^{N1}{\displaystyle \frac{d\mathrm{\Gamma }}{dx}}`$ (96)
$`=`$ $`\mathrm{\Gamma }^0C_V^2(\stackrel{~}{\mu })\left\{1{\displaystyle \frac{\alpha _sC_F}{4\pi }}\left[2\mathrm{log}^2{\displaystyle \frac{N}{n_0}}7\mathrm{log}{\displaystyle \frac{N}{n_0}}\mathrm{log}^2{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}5\mathrm{log}{\displaystyle \frac{\stackrel{~}{\mu }^2}{m_b^2}}\right]\right\}+\mathrm{},`$ (97)
where $`n_0=e^{\gamma _E}`$. This needs to be compared to the one-loop matrix element of $`O(N;\mu )`$, which can be obtained from (74):
$$O(N;\stackrel{~}{\mu })=_0^1𝑑yy^{N1}O(y;\stackrel{~}{\mu })=1\frac{\alpha _sC_F}{4\pi }\left[4\mathrm{log}^2\frac{\stackrel{~}{\mu }N}{m_bn_0}4\mathrm{log}\frac{\stackrel{~}{\mu }N}{m_bn_0}\right]+\mathrm{}.$$
(98)
The one loop matching coefficient is now easily determined using (93), (96) and (98) and we find
$$C^{(1)}(N;\stackrel{~}{\mu })=\frac{\alpha _sC_F}{4\pi }\left[2\mathrm{log}^2\frac{\stackrel{~}{\mu }^2N}{m_b^2n_0}+3\mathrm{log}\frac{\stackrel{~}{\mu }^2N}{m_b^2n_0}\right]+\mathrm{}.$$
(99)
At the matching scale $`\stackrel{~}{\mu }=m_b\sqrt{n_0/N}`$ all logarithms in this matching coefficient vanish. Furthermore, from (98) it is clear that the matrix element $`O(N;\stackrel{~}{\mu })`$ contains no large logarithms of $`N`$ at the scale $`\stackrel{~}{\mu }=m_bn_0/N`$. Thus we run in the collinear-soft theory from $`m_b`$ to $`m_b\sqrt{n_0/N}`$, perform the OPE, and run $`C(N;\stackrel{~}{\mu })`$ from $`m_b\sqrt{n_0/N}`$ to $`m_bn_0/N`$.
The running of the coefficient $`C_V`$ in the collinear-soft theory from the scale $`m_b`$ to the scale $`m_b\sqrt{n_0/N}`$ is obtained by setting $`\stackrel{~}{\mu }=m_b\sqrt{n_0/N}`$ in (48). The running in LEET is determined by the RGE for $`C(N;\stackrel{~}{\mu })`$
$$\mu \frac{d}{d\mu }C(N;\stackrel{~}{\mu })=\gamma (N;\stackrel{~}{\mu })C(N;\stackrel{~}{\mu }),$$
(100)
where the anomalous dimension is given by
$$\gamma (N;\stackrel{~}{\mu })=_0^1𝑑zz^{N1}\gamma (z;\stackrel{~}{\mu })=\frac{\alpha _s(\stackrel{~}{\mu })C_F}{\pi }\left[12\mathrm{log}\left(\frac{\stackrel{~}{\mu }N}{m_bn_0}\right)\right].$$
(101)
The solution to this equation is
$$C(N;\frac{m_bn_0}{N})=C_V^2\left(m_b\sqrt{\frac{n_0}{N}}\right)\left(\frac{\alpha _s(m_b\frac{n_0}{N})}{\alpha _s(m_b\sqrt{\frac{n_0}{N}})}\right)^{\frac{2C_F}{\beta _0}\left(1+\frac{4\pi }{\beta _0\alpha _s}2\mathrm{log}\frac{N}{n_0}\right)}\left(\frac{n_0}{N}\right)^{\frac{2C_F}{\beta _0}}.$$
(102)
This sums perturbative logarithms of $`N`$ into the coefficient $`C(N)`$. We can then substitute this into (93) to obtain an expression for the resummed moments of the differential decay rate.
Using the result for $`C_V(\mu )`$ given in (48) and taking the matrix element of $`O(N;\mu )`$ between hadronic states, we find find the resummed expression for large photon energy moments of the decay $`BX_s\gamma `$
$$\mathrm{\Gamma }(N)=\mathrm{\Gamma }_0f(N;m_bn_0/N)\left(\frac{\alpha _s(m_b\sqrt{\frac{n_0}{N}})}{\alpha _s}\right)^{\frac{C_F}{\beta _0}\left(5\frac{8\pi }{\beta _0\alpha _s}\right)}\left(\frac{\alpha _s(m_b\frac{n_0}{N})}{\alpha _s(m_b\sqrt{\frac{n_0}{N}})}\right)^{\frac{2C_F}{\beta _0}\left(1+\frac{4\pi }{\beta _0\alpha _s}2\mathrm{log}\frac{N}{n_0}\right)}.$$
(103)
Logarithms are explicitly summed in this expression and only long distance physics is contained in the function $`f(N;m_bn_0/N)`$.
We can easily compare our results to those in the literature. A resummed expression for $`\mathrm{\Gamma }(N)`$ is given in Ref. :
$$\mathrm{\Gamma }(N)=\mathrm{\Gamma }_0f(N;m_b/N)\mathrm{exp}\left[_{n_0/N}^1\frac{dy}{y}\left(2_{m_by}^{m_b\sqrt{y}}\frac{d\mu }{\mu }\mathrm{\Gamma }_c(\mu )+\mathrm{\Gamma }(m_by)+\gamma (m_b\sqrt{y})\right)\right],$$
(104)
where, at one loop,
$$\mathrm{\Gamma }_c(\mu )=\frac{\alpha _s(\mu )C_F}{\pi },\mathrm{\Gamma }(\mu )=\frac{\alpha _s(\mu )C_F}{\pi },\gamma (\mu )=\frac{3\alpha _s(\mu )C_F}{4\pi }.$$
(105)
Note that the cusp anomalous dimension $`\mathrm{\Gamma }_c(\mu )`$ is the contribution to the anomalous dimension from the $`1/ϵ^2`$ counterterm. Using only the one loop cusp anomalous dimension, tree level matching, and the one loop running of $`\alpha _s`$, Eq. (104) resums leading logarithms and the same class of next-to-leading logarithms we resum in our calculation. Performing the integrals in the exponent we reproduce (103). Thus the approach presented here, based on an effective field theory, is in agreement with the factorization formalism approach for summing perturbative logarithms.
## V Conclusions
In the specific case of $`\overline{B}X_s\gamma `$ we have shown how Sudakov logarithms can be summed within an effective field theory framework. First we construct an intermediate theory which includes both collinear and soft degrees of freedom. By performing a one-loop calculation we show that this collinear-soft theory can be matched onto QCD at the scale $`m_b`$ without introducing logarithmic terms into the short-distance coefficient. In addition we determine the one-loop anomalous dimension and solve the RGE. Next we integrate out collinear modes at the scale $`m_b\sqrt{yx}`$ by switching to LEET. We perform an OPE in powers of $`(yx)`$ which leads to the appearance of a nonlocal operator where two vertices are separated along the light-cone. The matrix element of this operator between $`B`$ meson states is the structure function. We perform the OPE at one-loop in the collinear-soft theory and match onto the nonlocal operator in LEET. At the scale $`m_b\sqrt{yx}`$ no logarithmic terms are introduced into the short-distance coefficient.
In order to compare to the factorization formalism results in the literature we repeat our analysis for large moments of the decay rate. In this case we find that the collinear-soft theory matches onto LEET at the scale $`m_b\sqrt{n_0/N}`$, and that there are no large logarithms in the matrix element of the bilocal operator at the scale $`m_bn_0/N`$. Using the renormalization group equations in the collinear-soft theory we sum logarithms of $`N`$ between the scales $`m_b`$ and $`m_b\sqrt{n_0/N}`$. We then switch to LEET and sum logarithms of $`N`$ between the scales $`m_b\sqrt{n_0/N}`$ and $`m_bn_0/N`$. This sums all perturbative logarithms of $`N`$. We find that our result agrees with results presented in the literature. This gives us confidence that we have constructed the correct effective field theory.
Though we have presented this work entirely in the context of $`BX_s\gamma `$ our approach is general. It should be straightforward to apply the collinear-soft theory and LEET to other processes in which Sudakov logarithms arise. Furthermore, this approach could also be applied to exclusive decays, in which case one does not perform the final OPE onto LEET, but remains in the collinear-soft theory. This could be applied to recent results on factorization in nonleptonic decays , as well as LEET-based relations between form-factors in decays to highly energetic final states. Since these latter results depend only on the spin symmetry of LEET, which is also present in the collinear-soft theory, they should remain valid in the present approach.
## Acknowledgements
We thank Craig Burrell, Adam Leibovich, Aneesh Manohar, Tom Mehen, Dan Pirjol, Ira Rothstein and Iain Stewart for useful discussions. This work was supported in part by the Natural Sciences and Engineering Research Council of Canada and the Sloan Foundation.
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# 𝜏-function for analytic curves
(May 2000)
We review the concept of $`\tau `$-function for simple analytic curves. The $`\tau `$-function gives a formal solution to the 2D inverse potential problem and appears as the $`\tau `$-function of the integrable hierarchy which describes conformal maps of simply-connected domains bounded by analytic curves to the unit disk. The $`\tau `$-function also emerges in the context of topological gravity and enjoys an interpretation as a large $`N`$ limit of the normal matrix model.
1. Recently, it has been realized that conformal maps exhibit an integrable structure: conformal maps of compact simply connected domains bounded by analytic curves provide a solution to the dispersionless limit of the 2D Toda hierarchy. As is well known from the theory of solitons, solutions of an integrable hierarchy are represented by $`\tau `$-functions. The dispersionless limit of the $`\tau `$-function emerges as a natural object associated with the curves. In this paper we discuss the $`\tau `$-function for simple analytic curves and its connection to the inverse potential problem, area preserving diffiomorphisms, the Dirichlet boundary problem, and matrix models.
2. Inverse potential problem. Consider a closed analytic curve<sup>1</sup><sup>1</sup>1A closed analytic curve is the curve which can be parametrized by a function $`zx+iy=z(w)`$, analytic in a domain which includes the unit circle $`|w|=1`$ $`\gamma `$ in the complex plane and denote by $`D_+`$ and $`D_{}`$ the interior and exterior domains with respect to the curve. The point $`z=0`$ is assumed to be in $`D_+`$. Assume that the domain $`D_+`$ is filled homogeneously with electric charge, with a density which we set to be equal to 1. The potential $`\mathrm{\Phi }`$ created by the charge obeys the equation
$$_z_{\overline{z}}\mathrm{\Phi }(z,\overline{z})=\{\begin{array}{cc}1\hfill & \text{ if }z=x+iyD_+\hfill \\ 0\hfill & \text{ if }z=x+iyD_{}\hfill \end{array}$$
(1)
The potential $`\mathrm{\Phi }`$ can be written as an integral over the domain $`D^+`$:
$$\mathrm{\Phi }(z,\overline{z})=\frac{2}{\pi }_{D_+}d^2z^{}\mathrm{log}|zz^{}|$$
(2)
In the exterior domain $`D_{}`$, the potential is the harmonic function whose asymptotic expansion as $`z\mathrm{}`$ is given by
$$\mathrm{\Phi }^{}(z,\overline{z})=2t_0\mathrm{log}|z|+2e\underset{k>0}{}\frac{v_k}{k}z^k,$$
(3)
where
$$v_k=\frac{1}{\pi }_{D_+}z^kd^2z(k>0)$$
(4)
are the harmonic moments of the interior domain $`D_+`$ and
$$\pi t_0=_{D_+}d^2z$$
(5)
is its area. In the interior domain $`D_+`$, the potential (2) is equal to a function $`\mathrm{\Phi }^+`$, which is harmonic up to the term $`|z|^2`$. The expansion of this function around $`z=0`$ is
$$\mathrm{\Phi }^+(z,\overline{z})=|z|^2v_0+2e\underset{k>0}{}t_kz^k$$
(6)
Here
$$t_k=\frac{1}{\pi k}_D_{}z^kd^2z(k>0)$$
(7)
are the harmonic moments of the exterior domain $`D_{}`$ and
$$v_0=\frac{2}{\pi }_{D_+}\mathrm{log}|z|d^2z.$$
(8)
The two sets of moments (4) and (7) are related by the conditions that $`\mathrm{\Phi }^+=\mathrm{\Phi }^{}`$, $`_z\mathrm{\Phi }^+=_z\mathrm{\Phi }^{}`$ on the curve $`\gamma `$.
The inverse potential problem is to determine the form of the curve $`\gamma `$ given one of the functions $`\mathrm{\Phi }^+`$ or $`\mathrm{\Phi }^{}`$, i.e. given one of the infinite sets of moments. We will choose as independent variables the area $`\pi t_0`$ and the moments of the exterior $`t_k(k1)`$. Under certain conditions, they completely determine the form of the curve as well as the moments $`v_k(k0)`$ . More precisely, $`\{t_k\}_{k=0}^{\mathrm{}}`$ is a good set of local coordinates in the space of analytic curves. For simplicity we assume in this paper that only a finite number of $`t_k`$ are non-zero. In this case the series (6) is a polynomial in $`z`$, $`\overline{z}`$ and, therefore, it gives the function $`\mathrm{\Phi }^+`$ for $`zD_+`$. Note that $`t_0`$, $`v_0`$ are real quantities while all other moments are in general complex variables.
3. Variational principle. Consider the energy functional describing a charge with a density $`\rho (z,\overline{z})`$ in the background potential created by the homogeneously distributed charge with the density +1 inside the domain $`D_+`$ (1):
$$\{\rho \}=\frac{1}{\pi ^2}d^2zd^2z^{}\rho (z,\overline{z})\mathrm{log}|zz^{}|\rho (z^{},\overline{z}^{})\frac{1}{\pi }d^2z\rho (z,\overline{z})\mathrm{\Phi }(z,\overline{z}).$$
(9)
The first term is the 2D “Coulomb” energy of the charge while the second one is the energy due to the background charge. Clearly, the distribution of the charge neutralizing the background charge gives the minimum to the functional: $`\rho _0=1`$ inside the domain and $`\rho _0=0`$ outside. At the minimum the functional is equal to minus electrostatic energy $`E`$ of the background charge :
$$E=\underset{\rho }{\mathrm{min}}\{\rho \}=\frac{1}{\pi ^2}_{D_+}d^2z_{D_+}d^2z^{}\mathrm{log}|zz^{}|=\frac{1}{2\pi }_{D_+}d^2z\mathrm{\Phi }(z,\overline{z}).$$
(10)
Varying over $`\rho `$ and then setting $`\rho =1`$ inside the domain, we obtain eq.(6).
The first corollary of the variational principle is that the $`E`$ is a potential function for the moments. Eq. (6) suggests to treat $`v_0`$ and $`t_k`$ as independent variables, so moments of the interior, $`v_k`$, $`k1`$, and $`t_0`$ are functions of $`v_0`$ and $`t_k`$. Let us differentiate $`E`$ or $`\{\rho \}`$ at the extremum with respect to the parameters $`v_0`$, $`t_k`$. Since $`\rho _0`$ minimizes the functional, the derivative is equivalent to the partial derivative of $``$ at the fixed extremum $`\rho `$. This gives
$$\frac{E}{t_k}=v_k,\frac{E}{\overline{t}_k}=\overline{v}_k,\frac{E}{v_0}=t_0,$$
(11)
where the partial derivative with respect to $`t_k`$ is taken at fixed $`v_0`$ and $`t_j(j0,k)`$. Therefore the differential $`dE`$ reads
$$dE=\underset{k>0}{}(v_kdt_k+\overline{v}_kd\overline{t}_k)t_0dv_0.$$
(12)
Let us note that the variational principle may be formulated in a number of different ways. One particular variational principle is suggested by the matrix model discussed in the Sec.9. In this case one consider a charged liquid in the potential
$$V(z,\overline{z})=z\overline{z}+v_0\underset{k>0}{}\left(t_kz^k+\overline{t}_k\overline{z}^k\right)$$
(13)
defined everywhere on the plane and $`v_0`$ and $`t_k`$ are parameters. The energy of the charged liquid
$$\{\rho ,V\}=\frac{1}{\pi ^2}d^2zd^2z^{}\rho (z,\overline{z})\rho (z^{},\overline{z}^{})\mathrm{log}|zz^{}|+\frac{1}{\pi }d^2z\rho (z,\overline{z})V(z,\overline{z}).$$
(14)
reaches the minimum if the liquid forms a drop with the density $`\rho _0=1`$ bounded by the curve determined by parameters of the potential $`v_0`$ and $`t_k`$. For another version of the variational principle see .
4. $`\tau `$-function. It is more natural to treat the total charge $`t_0`$ rather than $`v_0`$ as an independent variable, i.e. to consider the variational principle at a fixed total charge $`t_0=\rho d^2z`$. This is achieved via the Legendre transformation. Let us introduce the function $`F=E+t_0v_0`$ , whose differential is
$$dF=\underset{k>0}{}(v_kdt_k+\overline{v}_kd\overline{t}_k)+v_0dt_0.$$
(15)
We define the $`\tau `$-function as $`\tau =e^F`$, so that
$$\mathrm{log}\tau =\frac{1}{2\pi }_{D_+}d^2z\mathrm{\Phi }(z,\overline{z})+t_0v_0=\frac{1}{\pi }_{D_+}\mathrm{log}\left|\frac{1}{z}\frac{1}{z^{}}\right|d^2zd^2z^{}.$$
(16)
The $`\tau `$-function is a real function of the moments $`\{t_0,t_1,t_2,\mathrm{}\}`$. Under the assumption that only a finite number of them are non-zero, we can substitute (6) into (16) and perform the term-wise integration. Taking into account that $`\frac{1}{\pi }_{D_+}|z|^2d^2z=\frac{1}{2}t_0^2+\frac{1}{2}_{k>0}k(t_kv_k+\overline{t}_k\overline{v}_k)`$ (a simple consequence of the Stokes formula), we get the expression for the $`\tau `$-function in terms of $`t_k`$ and $`v_k`$:
$$2\mathrm{log}\tau =\frac{1}{2}t_0^2+t_0v_0\frac{1}{2}\underset{k>0}{}(k2)(t_kv_k+\overline{t}_k\overline{v}_k).$$
(17)
Rephraising (11) we get the main property of the $`\tau `$-function, which has been used as its definition in Ref.
$$\frac{\mathrm{log}\tau }{t_k}=v_k,\frac{\mathrm{log}\tau }{\overline{t}_k}=\overline{v}_k,\frac{\mathrm{log}\tau }{t_0}=v_0$$
(18)
where the derivative with respect to $`t_k`$ is taken at fixed $`t_j`$ ($`jk`$).
Two immediate consequences of the very existence of the potential function are symmetry relations for the moments
$$\frac{v_k}{t_n}=\frac{v_n}{t_k},\frac{v_k}{\overline{t}_n}=\frac{\overline{v}_n}{t_k}$$
(19)
and the quasi-homogeneity condition for the $`\tau `$-function:
$$4\mathrm{log}\tau =t_0^2+2t_0\frac{\mathrm{log}\tau }{t_0}\underset{n>0}{}(n2)\left(t_n\frac{\mathrm{log}\tau }{t_n}+\overline{t}_n\frac{\mathrm{log}\tau }{\overline{t}_n}\right).$$
(20)
Apart from the term $`t_0^2`$, this formula reflects the scaling of moments as $`z\lambda z`$: $`t_k\lambda ^{2k}t_k`$ ($`k0`$), $`v_k\lambda ^{2+k}v_k`$ ($`k1`$).
As an illustration we present the $`\tau `$-function of ellipse . In this case only the first two moments $`t_1`$ and $`t_2`$ are nonzero<sup>2</sup><sup>2</sup>2The $`\tau `$-fuction for the ellipse (at $`t_1=0`$) appeared in Ref. as the limit of the Laughlin wave function or a planar limit of the free energy of normal matrix models, see Sec. 9:
$$\text{log}\tau =\frac{3}{4}t_0^2+\frac{1}{2}t_0^2\text{log}\left(\frac{t_0}{14|t_2|^2)}\right)+\frac{t_0}{14|t_2|^2}\left(|t_1|^2+t_1^2\overline{t}_2+\overline{t}_1^2t_2\right).$$
5. Schwarz function and generating function of the conformal map. Consider a univalent conformal map of the exterior domain $`D_{}`$ to the exterior of the unit disk and expand it in Laurent series:
$$w(z)=\frac{1}{r}z+\underset{j=0}{\overset{\mathrm{}}{}}p_jz^j,$$
(21)
where the coefficient $`r`$ is chosen to be real and positive. The series for the inverse map (from the exterior of the unit disk to $`D_{}`$) has a similar form:
$$z(w)=rw+\underset{j=0}{\overset{\mathrm{}}{}}u_jw^j.$$
(22)
Chosen $`w`$ on the unit circle, eq.(22) gives a parametrization of the curve. By the definition of an analytic curve, the map can be analytically continued to a strip-like neighborhood of the curve belonging to $`D_+`$. The continuation is given by the Riemann-Schwarz reflection principle (see e.g.):
$$w=(\overline{w}(S(z)))^1,$$
(23)
where $`S(z)`$ is the point reflected relative to the curve<sup>3</sup><sup>3</sup>3We use the notation: given an analytic function $`f(z)=_jf_jz^j`$, we set $`\overline{f}(z)=_j\overline{f}_jz^j`$.. Following , we call $`S(z)`$ the Schwarz function of the curve. Let us recall its construction. Write the equation for the curve $`F(x,y)=0`$ in complex coordinates, $`F(\frac{z+\overline{z}}{2},\frac{z\overline{z}}{2i})=0`$, and solve it with respect to $`\overline{z}`$. One gets the Schwarz function: $`\overline{z}=S(z)`$. The Schwarz function is analytic in a strip-like domain that includes the curve. On the curve the Schwarz function is equal to the complex conjugate argument. The main property of the Schwarz function is the obvious but important unitarity condition
$$\overline{S}(S(z))=z$$
(24)
(the inverse function coincides with the complex conjugate function). In terms of a conformal map the Schwarz function is
$$S(z)=rw^1(z)+\underset{j=0}{\overset{\mathrm{}}{}}\overline{u}_jw^j(z).$$
(25)
Using the Schwarz function one can write the moments of the exterior and the interior domains (4,7) as contour integrals <sup>4</sup><sup>4</sup>4This is due to a more general statement $`_{D_\pm }f(z)d^2z=\pm \frac{1}{2i}_\gamma f(z)S(z)𝑑z,`$ where $`f(z)`$ is an analytic function in the domain $`D_\pm `$.
$$t_n=\frac{1}{2\pi in}_\gamma z^nS(z)𝑑z,v_n=\frac{1}{2\pi i}_\gamma z^nS(z)𝑑z$$
(26)
Eq. (26) yields the Laurent expansion of the Schwarz function
$$S(z)=\underset{k=1}{\overset{\mathrm{}}{}}kt_kz^{k1}+\frac{t_0}{z}+\underset{k=1}{\overset{\mathrm{}}{}}v_kz^{k1}.$$
(27)
Now let us define the generating function $`\mathrm{\Omega }(z)`$, related to the Schwarz function by
$$S(z)=_z\mathrm{\Omega }(z).$$
(28)
The latter is given, according to (27), by the Laurent series
$$\mathrm{\Omega }(z)=\underset{k=1}{\overset{\mathrm{}}{}}t_kz^k\frac{1}{2}v_0+t_0\text{log}z\underset{k=1}{\overset{\mathrm{}}{}}\frac{v_k}{k}z^k$$
(29)
It can be represented as $`\mathrm{\Omega }(z)=\mathrm{\Omega }^{(+)}(z)+\mathrm{\Omega }^{()}(z)\frac{1}{2}v_0`$, where $`\mathrm{\Omega }^{(\pm )}(z)`$ are analytic in $`D_\pm `$ respectively:
$$\mathrm{\Omega }^{(+)}(z)=\frac{1}{\pi }_D_{}\text{log}\left(1\frac{z}{z^{}}\right)d^2z^{}=\underset{k=1}{\overset{\mathrm{}}{}}t_kz^k$$
(30)
$$\mathrm{\Omega }^{()}(z)=\frac{1}{\pi }_{D_+}\mathrm{log}(zz^{})d^2z^{}=t_0\mathrm{log}z\underset{k=1}{\overset{\mathrm{}}{}}\frac{v_k}{k}z^k$$
(31)
From (3,6) we see that $`\mathrm{\Phi }^{}(z,\overline{z})=2e\mathrm{\Omega }^{()}(z)`$ and $`\mathrm{\Phi }^+(z,\overline{z})=2e\mathrm{\Omega }^{(+)}(z)v_0|z|^2`$. Contrary to the potentials $`\mathrm{\Phi }^\pm `$, the analytical functions $`\mathrm{\Omega }^+`$ and $`\mathrm{\Omega }^{}`$ do not match each other on the curve. The discontinuity gives the value of the generating function restricted to the curve
$$\mathrm{\Omega }(z)=\frac{1}{2}|z|^2+2iA(z),z\gamma $$
(32)
where $`A(z)`$ is the area of the interior domain bound by the ray $`\phi =\text{arg}z`$ and the real axis. As a corollary, it is easy to show that variations of the $`\mathrm{\Omega }(z)`$ on the curve with respect to the real parameters $`t_0`$, $`et_k`$ and $`mt_k`$ are purely imaginary. This allows one to apply the Riemann-Schwarz reflection principle to analytical continuation of
$$H_k(z)=_{t_k}\mathrm{\Omega }(z),\overline{H}_k(z)=_{\overline{t}_k}\mathrm{\Omega }(z)$$
(33)
and to prove the fundamental relations
$$_{t_0}\mathrm{\Omega }(z)=\mathrm{log}w(z),$$
(34)
$$_{t_k}\mathrm{\Omega }(z)=\left(z^k(w)\right)_++\frac{1}{2}\left(z^k(w)\right)_0$$
(35)
$$_{\overline{t}_k}\mathrm{\Omega }(z)=\left(S^k(z(w))\right)_{}+\frac{1}{2}\left(S^k(z(w))\right)_0$$
(36)
The symbols $`(f(w))_\pm `$ mean a truncated Laurent series, where only terms with positive (negative) powers of $`w`$ are kept, while $`(f(w))_0`$ is the constant term ($`w^0`$) of the series. Note that the derivatives in eqs.(34-36) are taken at fixed $`z`$.
To prove (34), we first notice that
$$_{t_0}\mathrm{\Omega }(z(w))=\mathrm{log}z\frac{_{t_0}v_0}{2}+\text{negative powers in}z=\mathrm{log}wr\frac{_{t_0}v_0}{2}+\text{negative powers in}w.$$
Independently, one can show that $`_{t_0}v_0=2\text{log}r`$.
Then, using the Riemann-Schwarz reflection principle, we may write $`_{t_0}\mathrm{\Omega }(z(w))`$ also in the form $`_{t_0}\overline{\mathrm{\Omega }}(S(z(w))`$. Expanding the latter in $`S(z)`$ and then, using expansion of (25) in $`w`$, we have
$$_{t_0}\overline{\mathrm{\Omega }}(S(z(w))=\mathrm{log}S(z)\frac{_{t_0}v_0}{2}+\text{negative powers in}S(z)=\mathrm{log}w+\text{positive powers in}w.$$
Comparing both expansions, we conclude that $`_{t_0}\mathrm{\Omega }(z)=\mathrm{log}w(z)`$. Similar arguments are used in the proof of (35) and (36).
6. Dispersionless Hirota equation and the Dirichlet boundary problem. Using the representation (18) of the moments $`v_k`$ as derivatives of the $`\tau `$-function, one can express the conformal map $`w(z)`$ (34) through the $`\tau `$-function:
$$\text{log}w=\text{log}z_{t_0}\left(\frac{1}{2}_{t_0}+\underset{k1}{}\frac{z^k}{k}_{t_k}\right)\text{log}\tau .$$
(37)
With the help of the $`\tau `$-function, eqs.(35,36) can be similarly encoded as follows:
$$_z_\zeta \text{log}\left(w(z)w(\zeta )\right)=\frac{1}{(z\zeta )^2}+\left(\underset{k1}{}z^{k1}_{t_k}\right)\left(\underset{n1}{}\zeta ^{n1}_{t_n}\right)\text{log}\tau $$
(38)
$$_z_{\overline{\zeta }}\text{log}\left(w(z)\overline{w}(\overline{\zeta })1\right)=\left(\underset{k1}{}z^{k1}_{t_k}\right)\left(\underset{n1}{}\overline{\zeta }^{n1}_{\overline{t}_n}\right)\text{log}\tau $$
(39)
The derivation is similar to the one given in for the case of the KP hierarchy. Moreover, these equations in the integrated form are most conveniently written in terms of the differential operators
$$D(z)=\underset{k1}{}\frac{z^k}{k}_{t_k},\overline{D}(\overline{z})=\underset{k1}{}\frac{\overline{z}^k}{k}_{\overline{t}_k}$$
(40)
From (38,39) one obtains:
$$\text{log}\frac{w(z)w(\zeta )}{z\zeta }=\frac{1}{2}_{t_0}^2\text{log}\tau +D(z)D(\zeta )\text{log}\tau $$
(41)
$$\text{log}\left(1\frac{1}{w(z)\overline{w}(\overline{\zeta })}\right)=D(z)\overline{D}(\overline{\zeta })\text{log}\tau $$
(42)
Combining (37) and (41), one obtains the dispersionless Hirota equation (or the dispersionless Fay identity) for 2D Toda lattice hierarchy :
$$(z\zeta )e^{D(z)D(\zeta )\mathrm{log}\tau }=ze^{_{t_0}D(z)\mathrm{log}\tau }\zeta e^{_{t_0}D(\zeta )\mathrm{log}\tau }$$
(43)
Eq. (43), after being expanded in powers of $`z`$ and $`\zeta `$, generates an infinite set of relations between the second derivatives $`_{t_n}_{t_m}\text{log}\tau `$ of the $`\tau `$-function. Using (42) instead of (41), a similar equation for the mixed derivatives $`_{t_n}_{\overline{t}_m}\text{log}\tau `$ can be written:
$$1e^{D(z)\overline{D}(\overline{\zeta })\mathrm{log}\tau }=\frac{1}{z\overline{\zeta }}e^{_{t_0}(_{t_0}+D(z)+\overline{D}(\overline{\zeta }))\mathrm{log}\tau }$$
(44)
Let us conclude this section with two other forms of the dispersionless Hirota equation for the conformal map. They emphasize a relation between the Hirota equation and two fundamental objects of the classical analysis: the Green function of the Dirichlet problem<sup>5</sup><sup>5</sup>5This relation is pointed out to us by L. Takhtajan. and the Schwarz derivative.
The Green function of the Dirichlet boundary problem for the Laplace operator in $`D_{}`$ expressed through the conformal map $`w(z)`$ is:
$$G(z,\zeta )=\text{log}\left|\frac{w(z)w(\zeta )}{w(z)\overline{w}(\overline{\zeta })1}\right|$$
(45)
Combining (41) and (42), and using the notation (40), we represent the Green function as follows:
$$2G(z,\zeta )=2\text{log}|z^1\zeta ^1|+\left(_{t_0}+D(z)+\overline{D}(\overline{z})\right)\left(_{t_0}+D(\zeta )+\overline{D}(\overline{\zeta })\right)\text{log}\tau $$
(46)
This formula generalizes (37) since (46) becomes the real part of (37) as $`\zeta \mathrm{}`$. (As $`\zeta \mathrm{}`$, $`G(z,\zeta )\text{log}|w(z)|`$.) Note also that the real part of (37) can be written in the form
$$\mathrm{\Phi }(z,\overline{z})=2t_0\text{log}|z|+\left(D(z)+\overline{D}(\overline{z})\right)\text{log}\tau $$
where $`\mathrm{\Phi }`$ is the potential (2) ($`zD_{}`$).
The l.h.s. of eq.(41) generalizes the Schwarz derivative of the conformal map
$$T(z)\frac{w^{\prime \prime \prime }(z)}{w^{}(z)}\frac{3}{2}\left(\frac{w^{\prime \prime }(z)}{w^{}(z)}\right)^2=6\text{lim}_{_{z\zeta }}_z_\zeta \text{log}\frac{w(z)w(\zeta )}{z\zeta }$$
(47)
Taking the limit $`\zeta z`$ of both sides of (41), we get a relation between the Schwarz derivative and the $`\tau `$-function:
$$T(z)=6z^2\underset{k,n1}{}z^{kn}\frac{^2\text{log}\tau }{t_kt_n}$$
(48)
The latter can be used as an alternative definition of the $`\tau `$-function.
7. Integrable structure of conformal maps. Eqs.(34-36) allow one to say that the differential
$$d\mathrm{\Omega }=Sdz+\text{log}wdt_0+\underset{k=1}{\overset{\mathrm{}}{}}(H_kdt_k\overline{H}_kd\overline{t}_k)$$
(49)
generates the set of Hamiltonian equations for deformations of the curve due to variation of $`t_k`$:
$$_{t_k}S(z)=_zH_k(z),_{\overline{t}_k}S(z)=_z\overline{H}_k(z),$$
(50)
where we set $`H_0(z)=\mathrm{log}w(z)`$. The equations are consistent due to commutativity of the flows:
$$\left(_{t_j}H_k\right)_z=\left(_{t_k}H_j\right)_z=_{t_j}_{t_k}\mathrm{\Omega }(z)$$
(51)
Equations (50) are more transparent being written in terms of canonical variables. The differential $`d\mathrm{\Omega }`$ suggests that the pairs $`\mathrm{log}w,t_0`$ and $`z(w),S(z(w))`$ are canonical and establishes the symplectic structure for conformal maps. Indeed, treating $`w`$ as an independent variable, one rewrites eq. (34) as
$$\{z(w),S(z(w))\}=1$$
(52)
where the Poisson bracket $`\{,\}`$ is with respect to $`\mathrm{log}w`$ and the area $`t_0`$ is defined as
$$\{f,g\}=w\frac{f}{w}\frac{g}{t_0}w\frac{g}{w}\frac{f}{t_0}$$
(53)
where the derivatives with respect to $`t_0`$ are taken at fixed $`t_k`$ and $`w`$.
The other flows read
$$\frac{z(w)}{t_k}=\{H_k,z(w)\}$$
(54)
$$\frac{S(z(w))}{t_k}=\{H_k,S(z(w)\},$$
(55)
and similarly for the flows with respect to $`\overline{t}_k`$. Now the Hamiltonian functions $`H_k`$ and $`\overline{H}_k`$ are degree $`k`$ polynomials of $`w`$ and $`w^1`$ respectively.
The consistency conditions (50) now take the form of the zero-curvature conditions:
$$_{t_j}H_i_{t_i}H_j+\{H_i,H_j\}=0,$$
(56)
$$_{t_j}\overline{H}_i+_{\overline{t}_i}H_j+\{\overline{H}_i,H_j\}=0.$$
(57)
The infinite set of the Poisson-commutating flows form a Whitham integrable hierarchy . Eqs. (54,55) are the Lax-Sato equations for the hierarchy. They generate an infinite set of differential equations for the coefficients (potentials) $`u_j`$ of the inverse conformal map (22). The first equation of the hierarchy is
$$_{t_1\overline{t}_1}^2\varphi =_{t_0}\mathrm{exp}(_{t_0}\varphi ),_{t_0}\varphi =\mathrm{log}r^2.$$
(58)
The integrable hierarchy describing conformal maps is also known in the soliton literature as the dispersionless Toda lattice hierarchy, or SDiff(2) Toda hierarchy (see the next section)<sup>6</sup><sup>6</sup>6 A relation between conformal maps of slit domains and special solutions to equations of hydrodynamic type (Benney equations) was first observed by Gibbons and Tsarev . The algebra Sdiff(2) of area-preserving diffeomorphisms is the symmetry algebra of this hierarchy . Eqs. (54-57) describe infinitesimal deformations of the curve such that the area $`t_0`$ is kept fixed.
The integrable hierarchy possesses many solutions. The particular solution relevant to conformal maps is selected by the subsidiary condition (52). This condition, known as dispersionless string equation, has already appeared in the study of the $`c=1`$ topological gravity and in the large $`N`$ limit of a model of normal random matrices . The latter is discussed in Sec. 9.
8. Toda lattice hierarchy and its dispersionless limit. Below we review the two dimensional Toda lattice hierarchy and show that its dispersionless limit gives the equations describing conformal maps (35,36,54,55).
The 2D Toda hierarchy is defined by two Lax operators
$$L=r(t_0)e^{\mathrm{}\frac{}{t_0}}+\underset{k=0}{\overset{\mathrm{}}{}}u_k(t_0)e^{k\mathrm{}\frac{}{t_0}}$$
(59)
$$\overline{L}=e^{\mathrm{}\frac{}{t_0}}r(t_0)+\underset{k=0}{\overset{\mathrm{}}{}}e^{k\mathrm{}\frac{}{t_0}}\overline{u}_k(t_0)$$
(60)
acting in the space of functions of $`t_0`$ where the coefficients $`u_j`$ and $`\overline{u}_j`$ are functions of $`t_0`$ and also of two independent sets of parameters (“times”) $`t_k`$ and $`\overline{t}_k`$. Note that $`u_k`$ and $`\overline{u}_k`$ as well as $`t_k`$ and $`\overline{t}_k`$ in (59,60) are not necessarily complex conjugate to each other, although we choose them to be so.
The dependence of the coefficient $`u_k`$ and $`\overline{u}_k`$ on $`t_k`$ and $`\overline{t}_k`$ are given by the Lax-Sato equations:
$$\mathrm{}\frac{L}{t_k}=[H_k,L]$$
(61)
$$\mathrm{}\frac{L}{\overline{t}_k}=[L,\overline{H}_k]$$
(62)
and similar equations for $`\overline{L}`$. The flows are generated by
$$H_k=\left(L^k\right)_++\frac{1}{2}\left(L^k\right)_0$$
(63)
$$\overline{H}_k=\left(\overline{L}^k\right)_{}+\frac{1}{2}\left(\overline{L}^k\right)_0$$
(64)
where the symbol $`\left(L^k\right)_\pm `$ means positive (negative) parts of the series in the shift operator $`e^{\mathrm{}\frac{}{t_0}}`$. The first equation of the hierarchy is the Toda lattice equation
$$_{t_1\overline{t}_1}^2\varphi (t_0)=e^{\varphi (t_0+\mathrm{})\varphi (t_0)}e^{\varphi (t_0)\varphi (t_0\mathrm{})},$$
(65)
where $`r^2=e^{\varphi (t_0+\mathrm{})\varphi (t_0)}`$.
The spectrum of the Lax operator is determined by the linear problem $`L\mathrm{\Psi }=z\mathrm{\Psi }`$. The wave function $`\mathrm{\Psi }`$ is expressed through the $`\tau `$-function $`\tau _{\mathrm{}}`$ of the dispersionfull hierarchy (61, 62) by the following formula:
$$\mathrm{\Psi }(z;t_0,t_1,t_2,\mathrm{})=\tau _{\mathrm{}}^1(t_0,t_1,t_2,\mathrm{})z^{t_0/\mathrm{}}e^{\frac{1}{\mathrm{}}_{k>0}t_kz^k}e^{\mathrm{}_{k>0}\frac{z^k}{k}\frac{}{t_k}}\tau _{\mathrm{}}(t_0,t_1,t_2,\mathrm{})$$
(66)
Among many solutions of the hierarchy, one is of particular interest. It is selected by the string equation
$$[L,\overline{L}]=\mathrm{}$$
(67)
This solution is known to describe the normal matrix model at finite size of matrices .
The dispersionless limit of the Toda hierarchy is a formal semi-classical limit $`\mathrm{}0`$. To proceed we notice that the shift operator $`W=e^{\mathrm{}\frac{}{t_0}}`$ obeys the commutation relation $`[W,t_0]=\mathrm{}W`$. In the semiclassical limit it is supposed to be replaced by the canonical variable $`w`$ with the Poisson bracket $`\{\text{log}w,t_0\}=1`$. The Lax operator then becomes a $`c`$-valued function which is identified with the inverse conformal map $`z(w)`$ (22). Similarly, $`\overline{L}`$ is identified with $`S(z(w))`$. In their turn, the Lax-Sato equations (61,62) are identified with eqs.(54,55) for the conformal map. In the same fashion the dispersionless limit of the string equation (67) is identified with eq.(52). The semiclassical limits of the wave function and the $`\tau `$-function give the generating function $`\mathrm{\Omega }`$ and the dispersionless $`\tau `$-function: $`\mathrm{\Psi }e^{\mathrm{\Omega }/\mathrm{}}`$, $`\tau _{\mathrm{}}e^{(\mathrm{log}\tau )/\mathrm{}^2}`$. Similarly, eq. (43) is a semiclassical limit of the Hirota equation for the $`\tau `$-function of the 2D Toda hierarchy.
9. The $`\tau `$-function of the conformal map as large $`N`$ matrix integral. The integrable structure of conformal maps is identical to the one observed in a class of random matrix models related to noncritical string theories. Moreover, there exists a random matrix model whose large $`N`$ limit reproduces exactly the $`\tau `$-function for analytic curves.
Consider the partition function of the ensemble of normal random $`N\times N`$ matrices <sup>7</sup><sup>7</sup>7Earlier V. Kazakov pointed to us that the Lax equations (61,62) are generated by the Hermitian 2-matrix model with complex conjugated potentials. The latter and the normal matrix model have an identical $`1/N`$-expansion., with the potential (13):
$$\tau _{\mathrm{}}[t,\overline{t}]=𝑑M𝑑M^{}e^{\frac{1}{\mathrm{}}\mathrm{Tr}V(M,M^{})}$$
(68)
A matrix is called normal if it commutes with its Hermitian conjugated $`[M,M^{}]=0`$. Passing to the eigenvalues $`\mathrm{diag}(z_1,\mathrm{},z_N)`$ of the matrix $`M`$, one obtains the measure of the integral in a factorized form $`dMdM^{}_{i=1}^Ndz_id\overline{z}_i_{k<j}(z_kz_j)(\overline{z}_k\overline{z}_j)`$. Then the partition function is represents a two-dimensional Coulomb gas in the potential (13)
$$\tau _{\mathrm{}}[t,\overline{t}]=\underset{k=1}{\overset{N}{}}dz_kd\overline{z}_ke^{\frac{1}{\mathrm{}}V(z_k,\overline{z}_k)}\underset{i<j}{}e^{2\mathrm{log}|z_iz_j|}.$$
(69)
To proceed to the large $`N`$ limit one introduces a parameter $`t_0=\mathrm{}N`$ and expresses the integrand in terms of density of eigenvalues as $`e^{\mathrm{}^2\{\rho ,V\}}`$, where $`\{\rho ,V\}`$ is given by eq. (14). Then, the large $`N`$ ($`\mathrm{}0)`$ limit yields to the variational principle of Sec.3. In the large $`N`$ limit the eigenvalues of the matrix homogeneously fill the domain $`D_+`$ bound by the curve, characterized by the harmonic moments $`t_k`$ and the area $`t_0`$ and leads to the $`\tau `$-function defined by eq. (16). Other objects introduced in Secs. 3-7 can also be identified with expectation values of the matrix model. In particular the moments $`v_k`$ (eq.(4)) are
$$v_k=\mathrm{}\mathrm{Tr}M^k$$
and $`\mathrm{\Omega }^{}\frac{1}{2}v_0=\mathrm{}\mathrm{Tr}\mathrm{log}(zM)`$.
In order to identify the Lax operator, we follow . Introduce the basis of orthogonal polynomials $`P_n(z)=h_nz^n+\mathrm{}(n0),`$ by the orthonormality relations
$$m|nd^2z\overline{P_n(z)}e^{\frac{1}{\mathrm{}}V(z,\overline{z})}P_m(z)=\delta _{m,n}$$
(70)
The polynomials are uniquely defined by the potential $`V`$ up to phase factors. It is easy to see that the $`\tau `$-function is given by the product of the coefficients $`N!h_nh_{n1}\mathrm{}h_0`$ of the highest powers of the polynomials $`P_n(z)=h_nz^n+\mathrm{}`$. Then Lax operators $`L`$ and $`\overline{L}`$ appear as the operators $`m|z|n`$ and $`m|\overline{z}|n`$. Since $`zP_n(z)`$ can be expressed through polynomials of the degree not grater than $`n`$, one may represent $`m|z|n`$ and $`m|\overline{z}|n`$ in terms of shifts operators $`W=e^{\mathrm{}\frac{}{t_0}}`$ in the form of (59,60), where $`r(t_0=\mathrm{}n)=h_n/h_{n+1}`$.
Similar arguments allow one to identify the flows. Consider a variation of some operator $`m|O|n`$ under a variation of $`t_k`$. We have $`\mathrm{}_{t_k}m|O|n=m[|H_k,O]|n`$, where $`H_k=A_kA_{}^{}{}_{k}{}^{}`$ and $`m|A_k|n=m|_{t_k}|n.`$ Obviously $`H_k=L^k(W)+`$ negative powers of $`W`$. Choosing $`O`$ to be $`\overline{L}`$ (see (60)) which consists on $`W^1`$ and positive powers of $`W`$, one concludes that $`H_k`$ does not consists of negative powers of $`W`$. This brings us to eq.(63).
Finally, the operator $`D=m|\mathrm{}_z|n`$ is equal to
$$D=\overline{L}\underset{k1}{}kt_kL^{k1},$$
(71)
The Heisenberg relation $`[D,L]=\mathrm{}`$ prompts the string equation (67).
The matrix model also offers an effective method to derive eqs.(37-43) (see e.g. ).
10. $`\tau `$-function and spectral properties of the Dirichlet problem. This subject is under current study.
We thank M. Brodsky, V. Kazakov, S.P. Novikov and L.Takhtajan for valuable comments and interest to this work.
The work of I.K. is supported in part by European TMR contract ERBFMRXCT960012 and EC Contract FMRX-CT96-0012. The work of I.Kr. is supported in part by NSF grant DMS-98-02577. P.W. would like to thanks P.Bleher and A.Its for the hospitality in MSRI during the workshop on Random Matrices in spring 1999. I.Kr. and A.Z. have been partially supported by CRDF grant 6531. P.W. and A.Z. have been partially supported by grants NSF DMR 9971332 and MRSEC NSF DMR 9808595. The work of A.Z. was supported in part by grant INTAS-99-0590 and RFBR grant 00-02-16477. He also thanks for hospitality the Erwin Schrödinger Institute in Vienna, where this work was completed.
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# Relations among Heegner cycles on familes of abelian surfaces
## 1 Introduction
A conjecture of Beilinson and Bloch asserts that if $`𝐕`$ is a smooth, geometrically irreducible, projective variety defined over a number field $`F`$, then
$$rankCH_{hom}^p(𝐕_F)=ord_{s=p}L_F(H^{2p1}(𝐕),s).$$
Here $`CH_{hom}^p`$ denotes the group of codimension $`p`$ cycles homologically equivalent to zero modulo rational equivalence on the variety. In particular, this says that the rank of the Chow group is finite. In the case when $`p=1`$ the finiteness of the rank is the Mordell-Weil theorem. However very little is known about this conjecture in the case when $`p>1`$. The purpose of this paper is to make some progress towards this conjecture in the case of codimension 2 cycles on certain types of varieties.
The varieties in question are compactifications of modular families of abelian surfaces over modular or Shimura curves. Schoen \[Sc1\], in the case of modular curves and later Besser \[Be\], in the case of Shimura curves, constructed certain codimension 2 nullhomologous cycles supported in fibres over complex multiplication points on the modular curve. (In fact they showed that these cycles are in general not algebraically equivalent to 0 so they are actually in the so called Griffiths group). These ‘CM cycles’ are the analogues of Heegner points on modular curves and are similarly defined over number fields. Further, one can take certain traces of these cycles, analogous to Heegner divisors, to get cycles defined over fixed number fields $`F`$ and hence they give elements of the groups in question. The purpose of this paper is to construct relations of rational equivalence between these cycles in the hope that one can construct enough relations to prove that the groups are finite dimensional.
Our strategy is to use the localization sequence for higher Chow groups \[Bl1\]. From that one sees that to construct relations between codimension 2 cycles it is enough to construct elements of the group $`CH^2(𝒜_\eta ,1)`$ where $`𝒜_\eta `$ is the generic abelian surface of the family and then compute their boundary. It turns out that there are some very natural elements constructed by Collino \[Co\] and interestingly, their boundary can be expressed in terms of the CM cycles.
We then have the following:
###### Theorem 1.1.
Let $`𝐗(𝐃_\mathrm{𝟎},𝐍)`$ be a Shimura curve parametrising abelian surfaces with endomorphism ring an Eichler order (loc. cit. Section 2) of level $`N`$ in a division algebra of discriminant $`D_0`$. Let $`𝐖(𝐃_\mathrm{𝟎},𝐍)`$ denote (the non-singular compactification of) the universal abelian surface over $`𝐗(𝐃_\mathrm{𝟎},𝐍)`$. Let $`p,a,b`$ be the invariants which determine Hashi -moto’s model (ibid.) and let $`P`$ and $`Q`$ denote two $`2`$-torsion points in the generic fibre.
Then there are relations in $`CH^2(𝐖(𝐃_\mathrm{𝟎},𝐍))`$ of the form
$$\underset{[]}{}ϵ_{P,Q}([])\underset{r,s}{}\underset{d^2|\mathrm{\Delta }}{}d𝒵_{\mathrm{\Delta }/d^2,[],s_1/d,s_1/d}0$$
where $``$ runs through all even theta characteristics, $`r,s`$ such that $`\frac{n^2pr^2}{4}_0`$, $`\frac{4bDn^2s^2}{4}_0`$, $`n`$ is coprime with $`2k`$, $`r,s,n`$ mutually coprime, $`s_1=p(s/2)aDr)`$ and
$$\mathrm{\Delta }=\frac{ps^24aDrs+4bDr^24Dn^2}{4}$$
with $`d`$ running through all $`d^2|\mathrm{\Delta }`$ such that $`\mathrm{\Delta }/d^2`$ is still a discriminant.
Here $`𝒵_\mathrm{\Delta }`$ is a Heegner Cycle of discriminant $`\mathrm{\Delta }`$ (loc. cit Section 5 ), a codimension 2 cycle which is homologous to zero in $`CH^2(𝐖(𝐃_\mathrm{𝟎},𝐍))`$ and $`ϵ_{P,Q}([])`$ is a sign function which depends on the level and the points $`P`$ and $`Q`$ (loc. cit. Section 4).
Relations between Heegner points on modular curves have been studied before by many other authors. In Zagier \[Za\] and van der Geer \[Ge\], they construct relations between these points which are very similar to ours, the main difference being that our relations involve the level 2 structure while theirs do not.
As is known \[Za\],\[G-Z\], Heegner points are very closely related to coefficients of modular forms of weight $`\frac{3}{2}`$ and relations of rational equivalence between Heegner points imply relations between coefficients of these forms. Similarly, Heegner cycles are supposed to be related to coefficients of modular forms of weight $`\frac{5}{2}`$. Some evidence of this can be found in \[Sc2\],\[Zh\] and \[Ne\]. So relations of rational equivalence between these cycles should give rise to relations between coefficients of such modular forms.
More interestingly, recently Borcherds \[Bo\] has used his new constructions of automophic forms to construct relations between Heegner points on modular and Shimura curves and more generally, relations between special divisors on modular varieties, and prove that these special subvarieties are related to coefficients of modular forms. Starting from certain meromorphic modular forms he constructs other automorphic forms whose divisors are on such special subvarieties. It would be very interesting to see if his methods can be generalized to construct Collino type elements of the higher Chow groups and use them to prove that the Heegner cycles are related to coefficients of modular forms. At the moment it is not clear to me how this can be done, but a weaker statement, namely computation of the regulator of Collino’s element may be a little more tractable.
The outline of the paper is as follows: In Section 2 we describe the varieties and cycles in question. In Section 3 we introduce the higher Chow groups and the localization sequence that we use. In Section 4 we describe Collino’s construction of the elements of the higher Chow groups. In Section 5 we compute the boundary of the elements and relate them to the CM cycles. Finally, in section 6 we put them all together to get our main result and give some examples.
Acknowledgements: I would like to thank my advisor Spencer Bloch for suggesting the problem, Madhav Nori for refering me to Collino’s work, Patrick Brosnan for helping me get started and Najmuddin Fakruddin for his invaluable help and innumerable suggestions and corrections. I would also like to thank the hospitality of the Institute for Advanced Study where this paper was completed. Finally, I would like to thank Chad Schoen and the referee for their very useful comments.
## 2 Preliminaries
### 2.1 Eichler Orders
The varieties we consider are universal families of abelian surfaces whose endomorphism rings contain an order in a quaternionic division algebra. To describe precisely we use the theorems of Hashimoto \[Ha\] and Roberts \[Ro\]. Another standard reference is \[Vi\].
Let $`𝐁`$ be an indefinite quaternion algebra over $``$ with discriminant $`D_0`$. An Eichler Order of level $`D=D_0N`$ is an order $`𝒪`$ in $`𝐁`$ such that
$$𝒪_{\mathrm{}}\left(\begin{array}{cc}& \\ 0& \end{array}\right)mod\mathrm{}forall\mathrm{}|N$$
Let $`𝒪`$ be an Eichler order of level $`D=D_0N`$ where $`N`$ is a positive integer prime to $`D_0`$. Let $`𝒮`$ be the set of primes where the division algebra is ramified ( i.e those primes $`\mathrm{}`$ such that $`𝐁_{\mathrm{}}M_2(_{\mathrm{}})`$). Choose an auxiliary prime $`p`$ such that the Hilbert symbol $`(D,p)_{\mathrm{}}=1`$ if and only if $`\mathrm{}𝒮`$. Such a $`p`$ is guaranteed by Dirichlet’s theorem on primes in arithmetic progressions. Let $`a`$ and $`b`$ be such that $`a^2D+1=bp`$. Then we have the following theorem described in Hashimoto, \[Ha\], though parts of it existed in some form in earlier papers.
###### Theorem 2.1.
Let $`𝐁`$ be a quaternionic division algebra of discriminant $`D_0`$. Then:
1. $`𝐁ijij`$ where $`i^2=D`$, $`j^2=p`$ and $`ij=ji`$.
2. The order $`𝒪e_1e_2e_3e_4`$ where $`e_1=1,e_2=(1+j)/2,e_3=(aDj+ij)/p`$ and $`e_4=(i+ij)/2`$
3. There is a skew-symmetric pairing on $`𝐁`$ which is $``$ valued on $`𝒪`$ given by $`E(x,y)=tr(x\overline{y}i^1)`$ where $`\overline{y}`$ denotes conjugation. Further, the elements $`\eta _1=e_4\frac{p1}{2},\eta _2=aDe_1e_3,\eta _3=e_1`$ and $`\eta _4=e_2`$ are a symplectic basis for $`𝒪`$
Since $`𝐁`$ is indefinite, any two Eichler orders of level $`D`$ are conjugate \[B-D\], so there is no loss of generality in working with this model and we can fix the isomorphism of $`𝐁M_2()`$ by defining
$$\mathrm{\Phi }_{\mathrm{}}(i)=\left(\begin{array}{cc}0& 1\\ D& 0\end{array}\right),\mathrm{\Phi }_{\mathrm{}}(j)=\left(\begin{array}{cc}\sqrt{p}& 0\\ 0& \sqrt{p}\end{array}\right)$$
We will call this description of the quaternion algebra Hashimoto’s model.
Orientations
The above description is not rigid enough. So, following \[Ro\],\[B-D\] we rigidify the definition further by defining an Oriented Eichler Order of level $`D_0N`$ as follows. For each prime $`\mathrm{}`$ dividing $`D_0`$ there are two algebra homomorphisms,
$$𝔬_{\mathrm{}}:𝒪𝔽_{\mathrm{}}𝔽_{\mathrm{}}$$
$$\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)a\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)d$$
and similarly for a prime $`\mathrm{}`$ dividing $`N`$ there are two distinct algebra homomorphisms
$$𝔬_{\mathrm{}}:𝒪𝔽_{\mathrm{}}𝔽_\mathrm{}^2$$
$$\left(\begin{array}{cc}a& b\\ \mathrm{}b^\sigma & a^\sigma \end{array}\right)a\left(\begin{array}{cc}a& b\\ \mathrm{}b^\sigma & a^\sigma \end{array}\right)a^\sigma $$
where $`\sigma `$ is the non-trivial automorphism of $`𝔽_\mathrm{}^2`$. An oriented Eichler Order is an Eichler order along with a choice of one of these homomorphism $`𝔬_{\mathrm{}}`$, called a $`\mathrm{}`$-orientation , for each $`\mathrm{}`$ dividing $`D=D_0N`$.
### 2.2 Shimura Curves
With the model of $`𝐁`$,$`𝒪`$, the isomorphism $`\mathrm{\Phi }_{\mathrm{}}`$ and the orientation we can form the family of surfaces we are interested in. Let $`\mathrm{\Gamma }(1)=\mathrm{\Gamma }_𝒪(1)`$ be the group
$$\mathrm{\Gamma }_𝒪(1)=\{x𝒪|Nm(x)=1\}$$
where $`Nm`$ denotes the reduced norm. $`\mathrm{\Gamma }(1)`$ acts on the upper half plane $``$ through the embedding $`\mathrm{\Phi }_{\mathrm{}}`$. The quotient is an algebraic curve whose complex points represents a component of the moduli of abelian surfaces whose endomorphism rings contain $`𝒪`$. In this case the generic abelian surface has endomorphism ring actually equal to $`𝒪`$.
We further assume that the abelian surfaces have full level $`2k`$ structure for some $`k>1`$. This is because we will use the level 2 structure and we need further level structure to ensure that there is a universal family over the curve. The $`k`$-level structure will not play any other part in our calculations, so will be suppressed but not forgotten in the remaining part of this paper.
The universal family can be described as follows. Let
$$\mathrm{\Gamma }(2k)=\mathrm{\Gamma }_𝒪(2k)=\{\gamma \mathrm{\Gamma }(1)|\gamma Ker(\mathrm{\Gamma }(1)Aut(𝒪/2k𝒪))\}$$
The Jacobi Group $`\mathrm{\Gamma }(2k)𝒪`$ with multiplication defined by $`(\gamma ,x)(\mu ,y)=(\gamma \mu ,x\mu +y)`$ acts on $`\times ^2`$ via
$$(\left(\begin{array}{cc}a& b\\ c& b\end{array}\right),x)(\tau ,z_1,z_2)=(\frac{a\tau +b}{c\tau +d},(c\tau +d)^1(\left(\begin{array}{c}z_1\\ z_2\end{array}\right)+\mathrm{\Phi }_{\mathrm{}}(x)\left(\begin{array}{c}\tau \\ 1\end{array}\right)))$$
and the quotient is a complex threefold which is a family of abelian surfaces over the Shimura curve $`/\mathrm{\Gamma }(2k)`$.
The fibre over a point $`\tau `$ is
$$𝒜_\tau =^2/\mathrm{\Phi }_{\mathrm{}}(𝒪)\left(\begin{array}{c}\tau \\ 1\end{array}\right)$$
and $`𝒜_\tau 𝒜_\tau ^{}`$ if and only if $`\tau =\gamma (\tau ^{})`$ for some $`\gamma \mathrm{\Gamma }(1)`$, with their level structures coinciding if and only if $`\gamma \mathrm{\Gamma }(2k)`$.
Let $`𝐘=𝐘(𝐃_\mathrm{𝟎},𝐍)`$ denote the curve $`/\mathrm{\Gamma }(2k)`$. If $`D_0=1`$ the curve is not compact but it can be compactified by adding finitely many cusps. Let $`𝐗=𝐗(𝐃_\mathrm{𝟎},𝐍)`$ denote the compactified curve. Let $`𝐖=𝐖(𝐃_\mathrm{𝟎},𝐍)`$ denote the desigularisation of the self product of the universal family over $`𝐗(𝐃_\mathrm{𝟎},𝐍)`$ as described in \[Sc1\]. Let $`𝒲=𝒲(D_0,N)=𝒜_\eta `$ denote the generic fibre of this family. As remarked earlier, it is an abelian surface with endomorphism ring $`𝒪`$.
Polarizations
There is an involution on $`𝒪`$ defined by $`d^{}=i^1\overline{d}i`$ where $`i`$ is as in Hashimoto’s theorem. The skew symmetric pairing on $`^2`$ defined by
$$<\mathrm{\Phi }_{\mathrm{}}(x)\left(\begin{array}{c}\tau \\ 1\end{array}\right),\mathrm{\Phi }_{\mathrm{}}(y)\left(\begin{array}{c}\tau \\ 1\end{array}\right)>=E(x,y)$$
$`(1)`$
where $`E`$ is the form coming from part $`3`$ of theorem $`2.1`$, is a Riemann form for $`A_\tau `$. This is a principal polarization and $`d^{}`$ is the Rosati involution corresponding to it. A well known theorem \[Mu\] pg 190, says that the fixed points under this involution correspond to elements of the Néron-Severi group of the abelian surface. It is easy to see that the elements $`1`$,$`j`$ and $`k`$ are invariant under this involution and are clearly linearly independent, so from that one can see that the rank of the generic Néron-Severi group is $`3`$. A little more work gives the following, which is a simple corollary to Hashimoto’s theorem, but will be useful later on.
###### Corollary 2.2.
The elements $`e_1=1`$, $`e_2=\frac{1+j}{2}`$ and $`e_3=\frac{aDj+ij}{p}`$ give a basis for the Néron-Severi of the generic fibre. Further, for ample $`g`$, the cup product is given by $`<g,g>=2Nm(\varphi _g)`$, where $`\varphi _g`$ is the endomorphism corresponding to $`g`$. With that, the intersection matrix is
$$\begin{array}{cccc}& e_1& e_2& e_3\\ & & & \\ e_1& 2& 1& 0\\ e_2& 1& \frac{1p}{2}& aD\\ e_3& 0& aD& 2bD\end{array}$$
where $`a^2D+1=bp`$
###### Proof.
From the Riemann-Roch theorem, one sees that, for an ample divisor $`g`$,
$$<g,g>=2\chi ((g))and\chi ((g))^2=deg(\varphi _g)$$
where $`(g)`$ is the corresponding line bundle and $`\varphi _g`$ is the corresponding endomorphism.
This, coupled with the formula for the degree in terms of the norm
$$deg(\varphi )=Nm(\varphi )^2$$
where $`Nm`$ is the reduced norm, shows that
$$<g,g>=\pm 2Nm(\varphi _g)$$
The sign can be determined from the fact that $`<1,1>=2`$ where $`1`$ represents the class of the principal polarization.
From now on we pick curves in the generic fibre which represent these classes $`e_i`$ and we denote them by $`e_i`$ as well. This will not affect anything as all of our results will be modulo divisors homologous to zero in the fibres. We also fix the polarization to be the one in equation (1) and denote it by $`\mathrm{\Theta }`$.
### 2.3 Humbert Invariants
A lemma in \[Be2\], Lemma 4.1 shows that there is an equivalence of categories between the category of even lattices of rank n with a element $`\zeta `$ of norm 2 with the category of lattices of rank (n-1) such that the quadratic form represents only numbers $`0`$ or $`1`$ mod $`4`$. The correspondence is given by
$$(M,\zeta )N(M,\zeta )N(M_N,\zeta )$$
where
$``$ $`N(M,\zeta )`$ is the lattice such that $`2N(2)`$ is the orthogonal complement to $`\zeta `$ in $`2M+\zeta `$
$``$ $`M_N`$ is the sublattice of $`(N(2)\zeta )`$ containing $`N\zeta `$ and all elements of the form $`\frac{1}{2}(x+(x,x)\zeta )`$ with $`x`$ in $`N`$.
Here $`N(k)`$ denotes the lattice which has the same underlying $``$ module as $`N`$ but with the pairing multiplied by $`k`$.
The Néron-Severi lattice of an abelian surface is an even lattice. A principal polarization is an element of norm 2. So a choice of principal polarization allows us to work with a smaller lattice, namely the primitive Néron-Severi lattice.
$$NS(𝒜,\mathrm{\Theta })=\{vNS(𝒜)|<v,\mathrm{\Theta }>=0\}$$
with the pairing given by
$$(v,v)_\mathrm{\Theta }=2<v,v>$$
More generally, one can work with the entire $`H^2(𝒜.)`$ and the primitive cohomology, $`H^2(𝒜,\mathrm{\Theta })`$.
Classically, this was first studied by Humbert \[Hu\], and so we call the pairing $`(,)_\mathrm{\Theta }`$ the Humbert norm and if $`vNS(𝒜,\mathrm{\Theta })`$ then $`(v,v)_\mathrm{\Theta }=(v,v)`$ is the Humbert Invariant of v, H(v).
From now on unless otherwise mentioned, we will always work with the Humbert norm on a lattice. From the corollary above and the description of the correspondence above we see that the Humbert lattice for the generic Néron-Severi is generated by
$$\overline{e}_2:=\frac{1}{2}(2e_2<e_2,e_1>e_1)and\overline{e}_3:=\frac{1}{2}(2e_3<e_3,e_1>e_1)$$
with intersection matrix
$$\begin{array}{ccc}& \overline{e}_2& \overline{e}_3\\ & & \\ \overline{e}_2& p& 2aD\\ \overline{e}_3& 2aD& 4bD\end{array}$$
This matrix has determinant 4D.
Humbert studied the moduli of abelian surfaces whose Néron-Severi contains an element with a non-zero Humbert invariant, which are now called Humbert surfaces. This amounts to saying that the rank of the Néron-Severi is at least 2 and he realised that this was equivalent to saying that the endomorphism ring contains an order in a real quadratic field. One can easily see that, in general, the endomorphism algebra is simply the Clifford algebra of the Humbert lattice.
Humbert’s definitions were more analytic and quite different from the above defintions, which are due to Kani. For some computations, however, it is more useful to use the analytic definition, so for that reason we will give it here. The equivalence of the two can be found in Kani \[Ka\].
Let $`V^{}`$ be the set of integral, skew symmetric 4x4 matrices, considered as a subgroup of $`M_4()`$. Let $`J`$ be the subgroup generated by $`\left(\begin{array}{cc}0& I_2\\ I_2& 0\end{array}\right)`$ and let $`V`$=$`V^{}/J`$.
There is an action of $`Sp_4()`$ on $`V^{}`$ given by
$$vM^tvM$$
for $`M`$ in $`Sp_4()`$ which leaves $`J`$ invariant and hence descends to an action on $`V`$. A $`v`$ in $`V^{}`$ looks like
$$v=\left(\begin{array}{cccc}0& a& b_1& b_2\\ a& 0& b_3& b_4\\ b_1& b_3& 0& d\\ b_2& b_4& d& 0\end{array}\right)=\left(\begin{array}{cc}A& B\\ B^t& D\end{array}\right)$$
and one defines the Humbert Invariant to be
$$(v,v)=H(v)=(b_1+b_4)^24(b_1b_4b_3b_2ad)=Tr(B)^24det(B)+ad$$
$`H(v)`$ depends only on $`v`$ mod $`J`$ and $`v`$ is said to be primitive if it is not divisable by a natural number $`2`$. A lemma due to Humbert says that two primitive elements $`v_1`$ and $`v_2`$ are equivalent if and only if $`H(v_1)=H(v_2)`$.
Let $`\tau =\left(\begin{array}{cc}\tau _1& \tau _2\\ \tau _2& \tau _3\end{array}\right)`$ be a point in the Siegel upper half space $`_2`$. Every primitive $`v`$ in $`V`$ determines a subvariety given by $`\tau _2`$ such that
$$\left(\begin{array}{cc}\tau & I_2\end{array}\right)v\left(\begin{array}{c}\tau \\ I_2\end{array}\right)=0$$
$$\tau A\tau +\tau BB^t\tau +D=0$$
and if $`\tau `$ satisfies such an equation it is said to satisfy a singular relation of invariant H(v) and the subvariety is called the Humbert surface of invariant $`H(v)`$. $`V`$ is called the space of singular relations. Equivalently, a singular relation can be described as a relation between the coefficients of $`\tau `$ of the form
$$v=\alpha \tau _1+\beta \tau _2+\gamma \tau _3+\delta (\tau _2^2\tau _1\tau _3)+ϵ=0$$
with invariant $`H(v)=\beta ^24\alpha \gamma 4\delta ϵ`$. Sometimes it is more convenient to work with this description. We will denote $`v`$ by $`(\alpha ,\beta ,\gamma ,\delta ,ϵ)`$. In fact \[Ka\], Prop.5.4, there is a basis for the primitive part of $`H^2(𝒜,)`$ for which this vector gives the corresponding 2-form in the other description of the Humbert surface.
### 2.4 CM points and CM cycles
The cycles we are interested are codimension 2 cycles on these families of abelian surfaces. There are two basic types, horizontal cycles and vertical cycles. Horizontal cycles are roughly images of the Shimura curve under some section. Vertical cycles are those which are supported in fibres over some points. In most fibres, all such cycles will come by restriction of a divisor on universal family to the special fibre. However, there are some points in whose fibres there are extra elements in the Néron-Severi of the fibre, hence extra codimension 2 cycles in the family.
CM Points
As we have seen before, the rank of the generic Néron-Severi is 3, so when there is an extra cycle, the rank goes up to 4. If $`\tau `$ is a point in the moduli where the rank of the Néron-Severi of the fibre $`𝒜_\tau `$ is 4, then the abelian surface $`𝒜_\tau `$ is necessarily isogenous to a product of elliptic curves with complex multiplication by an the imaginary quadratic field $`K=(\tau )`$ and $`End_0(𝒜_\tau )M_2(K)`$. Such a point $`\tau `$ is called a CM point.
Determining a CM point is equivalent to determining an embedding of an imaginary quadratic field $`q:K𝐁`$ such that $`𝐁KM_2(K)`$ as from Prop. 9.4 of \[Sh\] there is a unique point $`\tau `$ in $``$ such that
$$\mathrm{\Phi }_{\mathrm{}}(q(K^\times ))=\{\gamma \mathrm{\Phi }_{\mathrm{}}(B^\times GL_2^+())|\gamma \tau =\tau \}$$
We can normalize these embeddings as in \[Sh\] (4.4.3). Once having normalized the embedding of $`K𝐁`$ one has also normalized the embedding of $`KM_2(K)=𝐁K`$. $`KEnd(𝒜_\tau )`$ is an order of discriminant $`\mathrm{\Delta }`$ and hence will be denoted by $`𝒪_\mathrm{\Delta }`$. Through the embedding $`q`$ and the maps $`𝔬_{\mathrm{}}`$ for $`\mathrm{}`$ dividing $`D`$, we have also chosen an orientation of the order $`𝒪_\mathrm{\Delta }`$.
CM cycles
Let $`\alpha `$ be a traceless element of $`𝒪_\mathrm{\Delta }`$ with $`\alpha ^2=\mathrm{\Delta }`$, so $`\alpha `$ is purely imaginary. The Rosati involution extends to an involution of $`𝐁K`$ and acts by
$$(d\alpha )^{}=d^{}\overline{\alpha }=d^{}\alpha $$
Observe that the element $`i\alpha `$ is fixed by the Rosati involution. We define a CM cycle class to be
$$𝒵_\tau =i\alpha $$
the class of this element in the Néron-Severi of $`𝒜_\tau `$. By construction $`𝒵_\tau `$ is not a generic class. Further, in this case one can see that the intersection form on the Néron-Severi lattice being thought of as Rosati fixed elements is given by $`<v,v>=2det(\varphi _v)`$. From that one can see that the CM cycle is orthogonal to the generic Néron-Severi and the Humbert norm is
$$(𝒵_\tau ,𝒵_\tau )=4D\mathrm{\Delta }=4Ddisc(𝒪_\mathrm{\Delta })$$
###### Remark 2.3.
This definition of the CM cycle almost agrees with the definitions in \[Sc1\] and \[Be\], where it is defined to be the minimal generator of the orthogonal complement of the generic Néron-Severi. The definitions are the same when $`\mathrm{\Delta }`$ is odd, but when $`\mathrm{\Delta }`$ is even, our cycle is twice the generator.
The CM cycles are the analogues of CM points on modular and Shimura curves and have many similar properties. They are also defined over number fields. Schoen, in the modular case and Besser, in the Shimura curve case, show that they are homologous to zero in $`CH^2(𝐖)`$ and so give rise to interesting elements in $`CH_{hom}^2(𝐖)`$. In fact he shows that in general they are non-trivial in the Griffiths group.
## 3 Higher Chow groups and the Localization sequence
In this section we introduce the higher Chow groups that we use and the Localization theorem. No proofs are given. The proofs of these results may be found in \[Bl1\].
### 3.1 The group $`CH^2(𝒜,1)`$
The group we use is $`CH^2(𝒳,1)`$ where $`𝒳`$ is surface. It is defined as a a certain subquotient of the group of codimension 2 cycles on $`𝒳\times 𝔸^1`$. A theorem of Bloch’s \[Bl1\] identifies it with the $`𝒦`$-cohomology group $`H^1(𝒳,𝒦_2)`$, where $`𝒦_2`$ is the sheaf coming from the presheaf given by
$$UK_2(U)$$
With this description and the Bloch-Gersten-Quillen resolution of the sheaf $`𝒦_2`$ one can see that an element of $`CH^2(𝒳,1)`$ is represented by a formal sum
$$(𝒞_i,f_i)$$
where $`𝒞_i`$ are curves on $`𝒳`$ and $`f_i`$ are functions on the $`C_i`$ such that they satisfy the cocycle condition
$$div(f_i)=0$$
In $`CH^2(𝒳,1)`$ there are certain elements coming from the product structure on Chow groups,
$$CH^1(𝒳,1)CH^1(𝒳,0)CH^2(𝒳,1)$$
$`(1)`$
A theorem of Bloch’s \[Bl1\] shows that
$$CH^1(𝒳,1)=\mathrm{\Gamma }(𝒳,𝒪_𝒳^{})CH^1(𝒳,0)=CH^1(𝒳)=Pic(𝒳)$$
These elements are of the form
$$(𝒞_i,f_i)$$
where $`𝒞_i`$ are divisors on $`𝒳`$ and $`f_i`$ are constant functions hence automatically satisfy the cocycle condition as $`div(f_i)=0`$. Such elements are called decomposable. More generally one can construct more such elements by looking at norms of elements of extensions of the base field. Let $`CH^2(𝒳,1)_{dec}`$ denote the subgroup of $`CH^2(𝒳,1)`$ generated by such cycles. The group
$$CH^2(𝒳,1)_{indec}=CH^2(𝒳,1)/CH^2(𝒳,1)_{dec}$$
is called the subgroup of indecomposable cycles and is sometimes non-trivial.
### 3.2 The Localization sequence
The localization sequence for higher Chow groups is an extension of the usual localization sequence for Chow groups to the left. Suppose $`𝒳`$ is a smooth projective variety and $`𝒴𝒳`$ is a closed divisor, then we have an exact sequence
$$\mathrm{}CH^2(𝒳,1)CH^2(𝒳\backslash 𝒴,1)\stackrel{}{}CH^1(𝒴,0)CH^2(𝒳,0)\mathrm{}$$
In particular, if $`𝒳`$ is a family of surfaces over a curve $`𝒵`$ and $`𝒴_z`$ is the fibre over a point $`z𝒵`$, then we have a sequence
$$\mathrm{}CH^2(𝒳,1)CH^2(𝒴_\eta ,1)\stackrel{}{}\underset{z𝒵}{}CH^1(𝒴_z)CH^2(𝒳)..$$
where $`𝒴_\eta `$ is the generic fibre of the family. So, from this our strategy becomes clear: to construct relations of rational equivalence between codimension 2 cycles supported in fibres of a family of surfaces it suffices to construct elements of the higher Chow group of the generic fibre and compute the boundary under the map $``$. If $`(𝒞_i,f_i)`$ is an element of $`CH^2(𝒴_\eta ,1)`$, so now the $`𝒞_i`$ are curves in the generic fibre, the boundary of the element is
$$((𝒞_i,f_i))=div(f_i)$$
where $`div(f_i)`$ denotes the divisor of $`f_i`$ being though of as a function on the closure $`\overline{𝒞}_i`$ of $`𝒞_i`$. This divisor may have vertical componenents. For example, if $`(𝒞,f)`$ is a decomposable element in $`CH^2(𝒴_\eta ,1)`$, where $`f`$ is now a function on the base,
$$((𝒞,f))=a_i𝒞_{x_i}$$
where $`div(f)=a_ix_i`$ and $`𝒞_{x_i}`$ denotes the restriction of $`𝒞`$ to the fibre over the point $`x_i`$. Note that in this manner one can only get relations between cycles obtained by the restriction of generic cycles to the various fibres. Since the cycles we are interested in are not of this type, one cannot hope to use decomposable elements for our purposes. However, on the generic abelian surface there exist indecomposable elements In the next section we describe these elements constructed by Collino \[Co\].
## 4 Collino’s elements of $`CH^2(𝒜,1)`$
### 4.1 Construction of the Element
In \[Co\], Collino constructs certain elements of $`CH^g(𝒜,1)`$ where $`𝒜`$ is the Jacobian of a hyperelliptic genus $`g`$ curve. In particular, if $`𝒜`$ is a simple, principally polarized abelian surface, it is the Jacobian of a smooth genus 2 curve, so one can use his construction. It is as follows:
Let $`𝒜=Jac(𝒞)`$ where $`𝒞`$ is a genus 2 curve which represents the principal polarization. Since $`𝒞`$ is hyperelliptic there is a function $`f:𝒞^1`$ such that $`divf=2(P)2(Q)`$ where $`P`$ and $`Q`$ are two ramification points.
There are maps $`i_P(x)=x(P)`$ and $`i_Q(x)=x(Q)`$ from $`𝒞Jac(𝒞)`$. Let $`𝒞_P`$ and $`𝒞_Q`$ denote the images of $`𝒞`$ under $`i_P`$ and $`i_Q`$ respectively and let $`f_P`$ and $`f_Q`$ denote the function $`f`$ being thought of as functions on $`𝒞_P`$ and $`𝒞_Q`$ respectively. Then
$$(C_P,f_P)+(C_Q,f_Q)$$
is an element of $`CH^2(𝒜,1)`$. This is because
$$div(f_P)=2(O)2(QP)anddiv(f_Q)=2(PQ)2(O)$$
but $`(PQ)=(QP)`$ in $`Jac(𝒞)`$, so
$$div(f_P)+div(f_Q)=2(O)2(QP)+2(PQ)2(O)=0.$$
Collino proves that this element is indecomposable on the generic abelian surface. Note that the exact choice of the function $`f`$ is not that important as any other choice of function with the same divisor would give the same class in the group of indecomposable cycles.
### 4.2 Computation of the Boundary
Let $`𝐗`$ as before be the Shimura curve with full level $`2k`$ structure and $`𝐖`$ and $`𝒲`$ be the universal family and the generic fibre respectively. Since we have full level 2 structure, by making a choice of two 2-torsion sections $`P`$ and $`Q`$, we can apply the construction in section $`4.1`$ to the generic fibre to get an element of $`CH^2(𝒲,1)=CH^2(𝒜_\eta ,1)`$. In this section we will compute the boundary of this element.
The boundary of this element of $`CH^2(𝒲,1)`$ is of the form
$$\underset{x}{}a_x𝒟_x$$
where $`𝒟_x`$ are codimension 1 cycles in the fibre over the point $`x`$. In many fibres, $`𝒟_x`$ is simply the restriction of the generic curves $`𝒞_P`$ and $`𝒞_Q`$ to the fibre, but in some cases it is more interesting. There are two such cases \[We\]. The first is when the point is a cusp, so the fibre is a degenerate abelian surface. In this case a theorem of Schoen \[Sc1\], essentially the Manin-Drinfeld principle, shows that these cycles are torsion in $`CH^2(𝐖)`$ so one can neglect them as we are only interested in the rational Chow group. The second case is when the curve breaks up into a sum of two elliptic curves intersecting at a point, such that the two ramification points lie on different components. Later on, in sections $`4.3`$ and $`5`$, we will describe precisely when this happens. In this case the cycle is not simply the restriction and involves the components in a non-trivial manner. To compute the boundary we do the following local computation.
###### Theorem 4.1.
Let $`\tau `$ be a point on $`𝐗`$ in whose fibre the genus 2 curve $`𝒞_P`$ degenerates into a sum of two elliptic curves $`_1`$ and $`_2`$ and such that $`O`$ lies on $`_1`$ and $`QP`$ lies on $`_2`$. Then
$$((C_P,f_P)+(𝒞_Q,f_Q))=2(_2_1)$$
up to the boundaries of decomposable elements and images of cycles homologous to 0 in the fibre.
###### Proof.
In a neighborhood of the point $`\tau `$ the boundary of $`(𝒞_P,f_P)`$ is of the form
$$((𝒞_P,f_P))=a_1+b_2+$$
where $`a`$ and $`b`$ are in $``$ and $``$ denotes the closure of the horizontal sections of $`div(f_P)`$. Locally, we can always find a decomposable element whose boundary is of the form $`b(_1+_2)`$ so subtracting this element from Collino’s element allows us to assume that the boundary is of the form
$$((𝒞_P,f_P))=a_1+$$
To show that the boundary is not simply the restriction of the closure of a divisor in the generic fibre amounts to showing that $`a0`$. To do this we intersect with $`_2`$ and use the fact that a function restricts to a divisor of degree 0 in a cycle not contained in its divisor. Intersecting with $`_2`$ gives
$$0=<(a_1+)._2>=a2$$
as by assumption, only $`QP`$ lies on $`_2`$ so $`(._2)=2`$
A similar computation for $`(𝒞_Q,f_Q)`$ shows that
$$((𝒞_Q,f_Q))=2_2^{}+^{}$$
where here $`C_Q=_1^{}+_2^{}`$ are translates of the $`_1`$ and $`_2`$ which split $`C_P`$ and $`^{}`$ is the closure of $`divf_Q`$.
The cocycle condition $`div(f_p)+div(f_Q)=0`$, gives $`+^{}=0`$, as $``$ and $`^{}`$ are the closures of these divisors. Adding the two and observing that $`(_1_2)(_1^{}_2^{})`$ is homologous to zero in the fibre (in fact it is torsion) we get our result.
###### Remark 4.2.
This calculation works only locally as, a priori, it is not clear that one can find a global decomposable element which has the required boundary. However, we will later give an argument that this is indeed the case.
### 4.3 Computing Signs
In this section we give a recipe for computing the signs of the boundary. As shown in the previous section, the sign depends on the position of the 2-torsion points $`O`$ and $`PQ`$ when the curve splits into a product of two elliptic curves. The best way to describe this is to use genus 2 theta functions with characteristics. A good reference for all the facts used here is \[Fr-Kr\].
Theta Functions
The degree 2 Theta function with characteristic $`[ϵ,ϵ^{}]`$ is defined to be
$$\mathrm{\Theta }[ϵ,ϵ^{}](z,\tau )=\underset{N^2}{}exp(2\pi i(\frac{1}{2}(N+\frac{ϵ}{2})\tau (N+\frac{ϵ}{2})^{tr})+(N+\frac{ϵ}{2})(z+\frac{ϵ^{}}{2})^{tr})$$
where $`\tau =\left(\begin{array}{cc}\tau _1& \tau _2\\ \tau _2& \tau _3\end{array}\right)`$ is a point on $`_2`$, $`z=[(z_1,z_2)]`$ is a point on $`^2`$ being thought of as a row matrix and $`[ϵ,ϵ^{}]=[(ϵ_1,ϵ_2),(ϵ_1^{},ϵ_2^{})]`$ is a pair of points in $`(/2)^2`$.
To a characteristic $`[ϵ,ϵ^{}]`$ we associate the $`2`$-torsion point $`I\frac{ϵ^{}}{2}^{tr}+\tau \frac{ϵ}{2}^{tr}`$ on the abelian surface $`𝒜_\tau `$ and call this the associated $`2`$-torsion point.
A theta function with characteristic is called odd or even depending on whether the corresponding function is an odd or even function. This is equivalent to $`ϵ_1ϵ_1^{}+ϵ_2ϵ_2^{}`$ being odd or even.
Up to a nowhere-vanishing factor, the theta function with characteristic is the same as the function that one would get by translating the theta function $`\mathrm{\Theta }[(0,0)(0,0)](z,\tau )`$ by the associated $`2`$-torsion point. This shows that there are six odd theta functions and ten even theta functions.
The functions on the moduli space given by
$$\mathrm{\Theta }[(ϵ_1,ϵ_2),(ϵ_1^{},ϵ_2^{})](\tau )=\mathrm{\Theta }[(ϵ_1,ϵ_2),(ϵ_1^{},ϵ_2^{})](0,\tau )$$
are called the thetanullwerte. Since odd theta functions, being odd functions, vanish identically at $`0`$ there are no odd thetanullwerte.
Let $`\mathrm{\Gamma }_2=Sp_4()`$ and $`\mathrm{\Gamma }_2(2)`$ be the principal congruence subgroup of level 2. From the transformation formula for the theta functions, one can see that the thetanullwerte descend to give modular forms of weight one-half on $`_2/\mathrm{\Gamma }_2(2)`$. A well known theorem, see for example \[Ge\], says that each thetanullwert vanishes precisely on one component of the moduli of products of elliptic curves.
The group $`\mathrm{\Gamma }_2/\mathrm{\Gamma }_2(2)Sp_4(/2)`$ acts on these theta functions or theta characteristics. It acts transitively on the set of odd theta characteristics and can be used \[Ig\] to get an isomorphism of $`Sp_4(/2)`$ with $`𝔖_6`$ . Let $`\mathrm{\Gamma }_2(o)`$ be the inverse image of $`𝔄_6`$ under this isomorphism and let $`\omega Sp_4()`$ be such that $`\omega \mathrm{\Gamma }_2(o)`$ and $`\omega ^2=I`$.
At a point $`\tau `$ where the polarization splits into a sum of two elliptic curves, $`𝒞=_1+_2`$, the existence of the curve $`_1`$ implies that the period $`\tau `$ satisfies a certain singular relation, $`v=(\alpha ,\beta ,\gamma ,\delta ,ϵ)`$. The curve $`_2`$ corresponds to the relation $`v`$. We can make a choice of $`_1`$ or $`_2`$ and look at the $`\mathrm{\Gamma }_2(o)`$ translates, namely
$$g_1g^{tr}\text{for all }g\text{ in }\mathrm{\Gamma }_2(o).$$
This allows us to make a uniform choice of elliptic curve $`_1`$, with respect to which we can compute the sign. The sign does not change within components.
To compute the sign, we first need to make a choice of an embedding of the genus 2 curve. In Farkas and Kra \[Fr-Kr\] chapter VII, they make a choice and with respect to their choice, the genus 2 curve $`𝒞_{P_1}`$ is given by the divisor odd theta function
$$\mathrm{\Theta }[(0,1),(0,1)](z,\tau )$$
and the six ramification points are the two torsion points associated to the characteristics
$$P_1=[(0,0),(0,0)]P_2=[(1,0),(0,0)]P_3=[(1,0),(1,1)]$$
$$P_4=[(1,1),(1,1)]P_5=[(1,1),(1,0)]P_6=[(0,0),(1,0)]$$
In the component of the moduli of products of elliptic curves given by $`\tau =\left(\begin{array}{cc}\tau _1& 0\\ 0& \tau _3\end{array}\right)`$, the odd theta function splits as
$$\mathrm{\Theta }[(0,1),(0,1)]((z_1,z_2),\tau )=\mathrm{\Theta }[(0,0)](z_1,\tau _1)\mathrm{\Theta }[(1,1)](z_2,\tau _3)$$
Since the zeroes of the elliptic theta functions are $`[(1,1)]=\frac{1}{2}(1+\tau _1)`$ and $`[(0,0)]=0`$ respectively, the elliptic curves meet at the 2-torsion point
$$R=[(1,0),(1,0)]$$
From this one can see that the points $`P_1,P_2`$ and $`P_6`$ lie one the elliptic curve $`z_1\times [(0,0)]`$ and the points $`P_3,P_4`$ and $`P_5`$ lie on the elliptic curve $`[(1,1)]\times z_2`$.
So, for example, if we choose $`P=P_1`$ and $`Q=P_3`$ and consider the element of $`CH^2(𝒲,1)`$ given by $`(𝒞_P,f_P)+(𝒞_Q,f_Q)`$, the boundary of this element would have a sign of $`1`$ on the component of the moduli where the curve $`𝒞_P`$ splits in to a sum of two elliptic curves meeting at the $`2`$-torsion point $`[(1,0),(1,0)]`$.
In general, to a component ( which corresponds to an even theta characteristic ) we can associate a $`(3,3)`$ configuration of the six ramification points, namely the first triple corresponding to the points lying on our choice of $`_1`$ and the second to the points lying on $`_2`$. To any pair of these ramification points we have a Collino element of $`CH^2(𝒲,1)`$ and the sign of the boundary on this component would depend on the which triple $`0`$ and $`PQ`$ belong to. If they belong to the same triple, then the sign is $`0`$ and if they belong to different triples, the sign is $`+1`$ or $`1`$ depending on whether $`0`$ or $`PQ`$ lies on $`_1`$.
The sign at a different component will depend on the action of $`\mathrm{\Gamma }_2(o)`$ on the $`(3,3)`$ configuration of the six ramification points. To understand this action, we observe that this can be translated in to the well known \[Ig\],\[G-H\] situation of the action of $`𝔄_6`$ on the $`(3,3)`$ configuration of the six odd theta characteristics. This is because the points associated to the six odd theta charateristics are nothing but the ramifications points of the genus 2 curve given by divisor of
$$\mathrm{\Theta }[(0,0),(0,0)](z,\tau )$$
As we used
$$\mathrm{\Theta }[(0,1),(0,1)](z,\tau )$$
to compute the action on the six ramification points we have, we simply translate our points by the point $`[(0,1),(0,1)]`$ which gives us the six odd theta characteristics, and then use the description of the action on them.
The table describes the correspondence between pairs of odd theta characteristics and even theta characteristics. We use the convention, as in \[Ig\] that the six odd theta chacteristics are
$$P_1=[(0,1),(0,1)]P_2=[(0,1),(1,1)]P_3=[(1,0),(1,0)]$$
$$P_4=[(1,0),(1,1)]P_5=[(1,1),(0,0)]P_6=[(1,1),(1,0)]$$
With that convention, the class $`(346,125)`$ corresponds to $`[(1,1),(1,1)]`$ The $`g`$ in the leftmost column represents the element of $`𝔄_6`$ which has to be applied to $`(346,125)`$ to obtain the pair on the right. In \[Ig\] he makes the isomophism of $`𝔖_6`$ with $`Sp_4(/2)`$ explicit, which along with some help from MAPLE, allows one to do the computation and make the following table.
Define the function
$$ϵ_{P,Q}(\tau )=\{\begin{array}{cc}1\hfill & \text{if }(O)\text{ lies on }(_1)_\tau \text{ and }(PQ)\text{ lies on }(_2)_\tau \hfill \\ 0\hfill & \text{if }(O)\text{ lies on }(_i)_\tau \text{ and }(PQ)\text{ lies on }(_i)_\tau \text{ for }i=1,2\hfill \\ 1\hfill & \text{if }(O)\text{ lies on }(_2)_\tau \text{ and }(PQ)\text{ lies on }(_1)_\tau \hfill \end{array}$$
From the above discussion, given a point $`\tau `$ on the component corresponding to $`[(1,1),(1,1)]`$ we can determine $`ϵ_{P,Q}(g\tau )`$ for $`g`$ in $`\mathrm{\Gamma }_2(o)`$.
### 4.4 Computing the relation
In this section, we restrict the computations done in the previous sections to Shimura curves.
From Hashimoto’s \[Ha\], Thm 3.5, explicit description of the Eichler order and Shimura curve one can explicitly describe the embedding $`\mathrm{\Psi }:_2`$ which is compatible with the embedding of the group of units $`\mathrm{\Gamma }`$ into $`\mathrm{\Gamma }_2(o)`$ :
###### Theorem 4.3.
The curve $`𝐗(𝐃_\mathrm{𝟎},𝐍)`$ satisfies singular relations of the form
$$\tau _1+\tau _2+\frac{p1}{4}\tau _3=0$$
$$2aD\tau _2+(\tau _2^2\tau _1\tau _3)+(a^2Db)D=0$$
giving rise to an embedding
$$\mathrm{\Psi }(\tau )=\frac{1}{p\tau }\left(\begin{array}{cc}\overline{\kappa }^2+\frac{(p1)aD}{2}\tau +D\kappa ^2\tau ^2& \overline{\kappa }(p1)aD\tau D\kappa \tau ^2\\ \overline{\kappa }(p1)aD\tau D\kappa \tau ^2& 12aD\tau +D\tau ^2\end{array}\right)$$
where $`\kappa =\frac{1+\sqrt{p}}{2}`$
From this one can see \[Ha\] that the $`\tau `$ in the image satisfy singular relations with invariants of the form $`pX^2+4aDXY+4bDY^2`$ for $`(X,Y)`$ coprime.
Let $`v`$ be the singular relation given by the existence of an elliptic curve of degree $`1`$, and let $`\mathrm{\Psi }(\tau )`$ be a point on $`\mathrm{\Psi }()`$ which satisfies this relation. Then $`\tau `$ satisfies an equation given by
$$\left(\begin{array}{cc}\mathrm{\Psi }(\tau )& I_2\end{array}\right)v\left(\begin{array}{c}\mathrm{\Psi }(\tau )\\ I_2\end{array}\right)=0$$
Since as functions of $`\tau `$, $`\tau _1`$,$`\tau _2`$, $`\tau _3`$ and $`\tau _2^2\tau _1\tau _3`$ are at worst quadratic when multiplied by $`p\tau `$ and since the singular relation between the $`\tau _i^{}s`$ is linear, the equation satisfied by $`\tau `$ is quadratic. Hence the point $`\tau `$ is imaginary quadratic so is a CM point on $``$.
We can look at the action of $`\mathrm{\Gamma }_2(o)`$ on $`v`$, namely
$$gv=g^{tr}vg\text{ for }g\text{ in }\mathrm{\Gamma }_2(o)\text{ }$$
The point $`g\mathrm{\Psi }(\tau )`$, which is given by the action of $`\mathrm{\Gamma }_2(o)`$ on $`_2`$ need not lie on $`\mathrm{\Psi }()`$. However, for some $`g`$ there exists a point $`\mathrm{\Psi }(\tau ^{})`$ which satisfies
$$\left(\begin{array}{cc}\mathrm{\Psi }(\tau ^{})& I_2\end{array}\right)g^{tr}vg\left(\begin{array}{c}\mathrm{\Psi }(\tau ^{})\\ I_2\end{array}\right)=0$$
If such a point exists it will be unique, and we will denote it by $`g(\tau )`$. This has the same level $`2`$ structure as $`g\tau `$, so $`ϵ_{P,Q}(g(\tau ))=ϵ_{P,Q}(g\tau )`$. Let $`\stackrel{~}{g(\tau )}`$ be the point $`\omega (g(\tau ))`$. Then $`ϵ_{P,Q}(\stackrel{~}{g(\tau )})=ϵ_{P,Q}(g(\tau ))`$
The relation we get is then a signed sum over the $`g(\tau )`$, and putting it all together gives us the following proposition.
###### Proposition 4.4.
Let $`(C_P,f_P)+(C_Q,f_Q)`$ be Collino’s element in the group $`CH^2(𝒜_\eta ,1)`$, where $`\eta `$ is the generic point of a Shimura curve $`𝐗`$ as before. Let $`\tau `$ be a point on the moduli where the curve $`C_P`$ splits as a sum of elliptic curves $`_1+_2`$, and make a choice of $`_1`$ as before. Assume further that $`O`$ lies on $`_1`$ and $`QP`$ lies on $`_2`$. Then one has a relation in $`CH^2(𝐖)`$ of the form
$$\underset{g\mathrm{\Gamma }_2(o)}{}ϵ_{P,Q}(g(\tau ))(2(_2_1)|_{g(\tau )}2(_2_1)|_{\stackrel{~}{g(\tau )}})0$$
up to the relations coming from decomposable elements and cycles homologous to zero in the special fibres, where $`ϵ_{P,Q}(g(\tau ))`$ can be computed from the table in section $`4.3`$.
### 4.5 Isogenies
In this section we will explain how we can use generic isogenies to get more relations between some cycles in $`CH^2(𝐖)`$.
Let $`g`$ be an element of $`Sp_4()`$. Let $`\tau `$ be a point. Then one has an induced isogeny
$$\mathrm{\Phi }_g:𝒜_{g\tau }𝒜_\tau $$
of some degree $`n^2`$. For a fixed degree, there are only finitely many classes of $`g`$ under the relation
$$g_1g_2\gamma Stab(\times )|g_1=\gamma g_2$$
We can use this information to compute the relations one would get by applying such isogenies to the original relation in Proposition $`4.4`$ . Equivalently, we could apply the isogenies to Collino’s element to get new elements and compute their boundaries.
In the previous section we have used the embedding $`\mathrm{\Psi }:_2`$. Let $`g`$ be in $`Sp_4()`$ of degree $`n^2`$ prime to $`2k`$. $`g\mathrm{\Psi }`$ gives another embedding of embedding of $`_2`$. Let $`\mathrm{\Gamma }_g`$ denote the group $`g\mathrm{\Gamma }g^1Sp_4()`$, where the action of $`g`$ is on the image of of $`\mathrm{\Gamma }`$ in $`Sp_4()`$. This acts on the new embedding of $``$.
Let $`𝐗_g=/\mathrm{\Gamma }_g`$ (or its compactification, if necessary). There is a map $`𝐗_g𝐗`$ given by $`\tau g^1\tau `$ and as above, a corresponding isogeny $`\mathrm{\Phi }_g`$ of the universal abelian surfaces over these curves. Let $`\zeta `$ denote Collino’s element in $`CH^2(𝒜_g,1)`$. We can compute relations resulting from this element in the universal family over $`𝐗_g`$, $`𝐖_g`$ and then push the relation down using the isogeny. This will give us some new relations in $`CH^2(𝐖)`$. A remark is necessary when the curve $`𝐗_g`$ is not compact as then it is not immediately clear that the isogeny extends to a morphism over the cusps. However, as the cycles supported in the cuspidal fibres are torsion as all these curves are quotients of the upper half plane by congruence subgroups \[Sl2\], one can always multiply the relation by a suitable number to kill those cycles, and then apply the isogeny. As the isogeny preserves the non-cuspidal points, there is no problem.
The relation that one obtains in $`CH^2(𝐖_g)`$ are supported in fibres where there is an elliptic curve of degree $`1`$. Let $`v`$ be the singular relation representing the elliptic curve of degree 1, $`_1`$. By translating by $`\mathrm{\Gamma }_2(o)`$, we can make a uniform choice of the elliptic curve $`_1`$ over the entire moduli of products of elliptic curves. Let $`\tau `$ be a point on $`𝐗_g`$ satisfying that singular relation. Then, as described above, the point $`g^1\tau `$ on $`𝐗`$ satisfies a relation corresponding to the existance of elliptic curve of degree $`n`$ in the fibre.
It is possible for $`g^1\tau `$ to lie on the image of some other point $`\tau ^{}`$ as
$$g^1\tau =\gamma g^1\tau ^{}\text{ for some }\gamma \text{ in }\mathrm{\Gamma }(2k)\text{ }g\gamma g^1\tau =\tau ^{}$$
However, we are looking at equivalence by the subgroup of finite index in $`g\mathrm{\Gamma }g^1`$ given by $`\mathrm{\Gamma }_g`$, and these points need not be $`\mathrm{\Gamma }_g`$ equivalent. But since $`g`$ gives rise to an isogeny of degree coprime with $`2k`$ the signs of $`_1_2`$ in all the points in the fibre over a point on $`𝐗`$ are the same.
Conversely, a theorem of Kani \[Ka\] asserts that if $`\tau _0`$ is a point on $`𝐗`$ where there is an elliptic curve of degree $`n`$ in its fibre, then $`\tau _0=g^1\tau `$ for some point $`\tau `$ on $`𝐗_g`$ where there is an elliptic curve of degree $`1`$ and $`g`$ a primitive isogeny of degree $`n^2`$.
Let $`g_1,\mathrm{}.,g_M`$ be representatives for the finitely many primitive isogeny classes of degree $`n^2`$. Then, since we have made a uniform choice of $`_1`$ over the entire moduli, we have a uniform choice on all the points $`\tau `$ on all the $`𝐗_{g_i}`$ which lie on the moduli of products. Since the isogenies are of degree prime to $`2k`$, they do not affect the level structure, and so one can define $`ϵ_{P,Q}(g^1\tau )=ϵ_{P,Q}(\tau )`$ for $`g`$ of degree $`n^2`$.
Applying the corresponding isogenies to the relations on $`𝐗_{g_i},i=1,\mathrm{}M`$ one has the following proposition, which is a generalization of the proposion in the previous section.
###### Proposition 4.5.
Let $`\tau _1`$ be a point on $`𝐗_{g_1}`$ , where $`g_1`$ is one of the representatives for the isogenies of degree $`n^2`$, and using $`\tau `$ make a uniform choice of elliptic curve $`_1`$ for the whole moduli space. Choose a point $`\tau _i`$ on each of the $`𝐗_{g_i}`$
Then, for each $`i`$, there are relations in $`CH^2(𝐖)`$ of the form
$$\underset{g\mathrm{\Gamma }_2(o)}{}ϵ_{P,Q}(g_i^1g(\tau _i))(2(_2_1)|_{g_i^1g(\tau _i)}2(_2_1)|_{\stackrel{~}{g_i^1g(\tau _i)}})0$$
where $`_1`$ and $`_2`$ are now elliptic curves of degree $`n`$ in the fibre over $`g_i^1g(\tau )`$.
###### Remark 4.6.
While it may happen that $`g_i\tau =g_j\tau ^{}`$ for some points $`\tau `$ and $`\tau ^{}`$, this can never happen on an embedding of a curve where there is no generic elliptic curve of degree $`n^2`$ as otherwise one would have too many elliptic curves in the fibre. This is because the classes of elliptic curves are linearly independent in the Néron-Severi. Both $`g_i`$ and $`g_j`$ would contribute 2 each. If generically there were no elliptic curves this would lead to a contradiction as there would be too many linearly independent elements.
Hence we get infinitely many relations between some cycles supported in points where there is an elliptic curve of degree prime to $`2k`$. Since generically, the Néron-Severi is of rank three, when this happens the Néron-Severi jumps to rank 4, and the abelian surface is necessarily isogenous to a product of isogenous CM elliptic curves, and there are CM cycles in those fibres.
In the next section we will use the relations above to get relations between the CM cycles.
## 5 Relations between CM cycles
### 5.1 Rewriting the relation
In this section we first describe the cycles $`_1_2`$ in terms of the CM cycles and show that we can modify Collino’s element by suitably chosen decomposable elements to get a relation only between them. Then we get a more explicit description of the points $`g(\tau )`$.
In all our computations we will work modulo cycles homologous to 0 in the fibre. From weight considerations, one can see that such cycles map to 0 in the second intermediate Jacobian, though they need not be 0 in the Chow group. Since the CM cycles themselves are defined modulo such cycles, this is not really a restriction. For simplicity, we will only work with the case of Collino’s original element.
If we fix the functions $`f_P`$ and $`f_Q`$, then the boundary of Collino’s element looks like
$$\underset{\tau }{}𝒟_\tau $$
where $`𝒟_\tau `$ is some element of the Néron-Severi group of the fibre over $`\tau `$. $`D_\tau `$ can be written in terms of the basis for the rational Néron-Severi given by $`e_1`$,$`e_2`$,$`e_3`$ and $`𝒵_\tau `$,
$$𝒟_\tau =b_\tau ^1e_1+b_\tau ^2e_2+b_\tau ^3e_3+c_\tau 𝒵_\tau $$
where $`𝒵_\tau `$ denotes the CM cycle if the point $`\tau `$ is a CM point and is $`0`$ otherwise
###### Lemma 5.1.
Let $`𝐖`$ be (the compactification of) the universal family of abelian surfaces as before and let $`𝒲`$ be the generic fibre. Recall that $`e_1`$,$`e_2`$ and $`e_3`$ generate the generic Néron-Severi. Then there are decomposable elements $`(e_i,f_{e_i})`$ for $`i=1,2,3`$ such that
$$(\underset{i}{}(e_i,f_{e_i}))=\underset{\tau }{}(\underset{i}{}b_\tau ^ie_i).$$
Hence there is a relation of the form
$$\underset{\tau }{}c_\tau 𝒵_\tau +\{hom\}0$$
in $`CH^2(𝐖)`$, where $`\{hom\}`$ denotes the image of cycles homologous to zero in the fibre.
###### Proof.
Let $`\pi _i:\overline{e_i}𝐗`$ be the maps from the closure of the $`e_i`$ to the base $`𝐗`$. Let $`M_i`$ be divisors in $`CH^1(𝐖)`$ be such that $`(M_i.e_j)=\delta _{ij}`$. Such $`M_i`$ exist as one can always find an orthonormal basis for the rational Néron-Severi. Intersecting $`M_i`$ with $`_\tau 𝒟_\tau `$ gives a relation in $`CH^2(\overline{e_i})`$ of the form
$$\underset{\tau }{}b_\tau ^i(\overline{e_i}M_i)|_\tau $$
The direct image
$$(\pi _i)_{}(\underset{\tau }{}(\overline{e_i}M_i)|_\tau )=b_\tau ^i\tau $$
is a rational equivalence of points on $`X`$ so is the divisor of a function, $`f_{e_i}`$.
These functions $`f_{e_i}`$ combined with the elements $`e_i`$ give the required decomposable elements. Subtracting these elements from Collino’s element gives a relation of the form
$$\underset{\tau }{}c_\tau 𝒵_\tau +\{hom\}0$$
###### Remark 5.2.
At the cuspidal points, the cycles that remain after subtracting off the boundary of these decomposable elements are orthogonal to the closure of the generic Néron-Severi and as mentioned before are torsion in the Chow group of the special fibre itself, so do not make a contribution after tensoring with the rationals.
One could also say things more intrinsically by using the fact that there is a projector in the ring of correspondences with rational coefficients on $`𝒜`$ which takes the varation of Hodge structure $`R^2\pi _{}()`$ onto $`Sym^2(R^1p_{}())`$ and the fact that the CM cycles lie in this part in the fibres over the CM points. We could then apply this correspondence to the element of $`CH^2(𝒜,1)`$ itself to get a relation involving only CM cycles.
As a result of this lemma one sees that since one gets a relation between CM cycles which all lie in $`CH_{hom}^2(𝐖)`$, we get a relation there and not merely in the Chow group.
We would now like to determine the points $`\tau `$ and the coefficients $`c_\tau `$. To determine $`c_\tau `$ we have to write the CM cycle in terms of the basis for the primitive Néron-Severi given by $`\overline{e}_i,i=2,3`$ and $`\overline{e}_4=\frac{1}{2}(2e_4<e_4.e_1>e_1)`$, where $`e_4=_1`$. Since $`<e_4,e_4>=0`$, one has $`(\overline{e}_4.\overline{e}_4)=1`$. Define $`r`$ and $`s`$ by
$$(\overline{e}_4.\overline{e}_2)=r(\overline{e}_4.\overline{e}_3)=s$$
so the intersection matrix looks like
$$\begin{array}{cccc}& \overline{e}_2& \overline{e}_3& \overline{e}_4\\ & & & \\ \overline{e}_2& p& 2aD& r\\ \overline{e}_3& 2aD& 4bD& s\\ \overline{e}_4& r& s& 1\end{array}$$
The idea now is to compute the CM cycle up to a multiple using the property that it is orthogonal to the generic Néron-Severi. Define $`d_\tau `$ to be the smallest multiple of the CM cycle $`𝒵_\tau `$ which lies in the $``$-span of $`\overline{e}_i,i=2,3,4`$. Then one has
###### Proposition 5.3.
Let $`d_\tau `$ be as above, then
$$d_\tau 𝒵_\tau =\frac{aDs2bDr}{2}\overline{e}_2+\frac{2aDrps}{2}\overline{e}_3+2D\overline{e}_4$$
up to a sign.
###### Proof.
Let
$$𝒵_\tau =x_2\overline{e}_2+x_3\overline{e}_3+x_4\overline{e}_4$$
Then $`(𝒵_\tau ,\overline{e}_i)=0`$ for $`i=2,3`$, so we get 2 equations
$$0=px_2+2aDx_3+rx_4$$
$$0=2aDx_2+4bDx_3+sx_4$$
Eliminating the variable $`x_2`$ gives
$$x_3=\frac{2aDrps}{4D}x_4$$
and similarly
$$x_2=\frac{2aDs4bDr}{4D}x_4$$
so the vector
$$(\frac{2aDs4bDr}{4D},\frac{2aDrps}{4D},1)$$
is orthogonal to the generic Néron-Severi hence a rational multiple of the CM cycle.
The smallest multiple of the vector such that it is an integral linear combination will be the $`d_\tau `$ times the CM cycle or half the CM cycle depending on the discriminant, at least up to a sign. Define $`b_{r,s}`$ as the multiple which gives it, hence it is either the smallest multiple or twice the smallest multiple.
From the description above, and the fact that $`s`$ is even, which we will see presently, it is clear that $`b_{r,s}|4D`$ and $`c_\tau =c_{r,s}=4D/b_{r,s}`$. The smallest multiple is determined by $`s_1=p(s/2)aDr`$, if $`s_1`$ is even, then it is $`D`$ otherwise it is $`2D`$.
The proposition will then follow from the following lemma
###### Lemma 5.4.
$`|b_{r,s}|=2D`$ always.
###### Proof.
Let $`v=(\alpha ,\beta ,\gamma ,\delta ,ϵ)`$ be the singular relation corresponding to the elliptic curve $`_1`$. From Hashimoto’s explicit description of the embedding, one can compute $`r`$ and $`s`$ to be
$$r=\beta 2\gamma \frac{1p}{2}\alpha $$
$$s=2aD\beta 2ϵ2D\delta (a^2Db)$$
Let
$$s_1=p(s/2)aDr=(\frac{(p1)}{2}aD\alpha (p1)aD\beta 2aD\gamma +((p1)a^2D)\delta +pϵ)$$
It turns out that $`\tau `$ satisfies an equation of discriminant $`(1/p)(s_1^24D(\frac{pr^2}{4}))`$. This is even if and only if $`s_1`$ is even. Hence $`b_{r,s}=2D`$ always, as if $`s_1`$ is even, then the smallest multiple is $`D`$ hence $`b_{r,s}=2D`$ is twice the smallest multiple. If $`s_1`$ is odd,then the smallest multiple is $`2D`$ and $`b_{r,s}`$ is the smallest multiple. This also says that the discriminant of $`\tau `$ is $`(1/p)(s_1^24D(\frac{pr^2}{4}))`$ up to a square factor.
This concludes the proof of the proposition.
Armed with this information we can compute discriminants of the orders that appear in the relation. They are the points on the moduli where the Néron-Severi is generated by the generic elements along with an elliptic curve of degree $`1`$.
###### Proposition 5.5.
The discriminants of the points $`\tau `$ where there is a elliptic curve of degree $`1`$ are
$$\mathrm{\Delta }/d^2=\frac{P(s,r)4D}{4d^2}=(p(s/2)^2aDrs+bDr^2D)/d^2$$
where $`r`$ and $`s`$, where $`r`$ and $`s`$ run through all possible numbers satisfying $`\frac{pr^2}{4}_0`$, $`\frac{4bDs^2}{4}_0`$ and $`d`$ is such that $`d^2|\frac{P(s,r)4D}{4}`$
###### Proof.
The idea is to compare the determinants of the intersection matrices of the two bases of the rational Néron-Severi. Using the basis coming from $`\overline{e}_1,\overline{e}_2`$ and the CM cycle $`d_\tau 𝒵_\tau `$, at a point of discriminant $`\mathrm{\Delta }_0`$ one gets the determinant of the intersection matrix to be
$$16D^2d_\tau ^2\mathrm{\Delta }_0$$
On the other hand, from the computation using the other basis given by $`\overline{e}_1,\overline{e}_2,\overline{e}_3`$ and the fact that the change of basis matrix has determinant $`b_{r,s}^2=4D^2`$, a simple calculation shows that the determinant is
$$(4DP(s,r))4D^2$$
where $`P(s,r)`$ is the quadratic form $`ps^24aDrs+4bDr^2`$.
Comparing the two shows that
$$\mathrm{\Delta }_0=\frac{P(s,r)4D}{4d_\tau ^2}=\mathrm{\Delta }/d_\tau ^2$$
where $`\mathrm{\Delta }=\frac{P(s,r)4D}{4}`$
If $`v_1`$ and $`v_2`$ are two elements of the primitive Néron-Severi, then the lattice generated by them is the Néron-Severi lattice of an abelian surface with multiplication by an Eichler order of discriminant $`\frac{v_1^2v_2^2(v_1.v_2)^2}{4}`$. This explains why $`\frac{pr^2}{4}`$ and $`\frac{4bDs^2}{4}`$ are in $``$.
Conversely, if $`r`$ and $`s`$ satisfy the conditions that $`\frac{pr^2}{4}`$ and $`\frac{4bDs^2}{4}`$ are in $``$ and $`d^2`$ divides $`\frac{P(s,r)4Dn^2}{4}`$ then one can see that the element
$$=\frac{1}{2D}(d𝒵_\tau \frac{aDs2bDr}{2}\overline{e}_2\frac{2aDrps}{2}\overline{e}_3)$$
where $`\tau `$ satisfies an equation of discriminant $`\frac{P(s,r)4D}{4d^2}`$, is a primitive integral element of invariant $`n^2`$, hence, from \[Ka\], one sees that it is an elliptic curve of degree $`1`$.
Since $`P(s,r)`$ is positive definite, the discriminant is negative for only finitely many values of $`s`$ and $`r`$.
To generalize Proposition $`5.5`$ for the relations we get from applying isogenies, we have to carry out the same computation using an elliptic curve of odd degree $`n`$ instead of $`1`$, and one gets that the discriminants that appear are
$$\mathrm{\Delta }=\frac{P(s,r)4Dn^2}{4d^2}=\frac{p(s/2)^2aDrs+bDr^2Dn^2}{d^2}$$
### 5.2 Heegner cycles
In this section we define the Heegner cycles which are sums of certain CM cycles with the same discriminant. As it turns out, our final result can be expressed in terms of these cycles. A reference for the facts used here is \[B-D\] or \[G-Z\].
Let $`𝒪`$ be, as before, the Eichler order in the quaternionic division algebra (or $`M_2(`$). Let $`K`$ be an imaginary quadratic field and let $`𝒪_\mathrm{\Delta }`$ be an order of discriminant $`\mathrm{\Delta }`$. If $`\tau `$ is a CM point, then $`\tau `$ determines an embedding of $`K𝐁`$ and $`\tau `$ is called a Heegner point of discriminant $`\mathrm{\Delta }`$ if
$$K𝒪=𝒪_\mathrm{\Delta }$$
Composing with the orientation on the Eichler order $`𝔬_{\mathrm{}}`$ gives an orientation on the Heegner point, namely a surjective linear map
$$\kappa _{\mathrm{}}:𝒪_\mathrm{\Delta }𝔽_{\mathrm{}}$$
if $`\mathrm{}|N`$ or
$$\kappa _{\mathrm{}}:𝒪_\mathrm{\Delta }𝔽_\mathrm{}^2$$
if $`\mathrm{}|D_0`$.
For a given orientation and level structure, it is well known \[G-Z\] that there are $`h_\mathrm{\Delta }`$ Heegner points of discriminant $`\mathrm{\Delta }`$, where $`h_\mathrm{\Delta }`$ is the class number of the order. Further, the Class group $`Pic(𝒪_\mathrm{\Delta })`$ acts on the points preserving the orientation. The Atkin-Lehner operator $`w_{\mathrm{}}`$, for $`\mathrm{}|D`$ also acts on the Heegner points by flipping the orientation at $`\mathrm{}`$ and also changing the ideal class. The action of the group $`\mathrm{\Omega }_DPic(𝒪_\mathrm{\Delta })`$ is transitive showing that for a fixed level there are $`2^th_\mathrm{\Delta }`$ Heegner points.
Let $`m`$ be a number such that
$$m\{\begin{array}{cc}asquaremod\mathrm{}if\mathrm{}|N\hfill & \\ anonsquaremod\mathrm{}if\mathrm{}|D_0\hfill & \end{array}$$
and let $`_{\mathrm{},m}`$ be the extension of $`_{\mathrm{}}`$ obtained by adding the roots of the equation $`x^2m=0`$. For $`\mathrm{}|N`$ this is simply $`_{\mathrm{}}`$ and for $`\mathrm{}|D_0`$ this is a quadratic extension. Let $`_{l,m}`$ be the ring of integers and $`𝒪_{\mathrm{\Delta },\mathrm{},m}=𝒪_\mathrm{\Delta }_{l,m}`$. Then
$$𝒪_{\mathrm{\Delta },\mathrm{},m}/\mathrm{}𝒪_{\mathrm{\Delta },\mathrm{},m}=\{\begin{array}{cc}𝔽_{\mathrm{}}𝔽_{\mathrm{}}\hfill & \text{if }\mathrm{}|N\hfill \\ 𝔽_\mathrm{}^2𝔽_\mathrm{}^2\hfill & \text{if }\mathrm{}|D_0\hfill \end{array}$$
and the different orientations correspond to the different canonical maps or equivalently the different primes lying over the primes dividing $`D`$.
From that one can see that the orientation classes are in bijection with solutions $`\mu `$ mod $`2D`$ of the equation
$$\mu ^2m\mathrm{\Delta }mod\mathrm{\hspace{0.17em}\hspace{0.17em}4}D$$
as $`\mu `$ will give a solution of
$$\mu ^2m\mathrm{\Delta }mod\mathrm{\hspace{0.17em}\hspace{0.17em}4}\mathrm{}$$
Since $`m^1`$ is a square in $`_{\mathrm{},m}`$, $`m^1=t_{\mathrm{}}^2`$ for some $`t_{\mathrm{}}_{\mathrm{},m}`$, so the ideal
$$𝔏=(\mathrm{},\frac{\mu t_{\mathrm{}}+\sqrt{\mathrm{\Delta }}}{2})$$
is a prime lying over $`\mathrm{}`$ in $`_{\mathrm{},m}`$. Conversely, given an oriented Heegner point, the kernels of the various orientation maps give rise to the different primes and hence a solution of the equation.
The choice of $`m`$ is not that important. Any number satisfying the conditions satisfied by $`m`$ will give the same result.
A Heegner point on $`𝐗`$ also comes with the data of a level $`2k`$ structure, and this has a projection onto the level 2 structure $``$, and we can further consider the class of $``$, $`[]`$ which is the even theta characteristic corresponding to it. We can make a uniform choice of an elliptic curves $`_1`$ and $`_2`$ of degree $`1`$, or more generally of degree $`n`$, for $`n`$ odd, over each of these classes. A Heegner Cycle is the CM cycle in the fibre over such a point. For a given level $`2`$ structure and orientation there are two Heegner cycles corresponding to which of the two pairs of three $`2`$ torsion points lies on $`_1`$. We define the sign of the Heegner cycle to be positive if the component of $`_2`$ in the direction of the CM cycle is positive. This depends on the configuration of the $`2`$-torsion points.
From class field theory, one knows that $`Pic(𝒪_\mathrm{\Delta })Gal(H_\mathrm{\Delta }/K)`$, where $`H_\mathrm{\Delta }`$ is the ring class field of $`𝒪_\mathrm{\Delta }`$. From this one can realise the action of $`Pic`$ as a Galois action, and one can see that the action on the Heegner cycles does not change the sign. The action of the Fricke involution $`w_D=_{\mathrm{}}w_{\mathrm{}}`$ is, up to an element of $`Pic`$, given by complex conjugation and it flips the orientation and the configuration of the pair of three $`2`$-torsion points and hence the sign.
We then can define a Heegner cycle of discriminant $`\mathrm{\Delta }`$, level $`[]`$ and orientation $`\mu `$ to be a sum
$$𝒵_{\mathrm{\Delta },[],\mu ,\mu }=\underset{[]}{}\underset{𝔞Pic(𝒪_\mathrm{\Delta })}{}𝒵_{𝔞,,\mu ,I}\underset{𝔞Pic(𝒪_\mathrm{\Delta })}{}𝒵_{𝔞,,\mu ,II}$$
where the $`I`$ and $`II`$ denote the possible configurations. This cycle is invariant under the action of $`Pic`$ and changes sign under the action of the Fricke involution so
$$𝒵_{\mathrm{\Delta },[],\mu ,\mu }=\underset{\sigma Gal(H/)}{}𝒵_{𝔞_0,,\mu _0,I_0}^\sigma $$
The cycle $`𝒵_{\mathrm{\Delta },[],\mu ,\mu }`$ is a different cycle which corresponds to the same configuration but opposite orientation. To reduce notation, and since both the configurations occur, we will suppress the $`I`$ and $`II`$.
For full level $`2k`$ structure we define the Heegner cycle to be the sum over all the points over the point determined by $`(𝔞,,\mu )`$ of the Heegner cycles, and we use the same notation to denote it.
In our situation, if $`r`$ and $`s`$ denote the intersection numbers of $`_1`$ with $`\overline{e}_2`$ and $`\overline{e}_3`$, then they determine the orientation of the Heegner cycle corresponding to it, which is of discriminant $`\mathrm{\Delta }=\frac{ps^24adrs+4bDr^24D}{4d^2}`$ for some $`d`$.
The prime $`p`$ is a number like $`m`$ above, namely a square mod $`\mathrm{}`$ for $`\mathrm{}|N`$ and a non-square mod $`\mathrm{}`$ for $`\mathrm{}|D_0`$. By ‘completing the square’ one can then see that $`s_1=p(s/2)aDr`$ satisfies an equation of the form
$$x^2=p\mathrm{\Delta }d^2mod\mathrm{\hspace{0.17em}\hspace{0.17em}4}D$$
and hence determines the orientation !
## 6 Main Result
In this section we state our main result. We will assume that we always have level $`2k`$ structure for some odd integer $`k`$ coprime to $`N`$.
###### Theorem 6.1.
Let $`𝐗(𝐃_\mathrm{𝟎},𝐍)`$ be a Shimura curve parametrising abelian surfaces with endomorphism ring an Eichler order of level $`N`$ in a division algebra of discriminant $`D_0`$. Let $`𝐖(𝐃_\mathrm{𝟎},𝐍)`$ denote (the non-singular compactification of) the universal abelian surface over $`𝐗(𝐃_\mathrm{𝟎},𝐍)`$. Let $`p,a,b`$ be the invariants which determine Hashimoto’s model and let $`P`$ and $`Q`$ denote two $`2`$-torsion points.
Then there are relations in $`CH_{hom}^2(𝐖(𝐃_\mathrm{𝟎},𝐍))`$ of the form
$$\underset{[]}{}ϵ_{P,Q}([])\underset{r,s}{}\underset{d^2|\mathrm{\Delta }}{}d𝒵_{\mathrm{\Delta }/d^2,[],s_1/d,s_1/d}0$$
where $``$ runs through all even theta characteristics and $`ϵ_{P,Q}([])`$ is the sign function as in Section $`4.3`$, $`r,s`$ such that $`\frac{n^2pr^2}{4}_0`$, $`\frac{4bDn^2s^2}{4}_0`$, $`n`$ is coprime with $`2k`$, $`r,s,n`$ mutually coprime, $`s_1=p(s/2)aDr)`$ and
$$\mathrm{\Delta }=\frac{ps^24aDrs+4bDr^24Dn^2}{4}$$
with $`d`$ running through all $`d^2|\mathrm{\Delta }`$ such that $`\mathrm{\Delta }/d^2`$ is still a discriminant.
###### Proof.
This is just a consequence of putting together the statements in the previous sections.
###### Remark 6.2.
For different choices of the prime $`p`$ and numbers $`a`$ and $`b`$ one could get possibly different relations. One way of interpreting the Hashimoto model is by observing that it gives an embedding of the Shimura curve in to the Siegel upper half space. Different choices of $`p`$,$`a`$ and $`b`$ correspond to possibly different embeddings and could result in different relations.
### 6.1 Examples
I. Suppose $`D_0=2.3`$, $`N=1`$ so $`D=D_0N=6`$ and $`n=1`$. The triple $`(p,a,b)=(5,2,5)`$ determines an embedding. The possible values for $`(r,s)`$ are
$$\{(1,0),(1,2),(1,4),(1,6),(1,8),(1,10)\}$$
Out of these only $`(1,4)`$ and $`(1,6)`$ give rise to negative values for $`\mathrm{\Delta }`$, which are $`4`$ and $`3`$ respectively.
Suppose the two torsion points are the ones associated to
$$P=[(0,0),(0,0)]\text{ and }Q=[(1,0),(0,0)]$$
so they correspond to the odd characteristics $`[(0,1),(0,1)]`$ and $`[(1,1),(1,0)]`$ respectively. Then from the table one can read off the signs at the various components of the moduli of products corrsponding to the even theta characteristics. In this case, for example, the signs are $`1`$ at $`[(1,1),(1,1)],[(0,1)(0,0)]`$ and $`[(0,0),(1,0)]`$, $`1`$ at $`[(1,0),(0,0)],[(0,1),(1,0)]`$ and $`[(0,0),(1,1)]`$ and $`0`$ elsewhere.
So one has a relation in $`CH^2(𝐖(\mathrm{𝟐},\mathrm{𝟑}))`$ of the form
$$𝒵_{3,[1000],3,3}+𝒵_{4,[1000],2,2}+𝒵_{3,[0110],3,3}+𝒵_{4,[0110],2,2}$$
$$+𝒵_{3,[0011],3,3}+𝒵_{4,[0011],2,2}𝒵_{3,[1111],3,3}𝒵_{4,[1111],2,2}$$
$$𝒵_{3,[0100],3,3}𝒵_{4,[0100],2,2}𝒵_{3,[0010],3,3}𝒵_{4,[0010],2,2}$$
$$+𝒵_{3,[1000],3,3}+𝒵_{4,[1000],2,2}+𝒵_{3,[0110],3,3}+𝒵_{4,[0110],2,2}$$
$$+𝒵_{3,[0011],3,3}+𝒵_{4,[0011],2,2}𝒵_{3,[1111],3,3}𝒵_{4,[1111],2,2}$$
$$𝒵_{3,[0100],3,3}𝒵_{4,[0100],2,2}𝒵_{3,[0010],3,3}𝒵_{4,[0010],2,2}=0$$
where for $`[abcd]`$ denotes the characteristic $`[(a,b),(c,d)]`$.
II. Different choices of $`(p,a,b)`$ leads to different embeddings of the modular curves and hence new relations. For example, if $`D=26`$, $`(p,a,b)=(5,2,21)`$, then there are relations between the $`\mathrm{\Delta }^{}s`$ given by $`(r,s,\mathrm{\Delta })=(1,18,11)`$, $`(1,20,20)`$ and $`(1,22,19)`$. However, if $`(p,a,b)`$ is $`(149,19,63)`$ then one has relations between $`(r,s,\mathrm{\Delta })`$ of the form $`(1,6,11)`$, $`(3,20,24)`$, $`(7,46,11)`$ and $`(9,60,8)`$. In particular, one may get a relation involving more than one Heegner point of the same discriminant.
> Ramesh Sreekantan
> Department of Mathematics
> Duke University
> Durham, NC 27708
> email:ramesh@math.duke.edu
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# Properties of Cosmic Shock Waves in Large Scale Structure FormationTo appear in The Astrophysical Journal
## 1 Introduction
“Cosmic shock waves”, formed in the course of large-scale structure formation, can contribute important roles in cosmology. They include external accretion shocks as well as merger and flow shocks internal to galaxy clusters. The pristine cosmic plasma accreting onto the large scale structure is deflected from the Hubble flow and first processed by accretion shocks (see e.g., Ryu & Kang 1997b). Evidence for their existence might be inferred from the observation of hot gas in the intracluster medium (ICM). In the commonly accepted paradigm for structure formation in the universe, gas accreting onto cosmic filaments and clusters of galaxies (GCs) has a typical bulk velocity up to $`\mathrm{a}\mathrm{few}10^3\mathrm{km}\mathrm{s}^1`$. This gas is then shock-heated to temperatures ranging from $`10^510^7`$ K in filaments and up to $`10^710^8`$ K in GCs (e.g., Kang et al. 1994a, KCOR94 hereafter; Cen & Ostriker 1994, CO94 hereafter; Cen & Ostriker 1999a). Merger shocks are produced during the mergers of sub-structures within a galaxy cluster and propagate through the hot ICM. In addition, during such a process, accretion shocks associated with the merging units also propagate through the ICM. Together with merging shocks they form a complex structure that can survive for long times inside the ICM after the end of the merger, because of the continuous gas inflow through filaments and sheets. We refer to these as flow shocks (see §3.1). There is now substantial observational evidence for temperature structure in clusters due to internal shock waves; these are appear to be mostly produced by merger events (e.g., Markevitch et al. 1998; Donnelly et al. 1999 and references therein) although recently some evidence might have appeared for flow shocks, too (Ensslin et al. (1998), §4).
Among other reasons for interest in cosmic shocks is their ability to efficiently accelerate particles to relativistic energies (cf. Blandford & Ostriker 1978, 1980; and Jones & Ellison 1991 for a recent review of this subject). In fact, relativistic cosmic-ray (CR) electrons have been observed in GCs through their synchrotron emission (e.g., Kim et al. 1989; Giovannini et al. 1993; Deiss et al. 1997). Extended sources of synchrotron radiation are commonly observed with spatial distribution similar to that of the thermal X-ray emission (see e.g., Liang 1999). Although the cooling time for such CR electrons is much shorter than the cluster ages, explicit signatures of particle aging are rare in the spectra of the observed sources. Since individual cluster galaxies are unlikely to replenish the ICM adequately with populations of relativistic particles, an efficient mechanism for extended particle acceleration is probably required to understand the CR electron replenishment.
Moreover, recently the EUVE satellite has revealed that many clusters possess an excess of extreme ultra-violet (EUV) radiation compared to what is expected from the hot, thermal X-ray emitting ICM (e.g., Lieu et al. 1996; Fabian 1996; Mittaz et al. 1997; Kaastra 1998). Further evidence for nonthermal activity in the ICM comes from detection of radiation in excess to thermal emission in the hard X-ray band above $`10`$ KeV (e.g., Henriksen 1998; Fusco-Femiano et al. 1999; Valinia et al. 1999; Sarazin 1999). The mechanism proposed for the origin of these components is the inverse-Compton (IC) scattering of cosmic microwave background photons by CR electrons, although it is not clear if the same electron population is responsible for producing both the EUV excess and the hard X-ray excess (Ensslin et al. 1999). Such detections suggest the possibility that nonthermal activities in the ICM are much higher than previously expected (Sarazin & Lieu 1998; Lieu et al. 1999).
Cosmic shock waves should be capable of accelerating CRs electrons responsible for the above emissions. However, at the same time, CR protons are produced. It is possible then, although not established yet, that if the above interpretation for the nonthermal radiation is correct, the CR protons produced at these shocks and accumulated throughout the cosmological evolution could provide a substantial fraction of the total pressure in GCs (Sarazin & Lieu 1998; Lieu et al. 1999). It is clear that if the CR pressure was ever comparable to the thermal pressure during the evolution of the universe, that would have a profound impact on cosmology. For instance, structure formation is heavily used as a probe for discriminating among cosmological models (e.g., Carlberg et al. 1997; Bahcall & Fan 1998) and hydrostatic equilibrium of the thermal ICM gas in the potential well of the total cluster mass is commonly assumed in order to derive GC masses (e.g., White et al. 1993; Evrard 1997). The presence of a nonthermal component obviously would alter the results in proportion to its relative importance. Furthermore some additional source of pressure is clearly required in GCs over that produced by adiabatic hydrodynamics both to produce the correct density profiles (e.g., Evrard 1990; Navarro, Frenk & White 1995) and to prevent catastrophic cooling flows (Suginohara & Ostriker 1999). Cosmic ray pressure in the inner parts of GCs may thus play a vital role in the hydrodynamic equilibrium of these systems.
Accretion shocks were also proposed as sites for acceleration of high energy CRs, protons and heavy nuclei, up to $`10^{18}10^{19}`$eV (Kang et al. 1996, 1997). In fact, given the large velocity of the accretion flows and the large size and long lifetimes of the associated shocks, such energies would be achievable through “cross-field” diffusion in perpendicular magnetohydrodynamic shocks. Magnetic fields in the ICM of GCs have been observed with strengths of the order of a few $`\times 0.1\mu `$G (e.g., Kim et al. 1989; Fusco-Femiano et al. 1999; Molendi et al. 1999). Outside GCs constraints on rotational measure from quasars impose an upper limit of $``$nG, based on the assumption of regularly alternating magnetic field (see e.g., Kronberg 1994). However, this limit can be shifted to higher values if a realistic distribution of magnetic field associated with cosmic structures is assumed. On such basis, Ryu et al. (1998) and Blasi et al. (1999) claimed a new upper limit $`\stackrel{<}{}`$1$`\mu `$G, at least along cosmic structures.
Cosmic shock waves could serve also as sites for the generation of weak seeds of magnetic field by the Biermann battery mechanism. It was proposed that these seeds could be amplified to strong magnetic field of up to $`\mu \mathrm{G}`$ in clusters if flows there can be described as Kolmogoroff turbulence (Kulsrud et al. 1997). However, further development into coherent magnetic field is unclear, since there is as yet no detailed theory capable of describing this process (see, e.g., Chandran 1997).
Additional roles which shock waves may play in cosmology have been explored by a number of authors (e.g., Ryu & Kang 1997a, RK97 hereafter; Quilis et al. 1998). RK97 compared analytical self-similar solutions for cluster formation in the Einstein-de Sitter universe (Bertschinger 1985) with one-dimensional numerical simulations in low density universes with/without a cosmological constant ($`0.1<\mathrm{\Omega }_M<1`$), where the properties of the accretion flow are related with the cluster’s mass, radius and temperature. The major conclusion was a possibly testable prediction about the difference in the accretion flow in different cosmological models. In particular, the accretion velocity onto clusters of a given mass or radius in low density universes is smaller by up to 45 % and 65 % respectively compared to that in the Einstein-de Sitter universe.
In the present paper we focus on the quantitative properties of large scale shocks produced by gas during the formation of cosmic structures. For this, the simulation data described in KCOR94 and CO94 have been used. The roles played by those shocks, especially with regard to CR acceleration and magnetic field generation and their consequences on cosmology, will be studied in future work. Details of our data analysis are described in §2.2. In particular, we computed the velocity, Mach number, radius of shocks and kinetic energy flow across them, which were not studied in previous works, and calculated their correlation with the cluster core temperature. The results are reported in §3. Finally, §4 concludes with a discussion.
## 2 Numerics
### 2.1 Simulations & Data
We study the properties of shock waves associated with the large scale structure of the universe in numerical simulations. In particular we investigate the velocity and radius of shocks in terms of cluster’s X-ray temperature. The latter, in fact, is the most reliably reproducible quantity in numerical simulations (e.g., Kang et al. 1994b; Frenk et al. 1999) and has been recently measured with satellite observations (e.g., Arnaud 1994; Mushotzky & Scharf 1997).
Simulation data in two cosmological scenarios have been used: a standard CDM (SCDM) model and a CDM + $`\mathrm{\Lambda }`$ ($`\mathrm{\Lambda }`$CDM) model, respectively. The SCDM simulation was discussed by KCOR94 and is characterized by the following parameters: spectral index for the initial power spectrum of perturbations $`n`$ = 1, normalized Hubble constant $`h`$ H$`{}_{0}{}^{}/(100`$ $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$$`)=0.5`$, total mass density $`\mathrm{\Omega }_M=1`$, baryonic fraction $`\mathrm{\Omega }_b=0.06`$, and normalization $`\sigma _8=1.05`$. On the other hand, the $`\mathrm{\Lambda }`$CDM simulation was discussed by CO94 and featured the following parameters: $`n=1,h=0.6,\mathrm{\Omega }_M=0.45,\mathrm{\Omega }_\mathrm{\Lambda }=0.55(\mathrm{\Omega }_0=\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }=1),\mathrm{\Omega }_b=0.043`$ and $`\sigma _8=0.77`$. In both calculations a cubic region of size 85 $`h^1`$Mpc at the current epoch was simulated inside a computational box with $`270^3`$ cells and $`135^3`$ dark matter (DM) particles, allowing a spatial resolution of 0.315$`h^1`$Mpc. We refer to KCOR94 and CO94 for further details. The simulations were performed with the code described in Ryu et al. (1993). In addition to DM, the evolution of gas was followed with a grid-based hydrodynamic code.
### 2.2 Analysis Method
A crucial and lengthy part in the current analysis of the results is the extraction and interpretation of detailed information from the numerical data. Each simulation form thousands of clusters and a very rich system of associated shock waves. We outline our methods here.
#### 2.2.1 Cluster Identification and DM Related Properties
The first step to take is the identification of GCs in the data. This could be done either through the X-ray emissivity criterion detailed in KCOR94 or through the distribution of DM particles. The former is preferred when the derived quantities must be directly compared with observational tests. However, when studying dynamical properties of GCs such as accretion velocity onto them, it would be better to use DM particles for identification, because it is mainly the gravitational contribution of this component that determines those properties. Therefore, for the identification of GCs we have adopted the DM-based “spherical over-density” method described in Lacey & Cole (1994) with a slight modification to speed up the calculation process. In particular, each DM particle is placed inside a cell (of physical size \[0.315$`h^1`$Mpc\]<sup>3</sup>) of the full 270<sup>3</sup> grid according to its position. If the number of particles inside such a cell exceeds a given threshold $`N_{threshold}`$, then, for each particle location, we define a local number density $`n_{\mathrm{}}=3(N+1)/(4\pi R_N^3)`$, where $`R_N`$ is the distance to the $`N`$th nearest neighbor. We then take the highest density particle locations as the candidate centers of clusters. The “candidate” particle locations are then sorted by density and a sphere is grown around each one of them, with the radius being increased until the mean density decreases to a value $`\delta \overline{\rho }`$ where $`\delta `$ is a selected parameter and $`\overline{\rho }=\mathrm{\Omega }_M\rho _{crit}`$ is the mean background density of the universe. The center of mass (CM) of the particles inside the sphere is calculated and taken as a new center and the overall process iterated until center corrections are smaller than $`ϵ/n_{\mathrm{}}^{1/3}`$, where $`ϵ`$ is a small parameter. Once the cluster radii, $`R_\delta `$, are thus defined, we reject from our list the smaller of any two clusters whose spatial separation is shorter than $`3/4`$ the sum of their radii. For our analysis we choose $`N_{threshold}=50`$; as in Lacey & Cole (1994) we set $`N=10`$ and $`ϵ=0.1`$. We chose $`\delta =80`$, since then $`\delta \rho `$ is about the average density inside the first caustic around clusters in SCDM (Bertschinger 1985; RK97). So the cluster radius is defined by $`R_{80}`$.
We point out for comparison that we have also identified GCs with the X-ray luminosity method outlined in KCOR94. In general we have found the relations between various physical quantities are not affected significantly by the method used. However, since the criteria are different, the samples considered are not identical, with differences depending on the thresholds used.
#### 2.2.2 Cluster Baryonic Matter Related Quantities
Once the cluster centers and radii, $`R_{80}`$, have been defined we turn to determine their X-ray based temperature, $`T_x`$, and luminosity $`L_x`$; in addition we compute the shock radius and velocity, $`R_s`$ and $`v_s`$ respectively, around clusters. The cluster temperature $`T_x`$ is computed by averaging over the cluster core region defined by $`r<R_{avg}=`$ 0.5 $`h^1`$Mpc. The X-ray luminosity $`L_x`$, however, is averaged over a volume of $`r<R_{avg}=`$ 1.0 $`h^1`$Mpc, since bremsstrahlung emissivity $`j_{ff}(\nu )`$ (KCOR94, Eq. 2) is substantial up to a distance of $`1`$ $`h^1`$Mpc. Operationally, the average is contributed from all the the cells falling inside a sphere of volume $`V`$ (around the cluster), with a weight function, $`w`$, given by the cell intersection with $`V`$. So, for example, $`w=1`$ for central cells and $`w1`$ for the marginal ones. This detail is important, because of the rapid drop of the thermodynamic variables with distance from the cluster center. Altering the “natural” weight of the central cells, unless motivated by a corrective purpose, introduces a source of error in the calculated quantities. While our temperature metrics should be fairly accurate, we expect, on the basis of resolution studies (Cen & Ostriker 1999b), that the X-ray luminosities are systematically underestimated while the rank order of the luminosities will be correct.
Once the cluster luminosity has been determined, we further select our sample based upon the criterion $`L_xL_{ff}=10^{41}\text{erg s}^1`$. We note that given the different numerical schemes for the evolution of the baryonic and DM related variables, the position of the X-ray luminosity peak and the DM-based CM of the same cluster may not correspond exactly. The shift is limited to one grid cell for the high luminosity clusters, but can amount a few cells for the faint ones. Our combined choice of $`N_{threshold}`$ and of a limiting luminosity assures that the luminosity peak and DM-based CM do not differ more than 1 grid cell.
Finally, in order to determine the shock radius, $`R_s`$, and velocity, $`v_s`$, we have labeled those cells where shocks are located. Shocks have been identified as compression regions ($`𝐯<0`$) meeting the following requirements for pressure ($`P`$) and velocity ($`v`$) jumps (index 1 and 2 refer to pre-shock and post-shock quantities):
$`{\displaystyle \frac{P_2P_1}{P_1}}`$ $`>`$ $`1.5`$ (1)
$`v_1v_2=2{\displaystyle \frac{c_1}{\gamma +1}}{\displaystyle \frac{M_1^21}{M_1}}`$ $`>`$ $`0.87c_1,`$ (2)
where $`c_1`$ is the pre-shock sound speed. These criteria select shocks with Mach number $`M_1>\sqrt{3}`$. However, once shocks have been identified with this method, we adopt the temperature jump (across the shock) instead of the pressure in order to calculate their properties. In fact, when a shock is only a few computational cells apart from a cluster center, compression in excess to that produced by the shock wave can result from the gravitational potential of the cluster itself. Given finite numerical resolution, it becomes then impossible to separate the two effects (due to shock and gravitational compression). However, since the cluster core is thought to be rather isothermal (Evrard et al. 1996), adiabatic gravitational compression affects density, pressure and velocity substantially more than the temperature. Therefore, for cosmological simulations, shock properties based on the temperature jump are the most reliable. Our tests clearly support this point.
The condition enforced by Eq. 2 does not add any physical property to the selected shocks, with respect to Eq. 1. However, it is of great advantage in our numerical effort to rule out compression regions simply due to the hydrostatic equilibrium inside the cluster gravitational potential, and otherwise detectable by Eq. 1 as shocked cells. This scheme has proven successful in tests with shocks of known hydrodynamic properties for several different applications. Once the shocks have been properly identified, the shock radius, $`R_s`$, is calculated by taking the average of the 6 distances to the first “shocked” cells encountered along the coordinate directions ($`x,y`$ and $`z`$) from the cluster center. In addition, the normal component of the pre-shock velocity is calculated using the hydrodynamic shock-jump conditions (e.g., Landau & Lifshitz 1997). We note that the normal component of the shock velocity $`v_{n1}`$ and not $`v_1`$ is really the quantity of interest here. Therefore in the following this is the quantity we refer to as the shock velocity $`v_s`$. In fact it determines the shock Mach number and flux of kinetic energy available for particle acceleration, the parameters most relevant for production of CR populations.
## 3 Results
### 3.1 Morphological Structure of Cosmic Shock Waves
Figs. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal and Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal illustrate typical structures found in the simulations. They represent slices of 0.315$`h^1`$Mpc thickness showing bremsstrahlung emission (gray-scale) with superposed contours of the shock compression factor ($`v`$). In both cases shocks confine filamentary structures where the temperature can vary from $`10^4`$ K to $`10^7`$ K. The most interesting feature, however, is the complexity of shocks around GCs. In both cosmological models, in fact, the accretion flow develops multiple shocks extending over a region of $`510`$ $`h^1`$Mpc size around the GCs. This is a generic feature of accretion flows there and it is directly related with the hierarchical process for cluster formation. In fact the merging of two sub-structures into a single unit produces at least two main effects related to shocks. First, shock waves of low Mach number are generated in the collision of the clusters’ ICMs, which are commonly referred as merger shocks. But also, part of the accretion shocks previously associated with each sub-structure end up propagating through the ICM of the newly formed structure, reaching deep inside the cluster core. These shocks, which we refer to as flow shocks, are subsequently “fed” by residual gas motions in the ICM and ongoing gas inflow accreting along filaments and sheets. Their presence, in addition to outer accretion shocks, provides additional heating of the ICM and makes its thermal structure not quite uniform over a region of several $`h^1`$Mpc of size. As already pointed out in §1, Markevitch et al. (1998 and also references therein) and Donnelly et al. (1999 and also references therein) showed evidences for significant temperature structures inside clusters. They attribute them to the presence of shocks associated with merger events. However, in some cases the merging is only inferred from the temperature map; also, in other cases temperature asymmetry is observed when the merger is just beginning to take place and therefore has not been able to affect the cluster temperature structure yet. In the future cluster shocks could be identified independently of the presence of merging processes. Some of this evidence may actually already be available (Ensslin et al. 1998, see §4 for further discussion).
A three dimensional perspective of shocks around GCs in SCDM is offered in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal. This is a combination of shock-strength isosurfaces (a) and volume rendering of bremsstrahlung emission (b), for a portion of the computational domain of size $`30\times 40\times 30`$ ($`h^1`$Mpc)<sup>3</sup> at $`z=0`$. With the help of the bremsstrahlung emission, which identify the GCs, one is allowed to locate shocks in their appropriate cosmological context. This image reveals a further degree of complexity of real cosmic shocks with respect to the two-dimensional slices above. Namely, in addition to being multiple associations with individual clusters, such shocks are also largely connected topologically with neighboring structures. Their shapes, far from spherical, extend over tens of Mpc forming a continuum that envelops all nearby clusters.
In conclusion, the hierarchical process for structure formation produces an extremely complex shock structure around clusters and groups of clusters (or superclusters). These shock waves are neither spherical nor identifiable by a simple surface. Indeed they intersect each other, forming nested shock surfaces, penetrating deep inside the ICM of individual clusters.
### 3.2 Physical Relations for Accretion Shock Waves
In this section we aim for a more quantitative description of accretion flows and shocks. Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal includes four histograms showing the distribution, at redshift zero, of GCs (N\[GC\]) with respect to the Mach number ($`M`$) of the associated “external” accretion shocks (top) and “internal” shocks (bottom). The latter have been selected inside regions of radius 0.5 $`h^1`$Mpc from the cluster centers. The results are plotted on the left (right) for the SCDM ($`\mathrm{\Lambda }`$CDM) model. Mach numbers associated with external accretion shocks are remarkably large, ranging from $`10`$ up to a few $`\times 10^3`$ in both cosmologies. We note that this is partly because photo-heating of pre-shock gas by metagalactic radiation field was not included. In addition, feedback processes from massive stars in galaxies may raise the temperature in the outskirts of clusters further. So without these sources of heating the pre-shock gas stays cool with sound speed less than $`1\mathrm{km}\mathrm{s}^1`$. If photo-heating was properly included, however, the pre-shock gas temperature would have been raised to circa $`10^4`$ K, with corresponding sound speed $`30\mathrm{km}\mathrm{s}^1`$ (see e.g., Ostriker & Cen 1996; Cen & Ostriker 1999a). Hence, the Mach numbers should be significantly smaller than those in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal. Yet they would still reach up to $`100`$ for external accretion shocks. On the other hand, the Mach numbers associated with internal shocks are much smaller, mostly in the range $`310`$ and peaking about 5. It thus emerges that, because internal shocks propagate through a significantly hotter medium with typically $`T10^610^7`$ K, the velocities associated with both external and internal shocks are comparable.
Among the characteristics of accretion flows onto GCs the most relevant quantities are the shock velocity, $`v_s`$, and radius, $`R_s`$. These are plotted as a function of each cluster’s X-ray temperature in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal, again on the left (right) for SCDM ($`\mathrm{\Lambda }`$CDM). According to the self-similar solution of one-dimensional spherical accretion in the SCDM universe, the cluster is confined by an accretion shock at $`R_s`$ and the temperature of the cluster gas increases inward (Bertschinger 1985; RK97). By choosing the gas temperature at $`r=0.3R_s`$ as a representative value for the cluster’s temperature $`T_x`$, the following relations of $`v_s`$ versus $`T_x`$ and $`R_s`$ versus $`T_x`$ are expected:
$`v_s=1.75\times 10^3\mathrm{km}\mathrm{s}^1\left({\displaystyle \frac{T_x}{7.8\times 10^7\text{K}}}\right)^{\frac{1}{2}}`$ (3)
$`R_s=2.12h^1\text{Mpc}\left({\displaystyle \frac{T_x}{7.8\times 10^7\text{K}}}\right)^{\frac{1}{2}}.`$ (4)
By fitting our data for $`v_s`$ and $`R_s`$ to a function of the form
$$f(T_x)=K\left(\frac{T_x}{7.8\times 10^7\text{K}}\right)^\alpha ,$$
(5)
as suggested by Eqs. 3 and 4, we have obtained the values for the coefficients $`K_{v_s},\alpha _{v_s}`$ and $`K_{R_s},\alpha _{R_s}`$ respectively, that are reported in Table 1. Note that $`\alpha _{v_s}=\alpha _{R_s}=0.5`$ is expected for both SCDM and $`\mathrm{\Lambda }`$CDM from scaling relations. Clearly the expected trends are reproduced only for the shock velocity, whereas strong deviations from Eq. 4 appear in the plots for $`R_s`$ in both cosmologies (Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal bottom panels). The most dramatic are the differences for $`\alpha _{R_s}`$. The small value reported in Table 1 with respect to that in Eq. 4 suggests $`R_s`$ is almost independent of the type of cluster. Although numerical errors both in the simulations and in the data analysis must be considered, such discrepancies are probably true and related to the actual structure of the flows. As already described in the previous section, the formation process of GCs imprints complex, irreducible three-dimensional shock structures which make it difficult to describe each GC with a single shock radius.
On the other hand, as long as the thermalization of the accretion kinetic energy takes place, the $`v_s`$ versus $`T_x`$ relation is less affected by the flow structure, with $`\alpha _{v_s}`$ very close to 0.5 for both cosmologies. $`K_{v_s}`$, is consistent with the predictions of self-similar solutions for SCDM as well as for $`\mathrm{\Lambda }`$CDM due to the same reason. We point out that $`v_s`$ is the normal component of the accreting gas in the shock rest frame. This is the component that undergoes dissipation originating the postshock gas temperature. However, the three-dimensional accretion flow also possesses transverse velocity components that are not thermalized across the shock and that can generate turbulent motions inside the IGM of GCs. In this regard, Eulerian, uniform grid-based schemes may be among the best choices to capture this component of the flow in term of a balance between resolution and computational performance. In fact, on the one hand in Smoothed Particle Hydrodynamic methods turbulence can be suppressed by excessive viscosity. On the other hand, the higher computational cost paid by the advantage of having a higher resolution with Adaptive Mesh Refinement techniques is not completely satisfactory because the small scales of the turbulent component are not fully generated.
We have also calculated the flux of kinetic energy across shocks defined as
$$\mathrm{\Phi }_{E_k}=\frac{1}{2}\rho v_s^2R_s^2v_s.$$
(6)
This quantity is shown in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal as a function of $`T_x`$ for the SCDM (left) and $`\mathrm{\Lambda }`$CDM (right) models. It is of interest in relation to CR acceleration, because it represents the amount of power available for conversion into supra-thermal particles. It provides a large amount of power of the same order of the X-ray luminosity emitted from a central region of 0.5 $`h^1`$Mpc (see KCOR94 and CO94 for the amount of the cluster’s X-ray luminosity). It is clear that if a modest fraction of this inflowing kinetic energy can be converted into CRs, these may become important sources of emissions (by synchrotron and inverse-Compton scattering) and even play a dynamical role through CR pressure. $`\mathrm{\Phi }_{E_k}`$ was fitted by a power law of the form
$$\mathrm{\Phi }_{E_k}=K_\mathrm{\Phi }\left(\frac{T_x}{7\times 10^7\text{K}}\right)^{\alpha _\mathrm{\Phi }}.$$
(7)
The best fit parameters are reported in Table 1. As we can see, the normalization factor, $`K_\mathrm{\Phi }`$, is larger for the $`\mathrm{\Lambda }`$CDM case than for the SCDM one. Moreover, both slopes are larger than the values implied from the slopes of $`v_s`$ and $`R_s`$ in combination with Eq. 6. This is probably due to an additional dependence of the accreting gas density on the cluster temperature. In fact, $`\mathrm{\Phi }_{E_k}`$ is calculated at the first shock cells encountered along the coordinate axis from the GC center (and then averaged over the accretion surface). Since such cells could well be inside “external” shocks, as illustrated in Figs. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal and Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal, the density of the accreting gas depends on the properties of the cluster environment. In particular we know that: (1) for a given cosmological model such density is higher around larger (higher temperature) clusters, implying a steeper increase of $`\mathrm{\Phi }_{E_k}`$ than for a temperature independent density value. Also, (2) for the same cluster temperature, the corresponding cluster gas density and, therefore, $`K_\mathrm{\Phi }`$ are larger for a $`\mathrm{\Lambda }`$CDM model than in a SCDM one. This second point is related to the fact that in general clusters in $`\mathrm{\Lambda }`$CDM have temperature smaller that those in SCDM (e.g., Figs. 5 and 6; see also KCOR94, CO94).
It turns out that $`\mathrm{\Phi }_{E_k}`$ is a steep function of cluster temperature, spanning several orders of magnitude in the temperature range of the identified clusters. This means that if CR acceleration mechanism at shocks around GCs possesses an injection mechanism and an efficiency independent of the cluster properties (e.g., mass and temperature), then we would expect hotter clusters to store a relatively larger amount of nonthermal energy in the form of relativistic particles. Such trend has been observed already: for instance, Liang (1999) reported a positive correlation between the radio emission and the X-ray temperature in GCs.
### 3.3 Evolutionary Trends
In general, complex shock structures are already present at high redshift. At $`z=510`$ shock waves are well formed and have already developed connections with neighbor clusters or protoclusters. The strength of the shocks is largest around the most massive objects, yet far from uniform. As the evolution advances, mergers occur on all scales, affecting all types of structures including shocks. As a result, at $`z=0`$, many filamentary structures have coalesced into larger ones. In addition the shocks associated with them have become stronger and more uniform, due to the increased amount of matter onto which the gas is being accreted.
In order to study the characteristics of shock evolution we define the following quantity:
$$S(z,M)=\frac{1}{N_{tot}(1+z)dx}_{\mathrm{}}^M\frac{d}{dM^{}}(N_{shock}[z,>M])dM^{}$$
(8)
Here $`N_{shock}(z,M)`$ is the number of cells hosting a shock of Mach number greater than $`M`$ at redshift $`z`$, $`N_{tot}`$ the total number of cells in the computational box and $`(1+z)dx`$ is the comoving cell size. Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal shows the evolution of $`S(z,1)`$, representing the inverse of the average comoving distance between shocks of any Mach number, as a function of redshift ($`z`$) for SCDM (open circles) and $`\mathrm{\Lambda }`$CDM (solid circles). This quantity also represents the ratio of shock surface over space volume. The following characteristics stand out from this plot. First, it is clear that cosmic plasma was populated with many more shocks in the past than nowadays. Second, while $`S(z,1)`$ peaks at $`z4.6`$ in the $`\mathrm{\Lambda }`$CDM model, we can only say that it peaks at $`z5`$ in the SCDM model, given our limited data-set. Also, in recent epochs $`S(z,1)`$ is substantially larger in the $`\mathrm{\Lambda }`$CDM model than in the SCDM model. Finally, the evolution is smooth in the $`\mathrm{\Lambda }`$CDM case, but shows abrupt transitions in the SCDM one.
As for the first point made above, the larger area of shock surfaces is due to the extremely filamentary structure of the universe at higher redshifts. Filaments are confined by accretion shocks. As already pointed out in §3.1, as structure formation progresses, filaments coalesce, therefore growing thicker and rarer. Although the size of their associated shocks increases, their reduced population plays the dominant role determining overall a decrease of the total area of shock surfaces in a comoving volume.
The second point then tells us that although such a process takes place in both cosmologies, there are many more filaments today in the $`\mathrm{\Lambda }`$CDM than in the SCDM scenario. This is clearly illustrated in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal, where we show two-dimensional slices, with thickness 0.315$`h^1`$Mpc, of the temperature structure of the universe at two different redshifts, $`z=3`$ (a) and $`z=0`$ (b), for SCDM (left) and $`\mathrm{\Lambda }`$CDM (right) models. There we can see that at $`z=3`$ the two universes look quite similar, with comparable amounts of filaments, in accord with Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal. However, at later times, e.g., $`z=0`$, filaments in the SCDM model are fatter and rarer as compared to the $`\mathrm{\Lambda }`$CDM case. This finding is a reflection of the different initial conditions in the two models. In particular, the larger amplitude of the primordial perturbations used in the SCDM, $`\sigma _8=1.05`$, gives rise, at current epoch, to a more clustered but less filamentary structure in this model than in $`\mathrm{\Lambda }`$CDM (cf. KCOR94; CO94).
Turning to the final point, sudden reductions of $`S(z,1)`$ are located at $`z3,0.7`$ and 0.2 in SCDM. These reflect the occurrences of a higher rate of merging processes at those particular epochs.
Further details related to the time evolution of shocks are illustrated in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal, where we plot as a function of the Mach number the quantity
$$W(z,M)=\frac{d}{dz}S(z,M).$$
(9)
$`W(z,M)`$ expresses the negative of the rate of formation of shocks with Mach number greater than $`M`$ at a particular epoch $`z`$. Thus, from Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal we can see that at early times ($`z=4.5`$) shocks are forming with Mach number in the range between circa 10 and a few $`\times 10^2`$ for both $`\mathrm{\Lambda }`$CDM and SCDM cases. No shocks exist with $`M>`$ a few $`\times 10^2`$ and for $`M<10`$ shocks are being depleted due to merger events. Later on, at $`z=1.25`$, shocks start forming in the range $`10^2M10^3`$. In accord with Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal the total rate of shock formation ($`W(z,1)`$) is always negative. As already pointed out, the numerous weak filament shocks are replaced by stronger but rarer shocks. The last two panels ($`z=0.6`$ and $`z=0.05`$ respectively) show the smooth shock formation evolution in the $`\mathrm{\Lambda }`$CDM . A small amount of shock formation still occurs for high Mach numbers but overall the shock population is decreasing at an increasing rate. The situation is more complex in the SCDM case. At $`z=0.6`$, identified above as an epoch of high merging rate, shocks are reduced at any Mach number. At $`z=0.05`$ the largest Mach number shocks ($`M5\times 10^2`$) are depleted indicating merging of the most massive objects.
In order to assess the relative importance of shocks at different epochs in terms of CR contribution, in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal we plot the adimensional quantity
$$F(z)=\frac{(\mathrm{\Phi }_{E_k}[z])_{shock}}{(\rho _cH_0^3\lambda _{NL}^5)_{z=0}}$$
(10)
where $`\mathrm{\Phi }_{E_k}`$ is the total flux of kinetic energy through shocks (of any Mach number), $`\rho _c`$ is the critical density, $`H_0^3`$ the Hubble constant and $`\lambda _{NL}`$ the non-linear perturbation wavelength set to 50 $`h^1`$Mpc. Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal shows that the flux of kinetic energy through shocks today has increased by a few order of magnitudes with respect to early epochs, say $`z5`$. Thus, today’s shocks retain more kinetic energy than ever. However, the time integrated flux of kinetic energy through shocks (e.g., since $`z=5`$) is much larger than the thermal energy content at $`z=0`$, which, according to our data, has mostly been produced after $`z1.5`$. This is not due to low thermalization efficiency of shocks. Rather, that is because, although shocks form at much higher redshifts, the thermal energy they produce undergoes severe adiabatic losses due to cosmological expansion. On the other hand, after roughly $`z1.5`$, such thermal energy is retained inside well formed structures such as clusters and filaments.
In order to identify the characteristics of the most relevant shocks, i.e., those that process most of the gas, we calculate
$$Y(M)=\frac{1}{E_{th}(z=0)}_{t(z=1.5)}^{t(z=0)}\frac{d\mathrm{\Phi }_{E_k}(M)}{dLogM}𝑑t^{},$$
(11)
where the extremes of integration have been chosen on the basis of the arguments in the previous paragraph. Here, $`\mathrm{\Phi }_{E_k}(M)`$ is the kinetic energy flux through shocks with Mach number between Log$`(M)`$ and Log$`(M+`$d$`M)`$, and $`E_{th}(z=0)`$ is the total thermal energy inside the computational box at $`z=0`$. This quantity is plotted in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal, as a function of Mach number $`M`$, in the right (left) panel for the SCDM ($`\mathrm{\Lambda }`$CDM ) model.
We can see that most of the flux of kinetic energy occurs through shocks with (“low”) Mach numbers around 4 which correspond to our internal shocks. In fact, although most of the shocks have Mach numbers much larger than that (see histograms of Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal), low Mach number shocks are typically located inside much denser regions (formed structure) and therefore process much more matter and kinetic energy than on average (as already pointed out in §3.2 the gas velocities for internal and external shocks are comparable). This depiction is not just the result of the integration. A more detailed analysis of plots of $`dY(M)/dt`$, describing the flux of kinetic energy as a function of $`M`$, for different redshifts (not shown here) shows that the the qualitative features of the curve in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal are common at any $`z`$ and therefore the low Mach number shocks are always responsible for most of the processing of the kinetic energy of the peculiar motions.
Finally, integration of the area underneath each curve in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal represents the total kinetic energy passed through shocks since $`z=1.5`$ divided by the thermal energy at $`z=0`$. Its value is $`17`$ for the SCDM and $`13`$ for the $`\mathrm{\Lambda }`$CDM model respectively. If a fraction $`ϵ10^2`$ of such energy is transferred to CR protons, then the energy stored up in CRs today should amounts to about 15 % of the thermal energy inside formed structure. This is only a rough estimate which needs to be refined by more accurate calculations.
## 4 Discussion
We have studied the properties of “cosmic shock waves” associated with the large scale structure of the universe in two different cosmological scenarios, namely SCDM and $`\mathrm{\Lambda }`$CDM. Such shocks reveal remarkable properties. In fact, hierarchical formation histories of GCs produce highly complex flows and shock structures, which extend over scales of several Mpc. In addition to accretion shocks (responsible for heating infalling gas) and merger shocks, flow shocks also appear and propagate through the thermalized ICM, providing extra gas heating. It turns out that the morphology of shocks associated with a large scale structure is complex and irreducibly three-dimensional and spherical shapes are inadequate to their description. Only for the external accretion shocks, located far away from the cluster core, some form regularity is recovered. This is an important issue especially in perspective of those missions with the next generation of high resolution X-ray telescopes (Chandra and XMM) which are planning to detect shocks in the ICM. It is worth mentioning that Ensslin et al. work (Ensslin et al. 1998) might already provide observational evidence for the presence of flow shocks in cosmic structure. Their conclusions are based on the assumption that the observed radio emission is due to particles currently accelerated at shocks there and injected from a “radio relic”, a remnant previously associated with some radio galaxies. For example, for 1253+275, they find a pre-shock gas temperature $`T0.51`$ KeV which shows that this is not the case of an accretion shock but that of either a merger or a flow shock (propagating through the ICM). Since there is no evidence for a merging process in 1253+275, it must be a flow shock. In addition, from their reported data on the pressure jump we infer $`M34.2`$, well in the expected range for the internal shocks shown in Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal (bottom panel).
Cosmic shocks are also ideal sites for particle acceleration. We have shown already in §3.2 that cosmic shocks provide enough power to produce copious CRs. The details of the produced populations will depend on the injection mechanism and scattering agent, i.e., the magnetic field and the diffusion properties. Here, we further stress the important role of merger and flow shocks. These shocks may be responsible, not only for acceleration of CRs out of the thermal pool of the ICM, but also for the re-acceleration of CRs produced at accretion shocks and/or ejecta from radio-galaxies, AGNs or normal galaxies. In addition, they could be crucial in terms of acceleration and transport of ultra high energy CRs, because the scattering mean free path of these particles is of the same order as the typical separation between accretion/merger/flow shocks in the ICM. As already pointed out at the beginning of this section, there seems to be solid foundation for the existence of such shocks.
As described in §1, the presence of relativistic CR electrons in GCs has been inferred through observations of diffuse synchrotron radiation from radio halos. In addition, there is evidence for excess of radiation in both EUV and hard X-ray bands, with respect to thermal emission. Although still of controversial interpretation, such excesses are probably due to IC emission of CR electrons scattering off cosmic microwave background photons. Published studies, however, reveal that for an accurate interpretation of the constraints from the combined non-thermal emission components, it is crucial to have a detailed depiction of the relative distribution of particles and magnetic fields (Ensslin et al. (1999)).
The proton component of CRs has not been directly observed. Nonetheless, given the estimates for the CR electron component, Lieu et al. (1999) concluded that their contribution in terms of dynamical pressure in GCs could be comparable to the thermal gas. This is consistent with the estimate inferred in the previous section and, as already pointed out in the introduction, has important consequences for cosmology. We point out that cosmic shock waves have existed ever since nonlinear structure formation was initiated at high redshift. This was shown through Fig. Properties of Cosmic Shock Waves in Large Scale Structure Formation<sup>5</sup><sup>5</sup>affiliation: To appear in The Astrophysical Journal in §3.3. Therefore, the importance of nonthermal activities in the cosmic plasma traces back to early epochs. Collisions of CR proton in the ICM generate, however, a flux of gamma ray photons through the production and subsequent decay of neutral pions; as pointed out by Blasi (1999), such gamma ray flux seems to be only marginally compatible with the upper limits measured by EGRET for Coma and Virgo clusters. But, again, the spatial and spectral distribution of CR, both depending on the overall cosmological history of these particles, play a crucial role in the determination of the expected gamma ray flux. In any case, the advent of the new generation of $`\gamma `$-ray facilities (GLASS, VERITAS) characterized by a much higher sensitivity (cf. Blasi 1999 for more details) will definitely settle the issue.
From this depiction it emerges the importance and the necessity to understand the role of CRs in cosmology. For this purpose we are developing numerical tools in order to treat consistently magnetic fields and CRs in numerical simulations of structure formation. Such tools, in fact, will allow us to follow explicitly the evolution of the magnetic field as well as the acceleration and transport of CRs. With this information we will be able to carry out very useful comparisons between numerical and observational results in various bands of the electromagnetic spectrum.
FM was supported in part by a Doctoral Dissertation Fellowship at the Uiversity of Minnesota. FM and TWJ were supported in part by NSF grants AST9616964 and INT9511654, NASA grant NAGS-5055 and by the Minnesota Supercomputing Institute. DR and HK were supported in part by grant 1999-2-113-001-5 from the interdisciplinary Research Program of the KOSEF. RC and JPO were supported in part by NSF grants AST-9803137 and ASC-9740300.
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# COMPTON SCATTERING IN ULTRA-STRONG MAGNETIC FIELDS: NUMERICAL AND ANALYTICAL BEHAVIOR IN THE RELATIVISTIC REGIME
## 1 INTRODUCTION
Recent observations are providing evidence for the existence of isolated neutron stars having ultra-strong magnetic fields. Assuming that the spin-down of isolated neutron stars is a result of electromagnetic dipole radiation, the measured period and the period derivative give the strength of the surface magnetic field as $`B_o=6.4\times 10^{19}(P\dot{P})^{1/2}`$ (Shapiro & Teukolsky 1983 and Usov & Melrose 1995). Typical radio pulsars have period and period derivative distributions (from the Princeton Pulsar Catalogue: Taylor, Manchester, & Lyne 1993) suggesting magnetic field strengths between $`10^{11}`$ and $`10^{13}G`$, with about two dozen pulsars having spin-down fields greater than $`10^{13}G`$; the highest to date is $`B_o=1.1\times 10^{14}G`$, a product of the Parkes multi-beam survey (Camilo et al. 1999). Although there have not been any radio pulsars detected with a magnetic field much exceeding $`B_o=10^{14}G`$ (perhaps with the exception of the unconfirmed observation of SGR 1900+14; Shitov 1999; Shitov, Pugachev & Kutuzov 2000), growing evidence for a new class of isolated neutron stars with ultra-strong magnetic fields ($`B_o>10^{14}G`$) has come from the observations of Soft Gamma-Ray Repeaters (SGRs) and Anomalous X-ray pulsars (AXPs). The five known SGRs are transient sources that undergo repeated outbursts of gamma rays and all have been associated with young ($`t<10^5yr.`$) supernova remnants. Last year, Kouveliotou et al. (1998a) detected a 7.47 s period in quiescent emission of SGR1806-20 in the X-ray band. Recently, Hurley et al. (1999) detected a periodicity of 5.16 s for SGR1900+14 during a giant burst having an energy of $`10^{45}erg`$. Kouveliotou et al. (1999) also observed the same period in the quiescent X-ray emission. Assuming dipole radiation torques, the measured period derivatives imply surface magnetic fields between $`10^{14}10^{15}G`$, well above the quantum critical field, $`B_{\mathrm{cr}}=4.4\times 10^{13}G`$.
Evidence for ultra-strong magnetic fields has also come from observations of AXPs, a group of six or seven X-ray pulsars with supersecond periods that exhibit anomalous characteristics in comparison to the properties of accreting X-ray pulsars. The lack of optical counterparts and orbital Doppler shifts (Steinle et al. 1987, Mereghetti et al. 1992 and Mereghetti & Stella 1995) suggest that these objects are isolated pulsars. Their more-or-less steady spin-down and young characteristic ages, $`t<10^5yr.`$ (Vasisht & Gotthelf 1997), support this assertion. Several AXPs have been associated with young supernova remnants, also suggesting neutron star origin. The AXPs are bright X-ray sources with luminosities, $`L_X10^{35}erg/s`$, far exceeding their spin-down luminosity. This energetics issue has motivated Thompson & Duncan (1996) and Kulkarni & Thompson (1998) to suggest that, unlike rotation-powered pulsars, the X-ray and particle emission in AXPs is powered by a decay of the magnetic field in the stellar interior.
Various studies indicate that inverse Compton scattering (ICS) plays a significant role in the magnetospheric physics of strongly magnetized neutron stars. Relativistic electrons accelerated above the polar cap can Compton scatter off thermal radiation from the neutron star surface, producing high energy gamma rays that can power pair cascades. Daugherty & Harding (1989) found that in the presence of a strong magnetic field, resonant scattering greatly increases electron energy losses over those of non-resonant scattering, making Compton scattering efficient even at lower temperatures. Recent studies have indicated the importance of resonant Compton scattering over curvature radiation in pulsar polar cap acceleration models. Luo (1996) and Zhang et al. (1997) observed that if the polar cap temperature and the magnetic field are sufficiently high, the thickness of the accelerating gap is limited more efficiently by pairs from Compton scattered photons than by pairs from curvature radiation photons. Harding & Muslimov (1998) considered ICS by the trapped, back flowing positrons and found that the pairs from the ICS photons may cause surface acceleration gaps to be unstable, forcing them to higher altitudes. These studies indicate the evolving, critical role of Compton scattering in polar cap models.
Compton scattering is also very important in SGR radiation models. The highly super-Eddington luminosities of the bursts ensures high densities of both photons and particles, so that scattering will be a critical factor. Paczynski (1992) proposed that the lower scattering cross section below the resonance in the strong magnetic field (for photons in the perpendicular polarization mode) could allow super-Eddington luminosities. However, Miller et al. (1995) argued that scattering into the parallel mode keeps the radiation pressure high, and thus the effective Eddington luminosity lower, in hydrostatic atmospheres. In Thompson & Duncan’s (1995) model for the radiation from SGR bursts, Compton scattering plays a critical role in establishing equilibrium between pairs and photons and in the spectral formation. Compton scattering may also be important in the photon splitting cascade model for SGR burst emission (Baring 1995, Harding, Baring & Gonthier 1996 and Harding, Baring & Gonthier 1997). The issues raised and discussed by these papers all depend critically on the polarization state and the angular distribution of the photons involved in a scattering event.
The full QED expressions for the relativistic, magnetic cross section of Compton scattering were derived separately by Daugherty and Harding (1986, hereafter DH86), and Bussard et al. (1986) and discussed by Meszaros (1992). Because of the complexity of the expressions, they have been applied to the study of relativistic Compton scattering in high magnetic fields only for the case of a one-dimensional, thermal electron distribution. These studies (Alexander & Meszaros 1989, 1991, Harding & Daugherty 1991, Araya & Harding 1996, 1999) focussed on models of cyclotron line formation in accreting neutron star atmospheres and gamma-ray bursts. The inverse Compton scattering models for pulsars and SGRs described in the preceeding paragraphs involve non-thermal, highly relativistic electrons. Such models have not to date incorporated the QED scattering cross sections, because the larger length scales require treatment of field inhomogeneity (unlike the cyclotron line scattering models which assume homogeneous fields). Consequently, these inverse Compton scattering studies have had to approximate the scattering rates by a combination of the Klein-Nishina cross section for non-resonant Compton scattering and the non-relativistic (Thomson) limit (Canuto, Lodenquai, & Ruderman 1971, Blandford & Scharlemann 1976 and Herold 1979) for resonant Compton scattering. As a result, they do not include the quantum relativistic effects of a strong magnetic field ($`B>0.1`$, where here and throughout the paper $`B`$ is given in units of the critical field, $`B_{\mathrm{cr}}=m_e^2c^3/e\mathrm{}=4.414\times 10^{13}`$Gauss).
It is the purpose of this paper to explicitly present the major features of Compton scattering in strong sub-critical and super-critical magnetic fields, providing a development of simplified expressions for the magnetic scattering cross section, to facilitate applications to environments near the surfaces of pulsars, SGRs or AXPs. We extend the work of DH86 by obtaining expressions suitable for rapid computation of the polarized, differential and integrated cross sections, applicable to Compton scattering of highly relativistic electrons moving along the magnetic field. In this case, broadly applicable in astrophysical problems, photons propagate along the field in the electron rest frame. In this specialized axisymmetric case, there is a degeneracy between polarization transitions $``$ and $``$, with a similar identity of cross sections for the modes $``$ and $``$. Below the cyclotron fundamental, mostly $``$-photons are produced in scatterings, a situation that also pertains above this resonance for sub-critical fields. However, an interesting discovery of this paper is that for super-critical fields, the reverse situation arises above the cyclotron fundamental, and a preponderance of photons of parallel polarization results from scatterings. We derive an analytic approximation that describes well the integrated cross section for Compton scattering in both sub-critical and super-critical magnetic fields. The effects of the strong magnetic fields on the angle-integrated cross sections and angular distributions of scattered photons is also studied, noting in particular significant differences for large scattering angles between the high energy magnetic forms and the non-magnetic Klein-Nishina results. Finally, we discuss the impact of this field-dependent cross section on the simulation of acceleration and cascade processes in pulsars and SGR sources.
## 2 QED COMPTON SCATTERING CROSS SECTION
The present study follows the development of DH86, applying their work to a particular case for scattering of relativistic electrons. The expressions developed in this paper use the Johnson & Lippmann (1949) electron wavefunctions. Graziani (1993) and Graziani, Harding & Sina (1995) have indicated that the choice of these wavefunctions do not diagonalize the magnetic moment operator parallel to the external field resulting in some limitations and inaccuracies when spin-dependent cross sections are used. Therefore, in this work, we present results that are spin-averaged. We will further discuss this issue in section 4. The differential cross section in the rest frame of the electron is given in equation (11) of DH86, with the denominator term later corrected in Harding & Daugherty (1991), by the expression
$$\frac{d\sigma }{d\mathrm{\Omega }^{}}=\frac{3\sigma _\mathrm{T}}{16\pi }\frac{\omega ^{}e^{(\omega ^2\mathrm{sin}^2\theta ^{}+\omega ^2\mathrm{sin}^2\theta )/2B}}{\omega \left(1+\omega \omega ^{}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\mathrm{cos}\theta ^{}\right)\left(2+\omega \omega ^{}\right)}\left|\underset{n=0}{\overset{\mathrm{}}{}}\left[F_n^{(1)}e^{i\mathrm{\Phi }_1}+F_n^{(2)}e^{i\mathrm{\Phi }_2}\right]\right|^2.$$
(1)
While photon-electron interactions may excite electrons above the ground state in this case (to quantum numbers $`\mathrm{}0`$), the initial electron is assumed to be in its ground state with $`m=0`$ Landau state and its spin anti-parallel to the magnetic field. This assumption is valid as the synchro-cyclotron decay of excited states occurs during an extremely short time-scale. The incident and scattered photon energies denoted by, $`\omega `$ and $`\omega ^{}`$, respectively are in units of the electron rest mass energy, $`m_ec^2`$. The incident and scattered angles are denoted by $`\theta `$ and $`\theta ^{}`$, respectively, with respect to the z-axis determined by the direction of the magnetic field. The common phase factors are given by
$$\mathrm{\Phi }_1=\mathrm{\Phi }_2=\frac{\omega \omega ^{}\mathrm{sin}\theta \mathrm{sin}\theta ^{}}{2B}\mathrm{sin}(\varphi \varphi ^{}).$$
(2)
The sum in equation (1) is over the intermediate Landau states, labelled by quantum number $`n`$. The $`F`$ terms and phase factors are associated with the two different Feynman diagrams shown in Figure 1, and are listed explicitly in the Appendix. The $`F`$ terms of the second Feynman diagram can be obtained from those of the first diagram through the crossing-symmetry replacements
$$\omega \omega ^{},kk^{},\beta \beta ^{},\text{and}\epsilon \epsilon ^{}.$$
(3)
For each Feynman diagram, each $`F`$ expression has a spin no-flip and a spin flip $`F`$ term associated with it, due to the spin degeneracy of the final states. Notation used in defining the $`F`$ terms and appearing elsewhere in this paper includes $`\epsilon `$, which represents the photon polarization components defining two orthogonal linear polarization vectors as given in Daugherty & Bussard (1980):
$`\epsilon ^{}`$ $`=`$ $`\mathrm{cos}\theta \mathrm{cos}\varphi \widehat{x}\mathrm{cos}\theta \mathrm{sin}\varphi \widehat{y}+\mathrm{sin}\theta \widehat{z},`$
$`\epsilon _\pm ^{}`$ $`=`$ $`\epsilon _x^{}\pm i\epsilon _y^{}=\mathrm{cos}\theta e^{\pm i\varphi },`$
$`\epsilon _z^{}`$ $`=`$ $`\mathrm{sin}\theta ,`$
$`\epsilon ^{}`$ $`=`$ $`\mathrm{sin}\varphi \widehat{x}\mathrm{cos}\varphi \widehat{y},`$
$`\epsilon _\pm ^{}`$ $`=`$ $`\epsilon _x^{}\pm i\epsilon _y^{}=ie^{\pm i\varphi },`$
$`\epsilon _z^{}`$ $`=`$ $`0`$
and the “vertex” functions, which in the notation of DH86 have the form
$`\mathrm{\Lambda }_{\mathrm{},m}(\beta )=(i)^{GS}\left({\displaystyle \frac{S!}{G!}}\right)^{1/2}2^{(G+S)/2}(\beta ^{})^{\mathrm{}}\beta ^m\left({\displaystyle \frac{\left|\beta \right|^2}{2}}\right)^SL_S^{GS}\left({\displaystyle \frac{\left|\beta \right|^2}{2}}\right),`$ (5)
where
$`G=\mathrm{max}(\mathrm{},m),`$ $`S=\mathrm{min}(\mathrm{},m),`$
$`\beta =i{\displaystyle \frac{(k_x+ik_y)}{\sqrt{B}}},`$ $`\beta ^{}=i{\displaystyle \frac{(k_x^{}+ik_y^{})}{\sqrt{B}}}`$
and $`L_n^a(x)`$ are the associated Laguerre polynomials.
## 3 SCATTERING OF RELATIVISTIC ELECTRONS
We consider in this study the particular problem of scattering of photons from relativistic electrons, common to a variety of astrophysical phenomena. Such relativistic scattering leads to considerable simplification of the algebra. In this section we develop expressions for scattering of ultra-relativistic electrons with $`\gamma 1`$ moving parallel to the magnetic field lines. Generally, the photon may have any angle of incidence, $`\psi _i`$, in the laboratory frame with respect to the magnetic field. Due to the large $`\gamma `$’s, the laboratory angle, $`\psi _i`$, gets Lorentz contracted to $`\theta =\psi _i/2\gamma 0`$ degrees in the electron rest frame. The magnetic, nonrelativistic Thomson cross section (Canuto, Lodenquai & Ruderman 1971, Blandford & Scharlemann 1976, and Herold 1979) consists of two parts dependent on the incident photon angle as shown in the expression
$$\sigma _{\mathrm{Thomson}}=\sigma _\mathrm{T}\left[\mathrm{sin}^2\theta +\frac{\omega ^2}{2}(1+\mathrm{cos}^2\theta )\left(\frac{1}{(\omega B)^2}+\frac{1}{(\omega +B)^2}\right)\right].$$
(7)
Well below the resonance for angles, $`\theta \omega /B1`$, corresponding to $`\psi _i2\gamma \omega /B`$ in the laboratory frame, the $`\mathrm{sin}^2\theta `$ term is smaller than the resonant term. In the case of isotropic incident photons in the laboratory frame (where $`0<\psi _i<\pi `$), this constraint requires $`\omega \pi B/2\gamma `$, easily achieved with the large $`\gamma `$ of relativistic electrons considered throughout this paper. Under this assumption, the incident photon is parallel to the magnetic field lines and has no perpendicular momentum; hence
$`\theta `$ $`=`$ $`0,`$
$`k_{}`$ $`=`$ $`0,`$
$`\epsilon _z`$ $`=`$ $`0,`$ (8)
$`\mathrm{\Phi }_1=\mathrm{\Phi }_2`$ $`=`$ $`0,`$
$`\beta `$ $`=`$ $`0.`$
The coordinate system can be arbitrarily set so that the azimuthal angle, $`\varphi =0`$. Since the incident photon is parallel to the z-axis, there is no component of the polarization vector, $`\epsilon _z`$, along the z-axis, thereby, eliminating several terms in the $`F`$’s. The vertex functions associated with the incident photon become Kronecker delta functions, $`\mathrm{\Lambda }_{l,m}(0)=\delta _{lm}`$. This has the advantage of eliminating the sum over the intermediate states as only certain specific states contribute depending on the final Landau state, $`\mathrm{}`$. The vertex functions associated with the scattered photon can be written in terms of a single function using the following recursion relations
$`\mathrm{\Lambda }_{\mathrm{}1,0}(\beta ^{})`$ $`=`$ $`i{\displaystyle \frac{\sqrt{2\mathrm{}}}{\beta ^{}}}\mathrm{\Lambda }_{\mathrm{},0}(\beta ^{}),`$
$`\mathrm{\Lambda }_{\mathrm{},1}(\beta ^{})`$ $`=`$ $`{\displaystyle \frac{i\sqrt{2}}{\beta ^{}}}\left(\mathrm{}{\displaystyle \frac{\left|\beta ^{}\right|^2}{2}}\right)\mathrm{\Lambda }_{\mathrm{},0}(\beta ^{}),`$
$`\mathrm{\Lambda }_{\mathrm{}1,1}(\beta ^{})`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\mathrm{}}}{(\beta ^{})^2}}\left(\mathrm{}1{\displaystyle \frac{\left|\beta ^{}\right|^2}{2}}\right)\mathrm{\Lambda }_{\mathrm{},0}(\beta ^{}),`$ (9)
$`\mathrm{\Lambda }_{\mathrm{}+1,0}(\beta ^{})`$ $`=`$ $`i{\displaystyle \frac{\beta ^{}}{\sqrt{2(\mathrm{}+1)}}}\mathrm{\Lambda }_{\mathrm{},0}(\beta ^{}),`$
$`\mathrm{\Lambda }_{\mathrm{}2,0}(\beta ^{})`$ $`=`$ $`{\displaystyle \frac{2\sqrt{\mathrm{}(\mathrm{}1)}}{(\beta ^{})^2}}\mathrm{\Lambda }_{\mathrm{},0}(\beta ^{}).`$
The $`F`$ terms associated with each Feynman diagram can then be rewritten, with the common vertex function $`\mathrm{\Lambda }_{l,0}(\beta ^{})`$ being factored out of the scattering amplitudes to generate coefficients $`G_i`$ such that $`G_i\mathrm{\Lambda }_{\mathrm{},0}(\beta ^{})=F_{n,i}^{(1)}+F_{n,i}^{(2)}`$, where $`i`$ = flip and no-flip. Using the identity
$$\left|\mathrm{\Lambda }_{\mathrm{},0}(\beta ^{})\right|^2=\frac{\left|\beta ^{}\right|^2\mathrm{}}{2^{\mathrm{}}\mathrm{}!}=\frac{1}{\mathrm{}!}\left(\frac{\omega ^2\mathrm{sin}^2\theta ^{}}{2B}\right)^{\mathrm{}},$$
(10)
the differential cross section can be expressed compactly as
$$\frac{d\sigma _,}{d\mathrm{cos}\theta ^{}}=\frac{3\sigma _\mathrm{T}}{16\pi }\frac{\omega ^2e^{\omega ^2\mathrm{sin}^2\theta ^{}/2B}}{\omega (2+\omega \omega ^{})(\omega ^{}+\omega \omega ^{}(1\mathrm{cos}\theta ^{})\omega ^2\mathrm{sin}^2\theta ^{})}\frac{1}{\mathrm{}!}\left(\frac{\omega ^2\mathrm{sin}^2\theta ^{}}{2B}\right)^{\mathrm{}}G_,,$$
(11)
where
$$G_{}=\widehat{G}_{noflip}^{}+\widehat{G}_{flip}^{},G_{}=\widehat{G}_{noflip}^{}+\widehat{G}_{flip}^{},$$
(12)
and
$`\widehat{G}_{noflip}^{||}`$ $`=`$ $`{\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{noflip}^{||,||}\right|^2𝑑\varphi ^{}={\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{noflip}^{,||}\right|^2𝑑\varphi ^{}`$
$`=`$ $`2\pi \left\{\left[(B_1+B_3+B_7)\mathrm{cos}\theta ^{}(B_2+B_6)\mathrm{sin}\theta ^{}\right]^2+\left[B_4\mathrm{cos}\theta ^{}B_5\mathrm{sin}\theta ^{}\right]^2\right\},`$
$`\widehat{G}_{noflip}^{}`$ $`=`$ $`{\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{noflip}^{||,}\right|^2𝑑\varphi ^{}={\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{noflip}^,\right|^2𝑑\varphi ^{}`$
$`=`$ $`2\pi \left\{(B_1B_3B_7)^2+B_4^2\right\},`$
$`\widehat{G}_{flip}^{||}`$ $`=`$ $`{\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{flip}^{||,||}\right|^2𝑑\varphi ^{}={\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{flip}^{,||}\right|^2𝑑\varphi ^{}`$
$`=`$ $`2\pi \left\{\left[(C_1+C_3+C_7)\mathrm{cos}\theta ^{}(C_2+C_6)\mathrm{sin}\theta ^{}\right]^2+\left[C_4\mathrm{cos}\theta ^{}C_5\mathrm{sin}\theta ^{}\right]^2\right\},`$
$`\widehat{G}_{flip}^{}`$ $`=`$ $`{\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{flip}^{||,}\right|^2𝑑\varphi ^{}={\displaystyle \underset{0}{\overset{2\pi }{}}}\left|G_{flip}^,\right|^2𝑑\varphi ^{}`$
$`=`$ $`2\pi \left\{(C_1C_3C_7)^2+C_4^2\right\}.`$
In these developments, the $`\varphi ^{}`$ dependence and the imaginary terms are isolated in the polarization components and the phase exponentials, leading to elementary integrations over the azimuthal angle, $`\varphi ^{}`$. The $`B`$ and $`C`$ terms have the following forms:
$`B_1`$ $`=`$ $`{\displaystyle \frac{\left[2\omega \omega \omega ^{}(1\mathrm{cos}\theta ^{})\right]}{2(\omega B)}},`$
$`B_2`$ $`=`$ $`{\displaystyle \frac{\left[(\omega \omega ^{}\mathrm{cos}\theta ^{})\left(2lB\omega ^2\mathrm{sin}^2\theta ^{}\right)+2lB\omega \right]}{2\omega ^{}\mathrm{sin}\theta ^{}(\omega B)}},`$
$`B_3`$ $`=`$ $`{\displaystyle \frac{lB\left(2lB2B\omega ^2\mathrm{sin}^2\theta ^{}\right)}{\omega ^2\mathrm{sin}^2\theta ^{}(\omega B)}},`$
$`B_4`$ $`=`$ $`{\displaystyle \frac{\left[2\omega ^{}+\omega \omega ^{}(1\mathrm{cos}\theta ^{})\omega ^2\mathrm{sin}^2\theta ^{}\right]}{2\left[\omega \omega ^{}(1\mathrm{cos}\theta ^{})\omega B\right]}},`$
$`B_5`$ $`=`$ $`{\displaystyle \frac{(\omega \omega ^{}\mathrm{cos}\theta ^{})\omega ^{}\mathrm{sin}\theta ^{}}{2\left[\omega \omega ^{}(1\mathrm{cos}\theta ^{})\omega B\right]}},`$
$`B_6`$ $`=`$ $`{\displaystyle \frac{\mathrm{}B\mathrm{cos}\theta ^{}}{\mathrm{sin}\theta ^{}\left[\omega \omega ^{}(1\mathrm{cos}\theta ^{})\omega +B\right]}},`$
$`B_7`$ $`=`$ $`{\displaystyle \frac{2\mathrm{}(\mathrm{}1)B^2}{\omega ^2\mathrm{sin}^2\theta ^{}\left[\omega \omega ^{}(1\mathrm{cos}\theta ^{})\omega +B\right]}},`$
$`C_1`$ $`=`$ $`\sqrt{2\mathrm{}B}{\displaystyle \frac{\omega }{2(\omega B)}},`$
$`C_2`$ $`=`$ $`\sqrt{2lB}{\displaystyle \frac{\left[2\omega +2\omega ^2\omega \omega ^{}(1\mathrm{cos}\theta ^{})2lB+\omega ^2\mathrm{sin}^2\theta ^{}\right]}{2\omega ^{}\mathrm{sin}\theta ^{}(\omega B)}},`$
$`C_3`$ $`=`$ $`\sqrt{2lB}{\displaystyle \frac{\left(\omega \omega ^{}\mathrm{cos}\theta ^{}\right)\left(2lB2B\omega ^2\mathrm{sin}^2\theta ^{}\right)}{2\omega ^2\mathrm{sin}^2\theta ^{}(\omega B)}},`$
$`C_4`$ $`=`$ $`\sqrt{2\mathrm{}B}{\displaystyle \frac{\omega ^{}\mathrm{cos}\theta ^{}}{2\left[\omega ^{}\omega (1\mathrm{cos}\theta ^{})\omega B\right]}},`$
$`C_5`$ $`=`$ $`\sqrt{2\mathrm{}B}{\displaystyle \frac{\omega ^{}\mathrm{sin}\theta ^{}}{2\left[\omega ^{}\omega (1\mathrm{cos}\theta ^{})\omega B\right]}},`$
$`C_6`$ $`=`$ $`\sqrt{2\mathrm{}B}{\displaystyle \frac{\left[2\omega ^{}+\omega \omega ^{}(1\mathrm{cos}\theta ^{})\omega ^2\mathrm{sin}^2\theta ^{}\right]}{2\omega ^{}\mathrm{sin}\theta ^{}\left[\omega ^{}\omega (1\mathrm{cos}\theta ^{})\omega +B\right]}},`$
$`C_7`$ $`=`$ $`\sqrt{2\mathrm{}B}{\displaystyle \frac{(l1)B\left(\omega \omega ^{}\mathrm{cos}\theta ^{}\right)}{\omega ^2\mathrm{sin}^2\theta ^{}\left[\omega ^{}\omega (1\mathrm{cos}\theta ^{})\omega +B\right]}}.`$
Since this combination of expressions for the differential cross section is for the particular case of scattering of relativistic electrons with an incident photon angle of zero degrees in the electron rest frame, the resulting scattering rates are simple and do not possess the sum over Bessel functions as in Bussard et al. (1986). While there are four different possibilities for the scattering of polarized photons, ($``$, $``$, $``$, and $``$), for this special case under consideration, the expressions for the $``$, $``$ scattering have the same form, as well as those for the $``$, and $``$ scattering, as indicated above, resulting in two polarization channels in which the scattering process will produce either parallel or perpendicular polarized photons.
Observe also that the differential cross section is dependent on the final Landau state, $`\mathrm{}`$. Hence, in order to derive the complete contribution, a sum must be performed over all the contributing Landau states. Since the energy of the scattered photon may be expressed as
$$\omega ^{}=\frac{2(\omega \mathrm{}B)}{1+\omega (1\mathrm{cos}\theta ^{})+\left[\left(1+\omega (1\mathrm{cos}\theta ^{})\right)^22(\omega \mathrm{}B)\mathrm{sin}^2\theta ^{}\right]^{1/2}},$$
(15)
where $`\mathrm{}`$ is the final Landau state of the scattered electron, each final state has an energy threshold of $`\mathrm{}B`$. Therefore, the maximum contributing Landau state quantum number, $`\mathrm{}_{max}`$, to the cross section is pinned to the photon energy in cyclotron units: $`\omega /B1<\mathrm{}_{max}<\omega /B`$.
## 4 ANGLE-INTEGRATED CROSS SECTIONS
The differential cross section can be numerically integrated over $`\theta ^{}`$ using a Romberg integration scheme to obtain the energy dependent cross section. In Figure 2, we display the QED, exact angle-integrated cross section (solid curves) for the indicated magnetic fields, in units of $`B_{\mathrm{cr}}`$, as a function of the incident photon energy, $`\omega /B`$, in cyclotron energy units. We have averaged over the final spin of the electron and over the polarization of the scattered photon. For this particular case in the scattering of relativistic electrons, there is only one resonance occurring at the fundamental cyclotron resonance of $`\omega _B=B`$. We scale the photon energy by the cyclotron energy so that the resonance occurs at the same place independent of the magnetic field, $`B`$. For comparison, we also plot in the figures the nonrelativistic Thomson limit (dot-dashed curves) and the Klein-Nishina (dotted curves) predictions. The solid circles are the result of an approximation that is discussed later in section 6. For the exact calculation, we have summed over the all the contributing final Landau states.
Above the resonance, the exact cross section approaches the Klein-Nishina cross section. As expected for smaller fields, the convergence occurs at lower photon energies, as seen in the case of $`B=0.1`$. At this field strength, typical of radio pulsars, there are no significant deviations from the Thomson limit below the resonance and from the Klein-Nishina limit above the resonance. The main discrepancy occurs right above the resonance where the two limiting cases do not match the exact cross section. As the field strength increases, the exact cross section below the resonance drops significantly beneath the Thomson limit by over a factor of ten in the case of $`B=100`$ For scattering above the resonance at these high fields, there are deviations between the exact cross section and the Klein-Nishina cross section. However, as the energy increases, the exact cross section and the Klein-Nishina cross section appear to converge as seen in the cases of $`B=0.1`$ and $`B=1`$.
The trend as $`B`$ increases, evident in Figure 2, is for the magnitude of the cross section to drop at all energies, while the width of the resonance increases (for $`B1`$, when scaled in units of the cyclotron energy, this width actually declines). Since the resonance is formally divergent, the common practice (Xia et al. 1985, Daugherty & Harding 1989 and Dermer 1990) is to truncate it at $`\omega =B`$ by introducing a finite width $`\mathrm{\Gamma }`$ equal to the cyclotron decay width (inverse lifetime) for the $`\mathrm{}=10`$ transitions, corresponding to decay of an excited intermediate electron state. The procedure is to replace the resonant $`(\omega B)^2`$ denominator (e.g. see equation (LABEL:eq:dsig\_leq0\_approx) below) by $`\left[(wB)^2+G^2/4\right]`$. In the $`B1`$ regime, the cyclotron decay width assumes the well-known result $`\mathrm{\Gamma }=4\alpha B^2/3`$ in dimensionless units. When $`B1`$, Latal (1986) deduced that $`\mathrm{\Gamma }`$ is proportional to $`B^{1/2}`$, a dependence that can be inferred from the exact decay widths enunciated in equation (17) of Herold, Ruder & Wunner (1982) by noting that the angular distribution of radiation in cyclotron transitions is still quasi-isotropic (and not strongly beamed) for highly-supercritical fields. When substituted into the Lorentz profile prescription for truncating the resonance, these widths lead to the areas under the resonance (i.e. when integrating over $`\omega `$) being independent of $`B`$ in the magnetic Thomson regime of $`B1`$, and scaling as $`B^{1/2}`$ when $`B1`$; these results can be deduced using the $`\mathrm{}=0`$ approximation derived in equation (LABEL:eq:dsig\_leq0\_approx).
This area is an approximate measure of the importance of resonant Compton scattering for a particular situation. To see this, note that astrophysical problems generally have a broad energy distribution of beamed relativistic electrons interacting with not-so-highly-beamed low energy (soft) photons, so that a particular electron energy and soft photon angle (with respect to $`B`$) determines the value of $`\omega `$ in the electron rest frame. Summing over the electron energies and incoming photon angles amounts to a weighted integration of the area under the cross section. The weighting function is usually not very sensitive to $`\omega `$, so that the area under the curves gives a representative indication of the strength of resonant (as opposed to non-resonant) scattering provided the resonance is not too broad. Hence, it can be inferred that the resonant process is more important in supercritical fields than when $`B1`$. This conclusion stands even when it is noted that the magnetic Compton resonance is not truly cyclotronic in nature: the contribution of the $`\mathrm{}=1`$ transition to the right wing of the resonance plus the spread of parallel momenta introduce non-Lorentzian distortions to the resonance profile.
The photon polarization-dependent cross sections can also be easily obtained for the given integration of the $`G`$ terms shown above. In Figure 3, we present the QED Compton scattering cross section as a function of energy for magnetic fields of 0.1 and 10 times $`B_{\mathrm{cr}}`$. As mentioned earlier, for this particular case of scattering of relativistic electrons, there are two polarization scattering channels, in which the scattering leads to photons with parallel polarization (dashed curve) and with perpendicular polarization (dotted curve). The total cross section is shown as a solid curve. In the $`B=10`$ case, above the resonance, the scattering process preferably produces photons with parallel polarization, whereas below the resonance, the channel producing perpendicularly polarized photons dominates. This behavior, where perpendicular-polarized scattered photons dominate below the resonance and parallel-polarized scattered photons dominate above the resonance, is characteristic of the magnetic-relativistic cross section. In the nonrelativistic case, the perpendicular polarization channel will be three times larger than the parallel polarization channel, but has the same shape at all photon energies, as observed from equation (LABEL:eq:dsig\_leq0\_approx) letting $`\omega ^{}\omega `$ and integrating over $`\mathrm{cos}\theta ^{}`$. As can be seen in the Figure 3 for the $`B=0.1`$ case, $``$-polarization in the QED cross section dominates at low fields, thus the switching to $``$-polarization dominance at high fields is a relativistic effect. In the Klein-Nishina scattering there is no magnetic field and the initial photon polarization is important in determining the final photon polarization. In this case, there are three channels, $``$, $``$, and two degenerate switching channels $``$ and $``$. In the Klein-Nishina regime, right above the resonance, the $``$ channel dominates over the $``$ channel and the switching channels $``$ and $``$ which have the smallest contribution. Far above the resonance at high photon energies, the cross sections for the three channels merge as in the magnetic cross section in Figure 3.
In Figure 4, we show the contributions of the indicated final Landau states to the total cross section for the indicated field strength of 10 times $`B_{\mathrm{cr}}`$. The total cross section, represented by a thick-solid curve, is a result of summing over all contributing Landau states. Below the resonance, only the $`\mathrm{}=0`$ final state contributes (dotted curve) to the cross section due to the previously mentioned threshold associated with each $`\mathrm{}`$. The curve associated with $`\mathrm{}=1`$, having a similar shape as the $`\mathrm{}=0`$ curve, is plotted also as a dotted curve. The light solid curves represent a select group of the indicated higher final Landau states. As the photon energy increases, higher $`\mathrm{}`$ states may contribute more significantly than lower ones. For example, above a photon energy of 50, the $`\mathrm{}=10`$ state contributes more to the overall cross section than the $`\mathrm{}=0`$ or 1 states (dotted curves). Clearly for scattering above the resonance, many final states must be included for computational accuracy.
As mentioned earlier, the cross section of DH86 was derived using Johnson-Lippmann (JL) electron wavefunctions which do not correctly describe the spin-dependence of the S-matrix elements, but produce correct spin-averaged S-matrix elements. Thus, we have averaged over the initial and summed over final electron spin states in the expressions we derived. However, there is still a small error in the JL cross section at and above the cyclotron resonance, due to the spin-dependence of the intermediate states. This error is negligible for $`B<0.1`$ but grows with $`B`$ for $`B>0.1`$. We have evaluated this error through a numerical comparison of the JL cross section of DH86, derived in this paper that neglect the decay width of the intermediate states, with the cross section derived by Sina (1996) who used Sokolov-Ternov (ST) wavefunctions. For the case of the scattering of relativistic electrons with $`\psi _i=0`$, the two spin-averaged cross sections agree below the cyclotron energy, $`B`$, but do not quite agree for $`\omega B`$. We find that for $`B=0.1`$, there is a small error of 0.01% at $`\omega =B`$ in the spin-averaged cross section and for $`B=10`$, there is a somewhat larger error of 0.4% at $`\omega =B`$. The increasing error with increasing magnetic field is due to the intermediate states having non-zero momentum which increases the difference between JL and ST cross sections at higher $`B`$ fields. Thus the simplified expressions derived in this paper are accurate in their regions of validity. Furthermore, the cyclotron energy is high enough in supercritical fields that scattering above the resonance is less important than it is in subcritical fields. We plan to use the ST cross section in future derivation of simplified expressions for the scattering cross section above the cyclotron resonance.
## 5 SCATTERING TO $`\mathrm{}=0`$ FINAL STATES - BELOW THE RESONANCE
Due to the presence of the resonance in the cross section, the scattering process will try to select out resonant scattering, if the geometry and the kinematics permit. As the magnetic field increases, the photon energy, in the electron rest frame, required for resonant scattering (i.e., $`\omega =B`$) increases. If the source of photons is limited by blackbody temperatures and the field strength is large, the scattering will predominately occur much below the resonance. The cross sections, described here, are in the rest frame of the electron. If the electron is moving, the Lorentz-transformed photon energy in the electron rest frame is given by
$$E_{rest}=E_{lab}\gamma (1\beta \mathrm{cos}\psi _i),$$
(16)
where $`E_{rest}`$ and $`E_{lab}`$ are the energies of the photon in the rest and laboratory frames, respectively, and $`\psi _i`$ is the laboratory angle of the photon with respect to the electron direction. For small angles $`\psi _i0`$, where the photon is moving in the same direction as the electron, the photon energy is red shifted, $`E_{rest}E_{lab}/(2\gamma )`$. For large angles $`\psi _i\pi `$, where the photon and electron are colliding head on, the photon energy is blue shifted, $`E_{rest}2\gamma E_{lab}`$. For perpendicular scattering, the photon energy is also blue shifted, $`E_{rest}\gamma E_{lab}`$. In general, when the incident angle, $`\psi _i>\sqrt{2/\gamma }`$, the photon energy will be blue shifted. For relativistic electrons, the geometry of the interaction becomes important in determining whether the scattering takes place above or below the resonance.
If the scattering occurs below the resonance, the only contribution to the cross section is from the $`\mathrm{}=0`$ final state, since the final Landau state, $`\mathrm{}`$, state has an energy threshold of $`\mathrm{}B`$. Using the following identity valid for $`\mathrm{}=0`$,
$$\omega ^2\mathrm{sin}^2\theta ^{}=2\omega ^{}2\omega \omega \omega ^{}(1\mathrm{cos}\theta ^{}),$$
(17)
in the denominator of equation (11), the expression for the exact, QED, cross section of the $`\mathrm{}=0`$ final state has the following form
$$\frac{d\sigma }{d\mathrm{cos}\theta ^{}}=\frac{3\sigma _\mathrm{T}}{16\pi }\frac{\omega ^2}{\omega (2+\omega \omega ^{})(\xi \omega ^{})}e^{\omega ^2\mathrm{sin}^2\theta ^{}/2B}G,$$
(18)
where
$$\xi =2\omega \omega \omega ^{}(1\mathrm{cos}\theta ^{}).$$
(19)
All the $`C`$ terms associated with the spin flip drop out for the $`\mathrm{}=0`$ case, in addition the $`B_3`$, $`B_6`$ and $`B_7`$ terms are zero. The remaining nonzero $`B`$ terms simplify to the following
$`B_1`$ $`=`$ $`{\displaystyle \frac{\xi }{2(\omega B)}},`$
$`B_2`$ $`=`$ $`{\displaystyle \frac{\eta }{2(\omega B)}},`$
$`B_4`$ $`=`$ $`{\displaystyle \frac{\xi }{2(\omega \xi B)}},`$ (20)
$`B_5`$ $`=`$ $`{\displaystyle \frac{\eta }{2(\omega \xi B)}},`$
where
$$\eta =(\omega \omega ^{}\mathrm{cos}\theta ^{})\omega ^{}\mathrm{sin}\theta ^{}.$$
(21)
The exact polarization-dependent and averaged cross sections for $`\mathrm{}=0`$ can be expressed as
$`{\displaystyle \frac{d\sigma ^{||||}}{d\mathrm{cos}\theta ^{}}}`$ $`=`$ $`{\displaystyle \frac{d\sigma ^{||}}{d\mathrm{cos}\theta ^{}}}={\displaystyle \frac{3\sigma _\mathrm{T}}{32}}{\displaystyle \frac{\omega ^2e^{\omega ^2\mathrm{sin}^2\theta ^{}/2B}}{\omega (2+\omega \omega ^{})(\xi \omega ^{})}}`$
$`\times \left\{\left(\xi \mathrm{cos}\theta ^{}\eta \mathrm{sin}\theta ^{}\right)^2\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega \xi B)^2}}\right]\right\},`$
$`{\displaystyle \frac{d\sigma ^{||}}{d\mathrm{cos}\theta ^{}}}`$ $`=`$ $`{\displaystyle \frac{d\sigma ^{}}{d\mathrm{cos}\theta ^{}}}={\displaystyle \frac{3\sigma _\mathrm{T}}{32}}{\displaystyle \frac{\omega ^2e^{\omega ^2\mathrm{sin}^2\theta ^{}/2B}}{\omega (2+\omega \omega ^{})(\xi \omega ^{})}}`$ (23)
$`\times \left\{\xi ^2\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega \xi B)^2}}\right]\right\},`$
$`{\displaystyle \frac{d\sigma ^{ave}}{d\mathrm{cos}\theta ^{}}}`$ $`=`$ $`{\displaystyle \frac{3\sigma _\mathrm{T}}{32}}{\displaystyle \frac{\omega ^2e^{\omega ^2\mathrm{sin}^2\theta ^{}/2B}}{\omega (2+\omega \omega ^{})(\xi \omega ^{})}}`$
$`\times \left\{\left[\left(\xi \mathrm{cos}\theta ^{}\eta \mathrm{sin}\theta ^{}\right)^2+\xi ^2\right]\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega \xi B)^2}}\right]\right\}.`$
The average cross section is obtained by summing over the final and averaging over the initial photon polarizations.
## 6 APPROXIMATING THE $`\mathrm{}=0`$ FINAL STATES
An approximation to the exact $`\mathrm{}=0`$ differential cross section can be given by assuming that the scattering is significantly below the resonance, where $`\omega <B`$, and also $`\omega <1`$. We make this approximation by keeping only terms to first order in $`\omega `$ and $`\omega ^{}`$ in the two forms,
$`\xi `$ $`=`$ $`2\omega \omega \omega ^{}(1\mathrm{cos}\theta ^{})2\omega ,`$
$`\eta `$ $`=`$ $`(\omega \omega ^{}\mathrm{cos}\theta ^{})\omega ^{}\mathrm{sin}\theta ^{}0.`$ (24)
This leads to the following approximate expressions
$`{\displaystyle \frac{d\sigma ^{||||}}{d\mathrm{cos}\theta ^{}}}`$ $``$ $`{\displaystyle \frac{d\sigma ^{||}}{d\mathrm{cos}\theta ^{}}}={\displaystyle \frac{3\sigma _\mathrm{T}}{8}}{\displaystyle \frac{\omega \omega ^2\mathrm{cos}^2\theta ^{}}{(2\omega \omega ^{})}}\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega +B)^2}}\right],`$
$`{\displaystyle \frac{d\sigma ^{||}}{d\mathrm{cos}\theta ^{}}}`$ $``$ $`{\displaystyle \frac{d\sigma ^{}}{d\mathrm{cos}\theta ^{}}}={\displaystyle \frac{3\sigma _\mathrm{T}}{8}}{\displaystyle \frac{\omega \omega ^2}{(2\omega \omega ^{})}}\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega +B)^2}}\right],`$
$`{\displaystyle \frac{d\sigma ^{ave}}{d\mathrm{cos}\theta ^{}}}`$ $``$ $`{\displaystyle \frac{3\sigma _\mathrm{T}}{8}}{\displaystyle \frac{\omega \omega ^2(1+\mathrm{cos}^2\theta ^{})}{(2\omega \omega ^{})}}\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega +B)^2}}\right].`$
The nonrelativistic approximation would further lead to $`\omega ^{}=\omega `$, resulting in the nonrelativistic expressions found in Herold (1979).
These expressions, while much simpler than the previous ones, are still complex due to the functional form of $`\omega ^{}`$, having an angular dependence in its square root. Yet they are manageable and can be integrated over $`\mathrm{cos}\theta ^{}`$. After careful algebra, integration over $`\mathrm{cos}\theta ^{}`$ yields the following polarization-dependent and averaged, approximate cross sections
$`\sigma ^{||||}`$ $`=`$ $`\sigma ^{||}={\displaystyle \frac{3\sigma _\mathrm{T}}{16}}\left\{g(\omega )h(\omega )\right\}\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega +B)^2}}\right],`$
$`\sigma ^{||}`$ $`=`$ $`\sigma ^{}={\displaystyle \frac{3\sigma _\mathrm{T}}{16}}\left\{f(\omega )2\omega h(\omega )\right\}\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega +B)^2}}\right],`$ (26)
$`\sigma ^{ave}`$ $`=`$ $`{\displaystyle \frac{3\sigma _\mathrm{T}}{16}}\left\{g(\omega )+f(\omega )(1+2\omega )h(\omega )\right\}\left[{\displaystyle \frac{1}{(\omega B)^2}}+{\displaystyle \frac{1}{(\omega +B)^2}}\right].`$
where
$`g(\omega )`$ $`=`$ $`{\displaystyle \frac{\omega ^2(3+2\omega )+2\omega }{\sqrt{\omega (2+\omega )}}}\mathrm{ln}\left(1+\omega \sqrt{\omega (2+\omega )}\right)+{\displaystyle \frac{\omega }{2}}\mathrm{ln}(1+4\omega )`$
$`+\omega (1+2\omega )\mathrm{ln}(1+2\omega )+2\omega ,`$
$`f(\omega )`$ $`=`$ $`\omega ^2\mathrm{ln}(1+4\omega )+\omega (1+2\omega )\mathrm{ln}(1+2\omega ),`$ (27)
$`h(\omega )`$ $`=`$ $`\{\begin{array}{cc}\frac{\omega ^2}{\sqrt{\omega (2\omega )}}\mathrm{tan}^1\left(\frac{\sqrt{\omega (2\omega )}}{1+\omega }\right),\hfill & \text{for }\omega <2\text{,}\hfill \\ \frac{\omega ^2}{2\sqrt{\omega (\omega 2)}}\mathrm{ln}\left(\frac{\left(1+\omega +\sqrt{\omega (\omega 2)}\right)^2}{1+4\omega }\right),\hfill & \text{for }\omega >2\text{.}\hfill \end{array}`$ (30)
While one might expect $`h(\omega )`$ to become imaginary when $`\omega >2`$, given the $`(2\omega )`$ factor in the square-root terms in first expression for $`h(\omega )`$, the expression is completely general and remains real even when $`\omega >2`$. For the purpose of numerical calculations in a computer code, we introduce this second expression for $`h(\omega )`$ that can be coded using a natural logarithm when $`\omega >2`$.
Back in Figure 2, the solid circles represent the polarization-averaged cross section obtained from the above analytical approximation. The approximation is valid in the region below $`\omega <1`$ corresponding to $`\omega /B<1/B`$ along the photon energy axis in Figure 2. In this region, it agrees very well with the exact $`\mathrm{}=0`$ cross section. Above the region of validity, the approximation over estimates the exact $`\mathrm{}=0`$ cross section. However, the approximation does surprisingly well, when compared to the exact cross section, extrapolating above the region of validity. While the analytical approximation is a result of integrating the approximation to the exact $`\mathrm{}=0`$ differential cross section, it remains close to the total cross section at high $`\omega `$ above the resonance $`(\omega >B)`$ even for high field strengths.
## 7 ANGULAR DISTRIBUTIONS
We present in Figure 5, the differential cross section, $`d\sigma /d\theta ^{}`$, for a magnetic field of $`B=10B_{\mathrm{cr}}`$. In order to understand best the behavior of the angular distributions, we have plotted the differential cross section as a function of the logarithm of the scattered angle, $`\theta ^{}`$, of the photon in the electron rest frame. We sample the angular distributions beginning below the resonance, where only the $`\mathrm{}=0`$ final Landau state contributes, to high above the resonance where many Landau states contribute up to the threshold $`\mathrm{}_{max}`$: $`\omega /B1<\mathrm{}_{max}<\omega /B`$. The incident photon energies for each panel are indicated in units of the cyclotron energy, $`\omega _B=B`$. Also indicated are the final Landau states of the calculated angular distributions. As expected, the $`\mathrm{}=0`$ contribution is strong for all photon energies. As the photon energy increases high above the resonance (right panels of the figure), the angular distribution of the $`\mathrm{}=0`$ state reveals a dip at an energy-dependent angle in the forward direction. This dip is very steep, as indicated by the number of decades the distribution drops before approaching its minimum. Care must be exercised in this region, as one integrates the angular distributions. As the final Landau state increases, the angular distributions become more gaussian shaped, peaking at the same angle for a given photon energy as the minimum in the $`\mathrm{}=0`$ state. Above the resonance, the angular distributions evolve smoothly from a sharp minimum at low Landau states to a maximum at higher Landau states. Both the minima and maxima occur at invariably the same angle, $`\theta _o`$. This behavior is due to the functional form of the following factor in the differential cross section
$$f(\omega ,\theta ^{})=\frac{1}{\mathrm{}!}\left(\frac{\omega ^2\mathrm{sin}^2\theta ^{}}{2B}\right)^{\mathrm{}}\mathrm{exp}\left\{\frac{\omega ^2\mathrm{sin}^2\theta ^{}}{2B}\right\},$$
(31)
which controls the angular dependence of the differential cross section, $`d\sigma /d\mathrm{cos}\theta ^{}`$. The first part of this function arises from the vertex functions mentioned earlier in equation (10). The first derivative with respect to $`\theta ^{}`$ of this function goes to zero an angle, $`\theta _o`$, given by
$$\theta _o=\mathrm{tan}^1\left(\frac{\sqrt{1+2\omega }}{\omega }\right),$$
(32)
independent of $`\mathrm{}`$. Since $`f(\omega ,\theta ^{})`$ is only part of the overall differential cross section, this expression for $`\theta _o`$ is approximate. Yet, it predicts very well the angle at which the angular distributions experience minima and maxima with increasing Landau states. The scattered photon energy, $`\omega _o^{}`$, at which this peak in the angular distribution occurs is given by
$$\omega _o^{}=\frac{1+\omega }{1+2\omega }\left(1+2\omega \sqrt{(1+2\omega )(1+2\mathrm{}B)}\right).$$
(33)
The Landau state, $`\mathrm{}_s`$, at which the angular distribution evolves from a minimum to a maximum at a given photon energy, $`\omega `$, occurs when the second derivative of the above function, $`f(\omega ,\theta ^{})`$, with respect to $`\theta ^{}`$ is equal to zero and is given by the expression
$$\mathrm{}_s=\frac{\omega ^2}{2B(1+2\omega )}.$$
(34)
The energy of scattered photon at $`\mathrm{}_s`$ and at $`\theta _o`$ has the simple form
$$\omega _s^{}=\frac{(1+\omega )\omega }{1+2\omega }.$$
(35)
These expressions have served to guide the design of the algorithm that numerically integrates the angular distributions to obtain the integrated cross sections shown in Figure 2. For Landau states below, $`\mathrm{}_s`$, the angular distributions manifest the steep drop at $`\theta _o`$, therefore we integrate the angular distribution in two parts from $`\theta =0`$ to $`\theta _o`$ and from $`\theta _o`$ to $`\pi `$. When the Landau state is above, $`\mathrm{}`$s, the angular distributions are gaussian-shaped, and we integrate from $`\theta =0`$ to $`\pi `$.
The peak in the total angular distribution will also occur very near this angle $`\theta _o`$ as seen in Figure 6. Here we present the total angular distributions summed over all the contributing final Landau states represented by the solid curves. The dashed curves show the angular distribution for the $`\mathrm{}=0`$ state. The approximation to the exact $`\mathrm{}=0`$ differential cross section given in equation (LABEL:eq:dsig\_leq0\_approx) (averaged) is plotted as dot-dashed curves. Comparing the exact-summed angular distributions (solid curves) to the approximate distributions (dot-dashed) in Figure 6, one can see in Figure 2 that the approximation falls below the integrated cross sections because there is a deficiency in the approximation at large angles. As expected, the $`\mathrm{}=0`$ state is the dominant contribution to the angular distribution at photon energies below and right above the resonance. As the photon energy increases well beyond the resonance, the $`\mathrm{}=0`$ state contributes less significantly and higher Landau states become increasingly important contributors. Although the strength of large Landau states decreases rapidly as shown in Figure 5, they are numerous and contribute significantly when summed to obtain the overall angular distributions shown in Figure 6. The contribution of these higher Landau states occurs near $`\theta _o`$, where the total angular distributions peak. It is at $`\theta _o`$ where minima occur in the angular distributions of final states of $`\mathrm{}<\mathrm{}_s`$. Yet the approximation to the $`\mathrm{}=0`$ final Landau state peaks at approximately the angle $`\theta _o`$. The angular distributions of the exact, the Klein-Nishina, and the $`\mathrm{}=0`$ approximate cross sections all peak where $`\omega ^{}=\omega (\mathrm{at}\theta =\theta _o)`$. This is a result of the fact that the approximate angular distribution is governed by the kinematics. The scattered angle, $`\theta ^{}`$, is small at $`\theta _o`$, and the term $`2(\omega \mathrm{}B)\mathrm{sin}^2\theta ^{}`$ in the denominator of equation (15) is also small resulting in an expression very similar to the Klein-Nishina kinematics. At scattering angles below $`\theta _o`$, there is significant agreement between the exact and Klein-Nishina angular distributions. The small-angle, low recoil scatterings, where one would expect all to agree because of similar kinematics, does not probe the effects of the field.
Differences are seen at large angle, large recoil scattering, where the geometry of the magnetic scattering is impacted by Landau excitations. High above the resonance, magnetic effects become manifested in the backward direction. Comparisons with the Klein-Nishina angular distributions for large photon energies, suggest that at backward angles the Klein-Nishina cross section over estimates the exact, summed, angular distributions beyond an angle of about 30 degrees. At these backward angles, Landau states larger than zero contribute significantly. While the term $`2(\omega \mathrm{}B)\mathrm{sin}^2\theta ^{}`$ might be small at these angles, the $`\mathrm{}B`$ term in the numerator of $`\omega ^{}`$ in equation (15) becomes more important for larger $`\mathrm{}`$’s, and $`\omega ^{}`$ is significantly less than $`\omega `$, while $`\omega ^{}`$ in the Klein-Nishina kinematics remains close to $`\omega `$.
A quantity of interest in polar cap cascade models of highly magnetized gamma-ray pulsars is the mean value of the product $`\omega ^{}\mathrm{sin}\theta ^{}`$ achieved in resonant Compton scatterings. This product is a Lorentz invariant in transformations along the field lines, and represents the photon energy in the frame of reference where the upscattered photons move orthogonally to the local field. Hence, in conjunction with the value of $`B`$, this product principally controls the strength of strong-field photon attenuation processes such as pair creation $`\gamma e^+e^{}`$ and photon splitting $`\gamma \gamma \gamma `$ . The mean value of $`\omega ^{}\mathrm{sin}\theta ^{}`$ in scatterings is therefore extremely informative for pulsar cascade modelers, and accordingly is plotted as a function of incident photon energy for different $`B`$ in Figure 7. The average was formed by weighting the differential cross sections such as those in Figure 5 with $`\omega ^{}\mathrm{sin}\theta ^{}`$ using equation (15), summing over quantum numbers $`\mathrm{}`$ and integrating over $`\theta ^{}`$, and then dividing by the total cross section (see Figure 2). The resulting curves exhibit a generally increasing function of $`\omega `$, with structure at the cyclotron resonances that becomes prominent in critical and supercritical fields due to the enhanced importance of $`\mathrm{}>0`$ (excited final state) scatterings above the fundamental; for a given $`\theta ^{}`$, higher $`\mathrm{}`$ values imply lower final photon energies $`\omega ^{}`$ (see equation ). The most salient property of Figure 7 is that the criterion for the scattered photons to generally rise above pair threshold is largely insensitive to the value of $`B`$, and is contingent upon the initial photon energy exceeding about 5-10 MeV in the electron rest frame.
Other general properties of $`\omega ^{}\mathrm{sin}\theta ^{}`$ in Figure 7 can be understood as follows. Obviously there is naturally no expectation that the behavior of Figure 7 at low $`B`$ should mimic the non-magnetic scattering average at or below the cyclotron fundamental, because the total cross section does not approach the field-free Thomson limit when $`\omega B`$ (see Figure 2). In fact, contrary to such intuition, at energies well below the cyclotron fundamental, $`\omega ^{}\mathrm{sin}\theta ^{}`$ asymptotically approaches $`15\pi \omega /64`$ independent of $`B`$, derivable from the first line of equation (LABEL:eq:dsig\_leq0\_approx), a limit identical to the non-magnetic Thomson average. This ensues since, while the magnetic cross section has a different magnitude from the Thomson one, it possesses the same angular dependence (e.g. see equation (7.1b) of Rybicki & Lightman 1979), and $`\omega ^{}\omega `$ in this limit, independent of $`B`$. Departures from this low energy asymptote arise when $`B1`$. The analysis of the $`\omega B`$ case is less trivial. Since the Klein-Nishina cross section is reproduced in sub-critical fields (e.g. see Figure 2), it might be anticipated that the Klein-Nishina $`\omega ^{}\mathrm{sin}\theta ^{}_{\mathrm{KN}}`$ might result. This is realized, more or less, with the $`B=0.1`$ curve in Figure 7, which, if extrapolated, asymptotically approximates the non-magnetic Klein-Nishina average of
$$\omega ^{}\mathrm{sin}\theta ^{}_{\mathrm{KN}}\underset{\omega B}{}\frac{9\pi }{8}\frac{\sqrt{2\omega }}{2\mathrm{ln}(2\omega )+1}$$
(36)
when $`\omega 10^3`$. For higher $`B`$, deviations from pure Klein-Nishina behavior are more apparent, though the approximately $`\omega ^{1/2}`$ dependence of $`\omega ^{}\mathrm{sin}\theta ^{}`$ is generally maintained. The peak contribution to the cross section arises at scattered angles $`\theta ^{}1/\omega ^{1/2}`$ when $`\omega 1`$, as can be seen from equation (33), a consequence of kinematic constraints imposed by equation (15). In fact, this dependence of $`\theta ^{}`$ is similar to that for non-magnetic Klein-Nishina scattering, which possess similar (though not identical) kinematic restrictions. Furthermore, $`\theta ^{}`$ is proportional to $`\omega `$ in this regime (for example, see equations (33) and (35)), a dependence borne out by the Klein-Nishina cross section, though the constants of proportionality differ, and indeed are a weakly increasing function of $`B`$ in the magnetic case here. Hence, in summary, the quantum kinematics of magnetic Compton scattering are qualitatively similar to those of Klein-Nishina scattering, and control the behavior of $`\omega ^{}\mathrm{sin}\theta ^{}`$ when $`\omega B`$ and concomitant deviations from the non-magnetic case.
## 8 DISCUSSION
In this paper, we have extended the work of DH86 by exploring the regime of super-critical fields and by providing simplified and explicitly real expressions for the exact differential cross section for Compton scattering in the presence of strong magnetic fields. We have derived simple analytic approximations for both the differential and total cross sections, in the special case of scattering by highly relativistic electrons, important to a variety of astrophysical sources. These results will be very useful in studying the effects of Compton scattering in the ultra-strong magnetic fields believed to be present near stellar surfaces of SGRs and AXPs. They also provide much more accurate expressions for modeling Compton scattering in the fields, $`B>0.1`$, of many pulsars. From the comparison of the exact, angle-integrated cross section with the limiting cases of the non-relativistic Thomson and the Klein-Nishina cross sections (used in neutron star applications throughout the literature) depicted in Figure 2, we can draw the following conclusions about scattering in increasing high fields: (i) below the resonance, the exact cross section is depressed below the magnetic Thomson cross section (when $`\omega m_ec^2`$) differing by an order of magnitude or more for fields $`B>10`$; (ii) at the resonance, the exact cross section is dramatically reduced below the Thomson cross section, but the width of the resonance increases; (iii) far above the resonance, the exact cross section approaches the Klein-Nishina cross section, as expected for large photon energies, with the energy where the two merge an increasing function of $`B`$. The overall effect of the strong fields is to lower the cross section at all incident photon energies $`\omega m_ec^2`$, decreasing the electron scattering opacity.
Focusing on large photon energies above the resonance, the magnetic field has a smaller perturbation on the interaction than around the resonance and below, and the exact cross section tends toward the Klein-Nishina limit, a feature that is seen in Figure 2 for $`B10`$. However, the analytical demonstration that these two cross sections approach each other in this high energy domain requires an approximation to the sum over many Landau states of the scattered electron and will be treated in a later paper. By inspection of the angular distributions in Figure 6, one can see that the disagreement between the exact cross section and the Klein-Nishina cross section occurs for moderate to extreme backward scattering, where the interaction probes the dominant effects of the magnetic field. It is also important to note that as the photon energy increases, the number of final Landau states contributing to the overall cross section increases dramatically, as seen in Figure 4, and as expected from the $`\mathrm{}<\omega /B`$ kinematic constraint. Care must be exercised in adding the contributions, as the numerical value of the cross section becomes very small for large Landau states.
From Figure 2 it is clear that even in relatively low fields $`(B0.1)`$, neither the magnetic Thompson nor the Klein-Nishina cross sections provides an adequate approximation to the exact cross section in a region right above the resonance, making a better representation of the exact cross section important for astrophysical models. Since the $`\mathrm{}=0`$ scattering provides the sole contribution to the cross section below the resonance, we were motivated to use it as a basis for developing an analytic approximation that is strictly only valid for $`\omega 1`$ and $`B1`$. However, satisfyingly, the analytic approximation to the $`\mathrm{}=0`$ angle-integrated cross section can be extrapolated to $`\omega 1`$ domains, and is able to represent the exact cross section quite well even at super-critical fields $`(B10)`$. This approximation indeed provides a smooth match to the exact cross section below the resonance, through the resonance and then above it where the Klein-Nishina-like reductions take over.
Use of the more accurate numerical results or approximate expressions that we have given here for resonant scattering will generally diminish the effects attributed to resonant Compton scattering in astrophysical models that use a non-relativistic treatment. This is because the non-relativistic cross section over-estimates the exact cross section when extrapolated to the domains (i.e. $`\omega 1,B1`$) where relativistic effects are important or critical. Thus, we expect that the conditions for scattering to be significant in polar cap acceleration models will be more restricted: i.e., somewhat higher magnetic field strengths and soft photon densities will be required, than previously asserted, for Compton scattering to dominate over curvature radiation in the energy loss and pair production by primary particles (Sturner 1995, Harding & Muslimov 1998).
The cross section for Compton scattering is highly dependent on the photon polarization. Incident photons with either parallel or perpendicular polarization may switch polarization modes in a scattering. The calculations of this paper, as seen in Figure 3, show that the polarization-switching properties below the resonance are the same for the non-relativistic $`(B1)`$ and relativistic cases (i.e. more photons are scattered into the perpendicular mode). However for $`B1`$, the polarization-mode switching reverses above the resonance, so that scattering to the parallel mode is dominant. Such behavior contrasts that of the non-relativistic limit, where photons of perpendicular polarization are predominantly produced at all energies. In the case of the Klein-Nishina scattering where there is no magnetic field, the initial polarizations do matter in determining the final polarization resulting in three polarization channels, described earlier. Here too, the dominance of photons scattered to perpendicular polarization is manifested in the relativistic, non-magnetic Klein-Nishina case.
The polarization dependence of the scattering cross section in high fields will have significant implications for other polarization-dependent mechanisms such as pair production and, especially, photon splitting. As Baring & Harding (1998) noted, photon splitting could dominate pair production at supra-critical magnetic fields, thereby suppressing pair creation and possibly accounting for the radio quiescence of SGRs and AXPs. However, kinematic selection rules (Adler 1971; Shabad 1975) allow only one splitting mode to operate in the limit of weak dispersion, that in which photons with perpendicular polarization $`()`$ split into two photons, each with parallel polarization $`()`$. Under such restrictions, photon splitting would occur only once, and then pair production would take over as the dominant attenuating mechanism. It is possible that the dispersion characteristics of the ultra-strong field environment, or perhaps plasma properties present during outburst mode of SGRs, may permit the two other splitting modes allowed by CP (charge-parity) invariance to operate, providing parallel-mode photons the opportunity to split. Nevertheless, even if other modes do not become operational in high fields, Compton scattering below the resonance is able to convert the photons with parallel polarization into perpendicular polarization, refueling photon-splitting cascades. Optical depths for such scattering could be quite significant in SGRs during their high luminosity gamma-ray outbursts, provided that the photon field does not dominate the SGR energetics.
While the magnitude of the resonant cross section declines with increasing $`B`$, its width increases. For sub-critical fields, the two trends compensate each other to produce an area under the resonance curve (i.e. in Figure 2) insensitive to $`B`$. For supercritical fields, the area scales as $`B^{1/2}`$, as we have noted in Section 4. This area is an approximate measure of the importance of resonant Compton scattering for a particular situation. Hence it follows that, for a given soft photon field and electron population, resonant Compton scattering becomes more significant in magnetar-type fields. However, it also becomes more difficult to have photons with energies near the resonance.
Another quantity of interest to astrophysical modelers is the expectation of $`\omega ^{}\mathrm{sin}\theta ^{}`$ in scatterings. This is because this quantity, a Lorentz invariant in transformations along the field, represents the major controlling parameter (apart from $`B`$) for determining the rates of photon absorption processes such as pair creation and photon splitting in strong magnetic fields. When $`B5\times 10^{12}G`$, resonant Compton upscattering (i.e. the magnetic inverse Compton process) can be a major contributor to the gamma-ray emission of pulsars. Whether the Compton-upscattered photons produced by primary electrons can generate pairs, and therefore initiate pair cascades with steeper synchrotron radiation components, is contingent upon $`<\omega ^{}\mathrm{sin}\theta ^{}>`$ exceeding $`2m_ec^2`$. Figure 7 reveals that this criterion is roughly independent of $`B`$ for $`0.1<B<100`$, and that the necessary condition to spawn cascading is $`\omega 510MeV`$ in the electron rest frame. This translates into a particular electron energy and soft photon energy and angle with respect to $`B`$ in the observer’s frame that is readily identifiable for models. If the soft photons are quasi-isotropic thermal X-rays from the surface, then the primary electron Lorentz factors need to exceed $`\gamma 10^5(T/10^6K)^1`$ in order to satisfy this pair production criterion, independent of $`B`$. For attenuation by photon splitting, when $`B0.3`$ G, the energy at which splitting optical depths exceed unity can be below pair creation threshold by a decade or more (e.g. Baring 1991), so that the required value of $`<\omega ^{}\mathrm{sin}\theta ^{}>`$ can be much lower. Hence the target photon energies required for the resonant process to seed the splitting mechanism need only be around $`\omega 10keV1MeV`$ for $`B0.3`$ G (and approximately inversely proportional to $`B`$ in this instance), a much less stringent requirement than that for pair cascade initiation.
In conclusion, this paper has provided computations of resonant Compton scattering in a broad range of sub-critical and super-critical fields in the particular (but widely applicable) case of ultra-relativistic electrons moving along a uniform magnetic field. In doing so, we have simplified the exact QED differential cross section obtained in this specialization, and derived a compact analytic expression that approximates the cross section quite well at energies both below and above the resonance at the cyclotron fundamental, for fields $`B10`$. Such an approximation should prove extremely useful to astrophysical modelers interested in highly magnetized (normal and anomalous X-ray) pulsars and soft gamma repeaters. Significant deviations from the differential Klein-Nishina cross section were found for large scattering angles, though the total magnetic cross section was observed to approach the classic Klein-Nishina behavior at energies well above the resonance. Polarization properties of resonant scattering were also explored in detail, revealing that, as with the magnetic Thomson case, the scattered photons are predominantly of perpendicular polarization below the resonance. While this property persists above the resonance for sub-critical fields, a polarization-reversal arises in super-critical fields at $`\omega >B`$ so that parallel photons dominate the scattered photon population. Comprehension of such properties may be critical to model predictions of the emission from magnetars.
###### Acknowledgements.
We thank Ramin Sina for the use of his computer code that numerically calculates exact QED Compton scattering with Sokolov & Ternov electron spin states. We would like to express our sincere appreciation for the generous support of NASA under the Summer Faculty Fellowship Program, of the Michigan Space Grant Consortium, of the Research Corporation, and of the NSF under the REU program and through the grant NSF-9876670.
## Appendix A Appendix
Here the $`F_n^{(j)}`$ terms that contribute to the scattering amplitudes that appear in the general expression for differential cross section in Eq. (1) are listed, having been derived in the literature before (e.g. see DH86). Each of the $`F`$ terms consists of two parts due to the spin-degeneracy (above the ground state) of the final Landau states of the electron. For the first diagram in Figure 1, the electron spin no-flip and spin flip forms are given by
$`F_{n,noflip}^{(1)}={\displaystyle \frac{1}{(2\omega +\omega ^2\mathrm{sin}^2\theta 2nB)}}`$ (A6)
$`\times \left\{\begin{array}{c}\left[\omega (2+\omega \omega ^{})+\omega \mathrm{cos}\theta (\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{},n}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _z^{}\hfill \\ +\left[\omega (2+\omega \omega ^{})\omega \mathrm{cos}\theta (\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{},n1}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _+^{}\hfill \\ +\sqrt{2nB}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\left[\mathrm{\Lambda }_{\mathrm{},n}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _z^{}+\mathrm{\Lambda }_{\mathrm{},n1}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _+^{}\right]\hfill \\ +\sqrt{2\mathrm{}B}\left\{\begin{array}{c}\omega \mathrm{cos}\theta \left[\mathrm{\Lambda }_{\mathrm{}1,n}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _{}^{}+\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _z^{}\right]\hfill \\ +\sqrt{2nB}\left[\mathrm{\Lambda }_{\mathrm{}1,n}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _{}^{}\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _z^{}\right]\hfill \end{array}\right\}\hfill \end{array}\right\}`$
$`F_{n,flip}^{(1)}={\displaystyle \frac{1}{(2\omega +\omega ^2\mathrm{sin}^2\theta 2nB)}}`$ (A13)
$`\times \left\{\begin{array}{c}\left[\omega (2+\omega \omega ^{})+\omega \mathrm{cos}\theta (\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _z^{}\hfill \\ \left[\omega (2+\omega \omega ^{})\omega \mathrm{cos}\theta (\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{}1,n}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _{}^{}\hfill \\ +\sqrt{2nB}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\left[\mathrm{\Lambda }_{\mathrm{}1,n}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _{}^{}\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _z^{}\right]\hfill \\ \sqrt{2\mathrm{}B}\left\{\begin{array}{c}\omega \mathrm{cos}\theta \left[\mathrm{\Lambda }_{\mathrm{},n}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _z^{}\mathrm{\Lambda }_{\mathrm{},n1}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _+^{}\right]\hfill \\ +\sqrt{2nB}\left[\mathrm{\Lambda }_{\mathrm{},n}(\beta ^{})\mathrm{\Lambda }_{n1,0}(\beta )\epsilon _{}\epsilon _z^{}+\mathrm{\Lambda }_{\mathrm{},n1}(\beta ^{})\mathrm{\Lambda }_{n,0}(\beta )\epsilon _z\epsilon _+^{}\right]\hfill \end{array}\right\}\hfill \end{array}\right\},`$
where the $`\epsilon `$ represents the photon polarization components defining two orthogonal linear polarization vectors, as defined in Eq. (LABEL:eq:polarizdef), and the functions $`\mathrm{\Lambda }_{\mathrm{},m}(\beta )`$ are defined in Eq. (5) in terms of associated Laguerre polynomials. The $`F`$ terms for the second (exchange) Feynman diagram in Figure 1 are
$`F_{n,noflip}^{(2)}={\displaystyle \frac{1}{(\omega ^2\mathrm{sin}^2\theta ^{}2\omega ^{}2nB)}}`$ (A19)
$`\times \left\{\begin{array}{c}\left[\omega ^{}(2+\omega \omega ^{})+\omega ^{}\mathrm{cos}\theta ^{}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{},n}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _z\epsilon _z^{}\hfill \\ \left[\omega ^{}(2+\omega \omega ^{})\omega ^{}\mathrm{cos}\theta ^{}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{},n1}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _+\epsilon _{}^{}\hfill \\ +\sqrt{2nB}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\left[\mathrm{\Lambda }_{\mathrm{},n}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _z\epsilon _{}^{}+\mathrm{\Lambda }_{\mathrm{},n1}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _+\epsilon _z^{}\right]\hfill \\ +\sqrt{2\mathrm{}B}\left\{\begin{array}{c}\omega ^{}\mathrm{cos}\theta ^{}\left[\mathrm{\Lambda }_{\mathrm{}1,n}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _{}\epsilon _z^{}+\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _z\epsilon _{}^{}\right]\hfill \\ +\sqrt{2nB}\left[\mathrm{\Lambda }_{\mathrm{}1,n}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _{}\epsilon _{}^{}\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _z\epsilon _z^{}\right]\hfill \end{array}\right\}\hfill \end{array}\right\},`$
$`F_{n,flip}^{(2)}={\displaystyle \frac{1}{(\omega ^2\mathrm{sin}^2\theta ^{}2\omega ^{}2nB)}}`$ (A26)
$`\times \left\{\begin{array}{c}\left[\omega ^{}(2+\omega \omega ^{})+\omega ^{}\mathrm{cos}\theta ^{}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _z\epsilon _{}^{}\hfill \\ +\left[\omega ^{}(2+\omega \omega ^{})\omega ^{}\mathrm{cos}\theta ^{}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\right]\mathrm{\Lambda }_{\mathrm{}1,n}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _{}\epsilon _z^{}\hfill \\ +\sqrt{2nB}(\omega \mathrm{cos}\theta \omega ^{}\mathrm{cos}\theta ^{})\left[\mathrm{\Lambda }_{\mathrm{}1,n}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _{}\epsilon _{}^{}\mathrm{\Lambda }_{\mathrm{}1,n1}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _z\epsilon _z^{}\right]\hfill \\ \sqrt{2\mathrm{}B}\left\{\begin{array}{c}\omega ^{}\mathrm{cos}\theta ^{}\left[\mathrm{\Lambda }_{\mathrm{},n}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _z\epsilon _z^{}\mathrm{\Lambda }_{\mathrm{},n1}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _+\epsilon _{}^{}\right]\hfill \\ +\sqrt{2nB}\left[\mathrm{\Lambda }_{\mathrm{},n}(\beta )\mathrm{\Lambda }_{n1,0}(\beta ^{})\epsilon _z\epsilon _{}^{}+\mathrm{\Lambda }_{\mathrm{},n1}(\beta )\mathrm{\Lambda }_{n,0}(\beta ^{})\epsilon _+\epsilon _z^{}\right]\hfill \end{array}\right\}\hfill \end{array}\right\},`$
and can be obtained from Eq. (LABEL:eq:F1def) via the crossing symmetry in Eq. (3).
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# Introducing ZEUS-MP: A 3D, Parallel, Multiphysics Code for Astrophysical Fluid Dynamics
## 1 A Brief History of ZEUS
ZEUS is a family of codes for astrophysical fluid dynamics simulations developed at the Laboratory for Computational Astrophysics (LCA) of the National Center for Supercomputing Applications (NCSA) at the University of Illinois, Urbana-Champaign. The purpose of this paper is to announce the availability of the latest implementation: ZEUS-MP. Version 1.0 released 1/1/2000 is available from the LCA’s website lca.ncsa.uiuc.edu.
ZEUS has its roots in a 2D Eulerian hydro code developed by M. Norman for simulations of rotating protostellar collapse (?) while he was a student at the Lawrence Livermore National Laboratory. The hydrodynamics algorithm, which has changed little in subsequent versions, is based on a simple staggered-grid finite-difference scheme (??). Shock waves were captured within a few cells with a von Neumann-Richtmyer type artificial viscosity. A powerful and essential feature of the code, which has been retained in subsequent versions, was that the equations of self-gravitating hydrodynamics were solved on a moving Eulerian grid permitting accurate simulation over a range of scales in the collapsing protostar.
A significant improvement to the hydrodynamics algorithm came with the incorporation of the second order-accurate, monotonic advection scheme (?). This code version, innocuously called A2, was vectorized for the Cray-1 supercomputer at the Max-Planck-Institut für Astrophysik, and extensively applied to the simulation of extragalactic radio jets (?).
The first code called ZEUS was developed by David Clarke as a part of his PhD thesis on MHD jets (??) under Norman’s supervision. One of the principal challenges in numerical MHD simulations is satisfying the zero-divergence constraint on $`B`$. In our axisymmetric simulations, this was insured by evolving the toroidal component of the magnetic vector potential from which divergence-free poloidal field components can be derived, as well as evolving the toroidal magnetic field component directly. Third order-accurate monotonic advection was used for evolving $`A_\varphi `$ in order to improve the quality of the derived current densities.
The next development was a major rewrite and significant extension of ZEUS by James Stone as a part of his PhD thesis at the University of Illinois. The resulting code, named ZEUS-2D, solves the equations of self-gravitating radiation magnetohydrodynamics in 2D or 2-1/2D. Many new algorithms were developed and incorporated into ZEUS-2D including: (1) a covariant formulation, allowing simulations in various coordinate geometries; (2) a tensor artificial viscosity; (3) a new, more accurate MHD algorithm (MOC-CT) combining the Constrained Transport algorithm (?) with a Method Of Characteristics treatment for Alfvén waves; and (4) a variable tensor Eddington factor solution for the equations of radiation hydrodynamics. ZEUS-2D’s algorithms and tests are described in detail in a series of three papers (???)(the ZEUS Trilogy).
The MOC-CT algorithm for numerical MHD was specifically designed to be extensible to 3D, and work on a 3D version of ZEUS began in 1989 when David Clarke came to Illinois as Norman’s postdoc. Written for the Cray-2 supercomputer, ZEUS-3D physics options included hydrodynamics, MHD, self-gravity, and optically thin radiative cooling. Parallelization was done using Cray Autotasking compiler directives. Novel features of the code included the use of a custom source code pre-processor which handled a variety of source code transformations. Another useful feature of ZEUS-3D was an extensive set of inline graphics and diagnostic routines, as well as the ability to run in 1D and 2D mode.
With a grant from the National Science Foundation in 1992, the LCA was established with the purpose of disseminating ZEUS-2D, ZEUS-3D and the TITAN implicit adaptive-mesh radiation hydrodynamics code (?) to the international community. Currently, there are over 500 registered users of LCA codes in over 30 countries. Some recent applications of ZEUS include planetary nebulae (?), molecular cloud turbulence (?), and solar magnetic arcades (?).
Work was begun on ZEUS-MP in the fall of 1996 by Robert Fiedler and subsequently by John Hayes and James Bordner with support from the Department of Energy to explore algorithms for parallel radiation hydrodynamics simulations in 3D.
## 2 Why ZEUS-MP?
ZEUS-MP is a portable, parallel rewrite of ZEUS-3D. MP stands for: Multi-Physics, Massively-Parallel, and Message-Passing. 3D simulations are by their nature memory– and compute–intensive. The most powerful computers available today are parallel computers with hundreds to thousands of processors connected into a cluster. While some systems offer a shared memory view to the applications programmer, others, such as Beowulf clusters, do not. Thus, for portability sake we have assumed “shared nothing” and implemented ZEUS-MP as a SPMD (Single Program, Multiple Data) parallel code using the MPI message-passing library to affect interprocessor communication.
Figure 1 shows a block diagram of the major components of ZEUS-MP. ZEUS-MP is composed of an application layer (second row) and a libraries layer (third row), all resting on the MPI message passing library. A brief description follows.
There are four main physics modules: hydro, MHD, radiation transport, and self-gravity. The hydro and MHD modules are time-explicit, and thus require no linear algebra libraries for their solution. The hydro algorithm is a straight-forward 3D extension of the algorithm described in (?). The MHD algorithm is the MOC-CT algorithm described in (?) with the modifications described in (?) for enhanced stability in weakly magnetized, strongly sheared flows. At present, nonideal effects (viscosity, resistivity) are not included, however one can choose between an ideal gamma-law or isothermal equations of state. The radiation transport algorithm implements the time-implicit flux-limited diffusion algorithm of Stone (1994). The radiation and gas energy equations are solved as a coupled, implicit system, resulting in a large, sparse, banded system of linear equations which must be solved within an outer nonlinear Newton iteration. Two linear system solvers are built into ZEUS-MP: a conjugate gradient solver (CG/BiCG) with diagonal preconditioning, and a multigrid solver (MGMPI) (?). Problems involving self-gravity require the solution of the Poisson eqation. Two Poisson solvers are built into ZEUS-MP: a Fourier space solver for triply periodic cubic grids, and an elliptic finite difference solver for all other geometries and boundary conditions. The former utilizes the FFTw library developed at MIT, while the latter uses MGMPI.
As in earlier versions of ZEUS, the equations solved by ZEUS-MP are formulated on a covariant, moving Eulerian grid. Problems in Cartesian, cylindrical, and spherical polar coordinates can be run with a variety of boundary conditions and types (periodic, Dirichlet, Neumann). The linear system solvers are designed to handle all cases.
File I/O is done using NCSA’s HDF (Hierarchical Data File) standard (?), which is a widely adopted portable file format for scientific data.
## 3 Parallelism
ZEUS-MP utilizes domain decomposition (?) for parallelization, wherein the computational domain is subdivided into a number of equally-sized regions, each of which is assigned to a different processor for execution. Depending upon the problem size and the number of processors targeted, the user can specify a 1D “slab”, 2D “pencil”, or 3D “block” decomposition. A region is represented in processor memory as arrays of data storing the solution vector for a specific subdomain. The arrays are dimensioned so as to include two layers of buffer zones on each face of the block for the purpose of transferring boundary conditions from neighboring processors. Data transfer between neighboring blocks, as well as collective operations and global reductions are handled via MPI function calls.
ZEUS-MP performance has been optimized in two ways (?): (1) single processor performance, using a variety of standard cache optimization techniques; and (2) parallel performance, using asynchronous communication, wherein computation and communication is overlapped. Single-node performance on an MIPS R10000 processor is in the range of 100 MFlop/s. Parallel speedup results depend sensitively on the properties of the network hardware and software on the host computer. A collection of benchmark results can be found on the ZEUS-MP website zeus.ncsa.uiuc.edu/lca\_intro\_zeusmp.html. Figure 2 shows near-ideal scaling on an SGI/Cray Origin2000 for hydrodynamic and MHD tests of size $`256^3`$ and $`512^3`$ on up to 256 processors.
## 4 Source Code Availability
V1.0 is available now as a downloadable tar file from the ZEUS-MP website. Rudimentary documentation describing how to make the executable and run any of the five test problems is also available online. At present, the radiation module is not included; pending clarification of export restrictions, it will be made available in a later release.
Also available is a new 3D visualization tool called LCA Vision zeus.ncsa.uiuc.edu/$``$miksa/LCAVision1.0.html. Vision reads HDF files and provides a variety of visualization tools in an easy-to-use, menu-driven interface.
## 5 Sample Application
To illustrate ZEUS-MP’s capabilities, we present an application to magnetic star formation. The purpose is to understand the competition between turbulent, magnetic, and gravitational stresses in the formation of gravitationally bound cores in a turbulent molecular cloud (?). In particular, we want to explore whether the critical mass-to-flux ratio (?) is a good predictor of gravitational collapse in a cloud driven by turbulence.
The calculation is done in a triply periodic cube with $`256^3`$ cells distributed across 64 processors ($`4^3`$ cubic decomposition.) One begins with a uniform density gas filling the box threaded by a uniform magnetic field in the Z-direction. The gas is assumed to be isothermal. A turbulent velocity field is established by driving in Fourier space over a limited range of wavenumbers as described in (?). A statistical steady state is reached within a few dynamical times, and once it has, gravity is switched on. Thereafter, density peaks formed by colliding gas streams may become gravitationally bound and collapse. The problem is described by a few dimensionless parameters formed by combinations of the gas density, sound speed, magnetic field strength, rms velocity perturbation, and box length,: $`n_J=12`$, the number of thermal Jeans’ masses in the box; $`M/M_{cr}=1.1`$, the ratio of the box mass to the critical mass that can be supported by a static magnetic field; $`M_{turb}=5`$, the turbulent Mach number; and $`k_{drv}`$=2, the driving wavenumber in units of the inverse box length.
Fig. 3 shows isosurfaces of gas density at 1.82 freefall times after gravity was swithched on. The isolevel is eight times the initial uniform density. One can see evidence for filamentary and flattened condensations which are both aligned and perpendicular to the mean field direction. Gravitational collapse is already underway, as indicated by high densities in the centers of several filaments. We find, as did (?) using ZEUS-3D, that supersonic turbulence is not able to prevent gravitational collapse from occurring in magnetically supercritical clouds. Simulations at $`512^3`$ grid resolution and beyond are underway to investigate this further.
###### Acknowledgements.
I gratefully acknowledge my colleagues, past and present, who have materially contributed to ZEUS-MP: Jim Stone and David Clarke for their foundational contributions embodied in ZEUS-2D and ZEUS-3D; Robert Fiedler, who wrote the first hydrodynamics version of ZEUS-MP code and optimized it for cache-based parallel systems; John “radiation cowboy” Hayes, who implemented the implicit radiation diffusion module along with the CG/BiCG linear solver; Mordecai-Mark MacLow, who ported the HSMOC MHD algorithm from ZEUS-3D, which in turn was incorporated by Byung-Il Jun; Pakshing Li, who developed and tested both Poisson solvers and prepared the V1.0 release; and James Bordner, developer of the MGMPI multigrid solver and patient fixer of what we break. I thank my collaborators Fabian Heitsch, Mordecai MacLow and Pakshing Li for allowing me to show our unpublished results (Fig. 3). This work has been supported by contracts B324163 and B506131 from the Lawrence Livermore National Laboratory. I would like to thank Frank Graziani of B Division for his continued interest and support.
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# 1 Introduction
## 1 Introduction
The nonlinear sigma model (NLSM) with a Skyrme term has received a great deal of attention after the proposal by Adkins et al that this is a theory for weakly interacting mesons in the chiral limit resulting from the more fundamental theory for strong interactions, QCD, in the limit when the number of colors $`N_c`$ is taken very large. Using the collective coordinate method for the isospin rotation, the authors of have performed the semi-classical quantization of the model with the spectrum of static properties being accurate to within $`30\%`$ of the experimental data, an extraordinary result for such simple model. This model, however, presents some well known shortcomings. There is no compelling physical reason to include only the Skyrme term and not higher derivatives potentials. Also the quantum picture is known to be troubled by operator ordering ambiguities and the results following different approaches are not in complete agreement.
Technically the procedure in has reduced the field theory problem to that of a free point particle on a $`S^3`$ sphere. This theory is an example of a nonlinearly constrained second-class system. The quantization of such systems has been intensively studied in various contexts and the relevance of this problem for the quantization on curvilinear surfaces is well appreciate, both in the path integral and in the canonical approach. There has been different proposals in the sense of improving the physical description of the model. A new possible route, advocated in the present paper, would be to deform away from the intrinsic spherical symmetry. However there is no conclusive studies of nonlinear constraints of arbitrary geometries available in the literature. This gap makes it harder to implement this program since the exact shape of the constraint surface able to improve accuracy must come out from a numerical adjustment and is expected to be level dependent. All this points up to the necessity of a detailed study of quantization over arbitrary constraint surfaces.
It is the purpose of this work to illuminate the quantization of Skyrme’s collective rotational mode resulting from quadratic surfaces deviated from the spherical geometry. While most investigations are done towards understanding the quantum nature directly from the 2nd-class formulation, we reformulate the model as a gauge theory. We provide the gauge invariant reformulation of Skyrme’s model for a general nonlinearly constrained second-class surface. This is done through the iterative constraint conversion scheme developed by us a few years ago.
The technique that converses second-class constraints to first-class extending the original phase space with some auxiliary variables, was suggested by Faddeev and Shatashivilli with the addition of Wess-Zumino terms (WZ) on the original Hamiltonian. As far as we know the treatment of nonlinear systems as gauge theories appears already in Balachandran et al in and in a number of papers by Yamawaki et al.. Our presentation here is inspired by the idea of Kovner and Rosenstein (KR), using an analogy with QED to disclose a symmetry hidden in the NLSM. In we interpreted the KR symmetry as the gauge symmetry of a Wess-Zumino extended theory, with the geometrical second-class constraints converted to first-class. The hidden symmetry was shown to be a residual symmetry of the WZ orbits. The formalism developed in was used recently in to study nonlinear models restricted by a more general class of constraints. It provides a clear-cut geometrical interpretation for the Wess-Zumino gauge symmetry in the extended space as well as in the original configuration space.
In Section 2 we consider, with detailed care, the four dimensional SU(2) nonlinear sigma model with an stabilizing Skyrme term and review the semi-classical expansion of the collective rotational mode. Reduction to a nonlinear quantum mechanical model depending explicitly on the time-dependent collective variables satisfying a spherical constraint is performed. This will set the stage to allow the Skyrme’s model constraint to be deformed and then turned into a gauge theory in the extended WZ phase space. A systematical treatment of the constraint conversion is developed in the remaining of the paper and is distributed as follows. In section 3 the noninvariant aspect of the theory is reviewed and our notation is introduced. We consider a general setting and study the motion of a point particle in a N-dimensional Euclidean space moving freely on a nonlinear surface $`\mathrm{\Omega }(q)`$ embedded in the $`\mathrm{}^N`$. Special emphasis is given to the symmetric generating matrix that embraces the geometric information. In section 4 the second-class non-spherical model discussed in section 3 is transformed to a simple gauge invariant theory, written as the sum of the original plus the WZ term. To allow for a concrete calculation we adopt a perturbative point of view and consider small deformations away from the spherical symmetry. This introduces new parameters into the theory allowing for a better fit with experimental data. The last Section is reserved to discuss the physical meaning of our findings together with our final comments and conclusions.
## 2 The Skyrme Model Revisited
A few decades ago Skyrme proposed to describe baryons as topological solutions of the NLSM with an appropriate stabilizing term. The semi-classical quantization of the model was obtaining in separating the collective coordinate. Let us consider the SU(2) Skyrmion Lagrangian
$$L^{}=d^3x\left\{\frac{f_\pi ^2}{4}Tr\left(_\mu U^{}^\mu U\right)+\frac{1}{32e^2}Tr[U^{}^\mu U,U^{}^\nu U]^2\right\}$$
(1)
where $`f_\pi `$ is the pion decay constant and $`e`$ is a dimensionless parameter. $`U`$ is a $`SU(2)`$ matrix transforming as $`UAUB^1`$ under chiral $`SU(2)\times SU(2)`$, satisfying the boundary condition $`\underset{r\mathrm{}}{lim}U=I`$ so that the pion field vanishes as $`r`$ goes to infinity. There are soliton solutions described by the action (1) whose topological number are identified with the baryon number. To describe the static soliton we start with the ansatz $`U(r)=\mathrm{exp}\{i\stackrel{}{\tau }_a.\widehat{x}_af(r)\}`$ where $`\stackrel{}{\tau }_a`$ are Pauli matrices, $`\widehat{x}=\stackrel{}{x}/r`$ and $`\underset{r\mathrm{}}{lim}f(r)=0`$ and $`f(0)=\pi `$. Performing the collective semi-classical expansion in (1), where $`U(r,t)=A(t)U(r)A^{}(t)`$ and $`ASU(2)`$, we obtain after performing the space integral,
$$L^{}=M+Tr\left(_0A_0A^1\right).$$
(2)
M and $``$ are the soliton mass and the moment of inertia respectively which, in the hedgehog ansatz are given by
$$M=2\pi _0^{\mathrm{}}𝑑rr^2\left[f_\pi ^2\left(\left(\frac{df}{dr}\right)^2+2\frac{\mathrm{sin}^2f}{r^2}\right)+\frac{\mathrm{sin}^2f}{e^2r^2}\left(2\left(\frac{df}{dr}\right)^2+\frac{\mathrm{sin}^2f}{r^2}\right)\right]$$
(3)
and
$$=\frac{8\pi }{3}_0^{\mathrm{}}𝑑rr^2\mathrm{sin}^2f\left[f_\pi ^2+\frac{1}{e^2}\left(\left(\frac{df}{dr}\right)^2+\frac{\mathrm{sin}^2f}{r^2}\right)\right].$$
(4)
The matrix $`A`$ may be represented by $`A=a_0+i\stackrel{}{a}.\stackrel{}{\tau }`$ satisfying the spherical constraint
$$\mathrm{\Omega }=\underset{i=0}{\overset{3}{}}a_i^21.$$
(5)
In terms of these variables the Skyrmion Lagrangian (2) becomes
$$L^{}=M+2\underset{i=0}{\overset{3}{}}\dot{a}_{i}^{}{}_{}{}^{2}+\lambda \mathrm{\Omega },$$
(6)
with the spherical constraint (5) being implemented by the Lagrange multiplier $`\lambda `$. The Hamiltonian corresponding to (6) is
$$H^{}=M+\frac{1}{8}\underset{i=0}{\overset{3}{}}\mathrm{\Pi }_i^2\lambda \mathrm{\Omega },$$
(7)
with the canonical momenta,
$$\mathrm{\Pi }_i=4\dot{a}_i.$$
(8)
For simplicity we rescale the coordinates and introduce a new Lagrangian and Hamiltonian as
$`L`$ $`=`$ $`M+L^{}={\displaystyle \frac{1}{2}}\dot{q}_{i}^{}{}_{}{}^{2}+\overline{\lambda }\mathrm{\Omega },`$
$`H`$ $`=`$ $`M+H^{}={\displaystyle \frac{1}{2}}p_i^2\overline{\lambda }\mathrm{\Omega },`$ (9)
where
$`p_i`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{}}}\mathrm{\Pi }_i,`$
$`\dot{q}_i`$ $`=`$ $`2\sqrt{}\dot{a}_i,`$
$`\overline{\lambda }`$ $`=`$ $`4\lambda ,`$
$`\mathrm{\Omega }`$ $`=`$ $`q_i^2c;c=4.`$ (10)
We have reduced the Skyrmion problem to that of a non-relativistic unity mass particle constrained over a $`𝒮^3`$ sphere, a well known second-class problem. We are now ready to address the question of constraint conversion. A general construction based on the iterative formalism is presented in this paper. A particular emphasis is given to the extension of the spherical symmetry with dramatic consequences for the energy spectrum.
## 3 The Non-Spherical Second-Class Aspect
In this section we review the main features of a theory for a particle constrained to move on a 2nd-class nonlinear constraint surface. Let us consider the mechanical O(N) model which is a theory for a point particle with coordinates $`q_k(q_1,q_2,q_3,\mathrm{},q_N)`$ moving freely over a nonlinear surface,
$$\mathrm{\Omega }=\frac{1}{2}q_kT_{km}q_mc=0.$$
(11)
Here the surface generator symmetric matrix $`T_{km}`$ and the yet arbitrary constant $`c`$ characterize the surface. We are restricting our study to surfaces bilinear in the coordinates, whose exact structure is encoded in $`T_{km}`$. As so these matrices have a decisive role in determining the constraint structure of the model. This role will become clearer as we progress. The classical trajectory and its dynamics is governed by the Lagrangian,
$$=\frac{1}{2}\dot{q}^2\lambda \left(\frac{1}{2}qTqc\right),$$
(12)
with $`\lambda `$ being a Lagrange multiplier enforcing (11) as a constraint. To simplify the notation we omit the coordinate indices from now on unless to avoid confusion. The canonical analysis gives,
$$p_k=\dot{q}_k,$$
(13)
and the primary constraint
$$\varphi _1=p_\lambda 0.$$
(14)
The corresponding canonical Hamiltonian reads,
$$=\frac{1}{2}p_k^2+\lambda \left(\frac{1}{2}qTqc\right),$$
(15)
and the primary Hamiltonian is obtained from (15) enforcing the constraint (14) with a multiplier $`u`$,
$$_P=+u\varphi _1.$$
(16)
The nonlinear surface $`\mathrm{\Omega }(q)`$ reappears in the Dirac consistency chain as a secondary constraint,
$`\varphi _2`$ $`=`$ $`\{\varphi _1,_P\}`$ (17)
$`=`$ $`\mathrm{\Omega }.`$
The consistency algorithm demands the presence of a tertiary constraint,
$$\varphi _3=pTq$$
(18)
and a quaternary constraint as,
$$\varphi _4=pTp\lambda qT^2q.$$
(19)
where $`T_{km}^2=T_{kn}T_{nm}`$. Constraints $`\varphi _2`$ and $`\varphi _3`$ have a clear geometrical meaning. The secondary constraint enforces the particle coordinates $`q_k`$ to take values over the nonlinear surface. The tertiary constraint, called as transverse, enforces the particle velocity to take values over the tangent space of $`\mathrm{\Omega }(q)`$ at $`q_k`$. To disclose the meaning of $`\varphi _1`$ and $`\varphi _4`$ we compute the Poisson algebra that follows from the complete constraint set,
$`\{\varphi _1,\varphi _4\}`$ $`=`$ $`qT^2q,`$
$`\{\varphi _2,\varphi _3\}`$ $`=`$ $`qT^2q,`$
$`\{\varphi _2,\varphi _4\}`$ $`=`$ $`2qT^2p,`$ (20)
$`\{\varphi _3,\varphi _4\}`$ $`=`$ $`2pT^2p+2\lambda qT^3q.`$
There are no more constraints. From the time evolution consistency for the $`\varphi _4`$ constraint,
$$0(qT^2q)u+\{\varphi _4,\},$$
(21)
the multiplier $`u`$ is determined and the consistency chain stops. We may use $`\varphi _40`$ to determine the value of the multiplier $`\lambda `$ and then compute the Dirac brackets,
$`\{q_k,q_m\}^{}`$ $`=`$ $`0`$
$`\{q_k,p_m\}^{}`$ $`=`$ $`g_{km}`$
$`\{p_k,p_m\}^{}`$ $`=`$ $`h_{km}`$ (22)
where,
$`g_{km}`$ $`=`$ $`\delta _{km}t_k\mathrm{\Delta }^1t_m`$
$`h_{km}`$ $`=`$ $`t_k\mathrm{\Delta }^1_mt_np_nt_m\mathrm{\Delta }^1_kt_np_n.`$ (23)
Here the vector,
$$t_k=_k\mathrm{\Omega }$$
(24)
and
$$\mathrm{\Delta }=t_kt_k$$
(25)
are the basic geometric elements of the theory. The constraints $`\varphi _1`$ and $`\varphi _4`$ have no physical consequence. Their presence becomes necessary just to eliminate the multiplier sector and to give consistence to Dirac’s algorithm. In fact, an alternative approach where this sector is eliminate from the outset may be used. This completes the analysis of the 2nd-class aspect of the model.
## 4 Quantization of the Non-spherical first-class Skyrme Model
In this section we present our main result, i.e., the quantization of the deformed gauge invariant Skyrme model and the computation of the energy spectrum. The conversion of the second-class constraints to first-class by extension of the phase space was originally introduced to avoid the difficulties involving anomalies in chiral gauge theories, by removing the dynamical degree of freedom obstructing gauge symmetry. This scheme has also been used to covariantize the chiral boson constraint leading to a system with an infinite chain of 1st-class constraints. This formulation has been of some use in recent developments in the study of superstrings and dualities. The logical reasoning behind this approach is to maintain the iterative structure already present in Dirac’s formalism. In this sense, the iterative method converts the second-class constraints at their prompt appearance in the consistency chain.
To quantize the deformed Skyrme model and then compute the energy spectrum, it is necessary to know the non-spherical eigenfunctions corresponding to the Hamiltonian operator that now satisfies the new arbitrary geometry. Since it might become so complex, we consider to analyse the Skyrme model perturbatively deviated from the spherical geometry, using the spherical eigenfunctions as a zeroth-order input without any additional incoveniences. To this end we start proposing the perturbative non-spherical Skyrme model as,
$$=\frac{1}{2}\dot{q}_i^2\lambda (\frac{1}{2}q_i^2\frac{1}{2}\epsilon q_i\mathrm{\Delta }_{ij}q_jc),$$
(26)
where $`\epsilon `$ is a small parameter and the surface generator symmetric matrix $`T_{ij}`$ is assumed to be
$$T_{ij}=\delta _{ij}+\epsilon \mathrm{\Delta }_{ij},$$
(27)
with matrix $`\mathrm{\Delta }_{ij}`$ embracing the deviation from the spherical surface. The Hamiltonian computed through Legendre transformation is,
$$=\frac{1}{2}p_i^2+\lambda (\frac{1}{2}q_i^2\frac{1}{2}\epsilon q_i\mathrm{\Delta }_{ij}q_jc).$$
(28)
As in section 2, the model has four second-class constraints,
$`\varphi _1`$ $`=`$ $`p_\lambda ,`$
$`\varphi _2`$ $`=`$ $`{\displaystyle \frac{1}{2}}q_i^2{\displaystyle \frac{1}{2}}\epsilon q_i\mathrm{\Delta }_{ij}q_jc,`$
$`\varphi _3`$ $`=`$ $`q.p+\epsilon q.\mathrm{\Delta }.p,`$ (29)
$`\varphi _4`$ $`=`$ $`p^2\lambda ({\displaystyle \frac{1}{2}}q_i^2+{\displaystyle \frac{1}{2}}\epsilon q_i\mathrm{\Delta }_{ij}q_j)+\epsilon p.\mathrm{\Delta }.p\lambda (\epsilon q_i\mathrm{\Delta }_{ij}q_j+\epsilon ^2q_i\mathrm{\Delta }_{ij}q_j).`$
Recall that the first and the last constraints have no geometrical meaning. They are in fact artificial relations generated by the Dirac iterative process to preserve the original consistency presente in the dynamical phase space. Hence, these four second-class constraints could be reduced to only two, the geometrical ones ($`\varphi _2,\varphi _3`$). Computing the Lagrange multiplier $`\lambda `$ from the last constraint in (4), to first-order in $`\epsilon `$ gives,
$$\lambda =\frac{p^2}{2c}(1\frac{\epsilon }{2c}q.\mathrm{\Delta }.q+\epsilon \frac{p.\mathrm{\Delta }.p}{p^2}).$$
(30)
Bringing back this result into the Hamiltonian (28), it becomes,
$$=\frac{1}{4c}q_j^2p_i^2+\epsilon [\frac{p^2}{4c}(q.\mathrm{\Delta }.q)\frac{p^2}{8c^2}(q.\mathrm{\Delta }.q)q_j^2+\frac{p^2}{4c}(q.\mathrm{\Delta }.q)\frac{1}{2}p.\mathrm{\Delta }.p+\frac{\epsilon }{4c}(p.\mathrm{\Delta }.p)q_j^2],$$
(31)
with two second-class constraints, denoted subsequently as $`\chi _1=\varphi _2`$ and $`\chi _2=\varphi _3`$. At this stage the gauge invariant formulation of the second-class model with the introduction of the canonical WZ variables $`(\theta ,\pi _\theta )`$ starts. Firstly changing the nature of the geometrical constraint from second to first-class,
$`\chi _1\stackrel{~}{\chi }_1`$ $`=`$ $`\chi _1+\theta `$
$`\chi _2\stackrel{~}{\chi }_2`$ $`=`$ $`\chi _2\pi _\theta ,`$ (32)
so that the Poisson bracket between them becomes,
$$\{\stackrel{~}{\chi }_1,\stackrel{~}{\chi }_2\}=0.$$
(33)
After that, it is imperative to compute the new Hamiltonian that preserves the Dirac iterative process in the extended phase space,
$$\{\stackrel{~}{\chi }_1,\stackrel{~}{}\}=[\frac{q^2}{2c}+\frac{\epsilon }{2c}(q.\mathrm{\Delta }.q)]\stackrel{~}{\chi }_2.$$
(34)
Therefrom the Hamiltonian counter-term is found and, subsequently, the gauge invariant Hamiltonian,
$`\stackrel{~}{}={\displaystyle \frac{q^2p^2}{4c}}`$ $`+`$ $`\epsilon [{\displaystyle \frac{p^2}{4c}}(q.\mathrm{\Delta }.q){\displaystyle \frac{p^2}{8c^2}}(q.\mathrm{\Delta }.q)q_j^2+{\displaystyle \frac{p^2}{4c}}(q.\mathrm{\Delta }.q){\displaystyle \frac{1}{2}}p.\mathrm{\Delta }.p+{\displaystyle \frac{\epsilon }{4c}}(p.\mathrm{\Delta }.p)q_j^2]`$
$`{\displaystyle \frac{1}{4c}}(q^2`$ $`+`$ $`\epsilon (q.\mathrm{\Delta }.q))(q^2+2\epsilon (q.\mathrm{\Delta }.q))\pi _\theta ^2.`$ (35)
This concludes the gauge reformulation of the deformed Skyrme model.
To conclude our task, this gauge invariant model has to be adequately quantized. Thereby, the WZ symmetry is partially fixed as,
$$\psi =\theta .$$
(36)
Consequently, $`\stackrel{~}{\chi }_2`$ turns to second-class constraint,
$$\{\theta ,\stackrel{~}{\chi }_2\}=q^2+2\epsilon q.\mathrm{\Delta }.q$$
(37)
that allows to compute the corresponding Dirac brackets,
$`\{q_i,q_j\}^{}`$ $`=`$ $`\{p_i,p_j\}^{}=0,`$
$`\{q_i,p_j\}^{}`$ $`=`$ $`\delta _{ij}.`$ (38)
Notice that these brackets turn out to be canonical, illustrating the Maskawa-Nakajima theorem. The reduced Hamiltonian is rewritten as,
$$\stackrel{~}{}=\frac{q^2}{4c}p_iM_{ij}p_j+\epsilon [\frac{p^2}{2c}(q.\mathrm{\Delta }.q)\frac{p^2}{8c^2}(q.\mathrm{\Delta }.q)q^2+\frac{q^2}{4c}(p.\mathrm{\Delta }.p)\frac{1}{2}p.\mathrm{\Delta }.p(q.p)\frac{q.\mathrm{\Delta }.p}{2c}+\frac{(q.p)^2}{2cq^2}(q.\mathrm{\Delta }.q)]$$
(39)
where the matrix $`M_{ij}`$,
$$M_{ij}=\delta _{ij}\frac{q_iq_j}{q^2}$$
(40)
is singular. $`\stackrel{~}{\chi }_1|_{\theta =0}=\chi _1`$ remains a first-class constraint, identified as the “Gauss Law” generator of the remaining symmetry.
At this point we are interested to compute the energy spectrum of the model from the invariant Hamiltonian above and to compare it with the one given by ANW. To calculate the energy spectrum we will use the Dirac quantization method for first-class constraint systems. To this end we impose the first-class constraint in the spatial sector of the Hilbert space,
$$(\widehat{q}^2c)|phys>=0.$$
(41)
The momenta representation must reflect the constraint presence and its associated algebra. The first-class nature of the geometric constraint is represented as,
$$[\widehat{q}^2,\widehat{p}_i]=2ı\mathrm{}\widehat{q}_i.$$
(42)
This algebra is satisfied by the following commutator,
$$[\widehat{q}_i,\widehat{p}_j]=ı\mathrm{}\delta _{ij},$$
(43)
with the canonical representation for the momenta operator $`\widehat{p}_i`$,
$$p_i=ı\mathrm{}_i,$$
(44)
which satisfies the relations (42-43). The spectrum of the theory is obtained by evaluating the Hamiltonian meam value with typical eigenfunctions defined, for instance, $`|polyn>=\frac{1}{N(l)}(q_1+iq_2)^l`$ to give,
$$\widehat{\stackrel{~}{}}|polyn>=E_l|polyn>.$$
(45)
A direct calculation produces in this case
$$\widehat{\stackrel{~}{}}|polyn>=\frac{1}{4c}l(l+2)\left\{\left[1\epsilon \frac{q.\mathrm{\Delta }.q}{c}\right](q_1+iq_2)^l+2\epsilon (q_1+iq_2)^{l1}q_b\mathrm{\Delta }_{bc}_c(q_1+iq_2)\right\}.$$
(46)
Therefore, the consistence between (45) and (46) requires,
$$(q_1+iq_2)^{l1}q_b\mathrm{\Delta }_{bc}_c(q_1+iq_2)=(q_1+iq_2)^l.$$
(47)
For instance, a simple choice for the matrix $`\mathrm{\Delta }_{cb}`$, turning the surface slightly oblate,
$$\mathrm{\Delta }=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 0& 0& 0\end{array}\right)$$
(48)
produces the following eigenvalue,
$$E_l=\frac{1}{4c}l(l+2)\left[1+2\epsilon \left(1\frac{q_1^2+q_2^2}{2c}\right)\right].$$
(49)
Clearly, the extra term reflects the correction generated by the deformation away from the spherical symmetry. Note that this deformation matrix $`\mathrm{\Delta }_{ij}`$ has produced a new free parameter that may be used to better the fit energy spectrum to experimental data. Notice that for $`\epsilon >0`$ the above result improves the ANW estimate
$$E_l=E_l^{ANW}+\epsilon \mathrm{\Delta }E_l$$
(50)
Before closing this section we would like to mention an important feature related to the internal consistency of the formalism. If we choose the deformation matrix function $`\mathrm{\Delta }_{ij}=\delta _{ij}`$ (which automatically satisfy (47)) then such deformation corresponds only to a change in the radius of the sphere. A direct calculation from (46) shows that the spectrum remains unchanged. This is important since the radial direction was seen to be the gauge orbit \- spheres of different radii are physically equivalent. Only deformations on the geometry produce energy corrections.
## 5 Final Discussions
The quantization of the Skyrme model is a well known example of quantum mechanics on curved space. This comes from the presence of nonlinear second-class constraints but the passage to the quantum world is plagued by ambiguities. In general the use of the Poisson bracket structure is in conflict with the constraints. The use of the Dirac brackets helps solve this problem classically, but the new algebraic structure ends up being coordinate dependent posing a new problem in terms of ordering ambiguities at the quantum level. Different representations for the momentum operator brings additional terms into the quantum Hamiltonian associated to the zero-point motion. The quantum dynamics becomes dependent on the particular ordering adopted and is not unique.
To cure these sort of problems we propose to treat nonlinear second-class systems as gauge theories by converting the constraints into 1st-class. After conversion has been accomplished, the quantization may be implemented by the Dirac or Gupta-Bleuler procedures where the 1st-class gauge constraints select the physical Hilbert space. Since the Poisson brackets remain as the underlying algebraic structure, no ordering ambiguities affect the quantization process. The question that remains is how to convert efficiently the constraints so that different geometries could also be treated. This seems important to improve the quantization of the rotational modes of the Skyrme theory. The development of a formalism to answer this question is the main contribution of this work.
The development of this formalism was done in Section 4. We studied the model with the most general quadratic constraint and showed that a single WZ term produces the desired covariantization of the Skyrme Lagrangian, even for arbitrary geometries. The exact computation of the spectrum for an arbitrary geometry has shown to be a too difficult problem. We proposed to handle it in a systematic fashion by perturbative technique. A variety of deformations may then be examined through a proper choice of the matrix $`\mathrm{\Delta }_{ij}`$ satisfying the condition (47). The new parameters related to the geometry deformation modify the energy spectrum. Notice that when these parameters preserve the spherical symmetry the old results of ANW are recovered. It should be noticed that this deformation occurs in the internal space of constraints and, therefore, can not be compared directly with the results of the “deformed skyrmions” proposed earlier where the deformation was encoded in the profile function of the hedgehodge solution Our solution provides new free parameters related to the deformation which may be adjusted adequately to fit experimental results.
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# Conjecture on the Interlacing of Zeros in Complex Sturm-Liouville Problems
## I Introduction
The spectra of many classes of non-Hermitian $`𝒫𝒯`$-symmetric Hamiltonians are real and positive . It is believed that the positivity of the spectra is a consequence of $`𝒫𝒯`$ symmetry. Examples of heavily studied $`𝒫𝒯`$-symmetric Hamiltonians are
$$H=p^2(ix)^N(N2)$$
(1)
and
$$H=p^2x^4+2iax^3+(a^22b)x^2+2i(abJ)x(J\mathrm{integer},a^2+4b>K_{\mathrm{critical}}),$$
(2)
where $`K_{\mathrm{critical}}`$ grows with increasing $`J`$. For Hamiltonians like those in (1) and (2) the Schrödinger equations for the $`k`$th eigenfunction,
$$H\mathrm{\Psi }_k(x)=E_k\mathrm{\Psi }_k(x),$$
(3)
involve a complex potential and may require $`x`$ to be complex for the boundary conditions to be defined properly . Thus, the Schrödinger eigenvalue problems may be regarded as analytic extensions of Sturm-Liouville problems into the complex plane.
The eigenfunctions of a conventional self-adjoint Sturm-Liouville problem are complete. Completeness is the statement that a given function can be represented as a linear superposition of the eigenfunctions:
$$f(x)=\underset{n}{}a_n\mathrm{\Psi }_n(x).$$
(4)
It is necessary for the zeros of the eigenfunctions in the complete set to become dense on the interval in which the Sturm-Liouville problem is defined. If the zeros did not become dense, it would be impossible to represent a rapidly varying function . For conventional Sturm-Liouville problems one can prove that the zeros of successive eigenfunctions interlace, and this interlacing of the zeros ensures that the zeros become dense .
A major open mathematical question for $`𝒫𝒯`$-symmetric Hamiltonians is whether the eigenfunctions form a complete set. If the zeros of the eigenfunctions of complex Sturm-Liouville eigenvalue problems exhibit the property of interlacing, this provides heuristic evidence that the eigenfunctions might be complete. A proof of completeness would require that we identify the space in which they are complete, and we do not yet know how to do this. Nevertheless, if we can understand the distribution of the zeros of the eigenfunctions, we gain some insight into the question of completeness for eigenfunctions of $`𝒫𝒯`$-symmetric Sturm-Liouville problems.
## II Some Examples of Distributions of Zeros
We have studied three different complex $`𝒫𝒯`$-symmetric Hamiltonians. We find that in every case the qualitative features of the distribution in the complex plane of the zeros of the eigenfunctions are very similar: We observe a shifted interlacing of zeros. We believe that this pattern of zeros is universal.
Example 1: Large-$`N`$ limit of the $`(ix)^N`$ potential. The large-$`N`$ limit of a $`(ix)^N`$ potential is exactly solvable . In Fig. 1 the zeros of the $`14`$th and $`15`$th eigenfunctions are plotted and clearly exhibit a form of interlacing in the complex plane. For convenience, we have scaled the zeros by dividing by the magnitude of the turning points; we have then performed a linear transformation to fix the turning points at $`\pm 1`$. The appropriate scaling is
$$z=(xE^{1/N}+i)N/\pi .$$
(5)
As shown in Fig. 2, the zeros of the first 15 eigenfunctions interlace and appear to become dense in a narrow region surrounding an arch-shaped contour in the complex plane. This contour is the Stokes’ line that joins the turning points; that is, it is the path along which the phase in the WKB quantization condition is purely real (and thus the quantum-mechanical wave function is purely oscillatory). (It is interesting that this WKB path differs from the path that a classical particle follows in the complex plane as it oscillates between the turning points. The path that a classical particle follows is an inverted arch-shaped contour between the same two turning points.)
Example 2: Quasi-exactly-solvable $`x^4`$ potential. Next, consider the quasi-exactly-solvable potential in (2) with $`a=10`$ and $`b=2`$. The eigenfunctions of the Schrödinger equation have the form of an exponential multiplied by a polynomial. The zeros of this polynomial are easy to calculate numerically. For a given $`J`$ the polynomials in the eigenfunctions all have the same degree and, as a result, all of the eigenfunctions have the same number of zeros. However, for the $`k`$th wave function, $`Jk`$ zeros lie along the branch cut on the positive imaginary axis, and we consider these zeros to be irrelevant. For values of $`J`$ ranging from $`1`$ to $`21`$, the qualitative behavior is always the same. In Fig. 3 the results for $`J=21`$ are plotted and the relevant zeros again lie along the WKB paths in the complex plane. In this case the zeros are not contained in as narrow a region of the complex plane as for the $`(ix)^N`$ potential because the zeros have not been scaled as in Fig. 2. The zeros have an imaginary part that becomes more negative as $`k`$ increases and exhibit the complex version of interlacing.
In Fig. 4 the zeros are scaled so that $`z=x/|x_{\mathrm{TP}}|`$, where $`|x_{\mathrm{TP}}|`$ are the magnitudes of the classical turning points. (The classical turning points are the roots of $`x^4+20ix^3+96x^22ix=E`$.) The scaled zeros lie in a more compact region in the complex-$`z`$ plane than the zeros in Fig. 3. We believe the zeros become dense in this region. Notice that this arch-shaped region is broader than the corresponding region in Fig. 2 because the turning points do not all lie along the same polar angle; therefore, the scaling fixes the magnitudes but not the positions of the turning points. Potentials with various values of $`a`$ and $`b`$ were also investigated and similar results were obtained.
Example 3: $`ix^3`$ potential We obtain an $`ix^3`$ potential when we set $`N=3`$ in (1). Using Runge-Kutta techniques in the complex plane, we have plotted the level curves of the real and imaginary parts of the complex eigenfunctions. By finding the intersections of these level curves, we have determined the zeros of the eigenfunctions numerically. These zeros, which are shown in Fig. 5, lie along the Stokes’ lines of the WKB approximation. Again, the zeros exhibit the complex version of interlacing. They have an imaginary part that decreases as $`k`$ increases. In Fig. 6 the zeros are scaled by $`z=x/|x_{TP}|=xE^{1/3}`$, which fixes the magnitudes and positions of the turning points. This plot suggests that after the scaling the zeros become dense in the complex-$`z`$ plane. Once again, this plot suggests that the distribution of zeros in the complex plane is a universal property of complex Sturm-Liouville eigenvalue problems associated with $`𝒫𝒯`$-symmetric Hamiltonians.
Statement of the Conjecture
From our studies we observe that the unscaled zeros of the complex eigenfunctions do not become dense on a contour or in a narrow region of the complex-$`x`$ plane. In particular, for the $`ix^3`$ potential WKB theory predicts that $`E_kCk^{6/5}`$ ($`k\mathrm{}`$), where $`C`$ is a constant. Thus, the turning points behave like $`x_{TP}(iC)^{1/3}k^{2/5}`$ ($`k\mathrm{}`$) and
$$\frac{dx_{TP}}{dk}\frac{2}{5}(iC)^{1/3}k^{3/5}(k\mathrm{}).$$
(6)
Using Richardson extrapolation we have verified that the distance between zeros along the imaginary axis exhibits this $`k`$-dependence. Consequently, the distance from the contour along which the zeros of $`\mathrm{\Psi }_1`$ lie to the contour along which the zeros of $`\mathrm{\Psi }_k`$ for large $`k`$ lie is given by $`_{k=0}^{\mathrm{}}k^{3/5}`$ which is infinitely far away.
We have scaled the zeros by fixing the magnitudes of the turning points relative to a unit length. After this scaling is performed, the zeros appear to become dense in a narrow arch-shaped region in the scaled complex plane. If the zeros do become dense in this narrow region and exhibit the shifted complex version of interlacing, we conjecture that this behavior suggests that the eigenfunctions are complete in the scaled complex plane. Since we do not know what space within which to define completeness, we are unable to give a rigorous proof.
## ACKNOWLEDGMENTS
We wish to thank P. N. Meisinger for assistance in computer calculations. We also thank the U.S. Department of Energy for financial support.
FIGURE 2
FIGURE 3
FIGURE 4
FIGURE 5
FIGURE 6
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# Semi-Riemannian submersions from real and complex pseudo-hyperbolic spaces
## Introduction
The theory of Riemannian submersions was initiated by O’Neill and Gray . Presently, there is an extensive literature on the Riemannian submersions with different conditions imposed on the total space and on the fibres. A systematic exposition could be found in Besse’s book . Semi-Riemannian submersions were introduced by O’Neill in his book .
The class of harmonic Riemannian submersions, and in particular of those with totally geodesic fibres, is contained in the class of horizontally homothetic harmonic morphisms. For important results concerning the geometry of harmonic morphisms we refer to . Wood constructs examples of harmonic morphisms from Riemannian submersions with totally geodesic fibres by horizontally conformal deformation of the metric. Recently, Fuglede studied harmonic morphisms between semi-Riemannian manifolds (see ). In this paper we solve the classification problem of the semi-Riemannian submersions with totally geodesic fibres from real and complex pseudo-hyperbolic spaces.
R. Escobales , and A. Ranjan classified Riemannian submersions with totally geodesic fibres from a sphere $`S^n`$ and from a complex projective space $`P^n`$. M.A. Magid classified the semi-Riemannian submersions with totally geodesic fibres from an anti-de Sitter space onto a Riemannian manifold. In section §2 we classify the semi-Riemannian submersions with totally geodesic fibres from a pseudo-hyperbolic space onto a Riemannian manifold. Also we obtain the classification of the semi-Riemannian submersions with connected, complex, totally geodesic fibres from a complex pseudo-hyperbolic space onto a Riemannian manifold.
## 1. Preliminaries and examples
###### Definition 1.
Let $`(M,g)`$ be an $`m`$-dimensional connected semi-Riemannian manifold of index $`s`$ ($`0sm`$), let $`(B,g^{})`$ be an $`n`$-dimensional connected semi-Riemannian manifold of index $`s^{}s`$, ($`0s^{}n`$). A semi-Riemannian submersion (see ) is a smooth map $`\pi :MB`$ which is onto and satisfies the following three axioms:
* $`\pi _{}|_p`$ is onto for all $`pM`$;
* the fibres $`\pi ^1(b)`$, $`bB`$ are semi-Riemannian submanifolds of $`M`$;
* $`\pi _{}`$ preserves scalar products of vectors normal to fibres.
We shall always assume that the dimension of the fibres $`dimMdimB`$ is positive and the fibres are connected.
The tangent vectors to fibres are called vertical and those normal to fibres are called horizontal. We denote by $`𝒱`$ the vertical distribution and by $``$ the horizontal distribution.
B. O’Neill has characterized the geometry of a Riemannian submersion in terms of the tensor fields $`T`$, $`A`$ defined by
$`A_EF`$ $`=`$ $`h_{hE}vF+v_{hE}hF`$
$`T_EF`$ $`=`$ $`h_{vE}vF+v_{vE}hF`$
for every $`E`$, $`F`$ tangent vector fields to $`M`$. Here $``$ is the Levi-Civita connection of $`g`$, the symbols $`v`$ and $`h`$ are the orthogonal projections on $`𝒱`$ and $``$, respectively. The letters $`U`$, $`V`$ will always denote vertical vector fields, $`X`$, $`Y`$, $`Z`$ horizontal vector fields. Notice that $`T_UV`$ is the second fundamental form of each fibre and $`A_XY`$ is a natural obstruction to integrability of horizontal distribution $``$. The tensor $`A`$ is called O’Neill’s integrability tensor. For basic properties of Riemannian submersions and examples see , , . A vector field $`X`$ on $`M`$ is said to be basic if $`X`$ is horizontal and $`\pi `$related to a vector field $`X^{}`$ on $`B`$. Notice that every vector field $`X^{}`$ on $`B`$ has a unique horizontal lifting $`X`$ to $`M`$ and $`X`$ is basic. The following lemma is well known (see ).
###### Lemma 1.1.
We suppose X and Y are basic vector fields on M which are $`\pi `$-related to $`X^{}`$ and $`Y^{}`$. Then
* $`h[X,Y]`$ is basic and $`\pi `$-related to $`[X^{},Y^{}]`$;
* $`h_XY`$ is basic and $`\pi `$-related to $`_X^{}^{}Y^{}`$ , where $`^{}`$ is the Levi-Civita connection on B;
The O’Neill’s integrability tensor $`A`$ has the following properties (see or ).
###### Lemma 1.2.
Let $`X`$, $`Y`$ be horizontal vector fields and $`E`$, $`F`$ be vector fields on $`M`$. Then each of the following holds:
* $`A_XY=A_YX`$;
* $`A_{hE}F=A_EF`$;
* $`A_E`$ maps the horizontal subspace into the vertical one and the vertical subspace into the horizontal one;
* $`g(A_XE,F)=g(E,A_XF)`$;
* If moreover $`X`$ is basic then $`A_XV=h_VX`$ for every vertical vector field $`V`$;
* $`g((_YA)_XE,F)=g(E,(_YA)_XF)`$.
Let $`\widehat{g}`$ be the induced metric on fibre $`\pi ^1(\pi (p))`$, $`pM`$. We denote by $`R`$, $`R^{}`$, $`\widehat{R}`$ the Riemann tensors of the metrics $`g`$, $`g^{}`$, $`\widehat{g}`$ respectively.
The following equations, usually called O’Neill’s equations, characterize the geometry of a semi-Riemannian submersion (see , , ).
###### Proposition 1.3.
For every vertical vector fields $`U`$, $`V`$, $`W`$, $`W^{}`$ and for every horizontal vector fields $`X`$, $`Y`$, $`Z`$, $`Z^{}`$, we have the following formulae:
* $`R(U,V,W,W^{})=\widehat{R}(U,V,W,W^{})g(T_UW,T_VW^{})+g(T_VW,T_UW^{})`$,
* $`R(U,V,W,X)=g((_VT)_UW,X)g((_UT)_VW,X),`$
* $`R(X,U,Y,V)=g((_XT)_UV,Y)g(T_UX,T_VY)+g((_UA)_XY,V)+g(A_XU,A_YV)`$,
* $`R(U,V,X,Y)=g((_UA)_XY,V)g((_VA)_XY,U)+g(A_XU,A_YV)g(A_XV,A_YU)g(T_UX,T_VY)+g(T_VX,T_UY)`$,
* $`R(X,Y,Z,U)=g((_ZA)_XY,U)+g(A_XY,T_UZ)g(A_YZ,T_UX)g(A_ZX,T_UY)`$,
* $`R(X,Y,Z,Z^{})=R^{}(\pi _{}X,\pi _{}Y,\pi _{}Z,\pi _{}Z^{})2g(A_XY,A_ZZ^{})+g(A_YZ,A_XZ^{})g(A_XZ,A_YZ^{}).`$
Using O’Neill’s equations, we get the following lemma.
###### Lemma 1.4.
If $`\pi :(M,g)(B,g^{})`$ is a semi-Riemannian submersion with totally geodesic fibres then:
$`a)R(U,V,U,V)`$ $`=`$ $`\widehat{R}(U,V,U,V);`$
$`b)R(X,U,X,U)`$ $`=`$ $`g(A_XU,A_XU);`$
$`c)R(X,Y,X,Y)`$ $`=`$ $`R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Y)3g(A_XY,A_XY).`$
We recall the definitions of real and complex pseudo-hyperbolic spaces (see and ).
###### Definition 2.
Let $`<,>`$ be the symmetric bilinear form on $`^{m+1}`$ given by
$$<x,y>=\underset{i=0}{\overset{s}{}}x_iy_i+\underset{i=s+1}{\overset{m}{}}x_iy_i$$
for $`x=(x_0,\mathrm{},x_m),y=(y_0,\mathrm{},y_m)^{m+1}`$. For $`s>0`$ let
$`H_s^m=\{x^{m+1}|<x,x>=1\}`$ be the semi-Riemannian submanifold of
$$_{s+1}^{m+1}=(^{m+1},ds^2=dx^0dx^0\mathrm{}dx^sdx^s+dx^{s+1}dx^{s+1}+\mathrm{}+dx^mdx^m).$$
$`H_s^m`$ is called the $`m`$-dimensional (real) pseudo-hyperbolic space of index $`s`$. We notice that $`H_s^m`$ has constant sectional curvature $`1`$ and the curvature tensor is given by $`R(X,Y,X,Y)=g(X,X)g(Y,Y)+g(X,Y)^2`$. $`H_s^m`$ can be written as homogeneous space, namely we have $`H_s^m=SO(s+1,ms)/SO(s,ms)`$, $`H_{2s+1}^{2m+1}=SU(s+1,ms)/SU(s,ms)`$, $`H_{4s+3}^{4m+3}=Sp(s+1,ms)/Sp(s,ms)`$ (see ).
###### Definition 3.
Let $`(,)`$ be the hermitian scalar product on $`^{m+1}`$ given by
$$(z,w)=\underset{i=0}{\overset{s}{}}z_i\overline{w_i}+\underset{i=s+1}{\overset{m}{}}z_i\overline{w_i}$$
for $`z=(z_0,\mathrm{},z_m),w=(w_0,\mathrm{},w_m)^{m+1}`$. Let $`M`$ be the real hypersurface of $`^{m+1}`$ given by $`M=\{z^{m+1}|(z,z)=1\}`$ and endowed with the induced metric of
$$(^{m+1},ds^2=dz^0d\overline{z}^0\mathrm{}dz^sd\overline{z}^s+dz^{s+1}d\overline{z}^{s+1}+\mathrm{}+dz^md\overline{z}^m).$$
The natural action of $`S^1=\{e^{i\theta }|\theta \}`$ on $`^{m+1}`$ induces an action on $`M`$. Let $`H_s^m=M/S^1`$ endowed with the unique indefinite Kähler metric of index $`2s`$ such that the projection $`MM/S^1`$ becomes a semi-Riemannian submersion (see ). $`H_s^m`$ is called the complex pseudo-hyperbolic space. Notice that $`H_s^m`$ has constant holomorphic sectional curvature $`4`$ and the curvature tensor is given by $`R(X,Y,X,Y)=g(X,X)g(Y,Y)+g(X,Y)^23g(I_0X,Y)^2`$, where $`I_0`$ is the natural complex structure on $`H_s^m`$. $`H_s^m`$ is a homogeneous space, namely we have (see ) $`H_s^m=SU(s+1,ms)/S(U(1)U(s,ms))`$ and
$`H_{2s+1}^{2m+1}=Sp(s+1,ms)/U(1)Sp(s,ms)`$.
We denote by $`H^n(4)`$ the hyperbolic space with sectional curvature $`4`$, by $`H^n`$ the quaternionic hyperbolic space of real dimension $`4n`$ with quaternionic sectional curvature $`4`$.
Many explicit examples of semi-Riemannian submersions with totally geodesic fibres can be given following a standard construction (see for Riemannian case). Let $`G`$ be a Lie group and $`K`$, $`H`$ two compact Lie subgroups of $`G`$ with $`KH`$. Let $`\pi :G/KG/H`$ be the associated bundle with fibre $`H/K`$ to the $`H`$principal bundle $`p:GG/H`$. Let $`𝚐`$ be the Lie algebra of $`G`$ and $`𝚔𝚑`$ the corresponding Lie subalgebras of $`K`$ and $`H`$. We choose an $`Ad(H)`$invariant complement $`𝚖`$ to $`𝚑`$ in $`𝚐`$, and an $`Ad(K)`$invariant complement $`𝚙`$ to $`𝚔`$ in $`𝚑`$. An $`ad(H)`$-invariant nondegenerate bilinear symmetric form on $`𝚖`$ defines a $`G`$-invariant semi-Riemannian metric $`g^{}`$ on $`G/H`$ and an $`ad(K)`$-invariant nondegenerate bilinear symmetric form on $`𝚙`$ defines a $`H`$-invariant semi-Riemannian metric $`\widehat{g}`$ on $`H/K`$. The orthogonal direct sum for these nondegenerate bilinear symmetric forms on $`𝚙𝚖`$ defines a $`G`$-invariant semi-Riemannian metric $`g`$ on $`G/K`$. The following theorem is proved in .
###### Theorem 1.5.
The map $`\pi :(G/K,g)(G/H,g^{})`$ is a semi-Riemannian submersion with totally geodesic fibres.
Using this theorem we get the following examples.
###### Example 1.
Let $`G=SU(1,n)`$, $`H=S(U(1)U(n))`$, $`K=SU(n)`$. We have the semi-Riemannian submersion
$$H_1^{2n+1}=SU(1,n)/SU(n)H^n=SU(1,n)/S(U(1)U(n)).$$
###### Example 2.
Let $`G=Sp(1,n)`$, $`H=Sp(1)Sp(n)`$, $`K=Sp(n)`$. We get the semi-Riemannian submersion
$$H_3^{4n+3}=Sp(1,n)/Sp(n)H^n=Sp(1,n)/Sp(1)Sp(n).$$
###### Example 3.
Let $`G=Spin(1,8)`$, $`H=Spin(8)`$, $`K=Spin(7)`$. We have the semi-Riemannian submersion
$$H_7^{15}=Spin(1,8)/Spin(7)H^8(4)=Spin(1,8)/Spin(8).$$
###### Example 4.
Let $`G=Sp(1,n)`$, $`H=Sp(1)Sp(n)`$, $`K=U(1)Sp(n)`$. We obtain the semi-Riemannian submersion
$$H_1^{2n+1}=Sp(1,n)/U(1)Sp(n)H^n=Sp(1,n)/Sp(1)Sp(n).$$
###### Definition 4.
Two semi-Riemannian submersions $`\pi ,\pi ^{}:(M,g)(B,g^{})`$ are called equivalent if there is an isometry $`f`$ of $`M`$ which induces an isometry $`\stackrel{~}{f}`$ of $`B`$ so that $`\pi ^{}f=\stackrel{~}{f}\pi `$. In this case the pair $`(f,\stackrel{~}{f})`$ is called a bundle isometry.
We shall need the following theorem, which is the semi-Riemannian version of theorem 2.2 in .
###### Theorem 1.6.
Let $`\pi _1,\pi _2:MB`$ be semi-Riemannian submersions from a connected complete semi-Riemannian manifold onto a semi-Riemannian manifold. Assume the fibres of these submersions are connected and totally geodesic. Suppose $`f`$ is an isometry of $`M`$ which satisfies the following two properties at a given point $`pM`$:
* $`f_p:T_pMT_{f(p)}M`$ maps $`_{1p}`$ onto $`_{2f(p)}`$, where $`_i`$ denotes the horizontal distribution of $`\pi _i`$, $`i\{1,2\}`$;
* For every $`E`$, $`FT_pM`$, $`f_{}A_{1E}F=A_{2f_{}E}f_{}F`$, where $`A_i`$ are the integrability tensors associated with $`\pi _i`$.
Then $`f`$ induces an isometry $`\stackrel{~}{f}`$ of $`B`$ so that the pair $`(f,\stackrel{~}{f})`$ is a bundle isometry between $`\pi _1`$ and $`\pi _2`$. In particular, $`\pi _1`$ and $`\pi _2`$ are equivalent.
## 2. Semi-Riemannian submersions with totally geodesic fibres
###### Proposition 2.1.
If $`\pi :H_s^mB^n`$ is a semi-Riemannian submersion with totally geodesic fibres from an $`m`$-dimensional pseudo-hyperbolic space of index $`s`$ onto an $`n`$-dimensional Riemannian manifold then $`m=n+s`$, the induced metrics on fibres are negative definite and $`B`$ has negative sectional curvature.
###### Proof.
By lemma 1.4-b), we get $`g(A_XV,A_XV)=g(X,X)g(V,V)0`$ for every horizontal vector $`X`$ and for every vertical vector $`V`$. Therefore $`g(V,V)0`$ for every vertical vector $`V`$. By lemma 1.4-c), we have $`R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Y)=g^{}(\pi _{}X,\pi _{}X)g^{}(\pi _{}Y,\pi _{}Y)+g^{}(\pi _{}X,\pi _{}Y)^2+3g(A_XY,A_XY)<0`$ for every linearly independent horizontal vectors $`X`$ and $`Y`$. ∎
###### Proposition 2.2.
Let $`\pi :(M_s^{n+s},g)(B^n,g^{})`$ be a semi-Riemannian submersion from an $`(n+s)`$-dimensional semi-Riemannian manifold of index $`s1`$ onto an $`n`$-dimensional Riemannian manifold. We suppose $`M`$ is geodesically complete and simply connected. Then $`B`$ is complete and simply connected. If moreover $`B`$ has nonpositive curvature then the fibres are simply connected.
###### Proof.
Since $`M`$ is geodesically complete, the base space $`B`$ is complete.
Let $`\stackrel{~}{g}`$ be the Riemannian metric on $`M`$ defined by
$$\stackrel{~}{g}(E,F)=g(hE,hF)g(vE,vF)$$
for every $`E`$, $`F`$ vector fields on $`M`$. Since $`\stackrel{~}{g}`$ is a horizontally complete Riemannian metric (this means that any maximal horizontal geodesic is defined on the entire real line) and $`B`$ is a complete Riemannian manifold then $``$ is an Ehresmann connection for $`\pi `$ (see theorem 1 in ). By theorem 9.40 in , it follows $`\pi :MB`$ is a locally trivial fibration and we have an exact homotopy sequence
$$\mathrm{}\pi _2(B)\pi _1(fibre)\pi _1(M)\pi _1(B)0$$
Since $`M`$ is simply connected, we have $`\pi _1(B)=0`$.
If $`B`$ has nonpositive curvature, then $`\pi _2(B)=0`$ by theorem of Hadamard. It follows $`\pi _1(fibre)=0`$. ∎
###### Theorem 2.3.
If $`\pi :H_s^mB^n`$ is a semi-Riemannian submersion with totally geodesic fibres from a pseudo-hyperbolic space of index $`s>1`$ onto a Riemannian manifold then $`B`$ is a Riemannian symmetric space of rank one, noncompact and simply connected, any fibre is diffeomorphic to $`S^s`$ and $`s\{3,7\}`$.
###### Proof.
In order to prove that $`B`$ is a locally symmetric space we need to check that $`^{}R^{}0`$.
Let $`X_0^{}`$, $`X^{}`$, $`Y^{}`$, $`Z^{}`$ be vector fields on $`B`$ and let $`X_0`$, $`X`$, $`Y`$, $`Z`$ be the horizontal liftings of $`X_0^{}`$, $`X^{}`$, $`Y^{}`$, $`Z^{}`$ respectively. By definition of the covariant derivative we have
$`(_{X_0^{}}^{}R^{})(X^{},Y^{},Z^{})`$ $`=`$ $`_{X_0^{}}^{}R^{}(X^{},Y^{})Z^{}R^{}(_{X_0^{}}^{}X^{},Y^{})Z^{}`$
$`R^{}(X^{},_{X_0^{}}^{}Y^{})Z^{}R^{}(X^{},Y^{})_{X_0^{}}^{}Z^{}.`$
In order to prove that the curvature tensor $`R^{}`$ of the base space is parallel, we have to lift all vector fields in relation (2). By lemma 1.1, the horizontal liftings of $`_{X_0^{}}^{}X^{}`$, $`_{X_0^{}}^{}Y^{}`$ and $`_{X_0^{}}^{}Z^{}`$ are $`h_{X_0}X`$, $`h_{X_0}Y`$ and $`h_{X_0}Z`$, respectively.
We denote by $`R^h(X,Y)Z`$ the horizontal lifting of $`R^{}(X^{},Y^{})Z^{}`$. The convention for Riemann tensor used here is $`R(X,Y)=_X_Y_Y_X_{[X,Y]}`$. O’Neill’s equation $`vi)`$ gives us the following relation
$$R^h(X,Y)Z=h(R(X,Y)Z)+2A_ZA_XYA_XA_YZA_YA_ZX.$$
Using this relation we compute
$`(_{X_0}^{}R^{})(X^{},Y^{},Z^{})`$ $`=`$ $`\pi _{}[h_{X_0}(R^h(X,Y)Z)R^h(h_{X_0}X,Y)Z`$
$`R^h(X,h_{X_0}Y)ZR^h(X,Y)h_{X_0}Z]`$
$`=`$ $`\pi _{}[h_{X_0}h(R(X,Y)Z)hR(h_{X_0}X,Y)Z`$
$`hR(X,h_{X_0}Y)ZhR(X,Y)h_{X_0}Z`$
$`+2(h_{X_0}A_ZA_XYA_{h_{X_0}Z}A_XY`$
$`A_ZA_{h_{X_0}X}YA_ZA_Xh_{X_0}Y)`$
$`(h_{X_0}A_XA_YZA_{h_{X_0}X}A_YZ`$
$`A_XA_{h_{X_0}Y}ZA_XA_Yh_{X_0}Z)`$
$`(h_{X_0}A_YA_ZXA_{h_{X_0}Y}A_ZX`$
$`A_YA_{h_{X_0}Z}XA_YA_Zh_{X_0}X)].`$
Since $`H_s^m`$ has constant curvature, we have $`R(X,Y,Z,U)=0`$ for every vertical vector $`U`$ and for every horizontal vector fields $`X`$, $`Y`$, $`Z`$. This implies $`R(X,Y)Z`$ is horizontal and $`R(X,U)Y`$, $`R(U,X)Y`$, $`R(X,Y)U`$ are vertical. Hence
$`\pi _{}(_{X_0}h(R(X,Y)Z)hR(h_{X_0}X,Y)ZhR(X,h_{X_0}Y)ZhR(X,Y)h_{X_0}Z)=\pi _{}(_{X_0}R(X,Y)Z)\pi _{}(R(_{X_0}X,Y)ZR(v_{X_0}X,Y)Z)\pi _{}(R(X,Y)_{X_0}ZR(X,Y)v_{X_0}Z)=\pi _{}[(_{X_0}R)(X,Y,Z)]`$. Since $`H_s^m`$ has constant curvature, we get $`(_{X_0}R)(X,Y,Z)=0`$. So the sum of the first four terms in relation (2) is zero.
We have $`h_{X_0}A_ZA_XYA_{h_{X_0}Z}A_XYA_ZA_{h_{X_0}X}YA_ZA_Xh_{X_0}Y=`$
$`h((_{X_0}A)_Z(A_XY))A_Z(v(_{X_0}A)_XY)`$.
For the case of totally geodesic fibres, O’Neill’s equation $`v)`$ becomes
$$R(X,Y,Z,U)=g((_ZA)_XY,U).$$
By lemma 1.2 (f) and by the hypothesis of constant curvature total space we get
$$g((_ZA)_XU,Y)=g((_ZA)_XY,U)=0$$
for every horizontal vector fields $`X`$, $`Y`$, $`Z`$ and for every vertical vector field $`U`$. It follows $`h(_ZA)_XU=0`$ and $`v(_ZA)_XY=0`$ for every horizontal vector fields $`X`$, $`Y`$, $`Z`$ and for every vertical vector field $`U`$. Therefore $`h((_{X_0}A)_Z(A_XY))=0`$ and $`v((_{X_0}A)_XY)=0`$. This implies
$$h_{X_0}A_ZA_XYA_{h_{X_0}Z}A_XYA_ZA_{h_{X_0}X}YA_ZA_Xh_{X_0}Y=0.$$
By circular permutations of $`(X,Y,Z)`$ in the last relation we get
$$h_{X_0}A_XA_YZA_{h_{X_0}X}A_YZA_XA_{h_{X_0}Y}ZA_XA_Yh_{X_0}Z=0,$$
$$h_{X_0}A_YA_ZXA_{h_{X_0}Y}A_ZXA_YA_{h_{X_0}Z}XA_YA_Zh_{X_0}X=0.$$
So the sum of all terms in relation (2) is zero.
We proved that $`(_{X_0^{}}^{}R^{})(X^{},Y^{},Z^{})=0`$ for every vector fields $`X_0^{}`$, $`X^{}`$, $`Y^{}`$, $`Z^{}`$, so $`B`$ is a locally symmetric space. By proposition 2.2, $`B`$ is simply connected and complete. Therefore $`B`$ is a Riemannian symmetric space. By proposition 2.1, $`B`$ has negative sectional curvature. Hence $`B`$ is a noncompact Riemannian symmetric space of rank one.
Let $`bB`$. Since $`\pi ^1(b)`$ is a totally geodesic submanifold of a geodesically complete manifold, $`\pi ^1(b)`$ is itself geodesically complete. Since $`R^{}(X^{},Y^{},X^{},Y^{})0`$ for every $`X^{}`$, $`Y^{}`$ tangent vectors to $`B`$, we have $`\pi _1(fibre)=0`$, by proposition 2.2. Since $`(\pi ^1(b),\widehat{g})`$ is a complete, simply connected semi-Riemannian manifold of dimension $`r`$ and of index $`r`$ and with constant sectional curvature $`1`$, it follows $`(\pi ^1(b),\widehat{g})`$ is isometric to $`H_s^s`$ (see Proposition 23 from page 227 in ). Hence any fibre is diffeomorphic to $`S^s`$.
We shall prove below that the tangent bundle of any fibre is trivial. From a well known result of Adams it follows that $`s\{1,3,7\}`$.
###### Lemma 2.4.
The tangent bundle of any fibre is trivial.
###### Proof of lemma 2.4.
Since $`g(A_XV,A_XV)=g(X,X)g(V,V)`$, we have that
$`A_X:𝒱`$, $`VA_XV`$ is an injective map and $`dim𝒱dim`$, if $`g(X,X)0`$.
For any horizontal vector field $`X`$, we denote by $`A_X^{}:𝒱`$ the map given by $`A_X^{}(Y)=A_XY`$. By O’Neill’s equation $`iv)`$, we have $`g(A_XV,A_XW)=g(X,X)g(V,W)`$ for every vertical vector fields $`V`$ and $`W`$. Hence, by lemma 1.2 $`(d)`$, we get $`A_X^{}A_XV=g(X,X)V`$ for every vertical vector field $`V`$. If $`g(X,X)0`$ anywhere then $`A_X^{}`$ is surjective and hence $`dim𝒱=dimdim\mathrm{ker}A_X^{}`$. By lemma 1.2 $`(d)`$, we have $`A_XX=0`$. This implies $`dim\mathrm{ker}A_X^{}1`$.
Let $`bB`$ and $`xT_bB`$ with $`g(x,x)=1`$. We denote by $`X`$ the horizontal lifting along the fibre $`\pi ^1(b)`$ of the vector $`x`$. Let $`p`$ an arbitrary point in $`\pi ^1(b)`$ and let $`\{X(p),y_1,\mathrm{},y_l\}`$ be an orthonormal basis of the vector space $`\mathrm{ker}A_{X(p)}^{}`$. Since $`\pi _p`$ sends isometrically $`_p`$ into $`T_bB`$ we have $`\{\pi _{}X(p),\pi _{}y_1,\mathrm{},\pi _{}y_l\}`$ is a linearly independent system which can be completed to a basis of $`T_bB`$ with a system of vectors $`\{x_{l+1},\mathrm{},x_{n1}\}`$. Let $`X,X_1,X_2,\mathrm{},X_{n1}`$ be the horizontal liftings along the fibre $`\pi ^1(b)`$ of $`x=\pi _{}X(p),\pi _{}y_1,\mathrm{},\pi _{}y_l,x_{l+1},\mathrm{},x_{n1}`$ respectively.
By lemma 1.4, we have for every $`q\pi ^1(b)`$ and for every $`i\{1,\mathrm{},l\}`$
$`3g(A_{X(q)}X_i(q),A_{X(q)}X_i(q))`$ $`=`$ $`R^{}(\pi _{}X(q),\pi _{}X_i(q),\pi _{}X(q),\pi _{}X_i(q))`$
$`R(X(q),X_i(q),X(q),X_i(q))`$
$`=`$ $`R^{}(x,\pi _{}y_i,x,\pi _{}y_i)`$
$`+g(X(q),X(q))g(X_i(q),X_i(q))g(X(q),X_i(q))^2`$
$`=`$ $`R^{}(x,\pi _{}y_i,x,\pi _{}y_i)`$
$`+g^{}(\pi _{}X(q),\pi _{}X(q))g^{}(\pi _{}X_i(q),\pi _{}X_i(q))`$
$`g^{}(\pi _{}X(q),\pi _{}X_i(q))^2`$
$`=`$ $`3g(A_{X(p)}X_i(p),A_{X(p)}X_i(p))`$
$`=`$ $`0.`$
Since the induced metrics on fibre $`\pi ^1(b)`$ are negative definite, we get $`A_{X(q)}X_i(q)=0`$.
By lemma 1.2 (a), we have $`A_{X(q)}X(q)=0`$. We proved that $`\{X(q),X_1(q),\mathrm{},X_l(q)\}\mathrm{ker}A_{X(q)}^{}`$. Since $`\pi _q`$ sends isometrically $`_q`$ into $`T_bB`$, we get $`\{X(q),X_1(q),\mathrm{},X_l(q)\}`$ is a basis of the vector space $`\mathrm{ker}A_{X(q)}^{}`$ for every point $`q\pi ^1(b)`$.
Let $`V_{l+1}=A_X^{}X_{l+1},\mathrm{},V_{n1}=A_X^{}X_{n1}`$ be tangent vector fields to the fibre $`\pi ^1(b)`$. We denote by $`Q_q`$ the vector subspace of $`_q`$ spanned by $`\{X_{l+1}(q),X_{l+2}(q),\mathrm{},X_{n1}(q)\}`$. Let $`\stackrel{~}{g}`$ be the Riemannian metric on $`\pi ^1(b)`$ given by $`\stackrel{~}{g}(V,W)=g(V,W)`$ for every $`V`$, $`W`$ vector fields tangent to $`\pi ^1(b)`$. Since $`dim𝒱_q=dimQ_q`$ and $`g(A_XV,A_XV)=\stackrel{~}{g}(V,V)`$, we get $`A_{X(q)}:(𝒱_q,\stackrel{~}{g})(Q_q,g)`$ is an isometry for every $`q\pi ^1(b)`$.
So $`\{V_{l+1},\mathrm{},V_{n1}\}`$ is a global frame for the tangent bundle of $`\pi ^1(b)`$. It follows the tangent bundle of the fibre $`\pi ^1(b)`$ is trivial. ∎
This ends the proof of theorem 2.3. ∎
By the classification of the Riemannian symmetric spaces of rank one of noncompact type, we have $`B`$ is isometric to one of the following spaces:
* $`H^n(c)`$ real hyperbolic space with constant sectional curvature $`c`$;
* $`H^k(c)`$ complex hyperbolic space with holomorphic sectional curvature $`c`$;
* $`H^k(c)`$ quaternionic hyperbolic space with quaternionic sectional curvature $`c`$;
* $`aH^2(c)`$ Cayley hyperbolic plane with Cayley sectional curvature $`c`$.
This will give us more information about the relation between the dimension of fibres and the geometry of base space.
###### Proposition 2.5.
Let $`\pi :H_s^{n+s}B^n`$ be a semi-Riemannian submersion with totally geodesic fibres.
* If $`s=3`$ then $`n=4k`$ and $`B^n`$ is isometric to $`H^k`$.
* If $`s=7`$ then we have one of the following situations:
+ $`n=8`$ and $`B^n`$ is isometric to $`H^8(4)`$; or
+ $`n=16`$ and $`B^n`$ is isometric to $`aH^2`$.
###### Proof.
Let $`Y`$, $`Z`$ be two linear independent horizontal vectors and let $`Y^{}=\pi _{}Y`$, $`Z^{}=\pi _{}Z`$. By proposition 2.1, the metric induced on fibres are negative definite. This implies $`g(A_ZY,A_ZY)0`$. By lemma 1.4, we get
$$K^{}(Z^{},Y^{})=\frac{R^{}(Z^{},Y^{},Z^{},Y^{})}{g^{}(Z^{},Z^{})g^{}(Y^{},Y^{})g^{}(Z^{},Y^{})^2}=1+\frac{3g(A_ZY,A_ZY)}{g(Z,Z)g(Y,Y)g(Z,Y)^2}1.$$
By Schwartz inequality applied to the positive definite scalar product induced on $``$, we have
$$g(A_ZY,A_ZY)=g(A_ZA_ZY,Y)\sqrt{g(A_ZA_ZY,A_ZA_ZY)}\sqrt{g(Y,Y)}.$$
By lemma 1.4, we get $`g(A_ZY,A_ZY)\sqrt{g(A_ZY,A_ZY)g(Z,Z)}\sqrt{g(Y,Y)}.`$ Thus $`g(A_ZY,A_ZY)g(Z,Z)g(Y,Y)`$. Therefore $`K^{}(Z^{},Y^{})=1+\frac{3g(A_ZY,A_ZY)}{g(Z,Z)g(Y,Y)}4`$ for every orthogonal vectors $`Z^{}`$ and $`Y^{}`$.
We proved that $`4K^{}1`$.
We shall prove that if the base space $`B`$ has constant curvature $`c`$ then $`c=4`$. It is sufficient to see that for any point $`bB`$ there is a $`2`$plane $`\alpha T_bB`$ such that $`K(\alpha )=4`$. We choose $`\alpha =\{\pi _{}Z,\pi _{}A_ZV\}`$ where $`Z`$ is a horizontal vector and $`V`$ is a vertical vector. By lemma 1.4, we have
$`R^{}(\pi _{}Z,\pi _{}A_ZV,\pi _{}Z,\pi _{}A_ZV)`$ $`=`$ $`R(Z,A_ZV,Z,A_ZV)`$
$`+3g(A_Z(A_ZV),A_Z(A_ZV)).`$
We notice that $`Z`$ and $`A_ZV`$ are orthogonal, because, by lemma 1.2, we have $`g(Z,A_ZV)=g(A_ZZ,V)=0`$.
By lemma 1.4, we have
$$g(A_XU,A_XU)=g(X,X)g(U,U)$$
for every horizontal vector $`X`$ and for every vertical vector $`U`$. By lemma 1.2 (d), we get $`g(A_XA_XU,U)=g(X,X)g(U,U)`$. Hence, by polarization, we find $`A_XA_XU=g(X,X)U`$ for every horizontal vector $`X`$ and for every vertical vector $`U`$. Therefore the relation (2) becomes
$`R^{}(\pi _{}Z,\pi _{}A_ZV,\pi _{}Z,\pi _{}A_ZV)`$ $`=`$ $`g(Z,Z)g(A_ZV,A_ZV)+3g(Z,Z)^2g(V,V)`$
$`=`$ $`4g(Z,Z)^2g(V,V)`$
$`=`$ $`4(g^{}(\pi _{}Z,\pi _{}Z)g^{}(\pi _{}A_ZV,\pi _{}A_ZV)g^{}(\pi _{}Z,\pi _{}A_ZV)^2).`$
Then $`K^{}(\pi _{}Z,\pi _{}A_ZV)=4`$. Therefore if the base space $`B`$ has constant curvature $`c`$ then $`c=4`$.
Let $`X`$ be a horizontal vector field. By lemma 1.4, $`Y\mathrm{ker}A_X^{}`$ if and only if
$$R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Y)=g^{}(\pi _{}X,\pi _{}X)g^{}(\pi _{}Y,\pi _{}Y)+g^{}(\pi _{}X,\pi _{}Y)^2.$$
For every $`X^{}T_{\pi (p)}B`$, we denote by
$$_X^{}=\{Y^{}T_{\pi (p)}B|R^{}(X^{},Y^{},X^{},Y^{})=g^{}(X^{},X^{})g^{}(Y^{},Y^{})+g^{}(X^{},Y^{})^2\}.$$
With this notation, $`\pi _{}(\mathrm{ker}A_{X(p)}^{})=_{\pi _{}X(p)}`$. Since $`\pi _{}`$ sends isometrically $`_p`$ into $`T_{\pi (p)}B`$, we have $`dimdim𝒱=dim\mathrm{ker}A_{X(p)}^{}=dim_{\pi _{}X(p)}`$
We compute $`dim_X^{}`$ from the geometry of $`B`$. We have the following possibilities for $`B`$:
###### Case 1.
$`B=H^k(4)`$.
The curvature tensor of hyperbolic space $`H^k(4)`$ is given by
$$R^{}(X^{},Y^{},X^{},Y^{})=4(g^{}(X^{},X^{})g^{}(Y^{},Y^{})g^{}(X^{},Y^{})^2).$$
We have $`_X^{}=\{\lambda X^{}|\lambda \}`$. Hence $`dim_X^{}=1`$. It follows $`dim=dim𝒱+1`$.
If $`s=3`$ then $`B^n`$ is isometric to $`H^4(4)`$, which falls in the case a), since $`H^4(4)`$ is isometric to $`H^1`$.
If $`s=7`$ then $`dim=8`$ and this is the case b ii).
###### Case 2.
$`B=H^k`$.
Let $`I_0`$ be the natural complex structure on $`H^k`$. The curvature tensor of complex hyperbolic space $`H^k`$ with $`4K^{}1`$ is given by
$$R^{}(X^{},Y^{},X^{},Y^{})=(g^{}(X^{},X^{})g^{}(Y^{},Y^{})g^{}(X^{},Y^{})^2+3g^{}(I_0X^{},Y^{})^2).$$
We get $`_X^{}=\{I_0X^{}\}^{}`$. So $`dim_X^{}=2k1=dim1`$. It follows $`dim𝒱=1`$.
###### Case 3.
$`B=H^k`$.
Let $`\{I_0,J_0,K_0\}`$ be local almost complex structures which rise to the quaternionic structure on $`H^k`$. The curvature tensor of the quaternionic hyperbolic space $`H^k`$ with $`4K^{}1`$ (see ) is given by
$`R^{}(X^{},Y^{},X^{},Y^{})`$ $`=`$ $`g^{}(X^{},X^{})g^{}(Y^{},Y^{})+g^{}(X^{},Y^{})^2`$
$`3g^{}(I_0X^{},Y^{})^23g^{}(J_0X^{},Y^{})^23g^{}(K_0X^{},Y^{})^2.`$
It follows that $`Y^{}_X^{}`$ if and only if $`g^{}(I_0X^{},Y^{})=g^{}(J_0X^{},Y^{})=g^{}(K_0X^{},Y^{})=0`$. Therefore $`_X^{}=\{I_0X^{},J_0X^{},K_0X^{}\}^{}`$. Hence $`dim_X^{}=4k3=dim3`$. We get $`dim𝒱=3`$.
###### Case 4.
$`B=aH^2`$.
Let $`\{I_0,J_0,K_0,M_0,M_0I_0,M_0J_0,M_0K_0\}`$ be local almost complex structures which rise to the Cayley structure on $`aH^2`$ Cayley hyperbolic plane. The curvature tensor of the Cayley plane $`aH^2`$ with $`4K^{}1`$ (see ) is given by
$`R^{}(X^{},Y^{},X^{},Y^{})`$ $`=`$ $`g^{}(X^{},X^{})g^{}(Y^{},Y^{})+g^{}(X^{},Y^{})^23g^{}(I_0X^{},Y^{})^2`$
$`3g^{}(J_0X^{},Y^{})^23g^{}(K_0X^{},Y^{})^23g^{}(M_0I_0X^{},Y^{})^2`$
$`3g^{}(M_0J_0X^{},Y^{})^23g^{}(M_0K_0X^{},Y^{})^2.`$
We get
$$_X^{}=\{I_0X^{},J_0X^{},K_0X^{},M_0X^{},M_0I_0X^{},M_0J_0X^{},M_0K_0X^{}\}^{}.$$
So $`dim_X^{}=dim7`$. It follows $`dim𝒱=7`$.
Summarizing all of the above, we obtain our main classification result.
###### Main Theorem 2.6.
Let $`\pi :H_s^mB`$ be a semi-Riemannian submersion with totally geodesic fibres from a pseudo-hyperbolic space onto a Riemannian manifold. Then the semi-Riemannian submersion $`\pi `$ is equivalent to one of the following canonical semi-Riemannian submersions, given by examples 1)-3):
* $`H_1^{2k+1}H^k`$,
* $`H_3^{4k+3}H^k`$,
* $`H_7^{15}H^8(4)`$.
###### Proof.
The index of the pseudo-hyperbolic space cannot be $`s=0`$. Indeed, by lemma 1.4, for $`s=0`$, we get $`0g(A_XV,A_XV)=g(X,X)g(V,V)0`$ for every horizontal vector $`X`$ and for every vertical vector $`V`$. But this is not possible.
By , any semi-Riemannian submersion with totally geodesic fibres, from a pseudo-hyperbolic space of index $`1`$ onto a Riemannian manifold is equivalent to the canonical semi-Riemannian submersion $`H_1^{2k+1}H^k`$.
It remains to study the case $`s>1`$. By theorem 2.3 and proposition 2.5, any semi-Riemannian submersion with totally geodesic fibres from a pseudo-hyperbolic space of index $`s>1`$ onto a Riemannian manifold is one of the following types:
* $`H_3^{4k+3}H^k`$, or
* $`H_7^{15}H^8(4)`$, or
* $`H_7^{23}aH^2`$.
In order to prove that any two semi-Riemannian submersions in one of the categories (1) or (2) are equivalent we shall modify Ranjan’s argument (see ) to our situation. In the category (3), we shall prove there are no such semi-Riemannian submersions with totally geodesic fibres.
First, we shall prove the uniqueness in the case $`H_3^{4k+3}H^k`$. Let $`pH_3^{4k+3}`$ and let $`𝒰:𝒱_pEnd(_p)`$ the map given by $`𝒰(v)(x)=A_xv`$ for every $`v𝒱_p`$ and for every $`x_p`$. We denote $`𝒰(v)`$ by $`A^v`$. It is trivial to see that $`A^v`$ is skew-symmetric (i.e. $`g(A^vx,y)=g(x,A^vy)`$). The O’Neill’s equation $`g(A_xv,A_xv)=g(x,x)g(v,v)`$ becomes $`g(A^vx,A^vx)=g(x,x)g(v,v)`$. This implies $`g(A^vA^vx,x)=g(x,x)g(v,v)`$. Hence, by polarization in $`x`$, we have $`g(A^vA^vx,y)=g(x,y)g(v,v)`$ for every $`y_p`$. So $`A^vA^vx=g(v,v)x`$. Again by polarization we get $`A^vA^w+A^wA^v=2g(v,w)Id`$. Let $`\stackrel{~}{g}`$ be the Riemannian metric given by $`\stackrel{~}{g}(v,w)=g(v,w)`$ for every $`v`$, $`w𝒱_p`$. It follows $`A^vA^w+A^wA^v=2\stackrel{~}{g}(v,w)Id_{𝒱_p}`$. This is the condition which allows us to extend $`𝒰`$ to a representation of the Clifford algebra $`Cl(𝒱_p,\stackrel{~}{g}_p)`$ of $`𝒱_p`$. We also denote by $`𝒰`$ the extension of $`𝒰`$. Since $`dim𝒱_p=3`$ and $`\stackrel{~}{g}_p`$ is positive definite, $`Cl(𝒱_p,\stackrel{~}{g}_p)`$ has at most two types of irreducible representations. We notice that $`_p`$ is a $`Cl(𝒱_p,\stackrel{~}{g}_p)`$-module which splits in simple modules of dimension $`4`$. The next step is to show that any two such simple modules in decomposition of $`_p`$ are equivalent. Let $`\{v_1,v_2,v_3\}`$ be an orthonormal basis of $`(𝒱_p,\stackrel{~}{g}_p)`$. Since the affiliation of a simple $`Cl(𝒱_p,\stackrel{~}{g}_p)`$-module to one of the two possible types is decided by the action of $`v_1v_2v_3`$, it is sufficient to check that $`A^{v_1}A^{v_2}A^{v_3}=Id_{𝒱_p}`$.
Consider the function $`xg(A^{v_1}A^{v_2}A^{v_3}x,x)`$ defined on the unit sphere in $`_p`$. We have $`g(A^{v_1}A^{v_2}A^{v_3}x,x)=g(A^{v_2}A^{v_3}x,A^{v_1}x)=g(A_xA_{A_xv_3}v_2,v_1)`$. A straightforward computation shows that $`A_xA_{A_xv_3}v_2`$ is orthogonal to $`v_2`$ and $`v_3`$. Hence $`A_xA_{A_xv_3}v_2`$ is a multiple of $`v_1`$.
By polarization of the relation $`A_xA_xv=g(x,x)v`$, we get $`A_xA_y+A_yA_x=2g(x,y)Id`$ for every horizontal vectors $`x`$ and $`y`$ . In particular, we have $`A_xA_{A_xv_3}v_2=A_{A_xv_3}A_xv_2+2g(x,A_xv_3)v_2=A_{A_xv_3}A_xv_2`$.
Let $`S`$ be the vector subspace of $`_p`$ spanned by $`\{x,A_xv_1,A_xv_2,A_xv_3\}`$. By lemma 1.4, we get $`K^{}(\pi _{}x,\pi _{}A_xv_i)=4`$ for all $`i\{1,2,3\}`$. By geometry of $`H^n`$, there exists a unique totally geodesic hyperbolic line $`H^1`$ passing through $`\pi (p)`$ such that $`T_{\pi (p)}H^1=\pi _{}S`$. Notice that for every orthonormal vectors $`y,zT_{\pi (p)}H^1`$, $`K^{}(y,z)=4`$. In particular we have $`K^{}(\pi _{}A_xv_2,\pi _{}A_xv_3)=4`$. Hence $`g(A_{A_xv_3}A_xv_2,A_{A_xv_3}A_xv_2)=1`$. It follows that $`A_xA_{A_xv_3}v_2=\pm v_1`$. Hence $`g(A^{v_1}A^{v_2}A^{v_3}x,x)=\pm 1`$ for all unit vectors $`x`$. Since the function $`xg(A^{v_1}A^{v_2}A^{v_3}x,x)`$ defined on the unit sphere in $`_p`$ is continuous, we get either
* $`g(A^{v_1}A^{v_2}A^{v_3}x,x)=1`$ for any unit horizontal vector $`x`$, or
* $`g(A^{v_1}A^{v_2}A^{v_3}x,x)=1`$ for any unit horizontal vector $`x`$.
We may assume the case (i) holds.
If the case (ii) is happen, we replace the orthonormal basis $`\{v_1,v_2,v_3\}`$ of $`(𝒱_p,\stackrel{~}{g}_p)`$ with the orthonormal basis $`\{v_1,v_2,v_3\}`$. So for this new basis we are in the case (i).
Since $`A^{v_1}A^{v_2}A^{v_3}`$ is an isometry, we have $`g(A^{v_1}A^{v_2}A^{v_3}x,A^{v_1}A^{v_2}A^{v_3}x)g(x,x)=g(x,x)^2=1=g(A^{v_1}A^{v_2}A^{v_3}x,x)^2`$ for all unit horizontal vectors $`x`$.
So the Schwartz inequality for the scalar product $`g|__p`$
$$g(A^{v_1}A^{v_2}A^{v_3}x,x)^2g(A^{v_1}A^{v_2}A^{v_3}x,A^{v_1}A^{v_2}A^{v_3}x)g(x,x)$$
becomes equality. It follows that $`A^{v_1}A^{v_2}A^{v_3}x=\lambda x`$ for some $`\lambda `$. Because $`A^{v_1}A^{v_2}A^{v_3}`$ is an isometry and we assumed the case (i), it follows $`\lambda =1`$. We proved that $`A^{v_1}A^{v_2}A^{v_3}x=x`$ for all unit horizontal vectors $`x`$. Obviously, $`A^{v_1}A^{v_2}A^{v_3}x=x`$ for all $`x_p`$.
Let $`\pi ^{}:H_3^{4k+3}H^k`$ be another semi-Riemannian submersion with totally geodesic fibres. For an arbitrary chosen point $`qH_3^{4k+3}`$, we consider horizontal and vertical subspaces $`_q^{}`$ and $`𝒱_q^{}`$. Let $`\{v_1^{},v_2^{},v_3^{}\}`$ be an orthonormal basis in $`𝒱_q^{}`$ such that $`v_1^{}v_2^{}v_3^{}`$ acts on $`_q^{}`$ as $`Id`$. Let $`L_1:𝒱_q^{}𝒱_p`$ be the isometry given by $`L_1(v_i^{})=v_i`$ for all $`i\{1,2,3\}`$ and let $`Cl(L_1):Cl(𝒱_q^{})Cl(𝒱_p)`$ be the extension of $`L_1`$ to the Clifford algebras. The composition $`𝒰Cl(L_1):Cl(𝒱_q^{})End(_p)`$ makes $`_p`$ to be a $`Cl(𝒱_q^{})`$-module of dimension $`4k`$. Let $`_p=_1\mathrm{}_k`$ and $`_q^{}=_1^{}\mathrm{}_k^{}`$ be the decomposition of $`_p`$ and $`_q^{}`$ in simple $`Cl(𝒱_q^{})`$-modules, respectively. For each $`i`$ there is $`f_i:_i^{}_i`$ an equivalence of $`Cl(𝒱_q^{})`$-modules, which after a rescaling by a constant number is an isometry which preserves the O’Neill’s integrability tensors. Taking the direct sum of all these isometries, we obtain an isometry $`L_2:_q^{}_p`$ which preserves the O’Neill’s integrability tensors. Therefore $`L=L_1L_2:T_qH_3^{4k+3}T_pH_3^{4k+3}`$ is an isometry which maps $`_q^{}`$ onto $`_p`$ and $`A^{}`$ onto $`A`$. Since $`H_3^{4k+3}`$ is a simply connected complete symmetric space, there is an isometry $`f:H_3^{4k+3}H_3^{4k+3}`$ such that $`f(q)=p`$ and $`f_q=L`$ (see corollary 2.3.14 in ). Therefore, by theorem 1.6, we get $`\pi `$ and $`\pi ^{}`$ are equivalent.
Now, we shall prove that any two semi-Riemannian submersions $`\pi ,\pi ^{}:H_7^{15}H^8(4)`$ with totally geodesic fibres are equivalent. The proof is analogous to the case $`(1)`$, but it is easier.
Let $`p,qH_7^{15}`$ and let $`_p`$, $`𝒱_p`$ be the horizontal and vertical subspaces in $`T_pH_7^{15}`$ for $`\pi `$, let $`_q^{}`$, $`𝒱_q^{}`$ be the horizontal and vertical subspaces in $`T_qH_7^{15}`$ for $`\pi ^{}`$. Let $`\{v_1,\mathrm{},v_7\}`$ be an orthonormal basis of $`(𝒱_p,\stackrel{~}{g}_p)`$ and $`\{v_1^{},\mathrm{},v_7^{}\}`$ be an orthonormal basis of $`(𝒱_q^{},\stackrel{~}{g}_q)`$ such that $`A^{v_1}A^{v_2}\mathrm{}A^{v_7}=Id`$ and $`A^{v_1^{}}A^{v_2^{}}\mathrm{}A^{v_7^{}}=Id`$. Since $`dim𝒱_p=7`$, the irreducible $`Cl(𝒱_p,\stackrel{~}{g}_p)`$-modules are $`8`$-dimensional. Since $`dim_p=8`$, we get $`_p`$ is simple. Because $`A^{v_1}A^{v_2}\mathrm{}A^{v_7}=Id`$ and $`A^{v_1^{}}A^{v_2^{}}\mathrm{}A^{v_7^{}}=Id`$ we get $`_q^{}`$ and $`_p`$ are $`Cl(𝒱_q^{})`$-modules equivalent. Analogously to the case (1), we can construct an isometry $`L=L_1L_2:T_qH_7^{15}T_pH_7^{15}`$, which map $`_q^{}`$ onto $`_p`$ and $`A^{}`$ onto $`A`$. This produces an isometry $`f:H_7^{15}H_7^{15}`$ such that $`f(q)=p`$ and $`f_q=L`$ (see corollary 2.3.14 in ). Again by theorem 1.6, we get $`\pi `$ and $`\pi ^{}`$ are equivalent.
Now, we prove that there are no $`\pi :H_7^{23}aH^2`$ semi-Riemannian submersions with totally geodesic fibres. The proof is analogous to that of Ranjan (see proposition 5.1 in ).
$`_p`$ becomes a $`Cl(𝒱_p)`$-module by considering the extension of the map $`𝒰:𝒱_p\mathrm{End}(_p)`$, $`𝒰(V)(X)=A_XV`$ to the Clifford algebra $`Cl(𝒱_p)`$. Here $`Cl(𝒱_p,\stackrel{~}{g}_p)`$ denotes the Clifford algebra of $`(𝒱_p,\stackrel{~}{g}_p)`$, $`\stackrel{~}{g}(U,V)=g(U,V)`$ for every $`U`$, $`V𝒱_p`$. Since $`\stackrel{~}{g}_p`$ is positive definite, we have $`Cl(𝒱_p)(8)(8)`$. Hence, $`_p`$ splits into two $`8`$-dimensional irreducible $`Cl(𝒱_p)`$-modules. Since the induced metrics on fibres are negative definite we get $`\pi ^1(aH^1)`$ is totally geodesic in $`H_7^{23}`$ and isometric to $`H_7^{15}`$, by theorem 2.5 in . Here $`aH^1`$ denotes the Cayley hyperbolic line through $`\pi _{}X`$; We choose $`S`$ be the horizontal space of the restricted submersion $`\stackrel{~}{\pi }:H_7^{15}aH^1=H^8(4)`$. So for every $`X_p`$, $`g(X,X)0`$ we find an irreducible $`Cl(𝒱_p)`$-submodule $`S`$ of $`_p`$ passing through $`X`$. Since $`dim𝒱_p4`$, we get a contradiction.
R. Escobales classified Riemannian submersions from complex projective spaces under the assumption that the fibres are connected, complex, totally geodesic submanifolds. Using the main theorem 2.6, we obtain a classification of semi-Riemannian submersions from a complex pseudo-hyperbolic space onto a Riemannian manifold under the assumption that the fibres are connected, complex, totally geodesic submanifolds.
###### Proposition 2.7.
If $`\pi :H_s^mB^n`$ is a semi-Riemannian submersion with complex, connected, totally geodesic fibres then $`2m=n+2s`$, the induced metrics on fibres are negative definite and the fibres are diffeomorphic to $`P^s`$.
###### Proof.
We denote by $`J`$ the natural almost complex structure on $`H_s^m`$. By lemma 1.4, we have
a) $`\widehat{R}(U,V,U,V)=R(U,V,U,V)=(g(U,U)g(V,V)g(U,V)^2+3g(U,JV)^2)`$. Hence the fibres have constant holomorphic curvature $`4`$.
b) $`g(A_XU,A_XU)=(g(U,U)g(X,X)+3g(X,JU)^2)=g(U,U)g(X,X)`$, since the fibres are complex submanifolds. We obtain $`g(U,U)0`$ for every vertical vector field $`U`$.
c) $`R^{}(\pi _{}X,\pi _{}Y,\pi _{}X,\pi _{}Y)=R(X,Y,X,Y)+3g(A_XY,A_XY)=`$
$`(g(X,X)g(Y,Y)g(X,Y)^2+3g(X,JY)^2)+3g(A_XY,A_XY)0`$, since the induced metrics on fibres are negative definite. By proposition 2.2, it follows that the fibres are simply connected. Since the fibres are complete, simply connected, complex manifolds with constant holomorphic curvature $`4`$, we have that the fibres are isometric to $`H_s^s`$. ∎
###### Theorem 2.8.
If $`\pi :H_s^mB`$ is a semi-Riemannian submersion with connected, complex, totally geodesic fibres from a complex pseudo-hyperbolic space, then $`\pi `$ is, up to equivalence, the canonical semi-Riemannian submersion given by example $`4`$
$$H_1^{2k+1}H^k.$$
###### Proof.
Let $`\theta :H_{2s+1}^{2m+1}H_s^m`$ be the canonical semi-Riemannian submersion with totally geodesic fibres given in the definition 3 (see also or ). We have $`\stackrel{~}{\pi }=\pi \theta :H_{2s+1}^{2m+1}B`$ is a semi-Riemannian submersion with totally geodesic fibres, by theorem 2.5 in . Since the dimension of fibres of $`\stackrel{~}{\pi }`$ is greater than or equal to 2, we get, by main theorem 2.6, the following possible situations:
* $`m=2k+1`$, $`2s+1=3`$ and $`B`$ is isometric to $`H^k`$ or
* $`m=7`$, $`2s+1=7`$ and $`B`$ is isometric to $`H^8(4)`$.
First, we shall prove that any two semi-Riemannian submersions $`\pi ,\pi ^{}:H_1^{2k+1}H^k`$ with connected, complex, totally geodesic fibres are equivalent.
By proof of proposition 2.7, we have $`g(A_XU,A_XU)=g(U,U)g(X,X)`$. Let $`p,qH_1^{2k+1}`$ By proof of the main theorem, this implies $`A^vA^w+A^wA^v=2\stackrel{~}{g}(v,w)Id`$. The extension of $`𝒰:𝒱_pEnd(_p)`$ constructed in proof of the main theorem, to the Clifford algebra $`Cl(𝒱_p,\stackrel{~}{g}_p)`$ makes $`_p`$ a $`Cl(𝒱_p,\stackrel{~}{g}_p)`$-module which splits in $`k`$ irreducible modules of dimension $`4`$. By classification of irreducible representation for case $`dim𝒱_p=2`$ and $`\stackrel{~}{g}_p`$ positive definite, we have any two such irreducible $`Cl(𝒱_p,\stackrel{~}{g}_p)`$-modules are equivalent. Like in proof of the main theorem, we may construct an isometry $`L=L_1L_2:T_qH_1^{2k+1}T_pH_1^{2k+1}`$, which maps $`_q^{}`$ onto $`_p`$ and $`A^{}`$ onto $`A`$. This produces an isometry $`f:H_1^{2k+1}H_1^{2k+1}`$ with $`f(q)=p`$ and $`f_q=L`$ (see corollary 2.3.14 in ). Again by theorem 1.6, we get $`\pi `$ and $`\pi ^{}`$ are equivalent.
For the case $`ii)`$ we shall obtain that there are no $`\pi :H_3^7H^8(4)`$ semi-Riemannian submersions with complex, connected, totally geodesic fibres.
###### Proposition 2.9.
There are no $`\pi :H_3^7H^8(4)`$ semi-Riemannian submersions with connected, complex, totally geodesic fibres.
###### Proof.
The proof is based on Ranjan’s argument (see proof of main theorem in ). Here, we show how to modify Ranjan’s argument to our different situation.
Suppose there is $`\pi :H_3^7H^8(4)`$ a semi-Riemannian submersion with complex, connected, totally geodesic fibres. By main theorem 2.6, $`\stackrel{~}{\pi }=\pi \theta :H_7^{15}H^8(4)`$ is equivalent to the canonical semi-Riemannian submersion $`Spin(1,8)/Spin(7)Spin(1,8)/Spin(8)`$ given by example 3.
Let $`\sigma :Spin(1,8)SO(8,8)`$ be the spin representation of $`Spin(1,8)`$. $`Spin(1,8)`$ acts on $`H^8(4)`$ via double covering map $`Spin(1,8)SO(1,8)`$ and transitively on $`H_7^{15}_8^{16}`$. We denote by $`Cl^0(_1^9)`$ the even component of Clifford algebra $`Cl(_1^9)`$. Notice that $`Cl^0(_1^9)M(16,)`$, $`Cl(_1^9)M(16,)M(16,)`$ and the volume element $`\omega `$ in $`Cl(_1^9)`$ satisfies $`\omega ^2=1`$ (see ).
For any $`bH^8(4)`$, let $`G_b`$ be the isotropy group of $`b`$ in $`Spin(1,8)`$. If we restrict $`\sigma |_{G_b}`$ then $`\sigma |_{G_b}`$ breaks $`_8^{16}`$ into two $`\frac{1}{2}`$-spin representations. We will denote them by $`_\pm ^8`$. Hence $`_+^8H_7^{15}=\stackrel{~}{\pi }^1(b)`$. Let $`b^{}=\{x_1^9|<x,b>=0\}`$. We have $`Cl(b^{})Spin(1,8)=G_b`$, $`dimb^{}=8`$ and the following diagram is commutative
$$\begin{array}{ccc}G_b& & Cl^0(b^{})\\ & & & & \\ Spin(1,8)& & Cl^0(_1^9),\end{array}$$
where all arrows are standard inclusions. Let $`\{e_1,\mathrm{},e_8\}`$ be an orientated basis of $`b^{}`$. Then $`z^{}=e_1\mathrm{}e_8`$ lies in the centre of $`Cl^0(b^{})`$ and $`z^{}`$ acts by $`Id`$ on $`_+^8`$ and $`Id`$ on $`_{}^8`$. We have $`Cl(\sigma )(z^{})=\pm 1`$ on $`_\pm ^8`$.
Since $`_+^8H_7^{15}=\stackrel{~}{\pi }^1(b)`$, $`_+^8`$ is invariant under $`J`$ and so is $`_{}^8`$. Here $`J`$ denotes the natural complex structure on $`^{16}=^8`$. Hence $`Cl(\sigma )(z^{})`$ commutes with $`J`$. Let $`zCl(_1^9)`$ be the generator of the center of $`Cl(_1^9)`$. We have either $`z=e_1e_2\mathrm{}e_8b`$ or $`zb=e_1e_2\mathrm{}e_8=z^{}`$. Therefore $`Cl(\sigma )(zb)`$ commutes with $`J`$ for every $`bH^8(4)`$ and hence for every $`b^9`$.
Consider the linear map $`\alpha :^9M(16,)`$ given by $`bCl(\sigma )(zb)`$. It has the following properties:
* It factors through $`M(8,)M(16,)`$
* $`[Cl(\sigma )(zb)]^2=Cl(\sigma )((zb)^2)=Cl(\sigma )(|b|^2)=|b|^2Id.`$
Hence $`\alpha `$ extends to a homomorphism $`Cl(\alpha ):Cl(_1^9)M(8,)`$. But $`Cl(_1^9)M(16,)M(16,)`$ (see ). So the above homomorphism is impossible to exist. We get the required contradiction. ∎
This ends the proof of theorem 2.8.∎
###### Acknowledgments.
We would like to thank Dmitri Alekseevsky for the remarks and comments made to an earlier version of this work. The second author thanks to Stefano Marchiafava for useful discussions on this topic with the occasion of his visit to Rome in autumn of 1999.
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# Trapped condensates of atoms with dipole interactions
## I introduction
The recent success in atomic Bose-Einstein condensation (BEC) has stimulated great research activities into trapped quantum gases . To a remarkable degree, a single condensate wave function of the macroscopically occupied ground state, described by the nonlinear Schrödinger equation (NLSE) , captures all essential features of its coherence properties . In fact, one of the key diagnosis features for BEC, the reversed aspect ratio of a free expanding condensate, is described purely by the condensate wave function . In the standard treatment for the condensate wave function of interacting atoms, realistic inter-atomic potential $`V(\stackrel{}{R})`$ is often not directly used. Instead, a contact pseudo potential, $`u_0\delta (\stackrel{}{R})`$, obtained under the so-called shape independent approximation (SIA) is used. Such an idealization results in tremendous simplification, yet to date, SIA has worked remarkably well as verified by both theoretical calculations and experimental observations .
Currently available degenerate quantum gases are cold and dilute, the interaction is therefore dominated by s-wave collisions, described by a single atomic parameter: $`a_{\mathrm{sc}}`$, the s-wave scattering length, if the inter-atomic potential is isotropic and short ranged (decaying fast than $`1/R^3`$ asymptotically). The complete scattering amplitude is then isotropic and energy-independent, given by $`f(\stackrel{}{k},\stackrel{}{k}^{})=4\pi a_{\mathrm{sc}}`$ for collisions of incoming momentum $`\stackrel{}{k}`$ state scattering into $`\stackrel{}{k}^{}`$.
One of the attractive features of atomic degenerate gases lies at effective means for control of the atom-atom interaction . Indeed, very recently several groups have successfully implemented Feschbach resonance , thus enabling a control knob on $`a_{\mathrm{sc}}`$ through the changing of an external magnetic field. Other physical mechanisms also exist for modifying atom-atom interactions, e.g. the shape resonance as proposed in . In an external electric field, inter atomic potential is modified by the addition of an anisotropic (induced) dipole interaction.
Although anisotropically interacting fermi system has been an important area of study, e.g. liquid <sup>3</sup>He and d-wave paired high $`T_c`$ superconductors . Its bosonic counterpart has not been studied in great detail. In particular, we are not aware of any systematic approach for constructing an anisotropic pseudo potential .
For bosonic systems, another related topic is the condensate stability. Under the SIA, the scattering length takes a positive or negative value, corresponding to repulsive or attractive interactions. When $`a_{\mathrm{sc}}<0`$ occurs, self-interaction leads to a collapse of BEC in dimensions higher than 1 , thus the resulting condensate is limited by a critical number of particles . Anisotropic dipole interactions, on the other hand, are more complicated as both attractive and repulsive interactions arise along different directions. We note that several recent investigations have studied efforts of non-local interactions on condensate stability .
In this paper, we study properties of trapped BEC of atoms with dipole interactions , arising from either external electric field (induced) or permanent magnetic moments . We propose a practical method for constructing anisotropic pseudo potentials that can also be extended for investigation of polar molecular BEC . This paper is organized as following. We first briefly review the usual pseudo-potential approximation under the SIA. In Sec. II we describe and justify in detail a procedure for constructing effective low energy pseudo-potentials of anisotropic interactions. In Sec. III we provide our numerical procedure for solving the NLSE with anisotropic dipole interactions. Particular emphasis is put on the careful treatment of the singular origin of dipole interactions. We also present and discuss results from selected numerical calculations. To explain the stability region as well as the interesting aspect ratios observed from our numerical calculations, we perform in Sec. IV an analytic time dependent variational calculation. We compare the results obtained with direct numerical solutions of NLSE. Finally we conclude.
## II formulation
For $`N`$ trapped spinless bosonic atoms in a potential $`V_t(\stackrel{}{r})`$, the second quantized Hamiltonian is given by
$``$ $`=`$ $`{\displaystyle 𝑑\stackrel{}{r}\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})\left[\frac{\mathrm{}^2}{2M}^2+V_t(\stackrel{}{r})\mu \right]\widehat{\mathrm{\Psi }}(\stackrel{}{r})}`$ (1)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑\stackrel{}{r}𝑑\stackrel{}{r}^{}\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r}^{})V(\stackrel{}{r}\stackrel{}{r}^{})\widehat{\mathrm{\Psi }}(\stackrel{}{r}^{})\widehat{\mathrm{\Psi }}(\stackrel{}{r})},`$ (2)
where $`\widehat{\mathrm{\Psi }}(\stackrel{}{r})`$ and $`\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})`$ are atomic (bosonic) annihilation and creation fields. The chemical potential $`\mu `$ guarantees the atomic number $`\widehat{N}=𝑑\stackrel{}{r}\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})\widehat{\mathrm{\Psi }}(\stackrel{}{r})`$ conservation.
The bare potential $`V(\stackrel{}{R})`$ in (2) needs to be renormalized for a meaningful perturbation calculation . The usual treatment is based an effective interaction obtained by a resummation of certain classes of interaction diagrams . Physically the SIA can be viewed as a valid low energy and low density renormalization scheme, one simply replaces the bare potential $`V(\stackrel{}{R})`$ by a pseudo potential $`u_0\delta (\stackrel{}{R})`$ whose first order Born scattering amplitude reproduces the complete scattering amplitude ($`a_{\mathrm{sc}}`$). This gives $`u_0=4\pi \mathrm{}^2a_{\mathrm{sc}}/M`$.
When an electric field is introduced along the positive z axis, an additional term
$`V_E(\stackrel{}{R})`$ $`=u_2{\displaystyle \frac{Y_{20}(\widehat{R})}{R^3}},`$ (3)
appears in the atom-atom interaction , where $`u_2=4\sqrt{(\pi /5)}\alpha (0)\alpha ^{}(0)^2`$. $`\alpha (0)`$ is the atomic polarizability, and $``$ denotes the electric field strength. As was shown before , this modification results in a completely new form for the low-energy scattering amplitude
$`f(\stackrel{}{k},\stackrel{}{k}^{})|_{k=k^{}0}=4\pi {\displaystyle \underset{lm,l^{}m^{}}{}}t_{lm}^{l^{}m^{}}()Y_{lm}^{}(\widehat{k})Y_{l^{}m^{}}(\widehat{k}^{}),`$ (4)
with $`t_{lm}^{l^{}m^{}}()`$ the reduced T-matrix elements. They are all energy independent and act as generalized scattering lengths. The anisotropic $`V_E`$ causes the dependence on both incident and scattered directions: $`\widehat{k}`$ and $`\widehat{k}^{}=\widehat{R}`$.
We therefore propose a general (energy-independent) anisotropic pseudo potential constructed according to
$`V_{\mathrm{eff}}(\stackrel{}{R})=u_0\delta (\stackrel{}{R})+{\displaystyle \underset{l_1>0,m_1}{}}\gamma _{l_1m_1}{\displaystyle \frac{Y_{l_1m_1}(\widehat{R})}{R^3}},`$ (5)
whose first Born amplitude is
$`f_{\mathrm{Born}}(\stackrel{}{k},\stackrel{}{k}^{})=(4\pi )^2a_{\mathrm{sc}}Y_{00}^{}(\widehat{k})Y_{00}(\widehat{k}^{}){\displaystyle \frac{M}{4\pi \mathrm{}^2}}{\displaystyle \underset{l_1m_1}{}}\gamma _{l_1m_1}(4\pi )^2{\displaystyle \underset{lm}{}}{\displaystyle \underset{l^{}m^{}}{}}𝒯_{lm}^{l^{}m^{}}(l_1,m_1)Y_{lm}^{}(\widehat{k})Y_{l^{}m^{}}(\widehat{k}^{}),`$ (6)
with $`𝒯_{lm}^{l^{}m^{}}(l_1,m_1)=(i)^{l+l^{}}_l^l^{}I_{lm}^{l^{}m^{}}(l_1,m_1).`$ Both
$`I_{lm}^{l^{}m^{}}(l_1m_1)`$ $`=Y_{l^{}m^{}}|Y_{l_1m_1}|Y_{lm}`$ (12)
$`=(1)^m\sqrt{{\displaystyle \frac{(2l+1)(2l^{}+1)(2l_1+1)}{4\pi }}}\left(\begin{array}{ccc}l& l^{}& l_1\\ m& m^{}& m_1\end{array}\right)\left(\begin{array}{ccc}l& l^{}& l_1\\ 0& 0& 0\end{array}\right),`$
and
$`_l^l^{}`$ $`={\displaystyle _0^{\mathrm{}}}𝑑R{\displaystyle \frac{1}{R}}j_l(kR)j_l^{}(k^{}R)`$ (14)
$`={\displaystyle \frac{\pi }{8}}\eta ^l{\displaystyle \frac{\mathrm{\Gamma }(\frac{l+l^{}}{2})}{\mathrm{\Gamma }(\frac{3+l^{}l}{2})\mathrm{\Gamma }(l+\frac{3}{2})}}_2F_1({\displaystyle \frac{1l^{}+l}{2}},{\displaystyle \frac{l+l^{}}{2}},l+{\displaystyle \frac{3}{2}};\eta ^2),`$
can be computed analytically. The $`1/R^3`$ form in Eq. (5) assures all $`_l^l^{}`$ to be $`k=k^{}`$ independent (easily seen by a change of variable to $`x=kR`$ in the integral). Putting
$`f_{\mathrm{Born}}(\stackrel{}{k},\stackrel{}{k}^{})=f(\stackrel{}{k},\stackrel{}{k}^{}),`$ (15)
i.e. requiring the Bohn amplitude from the pseudo potential Eq. (5) to be the same as the numerically computed value $`f(\stackrel{}{k},\stackrel{}{k}^{})`$, one can solve for all $`\gamma _{l_1m_1}()`$ from known $`t_{lm}^{l^{}m^{}}()`$ . This reduces to a set of (under determined) linear equations
$`{\displaystyle \frac{M}{4\pi \mathrm{}^2}}{\displaystyle \underset{l_1m_1}{}}\gamma _{l_1m_1}(4\pi )𝒯_{lm}^{l^{}m^{}}(l_1,m_1)t_{lm}^{l^{}m^{}},`$ (16)
for all ($`lm`$) and ($`l^{}m^{}`$) with $`l,l^{}0`$, and separately $`a_{\mathrm{sc}}()=t_{00}^{00}()`$.
Considerable simplification arises further for bosons (fermions) as only even (odd) $`(l,l^{})`$ terms are needed to match in (16). Figure 1 displays the result of field dependent $`a_{\mathrm{sc}}()`$ <sup>41</sup>K atoms in the triplet electron spin state. Note the spikes of shape resonances. The Born amplitude for the dipole term $`V_E`$ is
$`f_{\mathrm{Born}}(\stackrel{}{k},\stackrel{}{k}^{})`$ $`=u_2{\displaystyle \frac{M}{4\pi \mathrm{}^2}}(4\pi )^2𝒯_{00}^{20}{\displaystyle \underset{lm,l^{}m^{}}{}}\overline{𝒯}_{lm}^{l^{}m^{}}Y_{lm}^{}(\widehat{k})Y_{l^{}m^{}}(\widehat{k}^{}),`$ (17)
with $`𝒯_{00}^{20}=0.023508`$. $`\overline{𝒯}_{lm}^{l^{}m^{}}=𝒯_{lm}^{l^{}m^{}}(2,0)/𝒯_{00}^{20}`$ are independent of electric field $``$ within the perturbative Born approximation as tabulated below.
We find that away from regions of shape resonances to be discussed elsewhere , all numerically computed $`t_{lm}^{l^{}m^{}}()`$ values, large enough to justify their inclusions, are actually all proportional to $`^2`$. Thus, we could rewrite $`f(\stackrel{}{k},\stackrel{}{k}^{})`$ as
$`f(\stackrel{}{k},\stackrel{}{k}^{})=(4\pi )t_{00}^{20}(){\displaystyle \underset{lm,l^{}m^{}}{}}\overline{t}_{lm}^{l^{}m^{}}Y_{lm}^{}(\widehat{k})Y_{l^{}m^{}}(\widehat{k}^{}),`$ (18)
with scaled quantities $`\overline{t}_{lm}^{l^{}m^{}}=t_{lm}^{l^{}m^{}}()/t_{00}^{20}()`$ now all being constants. We have since computed (numerically) for several alkali metal isotopes, our results for $`\overline{t}_{lm}^{l^{}m^{}}`$ are tabulated below.
The agreement between the first order Born approximation and the multi-channel scattering calculations is remarkable. We estimate the numerical scattering results to be accurate to a few percent (except for <sup>39</sup>K), independent of atoms being bosons (even $`l,l^{}`$) or fermions (odd $`l,l^{}`$). Only bosonic results are being considered in this paper. This is displayed by noticing the agreement between Table I ($``$ a few per cent) with Table II-VI. This interesting observation applies for all bosonic alkali triplet states: <sup>7</sup>Li, <sup>39,41</sup>K, and <sup>85,87</sup>Rb, for up to a field strength of $`3\times 10^6`$ (V/cm) computed by us. Physically, this implies the effect of $`V_E`$ is perturbative when $``$ remains small (in a.u.). For the convenience of further discussions, we tabulate polarizabilities of selected atoms in Table VII.
What is remarkable is the fact that $`𝒯_{00}^{20}()`$ and $`t_{00}^{20}()`$ also agree in absolute values except for a slight difference (1-6%). They are calculate below and tabulated in VIII.
$`u_2{\displaystyle \frac{M}{4\pi \mathrm{}^2}}(4\pi )^2𝒯_{00}^{20}`$ $`=16\pi \sqrt{{\displaystyle \frac{\pi }{5}}}𝒯_{00}^{20}\alpha ^2^2{\displaystyle \frac{M}{\mathrm{}^2}}`$ (20)
$`=1718\times \overline{M}\overline{\alpha }^2\overline{}^2a_0,`$
where all overlined quantities are in atomic units (a.u.).
An important parameter in our discussion is the ratio between $`u_2`$ and $`u_0`$. Since the results from first order Born approximation and the multichannel calculations are about the same, one can write this ratio as
$`c()={\displaystyle \frac{u_2}{u_0}}=\gamma ^2,`$ (21)
where $`\gamma 1.748\times 10^{17}\overline{\alpha }^2\overline{M}/\overline{a}_{\mathrm{sc}}`$ and $``$ is in unit of V/cm. The $`\gamma `$ values for selected atoma are tabulated in Table IX.
Most atoms possess permanent magnetic dipoles. With alkali metals, the magnetic dipole mainly originates from valance electron spin, typically measured in units of Bohr magneton. It is therefore interesting to compare electric dipole interactions with magnetic dipole interactions. The dipole interaction strength between atoms of a permanent magnetic dipole $`\mu `$ is
$`\mu ^2`$ $`=\overline{\mu }^2\mu _B^2`$ (23)
$`=1.331\times 10^5\overline{\mu }^2e^2a_0^2,`$
with $`\mu _B`$ the unit of Bohr magneton. For induced electric dipoles, the interaction strength is
$`\alpha ^2^2=\overline{\alpha }^2\overline{}^2e^2a_0^2.`$ (24)
A typical heavy alkali atom has $`\overline{\alpha }200`$, thus for which a 1 ($`\mu _B`$) magnetic moment corresponds to an effective electric field of $`3.3\times 10^5`$ (a.u.), or $`1.7\times 10^5`$ (V/cm). Atoms with large magnetic moments effectively simulates induced dipole interactions at a high value of equivalent electric field .
As can be concluded from comparing listed data in all tables, we can approximate Eq. (5) by keeping only the $`l_1=2`$, $`m_1=0`$ term to achieve a satisfactory level of accuracy. Thus away from shape resonances we take
$`V_{\mathrm{eff}}(\stackrel{}{R})=u_0\delta (\stackrel{}{R})+u_2Y_{20}(\widehat{R})/R^3,`$ (25)
where $`u_0=\frac{4\pi \mathrm{}^2}{M}a_{\mathrm{sc}}()`$ and $`u_2=c()u_0`$. The Hamiltonian (2) then becomes
$``$ $`=`$ $`{\displaystyle 𝑑\stackrel{}{r}\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})\left[\frac{\mathrm{}^2}{2M}^2+V_t(\stackrel{}{r})\mu \right]\widehat{\mathrm{\Psi }}(\stackrel{}{r})}`$ (26)
$`+`$ $`{\displaystyle \frac{u_0}{2}}{\displaystyle 𝑑\stackrel{}{r}\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})\widehat{\mathrm{\Psi }}(\stackrel{}{r})\widehat{\mathrm{\Psi }}(\stackrel{}{r})}`$ (27)
$`+`$ $`{\displaystyle \frac{u_2}{2}}{\displaystyle 𝑑\stackrel{}{r}𝑑\stackrel{}{r}^{}\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r})\widehat{\mathrm{\Psi }}^{}(\stackrel{}{r}^{})\frac{Y_{20}(\widehat{R})}{R^3}\widehat{\mathrm{\Psi }}(\stackrel{}{r}^{})\widehat{\mathrm{\Psi }}(\stackrel{}{r})},`$ (28)
with $`\stackrel{}{R}=\stackrel{}{r}\stackrel{}{r}^{}`$. The Heisenberg equation for $`\widehat{\mathrm{\Psi }}(\stackrel{}{r},t)`$ becomes nonlocal. At zero temperature the condensate wave function $`\psi (\stackrel{}{r},t)=\widehat{\mathrm{\Psi }}(\stackrel{}{r},t)`$ obeys the NLSE
$`i\mathrm{}{\displaystyle \frac{d}{dt}}\psi (\stackrel{}{r},t)=\left[{\displaystyle \frac{\mathrm{}^2}{2M}}^2+V_t(\stackrel{}{r})\mu +u_0|\psi (\stackrel{}{r},t)|^2+u_2{\displaystyle 𝑑\stackrel{}{r}^{}\frac{Y_{20}(\widehat{R})}{R^3}|\psi (\stackrel{}{r}^{},t)|^2}\right]\psi (\stackrel{}{r},t),`$ (29)
with $`\psi (\stackrel{}{r},t)`$ normalized to $`N`$ (the number of the atom in the condensate).
## III Numerical Studies
In this section we discuss the ground state properties of trapped condensates based on numerical solutions of NLSE (29). We start with a detailed analysis of our numerical procedure for handling the non-local dipole interaction .
### A The numerical procedure
We use steepest descent through a propagation of Eq. (29) in imaginary time ($`it`$) to find its ground state wave function. With an axial symmetric harmonic trap
$`V_t(\stackrel{}{r})={\displaystyle \frac{1}{2}}M\nu ^2(\lambda _x^2x^2+\lambda _y^2y^2+\lambda _z^2z^2),`$ (30)
we rescale Eq. (29) by introducing dimensionless units for length ($`a_\mathrm{t}=\sqrt{\mathrm{}/M\nu }`$), energy ($`\mathrm{}\nu `$), time ($`2i/\nu `$), and wave function ($`\sqrt{N/a_\mathrm{t}^3}`$). We than obtain
$`{\displaystyle \frac{d}{dt}}\psi (\stackrel{}{r},t)=\widehat{H}\psi (\stackrel{}{r},t),`$ (31)
with
$`\widehat{H}=^2+(\lambda _x^2x^2+\lambda _y^2y^2+\lambda _z^2z^2)2\mu +2(2\pi )^{3/2}P\left[|\psi (\stackrel{}{r},t)|^2c(){\displaystyle 𝑑\stackrel{}{r}^{}\frac{Y_{20}(\widehat{R})}{R^3}|\psi (\stackrel{}{r}^{},t)|^2}\right],`$ (32)
where $`P=\sqrt{2/\pi }Na_{\mathrm{sc}}/a_\mathrm{t}`$ and $`\psi (\stackrel{}{r},t)`$ is normalized to 1.
The ground state is found be starting with an arbitrary random wave funtion, and propagating Eq. (31) in $`t`$ until it’s stable (apart from its decreasing norm). In practice we chose an appropriate time step $`\mathrm{\Delta }t`$ and iterates the Eq. (31) according to
$`\psi (\stackrel{}{r},t+\mathrm{\Delta }t)=\psi (\stackrel{}{r},t)(\mathrm{\Delta }t)\widehat{H}\psi (\stackrel{}{r},t).`$ (33)
We renormalize $`\psi `$ to 1 after each iteration and adjust $`\mathrm{\Delta }t`$ to control the rate of convergence.
For a cylindrical symmetric trap ($`\lambda _x=\lambda _y=1,\lambda _z=\lambda `$), the ground state wave function also possesses the cylindrical symmetry. Therefore the non-local term simplifies to
$`{\displaystyle 𝑑\stackrel{}{r}^{}|\psi (\rho ^{},z^{})|^2\frac{Y_{20}(\widehat{R})}{R^3}}={\displaystyle 𝑑z^{}𝑑\rho ^{}𝒦(\rho ^{},\rho ,z^{}z)|\psi (\rho ^{},z^{})|^2},`$ (34)
with a kernel
$`𝒦(\rho ^{},\rho ,z^{}z)`$ $`=`$ $`\sqrt{{\displaystyle \frac{5}{\pi }}}{\displaystyle \frac{\rho ^{}}{[(\rho \rho ^{})^2+(z^{}z)^2]^2[(\rho +\rho ^{})^2+(z^{}z)^2]^{\frac{3}{2}}}}`$ (35)
$`([(\rho ^2\rho ^2)^22(\rho ^2+\rho ^2)(z^{}z)^23(z^{}z)^4]\mathrm{E}\left[{\displaystyle \frac{4\rho \rho ^{}}{(\rho +\rho ^{})^2+(z^{}z)^2}}\right]`$ (36)
$`+`$ $`[(\rho \rho ^{})^2+(z^{}z)^2](z^{}z)^2\mathrm{K}\left[{\displaystyle \frac{4\rho \rho ^{}}{(\rho +\rho ^{})^2+(z^{}z)^2}}\right]),`$ (37)
where $`\mathrm{E}[.]`$ and $`\mathrm{K}[.]`$ are standard Elliptical integrals.
We discretize the $`(\rho ,z)`$ plane into a two-dimensional grid of points such that wave function values at each point becomes a matrix. At each time step the matrix elements are updated according to (33). The derivatives in the Hamiltonian are evaluated by means of finite-difference methods. Typically, the ground state can be sufficiently well described using a grid of $`100\times 200`$ points in the range $`0<\rho <5`$ and $`5<z<5`$.
At first sight, one may naively underestimate the complication of Eq. (31) due to the non-local interaction term in (32). Several other groups have addressed non-local interactions previously . There is however, a significant numerical challenge with the dipole interaction, which is singular at the origin. To accurately represent its detailed structure, an enormously large grid is needed. Although a Fourier transform into momentum space could simplify the convolution operation of the nonlocal term. We found it hard to completely avoid the effort of the singularity this way by going to a momentum representation with a limited coarse grid . Physically, this singularity implies the presence of two different length scales for Eq. (31). We thus developed a numerical procedure with two overlaying grids: a coarse grid for the relatively smooth wave function and a much finer grid for computing the non-local dipole interaction kernel.
The kernel $`𝒦(\rho ^{},\rho ,z^{}z)`$ is divergent at $`\stackrel{}{r}=\stackrel{}{r}^{}`$, we thus define a cut-off radius $`R_c`$ such that $`𝒦(\rho ^{},\rho ,z^{}z)=0`$ whenever $`|\stackrel{}{r}\stackrel{}{r}^{}|<R_c`$. This cut off is chosen to be small enough that negligible errors result from the numerically represented kernel. Typically $`R_c50`$($`a_0`$) taken, which is much smaller than the wave function grid size. The rapid varying kernel $`𝒦(\rho ^{},\rho ,z^{}z)`$ is treated with a finer grid. Instead of directly integrating over $`𝒦(\rho ^{},\rho ,z^{}z)|\psi (\rho ^{},z^{})|^2`$ on the wave function grid, we first integrate the kernel separately over the fine grid around each of the wave function grid point. Such an integration is numerically intensive, but only needs to be performed once for each of the wave function grid point as the kernel is determined by the geometry of system. The integrated kernel values on the wave function grid remain the same for different traps and different number of atoms. Finally the non-local term is approximated by integrating over the wave function grid using the product of integrated kernel values and the wave function.
For a homogeneous distribution of aligned dipoles, the mean dipole interaction vanishes as $`Y_{20}(\widehat{R})`$ averages to zero upon integration over $`d\widehat{r}`$ or $`d\widehat{r}^{}`$. This property is maintained for our kernel Eq.(37) even though we have integrated over $`(\varphi \varphi ^{})`$ first. We have verified this by noting that the integration of $`𝒦(\rho ^{},\rho ,z^{}z)`$ over $`(\rho ,z)`$ and $`(\rho ^{},z^{})`$ does vanish.
For cylindrically symmetric traps, the wave function grid as well as the integrated kernel region is as illustrated in Fig. 2. An accurate representation of the integrated kernel requires a quadrature operation over a much finer grid for each of the shaded regions surrounding the wave function grid. As is shown in Fig. 2, there are three different types of shaded regions, labeled as 1, 2, and 3. Both types of 1 and 2 are boundary terms, which are not needed since they are respectively at $`z=\pm L_z`$ and $`\rho =0,L_\rho `$ where either the wave function $`\psi (\rho ^{},z^{})`$ varnishes or the integration measure $`\rho ^{}𝑑\rho ^{}`$ vanishes. Therefore, we need only to compute the integration of kernel over type 3 element by defining a much finer grid and use standard numerical quadrature techniques. When $`(\rho ,z)(\rho ^{},z^{})`$, the integration reduces to a two-dimensional one which can be easily performed. On the other hand, a three-dimensional integration is needed when $`(\rho ,z)=(\rho ^{},z^{})`$. In this case we have to carefully implement the cut-off radius $`R_c`$.
Figure 3 compares the kernel with the coarse grained integrated kernel. Note the significantly different vertical scale.
### B Vortex states
Simple vortex states with quantized circulations can also be considered by writing the condensate wave function in the form
$`\psi (\stackrel{}{r})=|\psi (\stackrel{}{r})|e^{in\varphi },`$ (38)
with $`\varphi `$ the azimuthal angle with respect to z-axis. The corresponding Eq. (31) for $`|\psi (\stackrel{}{r})|`$ is then modified by the addition of $`n^2/\rho ^2`$ to $`\widehat{H}`$ of Eq. (32).
### C Numerical Results
#### 1 Ground state wave function
The ground state properties of trapped condensates with dipole interactions were first discussed by us in . Our basic findings are: 1) condensates become elongated along the direction of external field while shrank in the orthogonal radial direction; 2) for given values of $`P`$ and $`\lambda `$, there exists a maximum $`c_M(P,\lambda )`$ of a threshold field strength beyond which condensate collapse occurs. Figure 4 displays numerically computed $`c_M`$ for $`\lambda =1`$ at several different $`P`$ values. For comparison we also show the results from variational calculations to be discussed later.
The condensate collapse is mainly due to attractive dipole interactions along the direction of external field. To minimize its total energy, condensed atoms prefer to align along the attractive direction (z-axis), while narrowing its width along the radially repulsive direction. The collapsing occurs when the radial width eventually approaches zero with increasing external field strength .
Figure 5 shows a vortex state wave function for $`n=1`$. The effects of dipole interactions are similar to the ground state. As $`c()`$ increases, the vortex state will also collapse. Because of the zero density inside the vortex core, $`c_M`$ in this case is much larger than for the ground state.
#### 2 Comparing numerical solutions with TFA
A useful approximation for the ground state solution of NLSE (29) is the Thomas-Fermi approximation (TFA). When the interaction between atoms are repulsive, i.e. with a positive scattering length, the condensate is expected to increase its size as compared to the single atom ground state in the trap. With more atoms, the larger the condensate size, eventually the spatial directives, consequently the kinetic energy term becomes negligible. In this limit, TFA is used to find the ground state wave function by neglecting the kinetic energy term in Eq. (32). With a SIA interaction term, the solution simply takes the shape of the inverted trap potential $`V_t(\stackrel{}{r})`$. The nonlocal dipole interaction term, however, prevents a simple analytic solution even with the TFA.
Typical solutions to Eq. (29) with and without TFA for $`P=5000`$, $`\lambda =1`$, and $`c()=0.2,0.6,0.7`$ are compared in Fig. 6.
## IV time-dependent variational analysis
The time-dependent variational approach can also be used to analyze solutions of (29) . We start by identifying a Lagrangian density $``$
$``$ $`=`$ $`{\displaystyle \frac{i}{2}}\mathrm{}\left[\psi (\stackrel{}{r}){\displaystyle \frac{\psi ^{}(\stackrel{}{r})}{t}}\psi ^{}(\stackrel{}{r}){\displaystyle \frac{\psi (\stackrel{}{r})}{t}}\right]`$ (39)
$`+`$ $`{\displaystyle \frac{\mathrm{}^2}{2M}}|\psi (\stackrel{}{r})|^2+V_\mathrm{t}(\stackrel{}{r})|\psi (\stackrel{}{r})|^2`$ (40)
$`+`$ $`{\displaystyle \frac{u_0}{2}}|\psi (\stackrel{}{r})|^4+{\displaystyle \frac{u_2}{2}}|\psi (\stackrel{}{r})|^2{\displaystyle 𝑑\stackrel{}{r}^{}\frac{Y_{20}(\widehat{R})}{R^3}|\psi (\stackrel{}{r}^{})|^2}.`$ (41)
The NLSE can then be found from a minimization of the action
$`S={\displaystyle 𝑑\stackrel{}{r}𝑑t}.`$ (42)
To simplify the variational calculation, we restrict $`\psi `$ to a convenient family of trial functions and study the time evolution of the parameters that define the family. A natural choice is a Gaussian-like function first used in
$`\psi (x,y,z,t)=A(t){\displaystyle \underset{\eta =x,y,z}{}}e^{[\eta \eta _0(t)]^2/2w_\eta ^2+i\eta \alpha _\eta (t)+i\eta ^2\beta _\eta (t)},`$ (43)
where $`A`$ (complex amplitude), $`w_\eta `$ (width), $`\alpha _\eta `$ (slope), $`\beta _\eta `$ (curvature radius)<sup>-1/2</sup>, and $`\eta _0`$ (center of cloud) are variational parameters. This approach, pioneered by Perez-Garcia et al. , has since been successfully used for many studies of trapped condensates, a more recent application attempted to explain the anomalous behavior in the finite temperature excitation experiment .
Our goal here is to find equations governing the variational parameters. To this aim, we insert (43) into (41) and calculate an effective Lagrangian $`L`$ by integrating the Lagrangian density over all space coordinates
$`L=={\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\stackrel{}{r},`$ (44)
to arrive at
$`L`$ $`={\displaystyle \frac{i\mathrm{}N}{2}}\left({\displaystyle \frac{\dot{A^{}}}{A^{}}}{\displaystyle \frac{\dot{A}}{A}}\right)`$ (47)
$`+{\displaystyle \frac{N}{2}}{\displaystyle \underset{\eta }{}}\left[\left(\mathrm{}\dot{\beta _\eta }+{\displaystyle \frac{2\mathrm{}^2\beta _\eta ^2}{M}}+{\displaystyle \frac{1}{2}}M\nu ^2\lambda _\eta ^2\right)(w_\eta ^2+2\eta _0^2)+\left(\mathrm{}\dot{\alpha _\eta }+{\displaystyle \frac{2\mathrm{}^2\alpha _\eta \beta _\eta }{M}}\right)2\eta +{\displaystyle \frac{\mathrm{}^2}{2Mw_\eta ^2}}+{\displaystyle \frac{\mathrm{}^2\alpha _\eta ^2}{M}}\right]`$
$`+{\displaystyle \frac{N^2}{2\sqrt{8}\pi ^{3/2}w_xw_yw_z}}\left[u_0+u_2\sqrt{{\displaystyle \frac{5}{16\pi }}}{\displaystyle 𝑑\stackrel{}{r}\mathrm{exp}\left\{\underset{\eta }{}\frac{\eta ^2}{2w_\eta ^2}\right\}\frac{3\mathrm{cos}^2\theta 1}{r^3}}\right],`$
where we have used atom number conservation
$`N=\pi ^{3/2}|A(t)|^2w_x(t)w_y(t)w_z(t)=\mathrm{const}.`$ (48)
At this point, the variation calculation effectively has been reduced to a finite dimensional problem, i.e., to solve the Lagrange equations
$`{\displaystyle \frac{d}{dt}}\left({\displaystyle \frac{L}{\dot{q_j}}}\right){\displaystyle \frac{L}{q_j}}=0,`$ (49)
with the notation
$`q_j\{w_x,w_y,w_z,A,A^{},x_0,y_0,z_0,\alpha _x,\alpha _y,\alpha _z,\beta _x,\beta _y,\beta _z\}.`$ (50)
We find equations for the center of the condensate
$`\ddot{\eta _0}+\lambda _\eta ^2\nu ^2\eta _0=0,`$ (51)
and the condensate widths satisfy
$`\ddot{w_\eta }+\lambda _\eta ^2\nu ^2w_\eta ={\displaystyle \frac{\mathrm{}^2}{M^2w_\eta ^3}}{\displaystyle \frac{N}{4\sqrt{2}\pi ^{3/2}M}}{\displaystyle \frac{}{w_\eta }}\left[{\displaystyle \frac{1}{w_xw_yw_z}}\left(u_0+u_2\sqrt{{\displaystyle \frac{5}{16\pi }}}{\displaystyle 𝑑\stackrel{}{r}\mathrm{exp}\left\{\underset{\eta }{}\frac{\eta ^2}{2w_\eta ^2}\right\}\frac{3\mathrm{cos}^2\theta 1}{r^3}}\right)\right].`$ (52)
The rest of the parameters can be obtained from
$`\beta _\eta ={\displaystyle \frac{M\dot{w_\eta }}{2\mathrm{}w_\eta }},`$ (53)
and
$`\alpha _\eta ={\displaystyle \frac{M}{\mathrm{}}}\left(\dot{\eta _0}{\displaystyle \frac{\eta _0\dot{w_\eta }}{w_\eta }}\right).`$ (54)
It is convenient to introduce new dimensionless variables $`\tau =\nu t`$ and $`w_\eta =a_\mathrm{t}v_\eta `$. We then arrive at
$`{\displaystyle \frac{d^2}{d\tau ^2}}v_\eta +\lambda _\eta ^2v_\eta ={\displaystyle \frac{1}{v_\eta ^3}}P{\displaystyle \frac{}{v_\eta }}\left[{\displaystyle \frac{1}{v_xv_yv_z}}\left(1\sqrt{{\displaystyle \frac{5}{16\pi }}}c(){\displaystyle 𝑑\stackrel{}{r}\mathrm{exp}\left\{\underset{\eta }{}\frac{\eta ^2}{2v_\eta ^2}\right\}\frac{3\mathrm{c}\mathrm{o}\mathrm{s}^2\theta 1}{r^3}}\right)\right],`$ (55)
where $`P=\sqrt{2/\pi }Na_{\mathrm{sc}}/a_\mathrm{t}`$. This equation describes the motion of a particle with coordinates $`(v_x,v_y,v_z)`$ in an effective potential
$`U(v_x,v_y,v_z)`$ $`={\displaystyle \frac{1}{2}}\left(\lambda _x^2v_x^2+\lambda _y^2v_y^2+\lambda _z^2v_z^2\right)+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{v_x^2}}+{\displaystyle \frac{1}{v_y^2}}+{\displaystyle \frac{1}{v_z^2}}\right)`$ (57)
$`+{\displaystyle \frac{P}{v_xv_yv_z}}\left[1\sqrt{{\displaystyle \frac{5}{16\pi }}}c(){\displaystyle 𝑑\stackrel{}{r}\mathrm{exp}\left\{\underset{\eta }{}\frac{\eta ^2}{2v_\eta ^2}\right\}\frac{3\mathrm{cos}^2\theta 1}{r^3}}\right].`$
For a cylindrically symmetric trap with $`\lambda _x=\lambda _y=1`$, $`\lambda _z=\lambda `$, all integrals can be performed analytically to yield
$`{\displaystyle \frac{d^2}{d\tau ^2}}v+v`$ $`={\displaystyle \frac{1}{v^3}}+{\displaystyle \frac{P}{v^3v_z}}\left[1c()f(\kappa )\right],`$ (58)
$`{\displaystyle \frac{d^2}{d\tau ^2}}v_z+\lambda ^2v_z`$ $`={\displaystyle \frac{1}{v_z^3}}+{\displaystyle \frac{P}{v^2v_z^2}}\left[1c()g(\kappa )\right],`$ (59)
with $`v_x=v_y=v`$, $`\kappa =v/v_z`$, and
$`f(\kappa )`$ $`={\displaystyle \frac{\sqrt{5\pi }}{6(1\kappa ^2)^2}}\left[4\kappa ^47\kappa ^2+2+9\kappa ^4H(\kappa )\right],`$ (60)
$`g(\kappa )`$ $`={\displaystyle \frac{\sqrt{5\pi }}{3(1\kappa ^2)^2}}\left[2\kappa ^4+10\kappa ^2+19\kappa ^2H(\kappa )\right].`$ (61)
$`H(\kappa )\mathrm{tanh}^1\sqrt{1\kappa ^2}/\sqrt{1\kappa ^2}`$. The equilibrium widths are then determined by
$`v_0`$ $`={\displaystyle \frac{1}{v_0^3}}+{\displaystyle \frac{P}{v_0^3v_{z0}}}\left[1c()f(\kappa _0)\right],`$ (62)
$`\lambda ^2v_{z0}`$ $`={\displaystyle \frac{1}{v_{z0}^3}}+{\displaystyle \frac{P}{v_0^2v_{z0}^2}}\left[1c()g(\kappa _0)\right].`$ (63)
Using $`v_{z0}`$ and $`\kappa _0`$, Eq. (63) can be rewritten as
$`\kappa _0v_{z0}`$ $`={\displaystyle \frac{1}{\kappa _0^3v_{z0}^3}}+{\displaystyle \frac{P}{\kappa _0^3v_{z0}^4}}\left[1c()f(\kappa _0)\right],`$ (64)
$`\lambda ^2v_{z0}`$ $`={\displaystyle \frac{1}{v_{z0}^3}}+{\displaystyle \frac{P}{\kappa _0^2v_{z0}^4}}\left[1c()g(\kappa _0)\right].`$ (65)
In following discussion, we consider only $`a_{\mathrm{sc}}()>0`$ case, which implies both $`P>0`$ and $`c()>0`$.
### A Equilibrium widths
First, for simplicity, we assume that our system satisfy Thomas-Fermi limit ($`P1`$), then we can safely ignore the kinetic term and rewrite Eq. (65) for $`\kappa _0`$ in the following form
$`\kappa _0^2\left[1c()g(\kappa _0)\right]=\lambda ^2\left[1c()f(\kappa _0)\right].`$ (66)
This equation can be solved graphically. From Fig. 7, we first note both $`f(\kappa )`$ and $`g(\kappa )`$ are monotonically decreasing functions bounded between $`\sqrt{5\pi }/3`$ and $`2\sqrt{5\pi }/3`$, also the inequality $`f(\kappa )>g(\kappa )`$ holds for all $`\kappa >0`$. Then for all $`\kappa \lambda `$, we have $`\kappa ^2\left[1c()g(\kappa )\right]>\lambda ^2\left[1c()f(\kappa )\right]`$, therefore, if $`\kappa _0`$ is a solution of Eq. (66), it must satisfy $`\kappa _0<\lambda `$. Meanwhile, when $`c()=0`$, the solution for $`\kappa `$ in Eq. (66) is $`\kappa _0=\lambda `$. This then proves that no matter what the initial field-free condensate aspect ratio is, the condensate always become more prolate along the electric field direction, i.e. approximately as illustrated in Fig. 8, it expands along the field direction but shrinks in the orthogonal direction. As will be discussed in more detail later, the total condensate volume actually shrinks with increasing fields because of the attractive dipole interaction.
This result, first explained by us in terms of the minimization of total energy, is different from the conclusion reached in Ref. based on a force argument. This interesting feature has also been independently verified by numerical solutions based on the FFT algorithm adopted by Goral et. al .
We can also rewrite Eq. (66) as
$`\kappa ^2\lambda ^2=c()h(\kappa ),`$ (67)
with $`h(\kappa )=\kappa ^2g(\kappa )\lambda ^2f(\kappa )`$. Figure 9 shows its graphical solutions ($`\lambda =1`$) at several different $`c()`$ values.
First, we note that $`h(\kappa =0)=\sqrt{5\pi }\lambda ^2/3`$. Thus, as long as $`c()<3/\sqrt{5\pi }`$, there will be one and only one root. This result may not look absolutely clear from the figure because of plotting constraints, but it can be seen clearly form Eq. (66). We also find that as $`c()`$ increases, there may exist one, two, or zero roots for $`\kappa `$. Once $`\kappa _0`$ is known, one can easily find solutions $`v_0`$ and $`v_{z0}`$, whose stability can also be checked straightforwardly. It turns out that, in all our calculations, whenever only one root for $`\kappa `$ occurs its corresponding solution for $`v_0`$ and $`v_{z0}`$ is always stable. If there are two roots for $`\kappa `$, the solution for $`v_0`$ and $`v_{z0}`$ corresponding to the smaller $`\kappa _0`$ root is a saddle point, thus always unstable, while the other is always stable.
We now consider the $`\lambda `$ dependence of the condensate property. Figure 10 shows the function $`h(\kappa )`$ at several different $`\lambda `$ values. We see that as $`\lambda `$ increases, the maximum value of function $`h(\kappa )`$, $`h_{\mathrm{max}}`$ also increases, and $`h_{\mathrm{max}}=0`$ at $`\lambda _c=5.170169`$. For those $`\lambda `$ values corresponding to a negative $`h_{\mathrm{max}}`$, no root of $`\kappa `$ is found if $`c()`$ becomes sufficiently large. But when $`\lambda \lambda _c`$, at least one root for $`\kappa `$ exists no matter how large a $`c()`$ is. This means the condensate cannot simply collapse even at these very large electric field strengths. Physically this implies the increased stability of condensate inside an electric field with increasing values of $`\lambda `$. A condensate of a pancake shape is more stable than one with a cigar shape. This can be simply understood from the following argument; The collapse of a condensate under electric fields is due to the alignment of atoms along the attractive z-axis direction. A larger $`\lambda `$ value prevents such alignment which increases both kinetic as well as trap potential energy, hence increases the stability. In Fig. 11, we display $`c_M`$ as a function of $`\lambda `$, which separates the stable and unstable regions. Rigorous numerically calculations, on the other hand, find that condensate can still collapse even when $`\lambda \lambda _c`$. For example, we found collapse occurs when $`P5000`$ with $`\lambda >7`$, \[$`c()>2.0`$\]. This difference maybe due to the choice of a simple Gaussian shaped variation function (43).
When the system is not in the Thomas-Fermi limit, we have to first solve for $`\kappa _0`$ from equation
$`\kappa _0^2\lambda ^2=c()[\kappa _0^2g(\kappa _0)\lambda ^2f(\kappa _0)]+\left\{{\displaystyle \frac{(\kappa _0^4\lambda _0^2)^4}{P^4\kappa _0^2}}\left[1\kappa _0^2c()(g(\kappa _0)\kappa _0^2f(\kappa _0))\right]\right\}^{1/5},`$ (68)
and then using
$$v_{z0}=P\frac{\kappa _0^2\left[1c()g(\kappa _0)\right]\lambda ^2\left[1c()f(\kappa _0)\right]}{\lambda ^2\kappa _0^4},$$
and $`v_0=\kappa _0v_{z0}`$ to find the equilibrium widths. In this case, we find possibilities for one, two, three, and four or no roots of $`\kappa _0`$ depending on values of $`P`$, $`c()`$, and $`\lambda `$. The stability conditions for these roots, however, are similar to the TFA case as discussed before—there is at most only one stable solution. In Fig. 12, we present $`\kappa _0`$, $`v_0`$, and $`v_{z0}`$ as functions of $`c()`$. We see that $`v_0`$ decreases with increasing $`c()`$, and the condensate collapses when $`v_0`$ goes to zero. Figure 12(c) displays the dependence of the condensate volume on $``$, where the volume is defined as the produce of its effective widths in three separate dimensions. In terms of absolute values, numerically calculated volume (averaged width) differs from the variational result by a few times due to the multiplying effects of three widths. The shrinking volume increases the condensate density, which in turn can significantly increase the three body loss process, providing a potentially useful diagnosis tool .
Figure 13 shows $`c_M`$ as a function of $`P`$ for different $`\lambda `$ values. We see that, for smaller $`P`$, a condensate with a small $`\lambda `$ can be more stable than condensates with larger $`\lambda `$; while for larger $`P`$, a condensate with a larger $`\lambda `$ is always more stable than condensates with smaller $`\lambda `$. This also confirmed by both variational (TFA) and numerical calculations.
After we first submitted this paper, a paper on the same topic become available . We therefore compare our results with those from in the next few figures. Figure 14 displays the change of aspect ratio of the ground state for two extreme values of $`\lambda =0.1,10`$ and for a small value of $`P=5`$, and can be directly compared with the Fig. 3 of . We note that essentially the same results were obtained as in presumably because their neglect of the s-wave interaction simply corresponds to $`P=0`$ of our more general results.
More interestingly, we show in Fig. 15 for a larger value of $`P=500`$. We find the aspect ratio now changes in the opposite direction with increasing dipole interaction strength. This reversal of aspect ratio with increasing values of $`P`$ (due to increasing in atom numbers or s-wave scattering length $`a_{\mathrm{sc}}`$) is due precisely to the physics of minimizing the total system free energy discussed earlier in the TFA. This phenomena was not observed in the simpler model of Ref. . We also find that $`\lambda _c`$ remains virtually independent of $`P`$ at the same value as in the TFA: 5.170169, consistent with Ref. .
### B Evolution of widths
The evolution of condensate widths are found by numerically integrating Eq. (59). Assuming initially $`c()=0`$, $`v(0)=v_0`$, $`v_z(0)=v_{z0}`$, and $`\dot{v}(0)=\dot{v}_z(0)=0`$, we can apply an electric field suddenly or slowly for $`t>0`$. Under stable conditions for the widths, we choose to apply electric field suddenly. Otherwise the electric is increased linearly within a ramp-up time, and then kept constant. First, for the stable case, Figure 16 shows condensate widths evolution up to $`t=30`$ (it has been calculated up to $`t=1000`$). When $`c()=0`$, the condensate remains at its initially equilibrium state, and the widths are unchanged with time. When $`c()0`$, condensate widths oscillate with time, and prolonged numerical propagation indicates that we always have $`v>0`$ and $`v_z>0`$, i.e. the condensate is stable.
We also see from these figures that the oscillation amplitudes increased with increasing $`c()`$. Then finally at some stage, we could arrive at $`v<0`$ or $`v_z<0`$, signaling the condensate collapse. Figure 17 indeed displays such cases when a linear ramp-on of the external electric field is applied.
### C Small amplitude shape oscillations
Once the equilibrium widths are found from numerically solving Eqs. (65), small amplitude oscillations can be studied by evaluating the matrix of the second order derivatives of the equivalent potential $`U(v_x,v_y,v_z)`$ Eq. (57). We find that it takes the following symmetric form
$`\left(\begin{array}{ccc}U_{11}& U_{12}& U_{13}\\ U_{12}& U_{11}& U_{13}\\ U_{13}& U_{13}& U_{33}\end{array}\right),`$ (72)
where $`U_{ij}=U_{ji}`$ due to nature of commuting derivative operations with different coordinates, and $`U_{11}=U_{22}`$ and $`U_{13}=U_{23}`$ due to the cylindrical symmetry. We find the oscillation frequencies to be
$`\nu _1`$ $`=\sqrt{U_{11}U_{12}},`$ (73)
$`\nu _{2,3}`$ $`={\displaystyle \frac{1}{\sqrt{2}}}\left[U_{11}+U_{12}+U_{33}\pm \sqrt{U_{11}^2+U_{12}^2+U_{33}^2+8U_{13}^2+2U_{11}U_{12}2U_{11}U_{33}2U_{12}U_{33}}\right]^{1/2},`$ (74)
where the expression for the matrix elements $`U_{ij}`$ are listed in Appendix A. Typical results and mode structure identifications are given in Fig. 18. We see that mode 1 and mode 3 are doubly degenerate when $`c()=0`$. This is due to the additional symmetry $`U_{11}=U_{33}`$ and $`U_{12}=U_{13}`$ for $`\lambda =1`$.
## V conclusion
In conclusion, we have performed a detailed study of trapped condensates with dipole interactions. We have developed a general scheme for constructing effective pseudo-potentials for anisotropic interactions , which guarantees the agreement between the first order Born scattering amplitude from the pseudo-potential and the complete scattering amplitude obtained from a multi-channel collision calculation. Our theory has been applied to the study of induced electric dipole interactions and can also be directly extended to magnetic dipole interactions as well as permanent electric dipole interactions of trapped molecules .
Finally we provide a reality check for prospects of experimental observations of the electric field induced interaction effects. Though the required fields are relatively high, there are evidences they can be created with current laboratory technology. In Ref. fields of upto $`2\times 10^5`$ (V/cm) were used to slow a molecular beam. Gould used fields upto $`4.6\times 10^5`$ (V/cm) in the measurement of atomic tensor polarizability, while Marrus et al. reported fields upto $`10^6`$ (V/cm). What is perhaps most encouraging is a recent experiment for cooling molecule beams with time-dependent (adiabatic from the view point of atomic internal dynamics) fields of upto $`10^7`$ (V/cm) . We also note that at the high fields being discussed in this paper, the tunneling ionization of atoms remain infinitesimally small .
This work is supported by the U.S. Office of Naval Research grant No. 14-97-1-0633 and by the NSF grant No. PHY-9722410. The computation of this work is partially supported by NSF through a grant for the ITAMP at Harvard University and Smithsonian Astrophysical Observatory.
## A U-matrix elements
After tedious calculations, we find that
$`U_{11}=1+{\displaystyle \frac{3}{v_0^4}}+{\displaystyle \frac{P}{v_0^4v_{z0}}}\left[2{\displaystyle \frac{\sqrt{5\pi }c()}{24(v_{z0}^2v_0^2)^3}}\left(32v_0^6+141v_0^4v_{z0}^254v_0^2v_{z0}^4+16v_{z0}^69(11v_0^2+4v_{z0}^2)v_0^4H(v_0/v_{z0})\right)\right],`$ (A1)
$`U_{33}=\lambda ^2+{\displaystyle \frac{3}{v_{z0}^4}}+{\displaystyle \frac{P}{v_0^2v_{z0}^3}}\left[2{\displaystyle \frac{\sqrt{5\pi }c()}{3(v_{z0}^2v_0^2)^3}}\left(4v_0^612v_0^4v_{z0}^2+51v_0^2v_{z0}^4+2v_{z0}^69(v_0^2+4v_{z0}^2)v_0^2v_{z0}^2H(v_0/v_{z0})\right)\right],`$ (A2)
$`U_{12}={\displaystyle \frac{P}{v_0^4v_{z0}}}\left[1{\displaystyle \frac{\sqrt{5\pi }c()}{24(v_{z0}^2v_0^2)^3}}\left(16v_0^6+51v_0^4v_{z0}^230v_0^2v_{z0}^4+8v_{z0}^645v_0^6H(v_0/v_{z0})\right)\right],`$ (A3)
and
$`U_{13}={\displaystyle \frac{P}{v_0^3v_{z0}^2}}\left[1{\displaystyle \frac{\sqrt{5\pi }c()}{6(v_{z0}^2v_0^2)^3}}\left(4v_0^636v_0^4v_{z0}^215v_0^2v_{z0}^4+2v_{z0}^6+45v_0^4v_{z0}^2H(v_0/v_{z0})\right)\right].`$ (A4)
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# DETECTION OF GRAVITATIONAL WAVES FROM INSPIRALING COMPACT BINARIES USING NON-RESTRICTED POST-NEWTONIAN APPROXIMATIONS
## 1 Introduction
and future space-based interferometers like LISA $`^\mathrm{?}`$.
In this paper, we consider the in-spiral phase of the coalescence of binary systems in circular orbit using post-Newtonian (PN) approximations to general relativity $`^\mathrm{?}`$.
The detection of in-spiral signals is carried out by cross-correlating the detector output with a discrete set of filters $`^{\mathrm{?},\mathrm{?}}`$, usually computed within the so-called restricted post-Newtonian approximation $`^\mathrm{?}`$: PN corrections are taken into account in the phase of the wave-form, whereas the amplitude is retained at the lowest Newtonian order. Thus, one discards all multipole components except the quadrupole one. Such simplification of the filter structure is believed to have negligible effects on the detection performances and is not expected to affect appreciably the statistical errors in the estimation of the source parameters.
Here we investigate the effects of the introduction of PN corrections also to the amplitude of the wave-form and discuss some implications for signal detection and parameter estimation. For sake of simplicity we will assume negligible spins and, as it is always the case for ground based experiments, circular orbits.
## 2 The wave-form
The signal produced at the output of an interferometric detector by a gravitational wave of polarization amplitudes $`h_+`$ and $`h_\times `$ can be written as
$$h(t)=F_+h_+(t)+F_\times h_\times (t),$$
(1)
where $`F_+`$ and $`F_\times `$ are the so-called beam pattern functions of the detector $`^\mathrm{?}`$; they depend on the location of the source in the sky $`(\theta ,\varphi )`$ and the polarization angle $`\psi `$. If we consider the in-spiral of a binary system of masses $`m_1`$ and $`m_2`$, $`h_+`$ and $`h_\times `$ read $`^\mathrm{?}`$
$$h_{+,\times }=\frac{2m\eta }{r}x\left\{H_{+,\times }^{(0)}+x^{1/2}H_{+,\times }^{(1/2)}+xH_{+,\times }^{(1)}+x^{3/2}H_{+,\times }^{(3/2)}+x^2H_{+,\times }^{(2)}+\mathrm{}\right\},$$
(2)
where $`x(m\omega )^{2/3}`$, $`\omega `$ is the system orbital frequency and $`r`$ the source distance; $`m=m_1+m_2`$, $`\mu =m_1m_2/m`$, $`\eta =\mu /m`$ and $`=\mu ^{3/5}m^{2/5}=m\eta ^{3/5}`$ are the total mass, the reduced mass, the symmetric mass ratio and the chirp mass, respectively.
The lower terms of the PN expansion for the plus and cross polarization are given by $`^\mathrm{?}`$
$`H_+^{(0)}`$ $`=`$ $`(1+c^2)\mathrm{cos}\mathrm{\Phi },`$ (3)
$`H_+^{(1/2)}`$ $`=`$ $`s{\displaystyle \frac{\sqrt{14\eta }}{8}}\left[(5+c^2)\mathrm{cos}({\displaystyle \frac{1}{2}}\mathrm{\Phi })9(1+c^2)\mathrm{cos}({\displaystyle \frac{3}{2}}\mathrm{\Phi })\right],`$ (4)
$`H_\times ^{(0)}`$ $`=`$ $`2c\mathrm{sin}\mathrm{\Phi },`$ (5)
$`H_\times ^{(1/2)}`$ $`=`$ $`{\displaystyle \frac{3}{4}}sc\sqrt{14\eta }\left[\mathrm{sin}({\displaystyle \frac{1}{2}}\mathrm{\Phi })3\mathrm{sin}({\displaystyle \frac{3}{2}}\mathrm{\Phi })\right],`$ (6)
where $`c=\mathrm{cos}\iota `$ and $`s=\mathrm{sin}\iota `$; $`\iota `$ is the angle between the direction of the source and the orbital angular momentum, and $`\mathrm{\Phi }`$ is twice the orbital phase.
If one considers only the first term in (2), corresponding to the Newtonian one, one gets the restricted PN approximation. Our goal is to study the full 2 PN wave-form; here we will present some preliminary results, where the 0.5 PN corrections to the amplitude are taken into account, keeping however the phase at 2 PN order. Although this is a simplified and, to some extent, arbitrary choice of signal, all new features of the wave-form are introduced, in particular more information about the masses and the position of the source.
Considering the amplitude through 0.5-PN order, the GW output at the detector can be written as
$$h(t)=\frac{2m\eta }{r}(m\pi F)^{2/3}\left\{h^{(0)}(t)+(m\pi F)^{1/3}h^{(1/2)}(t)\right\},$$
(7)
where $`F`$ is the quadrupole gravitational wave frequency, i.e., $`d\mathrm{\Phi }/dt=2\pi F`$, and
$$h^{(0)}(t)=\sqrt{F_+^2(1+c^2)^2+F_\times ^24c^2}\mathrm{cos}(\mathrm{\Phi }+\phi _{(0)}),$$
(8)
$$\phi _{(0)}=\mathrm{arctan}\left\{\frac{2cF_\times }{(1+c^2)F_+}\right\},$$
(9)
$`h^{(1/2)}(t)=s{\displaystyle \frac{\sqrt{14\eta }}{4}}`$ $`[\sqrt{F_+^2{\displaystyle \frac{(5+c^2)^2}{4}}+F_\times ^29c^2}\mathrm{cos}({\displaystyle \frac{1}{2}}\mathrm{\Phi }+\phi _{(1/2)})`$
$`+{\displaystyle \frac{9}{2}}\sqrt{F_+^2(1+c^2)^2+F_\times ^24c^2}\mathrm{cos}({\displaystyle \frac{3}{2}}\mathrm{\Phi }+\phi _{(0)})],`$
$$\phi _{(1/2)}=\mathrm{arctan}\left\{\frac{6cF_\times }{(5+c^2)F_+}\right\}.$$
(11)
The Fourier transform of $`h(t)`$, calculated using the stationary phase approximation $`^{\mathrm{?},\mathrm{?}}`$, reads:
$$\stackrel{~}{h}(\nu )=\sqrt{F_+^2(1+c^2)^2+F_\times ^24c^2}\sqrt{\frac{5\pi }{96}}\frac{(\pi \nu )^{7/6}}{r}^{5/6}\mathrm{exp}\left[i\left(2\pi \nu t_c+\mathrm{\Xi }(\nu )\phi _{(0)}\frac{\pi }{4}\right)\right]\mathrm{\Lambda }$$
(12)
where
$`\mathrm{\Lambda }=`$ $`1+{\displaystyle \frac{s}{4}}\sqrt{14\eta }\sqrt{{\displaystyle \frac{F_+^2(5+c^2)^2/4+F_\times ^29c^2}{F_+^2(1+c^2)^2+F_\times ^24c^2}}}(\pi m\nu )^{1/3}\left({\displaystyle \frac{1}{2}}\right)^{1/3}\times `$ (13)
$`\mathrm{exp}\left[i\left({\displaystyle \frac{1}{2}}\mathrm{\Xi }(2\nu )\mathrm{\Xi }(\nu )+\phi _{(0)}\phi _{(1/2)}\right)\right]`$
$`s{\displaystyle \frac{9}{8}}\sqrt{14\eta }(\pi m\nu )^{1/3}\left({\displaystyle \frac{3}{2}}\right)^{1/3}\mathrm{exp}\left[i\left({\displaystyle \frac{3}{2}}\mathrm{\Xi }\left({\displaystyle \frac{2}{3}}\nu \right)\mathrm{\Xi }(\nu )\right)\right],`$
and
$`\mathrm{\Xi }(\nu )=\varphi _c+{\displaystyle \frac{3}{4}}(8\pi \nu )^{5/3}`$ $`[1+{\displaystyle \frac{20}{9}}({\displaystyle \frac{743}{336}}+{\displaystyle \frac{11}{4}}\eta )(\pi m\nu )^{2/3}16\pi (\pi m\nu )`$
$`+10({\displaystyle \frac{3058673}{1016064}}+{\displaystyle \frac{5429}{1008}}\eta +{\displaystyle \frac{617}{144}}\eta ^2)(\pi m\nu )^{4/3}].`$
In the restricted PN approximation $`\mathrm{\Lambda }=1`$. Now it contains two additional contributions related to the 0.5 PN corrections to the amplitude. Notice the dependency of these two terms on $`\mathrm{sin}\iota `$ and $`\sqrt{14\eta }`$: the departure from the value $`\mathrm{\Lambda }=1`$ increases as $`\iota \pi /2`$ and $`\eta 0`$.
## 3 Formalism
We denote by $`h`$ the “true” GW signal and $`u(𝜽)`$ the family of templates, as a function of the parameter vector $`𝜽=(t_c,\varphi _c,𝝀)`$. The signal to noise ratio SNR, for optimal filtering, is defined as $`^\mathrm{?}`$
$$\mathrm{SNR}=\sqrt{\left(h|h\right)},$$
(15)
where $`(|)`$ denotes the usual inner product. The fraction of SNR obtained by cross-correlating a template $`u(𝜽)`$ with $`h`$ is given by the ambiguity function
$$𝒜(𝜽)=\frac{\left(h|u(𝜽)\right)}{\sqrt{\left(h|h\right)\left(u(𝜽)|u(𝜽)\right)}},$$
(16)
which depends on the choice of $`𝜽`$. The maximum of the ambiguity function over the whole parameter space is defined as the fitting factor $`^\mathrm{?}`$
$$FF=\begin{array}{c}\mathrm{max}\\ 𝜽\end{array}\frac{\left(h|u(𝜽)\right)}{\sqrt{\left(h|h\right)\left(u(𝜽)|u(𝜽)\right)}}.$$
(17)
The fitting factor is a measure of how well any chosen family of templates fits the signal $`h`$. The maximization of the ambiguity function over the extrinsic parameters $`\varphi _c`$ and $`t_c`$, phase and time of coalescence, is the so-called match
$$M(𝝀_1,𝝀_2)=\begin{array}{c}\mathrm{max}\\ \varphi _c,t_c\end{array}\frac{\left(h(𝝀_1)|u(𝝀_2)\right)e^{i(2\pi ft_c\varphi _c)}}{\sqrt{\left(h(𝝀_1)|h(𝝀_1)\right)\left(u(𝝀_2)|u(𝝀_2)\right)}}.$$
(18)
For the set up of the bank of filters, one sets a minimal match $`^{\mathrm{?},\mathrm{?}}`$ as the match between signal and template in the case where the signal is equidistant from all the nearest templates.
## 4 Results
In the results presented in the following, the signal $`h`$ is computed according to Eq. (13)-(15), whereas the template wave-forms are calculated within the usual restricted PN approximation, corresponding to $`\mathrm{\Lambda }=1`$ in Eq. (14). The noise curve is the one corresponding to the initial LIGO configuration $`^\mathrm{?}`$.
In table 1, we give the fitting factors and the ratios of SNR, $`=\sqrt{(u|u)/(h|h)}`$, for the pair of masses 0.1-10 $`M_{}`$ and 1.4-10 $`M_{}`$, for different orientation and polarization angles. We note that the fitting factor reaches a minimum for $`\iota =\pi /2`$ as expected, and monotonically increases as $`\iota 0`$ or $`\iota \pi `$; it depends rather weakly on $`\theta `$, $`\varphi `$ and $`\psi `$.
For a fix position in the sky and $`\iota =\pi /2`$, we calculate then the fitting factors for different mass pairs. The fitting factor varies from 0.87 to 1.0; $`FF`$ gets smaller as $`m`$ increases and/or $`\eta `$ decreases, see table 2.
It is now interesting to investigate the loss of SNR if one uses a restricted PN bank of filters to detect signals that include PN amplitude terms. The discrete mesh of filters is normally generated in such a way that any signal in the restricted PN plane produces a match $`M(u(𝝀_1),u(𝝀_2))`$ is always larger than a minimum value, usually set to 0.97. Assuming that the “true” GW signal, $`h`$, is includes PN corrections to the amplitude, so that it lies outside the template space, we want to quantify the match $`M(h(𝝀_1),u(𝝀_2))`$ between the PN signal and the nearest restricted PN template.
In table 3, we present the results obtained for a system of masses 0.1-10$`M_{}`$ and different choices of $`\iota `$. For each case we calculate the fitting factor between the signal and the family of templates; then, using different parameter values $`𝝀_1`$ and $`𝝀_2`$, we calculate the match in the restricted PN plane $`M(u(𝝀_1),u(𝝀_2))`$, and, for the same parameters, the match between the signal and the template $`M(h(𝝀_1),u(𝝀_2))`$.
What we observe is that the “real match” can be approximated as $`M(h(𝝀_1),u(𝝀_2))FF\times M(u(𝝀_1),u(𝝀_2))`$.
We turn now attention to the issue of estimating the source parameters. It is important to notice that for a waveform computed taking into account PN corrections to the amplitude, two additional parameters are involved; our choice corresponds to $`F_\times /F_+`$ and $`\mathrm{cos}\iota `$. We adopt here the standard variance-covariance matrix analysis $`^{\mathrm{?},\mathrm{?}}`$, although it can underestimate the statistical errors in the limit of low SNR $`^{\mathrm{?},\mathrm{?}}`$; such well-known problem is beyond the purposes of this work and for sake of simplicity we will remain into the usual frame of the computation of the Fisher information matrix.
In table 4, we compare the errors for SNR=10 between the restricted and the “non-restricted” PN approximation. The results refer to a system of 1.4-20$`M_{}`$ and $`\theta =\pi /6`$, $`\varphi =\pi /4`$, $`\psi =\pi /4`$, and $`\iota =\pi /3`$. The use of more accurate wave-forms leads to smaller errors (by roughly a factor $`2`$) in the determination of the masses; however the information about $`\iota `$ and $`F_\times /F_+`$ remain very poor.
In table 5 we provide the correlation coefficients for the same parameter choice. We notice that the correlation coefficients are also smaller and that the two new parameters have small correlations with the other ones. Notice that the amplitude is now correlated with the other parameters while it is not the case for the restricted PN case.
## References
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